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# Tunnelling current and emission spectrum of a single electron transistor under optical pumping
## I Introduction
The tunnelling current characteristics of a single electron transistor (SET) have been extensively studied and interesting physical phenomena, such as Coulomb blockade and Kondo effect have been found.<sup>1-7</sup> The main structure of an SET consists of three electrodes (source, drain and gate) and the active region consists of one or more quantum dots (QDs). Although the physics of Kondo effect is more profound than that of Coulomb blockade, the main device application of SET is based on the Coulomb blockade effect. The tunnelling current of an SET is sensitive to the size of the QD. The charging energy, $`\mathrm{\Delta }U`$, is typically larger than the energy level spacing, $`\mathrm{\Delta }E`$, caused by the quantum confinement in large QDs. Consequently, the tunnelling current displays a periodic Coulomb oscillation with respect to the gate voltage. This feature of periodic Coulomb oscillation is a consequence of homogeneity of $`\mathrm{\Delta }U`$ and negligible $`\mathrm{\Delta }E`$. However, such a Coulomb oscillation was observed only at very low temperature due to the small charging energy. Note that the charging energy must be greater than the thermal energy and tunnelling rate multiplied by $`\mathrm{}`$ in order for an SET to operate properly. For application in high temperature SETs, nanoscale QDs with large charging energy are desired.
To date two major kinds of nanoscale QDs have been used to construct SETs: the Si/Ge QDs and the InAs/GaAs self-assembled quantum dots (SAQDs). Silicon (Si) or germanium (Ge) QDs can be miniaturized to reach nanometer scale with advanced fabrication technology. Nanoscale Si/Ge SETs typically display non-periodic Coulomb oscillations at room temperatures,<sup>8-11</sup> as a result of the large charging energy ($`\mathrm{\Delta }U`$) and the unequally spaced values of $`\mathrm{\Delta }E`$, which are comparable to $`\mathrm{\Delta }U`$. Even though Si/Ge SETs can operate at room temperatures, it is difficult to gian full understanding of their tunnelling current characteristics because of the multi-valleyed nature of Si/Ge and the unknown Si/Ge-SiO<sub>2</sub> interface properties. Furthermore, the electron-hole recombination rate in Si/Ge QDs is not large enough for useful application in optoelectronics, since Si and Ge are indirect-band-gap semiconductors.<sup>12</sup>
The InAs/Gas SAQDs, on the other hand, are better understood, and because InAs is a direct-band-gap semiconductors, SETs constructed from InAs/Gas SAQDs can have useful application in optoelectronic devices such as single-photon generators and single-photon detectors. Recently, InAs/GaAs self-assembled quantum dot (SAQD) embedded in a p-n junction has been employed to generate single-photon emission for application in quantum cryptography.<sup>13-15</sup> Another way to achieve single-photon generation is to use isolated SAQD with optical pumping. Such studies have received a great deal of attention recently.<sup>16-21</sup> The spontaneous emission spectrum of a QD typically exhibits coexisting sharp emission peaks, which have been attributed to the electron-hole recombination in the exciton, trion and biexciton formed in the QD. The antibunching feature of the emitted single photon in this kind of device has been demonstrated.<sup>13,14</sup> In addition, gate electrodes have been used in some experiments to manipulate the emission spectrum.<sup>22-24</sup>.
Many theoretical efforts based on the Anderson model and Keldysh Green’s function method have been devoted to the studies of transport properties of SETs, including the dc and ac tunnelling current.<sup>5-7</sup> The studies on the ac tunnelling current of SET mainly focus on the photon-side band phenomena caused by the microwave pumping<sup>5,7</sup>. Only a handful theoretical studies considered the optical properties of SETs. Recently, we have studied the intraband transitions of SETs made of InAs/GaAs SAQDs for application as an infrared photodetector with operating wavelength near $`10\mu m`$.<sup>25,26</sup> On the other hand, the interband transitions in a QD have been adopted frequently to realize single-photon generation. Theoretical studies of the electronically driven single-photon generator (SPG) have been reported,<sup>27,28</sup> whereas theoretical studies on optically pumped SETs (including the coupling with leads) are still lacking. It is the purpose of this paper to provide a theoretical analysis for such a system. Furthermore, the measurement of exciton binding energy in a QD remains illusive, since in a typical photoluminescence measurement or the emission spectrum of a single QD, one can only observe the exciton recombination, but not the free electron-hole recombination. The exciton binding energy is an important physical quantity as it provides information about the effective dielectric constant and the charge distribution of electron and hole in the QD. In our current theoretical study, we will show that via the measurement of the tunnelling current of an SET made of a nanoscale QD under optical pumping one can determine the exciton binding energy of a QD unambiguously.
The schematic diagram for the SET system considered is shown in Fig. 1. Here a nanoscale QD is embedded in a n-i-n junction. This is different from the electrically driven SPG, where a QD is embedded in a p-n junction. Therefore, the physical process considered here is also quite diffrent. Using optical pumping to create electron-hole pairs in the SET, one can not only electrically manipulate the single-photon emission spectrum arising from exciton complexes in the QD, but also optically control the tunnelling current of the SET. In this paper, we will consider both the hole-assisted tunnelling current due to optical pumping and the electrode-controlled spontaneous emission spectrum.
Our calculation is based on the Keldysh Green’s function approach<sup>29</sup> within the Anderson model for a two-level system. A effective-mass model, which takes into account the anisotropy and inhomogeneity, is used to estimate the inter-particle Coulomb energies. We find that the optical excitation creates holes in the QD, which provide new channels (via the electron-hole interaction) for the electron to tunnel from the emitter to the collector. As a consequence, an electron can tunnel through the QD via four additional channels, characterized by the exciton, positive trion, negative trion, and biexciton states. Each addition channel can generate a new plateau (or oscillatory peaks) in the tunnelling current characteristics in addition to the typical plateau (or peaks) caused by the Coulomb blockade effect. This gives rise to a rich tunnelling current characteristics and it can be used to determine the exciton binding energy.
This paper is organized as the following. In section II, we derive the tunnelling current of SET under optical pumping. In section III, we calculate the polarization needed for describing the spontaneous emission spectrum. In section IV, we present calculations of the inter-particle Coulomb interactions within a simple but realistic effective-mass model. In section V, we discuss the results of our numerical calculations on the tunnelling current and spontaneous emission spectrum and demonstrate that the exciton binding energy can be extracted via the measurement of tunnelling current. Finally, a summary and concluding remarks are presented in section VI.
## II Tunnelling current
The system of interest is shown in Fig. 1, which consists of a single InAs/GaAs self-assembled quantum dot (SAQD) sandwiched between two n-doped GaAs leads. Electrons are allowed to tunnel from the left lead (emitter) to the right lead (collector) under the influence of an optical field. The Hamiltonian for the system is given by
$`H`$ $`=`$ $`{\displaystyle \underset{\sigma ,i=e,h}{}}E_id_{i,\sigma }^{}d_{i,\sigma }+{\displaystyle \underset{𝐤,\sigma ;\mathrm{}=L,R}{}}ϵ_𝐤c_{𝐤,\sigma ;\mathrm{}}^{}c_{𝐤,\sigma ;\mathrm{}}`$ (1)
$`+`$ $`{\displaystyle \underset{\mathrm{}=L,R;k,\sigma }{}}(V_{𝐤,\sigma ;\mathrm{}}c_{𝐤,\sigma ;\mathrm{}}^{}d_{e,\sigma }+V_{𝐤,\sigma ;\mathrm{}}^{}d_{e,\sigma }^{}c_{𝐤,\sigma ;\mathrm{}})`$
$`+`$ $`{\displaystyle \underset{𝐤,\sigma }{}}(\lambda _0e^{i\omega _0t}b_{e,𝐤,\sigma }b_{h,𝐤,\sigma }+\lambda _0^{}e^{\omega _0t}b_{h,𝐤,\sigma }^{}b_{e,𝐤,\sigma }^{})`$
$`+`$ $`{\displaystyle \underset{i,𝐤,𝐪,\sigma }{}}(g_{i,𝐪}Ab_{i,𝐤,\sigma }^{}d_{i,\sigma }+g_{i,𝐪}^{}A^{}d_{i,\sigma }^{}b_{i,𝐤,\sigma })`$
$`+`$ $`{\displaystyle \underset{𝐐}{}}(\lambda e^{i\omega t}a^{}d_{e,\sigma }d_{h,\sigma }+\lambda ^{}e^{i\omega t}ad_{h,\sigma }^{}d_{e,\sigma }^{}),`$
where the first term describes electrons in the QD. We assume that the quantum confinement effect is strong for the small InAs QD considered here. Therefore, the energy spacings between the ground state and the first excited state for electrons and holes, $`\mathrm{\Delta }E_e`$ and $`\mathrm{\Delta }E_h`$, are much larger than the thermal energy, $`k_BT`$, where $`k_B`$ and $`T`$ denote the Boltzmann constant and temperature. Only the ground state levels for electrons and holes, $`E_e`$ and $`E_h`$, are considered in the first term. The second term describes the kinetic energies of free carriers in the electrodes. Note that in the current setup, the gate electrode does not provide any carriers, but merely controls the energy levels of the QD. The third term describes the coupling between the QD and the leads. Note that only the conduction band of QD is coupled with the electrodes. The fourth term describes the interband optical pumping with a frequency $`\omega _0`$, which is in resonance with the transition energy for an electron-hole pair in the wetting layer. We treat the electromagnetic field as a semiclassical field. The fifth term describes the optical phonon assisted process for electrons in the wetting layer to relax into the InAs QD. The last term describe the coupling of the QD with the electromagnetic field of frequency $`\omega `$. $`\lambda \mu _r`$ is the Rabi frequency, where $`\mu _r=<f|𝐫|𝐢>`$ is the matrix element for the optical transition and $``$ is electric field per photon. $`a^{}(A^{})`$ and $`a`$ ($`A`$) denote the creation and annihilation operators of a photon (phonon), respectively.
Due to the large strain-induced splitting between the heavy-hole and light-hole bands for typical QDs, we only have to consider the heavy hole band (with $`J_z=\pm 3/2`$) and ignore its coupling with light-hole band caused by the QD potential. Thus, we can treat the heavy hole as a spin-1/2 particle with $`\sigma =,`$ representing $`J_z=\pm 3/2`$. This treatment is convenient for algebraic manipulations in the calculation of Green’s functions. Because the effect of inter-particle Coulomb interactions is significant in small semiconductor QDs, we take into account the electron Coulomb interactions and electron-hole Coulomb interactions by adding the following terms
$$H_I=\underset{i,\sigma }{}U_{i,i}d_{i,\sigma }^{}d_{i,\sigma }d_{i,\sigma }^{}d_{i,\sigma }+\underset{ij;\sigma ,\sigma ^{}}{}U_{i,j}d_{i,\sigma }^{}d_{i,\sigma }d_{j,\sigma }^{}d_{j,\sigma },$$
(2)
where $`d_{i,\sigma }^{}`$ and $`d_{i,\sigma }`$ are creation and annihilation operators for particles with spin $`\sigma `$ in the $`i`$th energy level of the QD. Once the Hamiltonian is constructed, the tunnelling current of SET can be calculated via the Keldysh Green’s function method.<sup>27</sup> We obtain the tunnelling current through a single dot (see appendix A)
$$J=\frac{2e}{\mathrm{}}\frac{dϵ}{2\pi }[f_L(ϵ\mu _L)f_R(ϵ\mu _R)]\frac{\mathrm{\Gamma }_L(ϵ)\mathrm{\Gamma }_R(ϵ)}{\mathrm{\Gamma }_L(ϵ)+\mathrm{\Gamma }_R(ϵ)}ImG_{e,\sigma }^r(ϵ),$$
(3)
where $`f_L(ϵ)`$ and $`f_R(ϵ)`$ are the Fermi distribution function for the source and drain electrodes, respectively. The chemical potential difference between these two electrodes is related to the applied bias via $`\mu _L\mu _R=eV_a`$. $`\mathrm{\Gamma }_L(ϵ)`$ and $`\mathrm{\Gamma }_R(ϵ)`$ denote the tunnelling rates from the QD to the left (source) and right (drain) electrodes, respectively. For simplicity, these tunnelling rates will be assumed energy and bias-independent. Therefore, the calculation of tunnelling current is entirely determined by the spectral function $`A=ImG_{e,\sigma }^r(ϵ)`$, which is the imaginary part of the retarded Green’s function $`G_{e,\sigma }^r(ϵ)`$. Eq. (3) is obtained provided that the condition $`\mathrm{\Gamma }_{L(R)}>\gamma _{e,c}`$ and $`\mathrm{\Gamma }_{L(R)}R_{eh}`$ are satisfied, where $`\gamma _{e,c}`$ is the captured rate of electrons. Because of the phonon bottleneck effect, we expect the rate for the electron to relax from the wetting layer to the QD is small. Thus, we have ignored the terms involving phonons in the calculation of the tunnelling current. Besides, photon current has also been neglected, since the electron-hole recombination rate, $`R_{eh}`$ is much smaller than the tunnelling rates.
The retarded Green’s function, $`G_{e,\sigma }^r(ϵ)`$, can be obtained by the equation of motion of $`G_{e,\sigma }^r(t)=i\theta (t)\{d_{e,\sigma }(t),d_{e,\sigma }^{}(0)\}`$, where $`\theta (t)`$ is a step function. The curly brackets represent the anti-commutator and the bracket $`\mathrm{}`$ denotes the thermal average. The Fourier transform of $`G_{e,\sigma }^r(t)`$ is given by
$$G_{e,\sigma }^r(ϵ)=_{\mathrm{}}^{\mathrm{}}𝑑tG_{e,\sigma }^r(t)e^{i(ϵ+i\eta )t}$$
(4)
with $`\eta `$ being a positive infinitesimal number. $`G_{e,\sigma }^r`$ is the full Green’s function, which includes all type of interactions. Given the conditions $`\mathrm{\Gamma }_{L(R)}\gamma _{e,c}`$ and $`\mathrm{\Gamma }_{L(R)}R_{eh}`$, we can drop the fifth and the last term of the Hamiltonian in the calculation of $`G_{e,\sigma }^r(ϵ)`$. To solve $`G_{e,\sigma }^r(ϵ)`$, we consider only the lowest order coupling between the electrodes and the QD. The equation of motion for $`G_{e,\sigma }^r(t)`$ leads to
$$\left(ϵE_e+i\frac{\mathrm{\Gamma }_e}{2}\right)G_{e,\sigma }^r\left(ϵ\right)=1+U_eG_{ee}^r\left(ϵ\right)U_{eh}\left(G_{eh,1}^r\left(ϵ\right)+G_{eh2}^r\left(ϵ\right)\right),$$
(5)
where the two particle Green’s functions, $`G_{ee}^r(ϵ)`$, $`G_{eh1}^r(ϵ)`$ and $`G_{eh,2}^r(ϵ)`$ arise from the particle correlation and they satisfy
$$\left(ϵ\left(E_e+U_e\right)+i\frac{\mathrm{\Gamma }_e}{2}\right)G_{ee}^r\left(ϵ\right)=N_{e,\sigma }U_{eh}\left(G_{ehe1}^r\left(ϵ\right)+G_{ehe2}^r\left(ϵ\right)\right),$$
(6)
$$\left(ϵ\left(E_eU_{eh}\right)+i\frac{\mathrm{\Gamma }_e}{2}\right)G_{eh1}^r\left(ϵ\right)=n_{h,\sigma }+U_eG_{ehe1}^r\left(ϵ\right)U_{eh}G_{ehh}^r\left(ϵ\right),$$
(7)
and
$$\left(ϵ\left(E_eU_{eh}\right)+i\frac{\mathrm{\Gamma }_e}{2}\right)G_{eh2}^r\left(ϵ\right)=n_{h,\sigma }+U_eG_{ehe2}^r\left(ϵ\right)U_{eh}G_{ehh}^r\left(ϵ\right).$$
(8)
$`N_{e,\sigma }`$, $`n_{h,\sigma }`$ and $`n_{h,\sigma }`$ in Eqs. (6)-(8) denote the steady-state electron and hole occupation numbers. The two-particle Green’s functions are coupled with the three-particle Green’s functions defined as $`G_{ehe1}^r(ϵ)=N_{e,\sigma }n_{h,\sigma }d_{e,\sigma },d_{e,\sigma }^{}`$, $`G_{ehe2}^r(ϵ)=N_{e,\sigma }n_{h,\sigma }d_{e,\sigma },d_{e,\sigma }^{}`$ and $`G_{ehh}^r(ϵ)=n_{h,\sigma }n_{h,\sigma }d_{e,\sigma },d_{e,\sigma }^{}`$. The equations of motion of the three-particle Green’s functions will lead to coupling with the four-particle Green’s functions (two electrons and two holes), where the hierarchy terminates. Thus, these three-particle Green’s functions can be expressed in the following closed form
$`G_{ehe1}^r(ϵ)=N_{e,\sigma }n_{h,\sigma }({\displaystyle \frac{1n_{h,\sigma }}{ϵ(E_e+U_eU_{eh})+i\frac{\mathrm{\Gamma }_e}{2}}}`$
$`+{\displaystyle \frac{n_{h,\sigma }}{ϵ(E_e+U_e2U_{eh})+i\frac{\mathrm{\Gamma }_e}{2}}}),`$ (9)
$`G_{ehe2}^r(ϵ)=N_{e,\sigma }n_{h,\sigma }({\displaystyle \frac{1n_{h,\sigma }}{ϵ(E_e+U_eU_{eh})+i\frac{\mathrm{\Gamma }_e}{2}}}`$
$`+{\displaystyle \frac{n_{h,\sigma }}{ϵ(E_e+U_e2U_{eh})+i\frac{\mathrm{\Gamma }_e}{2}}}),`$ (10)
and
$`G_{ehh}^r(ϵ)=n_{h,\sigma }n_{h,\sigma }({\displaystyle \frac{1N_{e,\sigma }}{ϵ(E_e2U_{eh})+i\frac{\mathrm{\Gamma }_e}{2}}}`$
$`+{\displaystyle \frac{N_{e,\sigma }}{ϵ(E_e2U_{eh}+U_e)+i\frac{\mathrm{\Gamma }_e}{2}}}).`$ (11)
Eqs. (9) and (10) describe the mixed amplitudes for the propagation of an electron either in the presence of another electron (with opposite spin) plus one hole, or another electron plus two holes. Eq. (11) describes the mixed amplitudes for the propagation of an electron either in the presence of two holes, or two holes plus another electron. Substituting Eqs. (9)-(11) into Eqs. (7) and (8), we obtain, after some algebras, the retarded Green’s function of Eq. (3)
$`G_{e,\sigma }^r\left(ϵ\right)=(1N_{e,\sigma })\{{\displaystyle \frac{1\left(n_{h,\sigma }+n_{h,\sigma }\right)+n_{h,\sigma }n_{h,\sigma }}{ϵE_e+i\frac{\mathrm{\Gamma }_e}{2}}}`$ (12)
$`+`$ $`{\displaystyle \frac{n_{h,\sigma }+n_{h,\sigma }2n_{h,\sigma }n_{h,\sigma }}{ϵE_e+U_{eh}+i\frac{\mathrm{\Gamma }_e}{2}}}+{\displaystyle \frac{n_{h,\sigma }n_{h,\sigma }}{ϵE_e+2U_{eh}+i\frac{\mathrm{\Gamma }_e}{2}}}\}`$
$`+`$ $`N_{e,\sigma }\{{\displaystyle \frac{1\left(n_{h,\sigma }+n_{h,\sigma }\right)+n_{h,\sigma }n_{h,\sigma }}{ϵE_eU_e+i\frac{\mathrm{\Gamma }_e}{2}}}`$
$`+`$ $`{\displaystyle \frac{n_{h,\sigma }+n_{h,\sigma }2n_{h,\sigma }n_{h,\sigma }}{ϵE_eU_e+U_{eh}+i\frac{\mathrm{\Gamma }_e}{2}}}`$
$`+`$ $`{\displaystyle \frac{n_{h,\sigma }n_{h,\sigma }}{ϵE_eU_e+2U_{eh}+i\frac{\mathrm{\Gamma }_e}{2}}}\}.`$
Here $`G_{e,\sigma }^r(ϵ)`$ contains an admixture of six possible configurations in which a given electron can propagate. These configurations are: empty state, one-hole state, two-hole state, one-electron state, one-electron plus one-hole state, and one-electron plus two-hole state. In Eq. (12), $`\mathrm{\Gamma }_e`$ is the electron tunnelling rate $`\mathrm{\Gamma }_e\mathrm{\Gamma }_L+\mathrm{\Gamma }_R`$. The particle correlation has not been included in the tunneling rates of Eq. (12). Such approximation has been known to be adequate for describing the Coulomb blockade effect, but not the Kondo effect. The electron occupation number of the QD can be solved self-consistently via the relation
$$N_{e,\sigma }=\frac{dϵ}{\pi }\frac{\mathrm{\Gamma }_Lf_L\left(ϵ\right)+\mathrm{\Gamma }_Rf_R\left(ϵ\right)}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}ImG_{e,\sigma }^r\left(ϵ\right).$$
(13)
$`N_{e,\sigma }`$ is limited to the region $`0N_{e,\sigma }1`$. Eq. (13) indicates that the electron occupation numbers of the QD, $`N_{e,\sigma }`$ and $`N_{e,\sigma }`$, are primarily determined by the tunnelling process. To obtain the electron and hole occupation numbers ($`n_{e,\sigma }=n_{e,\sigma }`$ and $`n_{h,\sigma }=n_{h,\sigma }`$) arising from the optical pumping, we solve the rate equations (see appendix B) and obtain
$$n_e=n_{e,\sigma }=n_{e,\sigma }=\frac{\gamma _{e,c}N_{e,𝐤}(1N_e)}{\gamma _{e,c}N_{e,𝐤}+R_{eh}n_h+\mathrm{\Gamma }_e}$$
(14)
and
$$n_h=n_{h,\sigma }=n_{h,\sigma }=\frac{\gamma _{h,c}N_{h,𝐤}}{\gamma _{h,c}N_{h,𝐤}+R_{eh}(n_e+N_e)+\mathrm{\Gamma }_{h,s}},$$
(15)
where $`\gamma _{e(h),c}`$ and $`N_{e(h),𝐤}`$ denote the captured rate for electrons (holes) from the wetting layer to the QD and the occupation number of electrons (holes) in the wetting layer. See Eqs. (37) and (38), $`N_{e(h),𝐤}`$ is a linear function of $`p_{exc}`$ for low pumping intensity and small $`\gamma _{e(h),c}`$. $`\mathrm{\Gamma }_{h,s}`$ denotes the non-radiative recombination rate for holes in the QD. Because $`\gamma _{e,c}N_{e,𝐤}/\mathrm{\Gamma }_e1`$, $`n_e`$ has been neglected in Eq. (12).
## III Spontaneous emission spectrum
When electrons and holes appear in the QD, their recombination leads to single-photon emission. The optical polarization of the emitted photon depends on the spin polarization of the electrons. An electron with spin +1/2 (-1/2) can recombine with a heavy hole of angular momentum +3/2 (-3/2) to emit a circularly ($`\sigma ^+or\sigma ^{}`$) polarized photon. In some experiments, up to four coexisting peaks (associated with exciton, negative trion, positive trion and biexciton) have been observed in the emission spectrum.<sup>16-18</sup> The emission process is described by the last term in the Hamiltonian given in Eq. (1). The bias-dependent emission spectrum can be obtained by finding the nondiagonal lesser Green’s function at nonequal time.<sup>7</sup> To simplify the problem, we assume that the applied bias is large enough to create the biexciton state (two electrons and two holes) in the QD.
With the above assumption, polarization $`P(\omega )=G_{eh}^<(\omega )`$ can be calculated by the equal-time Keldysh Green function, $`G_{eh}^<(t,t)=ia^{}(t)d_{h,\sigma }(t)d_{e,\sigma }(t)`$ (or density matrix method). Prior to the calculation of the polarization $`P(\omega )`$, a unitary transformation has been used to remove the phase of the optical transition term. The renormalized energy levels of the electron and hole become $`ϵ_e=E_e\frac{\omega }{2}`$ and $`ϵ_h=E_h\frac{\omega }{2}`$. Furthermore, the hopping terms between the leads and dot, $`V_{j=L,R;𝐤}(t)=V_{j=L,R;𝐤}exp^{\frac{(\pm i\omega t)}{2}}`$ become time dependent, in which the energy and time dependence of the coupling are factorized. This factorization leads to time-independent tunnelling rates and it simplifies the calculation. Solving the equation of motion for $`G_{eh}^<(t,t)=ia^{}(t)d_{h,\sigma }(t)d_{e,\sigma }(t)`$ and considering the steady-state condition, we obtain the lesser Green’s function
$$G_{eh}^<(\omega )=\frac{\lambda (U_eU_{eh})𝒫_1(\omega )(U_hU_{eh})𝒫_2(\omega )}{ϵ_e+ϵ_hU_{eh}+i\frac{\mathrm{\Gamma }}{2}},$$
(16)
where
$$=\mathrm{\Phi }|a^{}a|\mathrm{\Phi }(1N_{e,\sigma }n_{h,\sigma })N_{e,\sigma }n_{h,\sigma },$$
(17)
$`\mathrm{\Gamma }\mathrm{\Gamma }_e+\mathrm{\Gamma }_h`$, which leads to a broadening of the spectrum, and $`|\mathrm{\Phi }`$ denotes a photon state. $`𝒫_1(\omega )`$ and $`𝒫_2(\omega )`$ are given by
$`𝒫_1(\omega )=\lambda N_{e,\sigma }({\displaystyle \frac{1n_{h,\sigma }}{ϵ_e+ϵ_h+U_e2U_{eh}+i\frac{\mathrm{\Gamma }}{2}}}`$
$`+{\displaystyle \frac{n_{h,\sigma }}{ϵ_e+ϵ_h+U_e+U_h3U_{eh}+i\frac{\mathrm{\Gamma }}{2}}})`$ (18)
and
$`𝒫_2(\omega )=\lambda n_{h,\sigma }({\displaystyle \frac{1N_{e,\sigma }}{ϵ_e+ϵ_h+U_h2U_{eh}+i\frac{\mathrm{\Gamma }}{2}}}`$
$`+{\displaystyle \frac{N_{e,\sigma }}{ϵ_e+ϵ_h+U_e+U_h3U_{eh}+i\frac{\mathrm{\Gamma }}{2}}}).`$ (19)
The first term in the filling factor $``$ \[Eq. (17)\] arises from the stimulated process, which vanishes in the vacuum state. The second term, $`N_{e,\sigma }n_{h,\sigma }`$ is due to the spontaneous emission process. This is a quantum effect of the electromagnetic field. Expressions of Eqs. (18) and (19) are derived, respectively, from the equation of motions for $`ia^{}(t)N_{e,\sigma }(t)d_{h,\sigma }(t)d_{e,\sigma }(t)`$ and $`ia^{}(t)n_{h,\sigma }(t)d_{h,\sigma }(t)d_{e,\sigma }(t)`$, which are coupled with $`ia^{}(t)n_{h,\sigma }(t)N_{e,\sigma }(t)d_{h,\sigma }(t)d_{e,\sigma }(t)`$. The term $`ia^{}(t)N_{e,\sigma }(t)d_{e,\sigma }(t)d_{h,\sigma }(t)`$ describes an electron-screened electron-hole recombination process in a negative trion. Similarly, $`ia^{}(t)n_{h,\sigma }(t)d_{h,\sigma }(t)d_{e,\sigma }(t)`$, describes a hole-screened electron-hole recombination process in a positive trion, and $`ia^{}(t)n_{h,\sigma }(t)N_{e,\sigma }(t)d_{e,\sigma }(t)d_{h,\sigma }(t)`$ describes the biexciton-to-exciton transition process. In the photon vacuum state, we obtain
$`𝒫\left(\omega \right)/\left(N_{e,\sigma }n_{h,\sigma }\right)`$
$`=`$ $`\lambda ^2\{{\displaystyle \frac{\left(1N_{e,\sigma }\right)\left(1n_{h,\sigma }\right)}{E_gU_{eh}\omega +i\mathrm{\Gamma }/2}}+{\displaystyle \frac{\left(1n_{h,\sigma }\right)N_{e,\sigma }}{E_g+U_e2U_{eh}\omega +i\mathrm{\Gamma }/2}}`$
$`+`$ $`{\displaystyle \frac{\left(1N_{e,\sigma }\right)n_{h,\sigma }}{E_g+U_h2U_{eh}\omega +i\mathrm{\Gamma }/2}}`$
$`+`$ $`{\displaystyle \frac{N_{e,\sigma }n_{h,\sigma }}{E_g+U_h+U_e3U_{eh}\omega +i\mathrm{\Gamma }/2}}\}.`$
The spectrum of the polarization $`𝒫(\omega )`$ displays four peaks, which are attributed to the exciton $`X`$, negative trion $`X^{}`$, positive trion $`X^+`$ and biexciton $`X^2`$ (biexciton decaying to exciton). The corresponding peak positions occur at $`\omega =E_gU_{eh}`$, $`\omega =E_g+U_e2U_{eh}`$, $`\omega =E_g+U_h2U_{eh}`$ and $`\omega =E_g+U_e+U_h3U_{eh}`$, which are significantly influenced by the particle Coulomb interactions.
## IV Energy levels and inter-particle interactions
According to Eqs. (12) and (20), particle Coulomb interactions will significantly affect the tunnelling current of SET and the emission spectrum of single photons. To illustrate this effect, we apply our theory to a pyramid-shaped InAs/GaAs QD. First, we calculate the inter-particle Coulomb interactions using a simple but realistic effective-mass model. The electronic structures of InAs/GaAs QDs have been extensively studied by several groups with different methods.<sup>29-33</sup> The electron (hole)in the QD is described by the equation
$`[{\displaystyle \frac{\mathrm{}^2}{2m_{e(h)}^{}(\rho ,z)}}+V_{QD}^{e(h)}(\rho ,z)eFz]\psi _{e(h)}(𝐫)`$ (21)
$`=`$ $`E_{e(h)}\psi _{e(h)}(𝐫),`$
where $`m_e^{}(\rho ,z)`$ (a scalar) denotes the position-dependent electron effective mass, which has $`m_{eG}^{}=0.067m_e`$ for GaAs and $`m_{eI}^{}=0.04m_e`$ for the strained InAs in the QD. $`m_h^{}(\rho ,z)`$ denotes the position-dependent effective mass tensor for the hole. It is a fairly good approximation to describe $`m_h^{}(\rho ,z)`$ in InAs/GaAs QD as a diagonal tensor with the $`x`$ and $`y`$ components given by $`m_{t}^{}{}_{}{}^{1}=(\gamma _1+\gamma _2)/m_e`$ and the $`z`$ component given by $`m_{l}^{}{}_{}{}^{1}=(\gamma _12\gamma _2)/m_e`$. $`\gamma _1`$ and $`\gamma _2`$ are the Luttinger parameters. Their values for InAs and GaAs are taken from Ref. 34. $`V_{QD}^e(\rho ,z)`$ ($`V_{QD}^h(\rho ,z)`$) is approximated by a constant potential in the InAs region with value determined by the conduction-band (valence-band) offset and the deformation potential shift caused by the biaxial strain in the QD. These values have been determined by comparison with results obtained from a microscopic model calculation<sup>30</sup> and we have $`V_{QD}^e=0.5eV`$ and $`V_{QD}^h=0.32eV`$. The $`eFz`$ term in Eq.(21) arises from the applied voltage, where $`F`$ denotes the strength of the electric field. Using the eigenfunctions of Eq. (21), we calculate the inter-particle Coulomb interactions via
$$U_{i,j}=𝑑𝐫_1𝑑𝐫_2\frac{e^2[n_i(𝐫_1)n_j(𝐫_2)]}{ϵ_0|𝐫_1𝐫_2|},$$
(22)
where $`i(j)=e,h`$. $`n_i(𝐫_1)`$ denotes the charge density. $`ϵ_0`$ is the static dielectric constant of InAs. We have ignored the image force arising from the small difference of dielectric constant between InAs and GaAs. Although the Coulomb energies are different in different exciton complexes, their difference is rather small.<sup>27,28</sup> Therefore, only one set of Coulomb interaction parameters has been used in this study.
For the purpose of constructing the approximate wave functions, we place the system in a large confining cubic box with length L. Here we adopt $`L=40nm`$. The wave functions are expanded in a set of basis functions, which are chosen as sine waves
$`\psi _{nlm}(\rho ,\varphi ,z)`$ $`=`$ $`{\displaystyle \frac{\sqrt{8}}{\sqrt{L^3}}}\mathrm{sin}\left(k_lx\right)\mathrm{sin}\left(k_my\right)\mathrm{sin}\left(k_nz\right),`$ (23)
where $`k_n=n\pi /L`$,$`k_m=m\pi /L`$,$`k_{\mathrm{}}=\mathrm{}\pi /L`$. n, m and $`\mathrm{}`$ are positive integers. The expression of the matrix elements of the Hamiltonian of Eq. (21) can be readily obtained. In our calculation $`n=20`$, $`m=10`$, and $`\mathrm{}=10`$ are used in solving Eq. (21). Fig. 2 shows the lowest three energy levels of the pyramidal InAs/GaAs QD as a function of QD size. The ratio of height and base length is fixed at $`h/b=1/4`$, while $`h`$ varies from $`2.5`$ nm to $`6.5`$nm. Diagram (a) and (b) denote, respectively, the energy levels for electrons and holes. We can see that the energy level spacing between the ground state and the first excited state is much larger than $`k_BT=2meV`$, which will be considered throughout this article. Therefore, the lowest energy levels, $`E_e`$ and $`E_h`$, are adopted in the Hamiltonian of Eq. (1). The intralevel Coulomb interactions $`U_e`$ and $`U_h`$ and interlevel Coulomb interaction $`U_{eh}`$ are calculated using Eq. (22). Fig. 3 shows the inter-particle interactions as functions of QD size. The strengths of Coulomb interactions are inversely proportional to the QD size. However as the QD size decreases below a threshold value (around $`b=12nm`$), $`U_e`$ is significantly reduced due to the leak out of electron density for small QDs. These Coulomb interactions approach approximately the same value in the large QD limit. This indicates similar degree of localization for electron and hole in large QDs. We also note that $`U_{eh}`$ is smaller than $`U_e`$ in the large QD. This is due to the fact that in large QDs the degree of localization for the hole becomes similar to that for electron, while the anisotropic nature of hole wave function reduces $`U_{eh}`$. The repulsive Coulomb interactions, $`U_e`$ and $`U_h`$, are the origin of Coulomb blockade for electrons and holes, respectively. The attractive Coulomb interaction $`U_{eh}`$ gives rise to the binding of the exciton. To study the behavior of tunnelling current and bias-dependent spontaneous emission spectrum, we consider a particular pyramidal InAs/GaAs QD with base length $`b=13`$nm and height $`h=3.5nm`$. The other relevant parameters for this QD are $`E_{e,0}=0.14eV`$, $`E_{h,0}=0.125eV`$,$`U_e=16.1meV`$, $`U_{eh}=16.7meV`$ and $`U_h=18.5meV`$.
## V Results
In this section, we discuss our numerical results for the tunnelling current and spontaneous emission spectrum. For simplicity, we assume that the tunnelling rate $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.5meV`$ is bias-independent. We apply a bias voltage $`V_a`$ across the source-drain and $`V_g`$ across the gate-drain. With the applied voltages, the QD electron and hole energy levels will be shifted to $`E_e+\alpha eV_a\beta eV_g`$ and $`E_h+\alpha eV_a\beta eV_g`$, respectively. In our calculation, we assume $`\alpha =0.5`$ and $`\beta =0.7`$. Meanwhile the chemical potentials of the electrodes corresponding to a Fermi energy $`E_F=60meV`$ (relative to the conduction band minimum in the leads) are assumed to be $`70meV`$ below the energy level of $`E_e`$ at zero bias. Other parameters adopted are: $`\mathrm{\Gamma }_h=0.2meV`$ and $`R_{eh}=10\mu eV`$.
### V.1 Effects of optical pumping on the tunnelling current
Applying Eqs. (3), (13) and (15), we solve for the electron occupation number $`N_e=N_{e,\sigma }=N_{e,\sigma }`$ and the tunnelling current $`J`$. Fig. 4 shows the calculated results for $`N_e`$ and $`J`$ as functions of gate voltage with and without photo-excitation at zero temperature and $`V_a=2mV`$. Solid line and dashed line correspond to $`I=0`$ (no pump) and $`I=0.9`$ (with pump), respectively. Here $`I`$ is a dimensionless quantity, defined as $`I\gamma _{h,c}N_{h,𝐤}/\mathrm{\Gamma }_{h,s}`$ which is proportional to the pump power. The electron occupation number displays several plateaus, while the tunnelling current displays an oscillatory behavior. Both are typical behaviors due to the Coulomb blockade. The manipulation of gate voltage can tune the energy levels of the QD. For the no-pump case ($`I=0`$), carriers in the emitter electrode are allowed to tunnel into the QD as the gate voltage exceeds the threshold voltage, $`V_{g3}(100mV)`$. When the energy level of $`E_e+\alpha V_a\beta V_g`$ is below the emitter chemical potential, but $`E_e+\alpha V_a\beta V_g+U_e`$ above the emitter chemical potential, $`N_e`$ (for a given spin) reaches 0.5 and displays a plateau which is caused by the Coulomb blockade associated with $`U_e`$ (the charging energy of the electron ground state). At higher gate voltage, $`E_e+\alpha V_a\beta V_g+U_e`$ moves below the emitter chemical potential, and $`N_e`$ approaches 1. When electron-hole pairs are generated due to optical pumping in addition to the tunnelling electrons, various exciton complexes such as the exciton, positive and negative trions, and biexciton can be formed. These new channels allow carriers in the emitter electrode to tunnel into the QD at lower gate voltage. When the gate voltage is below $`V_{g2}`$, but above $`V_{g1}`$, only the positive trion state is below the emitter chemical potential, hence only a small plateau is seen. As the gate voltage increases, more exciton complex levels are depressed below the emitter chemical potential, and a series of plateaus appear. The value of $`N_e`$ depends on which exciton complex state is filled and on the pumping intensity.
The tunnelling current as a function of the gate voltage can be measured directly. We see that optical pumping leads to two additional peaks below the threshold voltage $`V_{g3}`$, which is caused by the electron tunnelling assisted by the presence of a hole in the QD. This interesting phenomenon was observed by Fujiwara et al. in an SET composed of one silicon QD and three electrodes.<sup>35</sup> The behavior of the photo-induced tunnelling current can be understood by analyzing the energy poles of the retarded Green’s function as given in Eq. (12). In Figure 4, the first peak of the dashed line is caused by a resonant tunnelling through the energy pole at $`ϵ=E_e2U_{eh}`$ (which corresponds to a positive trion state) with probability $`(1N_{e,\sigma })(n_{h,\sigma }n_{h,\sigma })`$. The second peak is caused by a pair of poles at $`ϵ=E_eU_{eh}`$ (the exciton state) and $`ϵ=E_e+U_e2U_{eh}`$ (the biexciton state) with probabilities $`(1N_{e,\sigma })(n_{h,\sigma }+n_{h,\sigma }2n_{h,\sigma }n_{h,\sigma })`$ and $`N_{e,\sigma }(n_{h,\sigma }n_{h,\sigma })`$. Since the magnitude of $`U_e`$ is very close to that of $`U_{eh}`$ for the present case, the corresponding two peaks merge into one. For other cases (e.g. smaller QDs), the difference between $`U_e`$ and $`U_{eh}`$ may be large enough for the peaks to be distinguished. The third peak is caused by another pair of poles at $`ϵ=E_e`$ (the single-electron state) and $`ϵ=E_e+U_eU_{eh}`$ (the negative trion state) with probabilities $`(1N_{e,\sigma })(1(n_{h,\sigma }+n_{h,\sigma })+n_{h,\sigma }n_{h,\sigma })`$ and $`N_{e,\sigma }(n_{h,\sigma }+n_{h,\sigma }2n_{h,\sigma }n_{h,\sigma })`$. The last peak locating near $`V_g=123mV`$ is due to the resonant tunnelling through the energy level at the pole $`ϵ=E_e+U_e`$ (the two-electron state) with probability $`N_{e,\sigma }(1(n_{h,\sigma }+n_{h,\sigma })+n_{h,\sigma }n_{h,\sigma })`$. Obviously, the relative strengths of these tunnelling current peaks are influenced by the hole occupation numbers. The gate voltage differences $`\mathrm{\Delta }V_{g21}=V_{g2}V_{g1}`$ and $`V_{g43}=V_{g4}V_{g3}`$ allow the determination of the electron-hole interaction and the electron-electron Coulomb repulsion, since $`\beta \mathrm{\Delta }V_{g21}=U_{eh}`$ and $`\beta V_{g43}=U_e`$. Consequently, such a measurement can be used to determine the exciton binding energy, which is otherwise not possible via typical photoluminescence or the electrically driven emission spectrum measurement. In such a measurement, it is necessary to distinguish $`U_{eh}`$ from $`U_e`$. Thus, we recommend using a smaller bias ($`V_a`$), which would make the tunnelling peaks sharper and the double-peak features at $`V_g=V_{g2}`$ and $`V_g=V_{g3}`$ better resolved.
Figure 5 shows the tunnelling current as a function of gate voltage for various strengths of photo-excitation power at zero temperature.According to the definition of $`\gamma _{e(h),c}=_𝐪|g_{e(h),𝐪}|^2Im1/(\gamma _{e(h),s}(\mathrm{\Omega }_𝐪+i\gamma _{e(h),s}))`$, the condition $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R\gamma _{e,c}`$ is still satisfied, even though we tune $`I=\gamma _{h,c}N_{h,𝐤}/\mathrm{\Gamma }_{h,s}`$ up to 10 in Fig. 5. $`I=10`$ can be regarded as increasing $`p_{exc}`$ by 10 times, but it is still in the weak pumping regime. We notice that increasing photo-excitation power tends to suppress the tunnelling current at high gate-voltage. For instance, the peak located at $`V_g=123mV`$ almost vanishes due to the much reduced probability for the pole at $`ϵ=E_e+U_e`$. Based on the results of Figs. 4 and 5, we see that the SET has potential application as an optically controlled switch. For practical applications, we must consider the behavior of the tunnelling current for an SET operated at finite temperatures. Figure 6 shows a comparison of the electron occupation number and tunnelling current at zero temperature and finite temperature ($`k_BT=2meV`$). We see that the plateaus become broadened and the magnitude of the tunnelling current peak is reduced as the temperature increases. Thus, a large charging energy is required for SETs to be operated at high temperatures. The tunnelling current as a function of applied bias ($`V_a`$) for various strengths of photo-excitation power at $`k_BT=2meV`$ and $`V_g=0`$ is shown in Fig. 7. The tunnelling current exhibits a staircase behavior. This is the well-known Coulomb blockade effect. The effect of the optical pumping is reduce the threshold voltage and increase the number of plateaus in the J-V characteristics. We also notice that a negative differential conductance occurs at high voltages when the resonance level of the QD is below the conduction band minimum of the emitter electrode.
### V.2 Effects of electrodes on the spontaneous emission spectrum
SET has been considered as a single-photon generater (SPG) for quantum information application such as quantum cryptography and teleportation.<sup>36</sup> Thus, it is of great interest to study the spontaneous single-photon emission spectrum of the SET. We previously reported the theoretical studies of the emission spectrum of a single QD embedded in a p-n junction.<sup>27,28</sup> Here, we consider the single-photon emission of an SET (a QD embedded in a n-i-n junction) under optical pumping. Since the magnitudes of inter-particle Coulomb interactions $`U_e`$, $`U_h`$ and $`U_{eh}`$ are fairly close, it will be difficult to resolve the peaks of exciton ( $`X`$), negative trion ($`X^{}`$), positive trion ($`X^+`$) and biexciton ($`X^2`$) for the case of large tunnelling rate as we considered above ($`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.5meV`$). Therefore, for this study we adopt a small tunnelling rate, $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0.1meV`$ in order to resolve these peaks in the emission spectrum. This can be achieved, for example, by increasing the barrier thickness between the QD and the conducting leads. Figure 8 shows the emission spectrum for various gate-voltages at fixed pump power ($`I=0.9`$). For $`V_g=100mV`$, the calculated spontaneous emission spectrum displays four peaks corresponding to $`X^{}`$, $`X`$, $`X^2`$ and $`X^+`$. The red shift of the $`X^{}`$ peak relative to the exciton peak agrees well with experimental observations.<sup>16,22,24,37</sup> The biexciton peak displaying a blue shift with respect to the exciton peak (showing an antibinding biexciton) is also consistent with the observation reported in ref.\[14,19-21\] Recently, studies of the binding and antibinding of biexcitons were reported in Ref. 21. Our calculation given by Eqs. (21) and (22) provides only the antibinding feature of biexciton. In Ref. 21 it is pointed out that the biexciton complex changes from antibinding to binding as the QD size increases. For QDs with dimension larger than the exciton Bohr radius (around 20nm), the correlation energy becomes significant. In this study, we have not taken into account the correlation energy. This is justified as long as we restrict ourselves to QDs with size less than 20nm. We see in this figure when the gate-voltage is increased up to $`V_g=125mV`$, only the $`X^{}`$ and $`X^2`$ peaks survive in the emission spectrum as $`N_e`$ approaches 1. Furthermore, we can adjust the intensity of excitation power to either enhance or suppress the strength of the biexciton peak. Thompson et al. demonstrated that the single photon generated by the biexciton state is significantly less in emission time than the single exciton state.<sup>19</sup> From this point of view, the biexciton state is favored for the application of single-photon generation. We find that one can increase the pump power in order to select the biexciton state as the desired source to produce the single-photon emission. This is demonstrated in Fig. 9, where the emission spectra for various strengths of excitation power at the fixed gate voltage $`V_g=125mV`$ are shown. We see that the negative trion saturates at $`I=1.9`$ and diminishes afterwards, while the biexciton peak increases quadratically with $`I`$ and becomes the dominant peak in the emission spectrum.
Finally, we discuss the emission spectrum of an isolated charge-neutral QD with $`\mathrm{\Gamma }_L=\mathrm{\Gamma }_R=0`$. For this case the rate equation of exciton number $`N_{X,\sigma }`$ can be readily obtained by solving the rate equation
$$\gamma _{X,c}N_k(1N_{X,\sigma })\mathrm{\Omega }\frac{d\omega }{\pi }\omega \rho (\omega )|ImP(\omega )|=0$$
(24)
with
$$P(\omega )=N_{X,\sigma }(\frac{1N_{X,\sigma }}{E_X\omega +i\mathrm{\Gamma }/2}+\frac{N_{X,\sigma }}{E_X+\mathrm{\Delta }E_c\omega +i\mathrm{\Gamma }/2}),$$
(25)
where $`\mathrm{\Omega }=4\pi n_r^3\mu _r^2/(6c^3\mathrm{}^3ϵ_0)`$ and $`n_r`$ is the refractive index of the system. $`\rho (\omega )=\omega ^2`$ arises from the density of states of photons. $`E_X`$ and $`\mathrm{\Delta }E_{X,c}`$ denote, respectively, the exciton energy and exciton-exciton correlation energy. $`\gamma _{X,c}`$ denotes the capture rate for excitons from the wetting layer to the QD. After the integration of $`\omega `$ in Eq. (24), we obtain
$$\mathrm{\Omega }N_{X,\sigma }((1N_{X,\sigma })E_X^3+N_{X,\sigma }(E_X+\mathrm{\Delta }E_c)^3)$$
(26)
for small $`\mathrm{\Gamma }`$, and
$$N_{X,}=N_{X,}=\frac{\gamma _{X,c}N_k}{\gamma _{X,c}N_k+\gamma _{X,r}},$$
(27)
where $`\gamma _{X,r}=\mathrm{\Omega }E_X^3`$ denotes the exciton decaying rate, and we have ignored $`\mathrm{\Delta }E_c`$, since it is small compared to $`E_X`$. According to Eq. (25), we see that the intensities of exciton and biexciton peaks are proportional to $`p_{exc}`$ and $`p_{exc}^2`$, if $`N_{X,\sigma }`$ is proportional to $`p_{exc}`$ (the excitation power). This is in very good agreement with many experimental reports.<sup>17,19-21,38</sup> Note that in some experiments, the pump energy is far above the QD band gap and the emission occurs after the electrons and holes are captured. This situation becomes similar to the open system considered in the SET with optical pumping. In this case, we can replace $`N_{e,\sigma }`$ in Eq. (20) by $`n_{e,\sigma }`$, the electron occupation due to optical pumping, and we can still prove that the intensities of $`X`$ and $`X^2`$ peaks as functions of the excitation power are described by Eqs. (14), (15) and (20). Thus, they exhibit linear and quadratic behavior as a function of $`p_{exc}`$, respectively. So, there is no qualitative difference between Eq. (20) and Eq. (24) regarding the behavior of $`X`$ and $`X^2`$, except that Eq. (20) also describes the possibility of emission due to recombination in trions. This provides an explanation to the experimental observation of Ref. 17, in which the pump energy is far above the QD band gap and both $`X^{}`$ and $`X^+`$ peaks (in addition to the $`X`$ and $`X^2`$ peaks) were observed in the emission spectrum.
## VI Summary
In this article we have studied the tunnelling current and emission spectrum of an SET under optical pumping theoretically. We apply our theoretical analysis to an SET, which consists of a single quantum dot embedded in a n-i-n junction. We use a simple but realistic effective mass model to calculate the inter-particle Coulomb interactions of pyramid-shaped InAs/GaAs QDs. These Coulomb interactions are essential for the study of optical and transport properties of QDs. The retarded Green’s function was calculated via a non-equilibrium Green’s function method. It is found that the holes in QD generated by optical pumping leads to new channels for electrons to tunnel from the emitter to the collector, which creates a variety of ways to control both the electrical and optical signals. The binding energy of exciton complexes as well as electron charging energy can be determined by examining the tunnelling current as a function of gate voltage $`V_g`$ with a small applied bias $`V_a`$. In addition, the emission spectrum of SET can be modified significantly by adjusting bias voltage, the gate voltage, or the pump intensity.
ACKNOWLEDGMENTS
This work was supported by National Science Council of Republic of China under Contract Nos. NSC 93-2215-E-008-014, NSC-93-2120-M-008-002, and NSC 93-2215-E-008-011.
## VII Appendix A
In the system described by Eq. (1), electromagnetic field of frequency $`\omega _0`$ is used to resonantly pump carriers into the InAs wetting layer. Therefore, electrons can be injected into the InAs QD via either the tunnelling process or carrier capture from the wetting layer. On the other hand, holes in InAs QD are provided only by the capture process. The dominant carrier capture process is the the optical-phonon assisted carrier relaxation as described by the fifth term in Eq. (1). Because of the phonon-bottleneck effect, the carrier capture rates for QDs are suppressed. Consequently, the condition $`\mathrm{\Gamma }_e\gamma _{e,c}`$ is usually satisfied, where $`\mathrm{\Gamma }_e`$ and $`\gamma _{e,c}`$ denote the electron tunnelling rate and electron capture rate, respectively. Based on the above condition, we can ignore the electron capture process in the calculation of the tunnelling current. The tunnelling current from the left electrode to the InAs QD can be calculated from the time evolution of the occupation number, $`N_L=_{𝐤,\sigma ;L}c_{𝐤,\sigma ;L}^{}c_{𝐤,\sigma ;L}`$.
Before the calculation of tunnelling current, an unitary transformation is employed to remove the phase of the optical pumping term. The renormalized energy levels of the electron and hole become $`ϵ_e=E_e\omega /2`$ and $`ϵ_h=E_h\omega /2`$. Furthermore, the hopping terms between the leads and the dot become time dependent, $`V_{𝐤,j}(t)=V_𝐤exp[(\pm i\omega t)/2]`$. Now, we can apply the formulas given in ref.. We obtain
$$J=\frac{2e}{\mathrm{}}\frac{dϵ}{2\pi }\mathrm{\Gamma }_L(ϵ)[\frac{1}{2}G_{e,\sigma }^<(ϵ)+f_L(ϵ)ImG_{e,\sigma }^r(ϵ)],$$
(28)
where $`\mathrm{\Gamma }_L(ϵ)=2\pi _𝐤|V_{L,𝐤}|^2\delta (ϵϵ_𝐤)`$, and $`f_L(ϵ)`$ is the Fermi distribution function. $`G_{e,\sigma }^<(ϵ)`$ and $`G_{e,\sigma }^r(ϵ)`$ are the Fourier transformation of the lesser Green function $`G_{e,\sigma }^<(t)=id_{e,\sigma }^{}(0)d_{e,\sigma }(t)`$ and retarded Green’s function $`G_{e,\sigma }^r(t)=i\theta (t)\{d_{e,\sigma }(t),d_{e,\sigma }^{}(0)\}`$. Applying Dyson’s equation, we obtain
$$G_{e,\sigma }^<(ϵ)=f^<(ϵ)[G_{e,\sigma }^r(ϵ)G_{e,\sigma }^a(ϵ)]$$
(29)
with
$$f^<(ϵ)=\frac{i_{𝐤,\mathrm{}}|V_{𝐤,\mathrm{}}|^2g_{e,𝐤,\sigma ;\mathrm{}}^<(ϵ)+i_Q|\lambda |^2𝒢_{h,\sigma }^<(ϵ)}{i_{𝐤,\mathrm{}}|V_{𝐤,\mathrm{}}|^2(g_{e,𝐤,\sigma ;\mathrm{}}^r(ϵ)g_{e,𝐤,\sigma ;\mathrm{}}^a(ϵ))i_Q|\lambda |^2Im𝒢_{h,\sigma }^r(ϵ)},$$
(30)
where $`g_{e,𝐤,\sigma ;\mathrm{}}^<(ϵ)=f_{\mathrm{}}(ϵ)(g_{e,𝐤,\sigma ;\mathrm{}}^r(ϵ)g_{e,𝐤,\sigma ;\mathrm{}}^a(ϵ))`$.$`(g_{\mathrm{}}^r(ϵ)g_{\mathrm{}}^a(ϵ))=i2\pi \delta (ϵϵ_𝐤)`$. The lesser Green’s function and retarded Green’s function for holes are denoted by $`𝒢_{h,\sigma }^<(ϵ)`$ and $`𝒢_{h,\sigma }^r(ϵ)`$. Due to the small electron-hole recombination rate $`R_{eh}1/ns`$, terms involving $`|\lambda |^2`$ can be droped. That is to ignore the photocurrent in Eq. (28). After rewriting $`f^<(ϵ)`$ as
$$f^<(ϵ)=\frac{\mathrm{\Gamma }_L(ϵ)f_L(ϵ)+\mathrm{\Gamma }_R(ϵ)f_R(ϵ)}{\mathrm{\Gamma }_L(ϵ)+\mathrm{\Gamma }_R(ϵ)},$$
(31)
and substituting it into Eq. (28), we obtain the tunnelling current
$$J=\frac{2e}{\mathrm{}}\frac{dϵ}{2\pi }\frac{\mathrm{\Gamma }_L(ϵ)\mathrm{\Gamma }_R(ϵ)}{\mathrm{\Gamma }_L(ϵ)+\mathrm{\Gamma }_R(ϵ)}[f_L(ϵ)f_R(ϵ)]ImG_{e,\sigma }^r(ϵ).$$
(32)
The expression of Eq. (32) was also obtained by Meir, Wingreen and Lee<sup>6</sup>. It is of no surprise that we reproduce their result, because the system considered here reduces to their case as the carrier relaxation process and electron-hole recombination process are both turned off. Note that the retarded Green’s function $`G_{e,\sigma }^r(ϵ)`$ includes the full effects due to the electron-electron Coulomb interaction and electron-hole Coulomb interaction and the effect due to the dot-electrodes coupling to the lowest order. Also note that the electrodes are coupled to the conduction band of the QD, but not to the valence band of the QD.
## VIII Appendix B
When light is used to resonantly create electrons and holes in the wetting layer, electrons and holes can relax to the InAs QD via optical phonon-assisted process. When electrons relax into InAs QD, they will tunnel out quickly since the electrons in the ground state of the InAs QD are coupled with the electrodes. However, holes in the InAs QD can decay only via recombination with electrons (radiation process) or via impurity scattering (nonradiation process). When holes are present in the InAs QD, we expect the tunnelling current to change significantly due to the large electron-hole Coulomb interactions \[see Eq. (12)\]. In this appendix, we give the derivations for Eqs. (14) and (15), which describe the electron and hole populations in the InAs QD under optical pumping. First, we introduce the electron and hole distribution functions in the wetting layer, $`N_{e,𝐤}=d_{e,𝐤}^{}(t)d_{e,𝐤}(t)`$ and $`N_{h,𝐤}=d_{h,𝐤}^{}(t)d_{h,𝐤}(t)`$. We use the equation-of-motion method to calculate $`N_{e,𝐤}=d_{e,𝐤}^{}(t)d_{e,𝐤}(t)`$ and $`N_{h,𝐤}=d_{h,𝐤}^{}(t)d_{h,𝐤}(t)`$. We obtain
$$\frac{d}{dt}N_{e,𝐤}=i\gamma _{e,s}N_{e,𝐤}+2Im\lambda _0^{}b_{e,𝐤}^{}b_{h,𝐤}^{}$$
$$2\underset{𝐪}{}Img_{e,𝐪}Ab_{e,𝐤}^{}d_{e,\sigma },$$
(33)
and
$$\frac{d}{dt}N_{h,𝐤}=i\gamma _{h,s}N_{h,𝐤}+2Im\lambda _0^{}b_{e,𝐤}^{}b_{h,𝐤}^{}$$
$$2\underset{𝐪}{}Img_{h,𝐪}Ab_{h,𝐤}^{}d_{h,\sigma },$$
(34)
where $`\gamma _{e(h),s}`$ denotes the carrier scattering rate due to mechanisms not considered in the Hamiltonian. $`\lambda _0^{}b_{e,𝐤}^{}b_{h,𝐤}^{}`$ represents the polarization for creating an electron-hole pair in the wetting layer. $`g_{e,𝐪}Ab_{e,𝐤}^{}d_{e,\sigma }`$ ($`g_{h,𝐪}Ab_{h,𝐤}^{}d_{h,\sigma }`$) are terms due to the process for transferring one electron (hole) from the QD to the wetting layer by absorbing one optical phonon. To terminate the equations of motion \[Eqs. (33) and (34)\], we need to solve $`\lambda _0^{}b_{e,𝐤}^{}b_{h,𝐤}^{}`$ and $`g_{e,𝐪}Ab_{e,𝐤}^{}d_{e,\sigma }`$ ($`g_{h,𝐪}Ab_{h,𝐤}^{}d_{h,\sigma }`$). The equation of motion for $`g_{e,𝐪}Ab_{e,𝐤}^{}d_{e,\sigma }`$ leads to
$$\frac{d}{dt}g_{e,𝐪}Ab_{e,𝐤}^{}d_{e,\sigma }=(i\mathrm{\Omega }_q\gamma _{e,s})Ab_{e,𝐤}^{}d_{e,\sigma }+|g_{e,𝐪}|^2$$
$$[(1+n_q)(1(n_e+N_e))N_{e,𝐤}n_q(n_e+N_e)(1N_{e,𝐤})],$$
(35)
where $`n_q`$ denotes the phonon distribution function, defined as $`n_q(\omega _q)=1/(exp(\omega _q/k_bT)1)`$. $`n_e=d_{e,\sigma }^{}(t)d_{e,\sigma }(t)_p`$ and $`N_e=d_{e,\sigma }^{}(t)d_{e,\sigma }(t)_t`$ denote the electron occupation number of the QD due to the pumping process and tunnelling process, respectively. $`\mathrm{\Omega }_q=(ϵ_𝐤E_e\omega _q)`$. Meanwhile, the polarization obeys the following equation
$$\frac{d}{dt}\lambda _0^{}b_{e,𝐤}^{}b_{h,𝐤}^{}=(i\mathrm{\Omega }_0+\gamma _{e,s}+\gamma _{h,s})\lambda _0^{}b_{e,𝐤}^{}b_{h,𝐤}^{}$$
$$i|\lambda _0|^2(1N_{e,𝐤}N_{h,𝐤}),$$
(36)
where $`\mathrm{\Omega }_0=ϵ_{e,𝐤}+ϵ_{h,𝐤}\omega _0`$. In the steady state, we obtain
$$N_{e,𝐤}=p_{exc}(1N_{e,𝐤}N_{h,𝐤})\gamma _{e,c}[(1+n_q)(1(n_e+N_e))$$
$$N_{e,𝐤}n_q(n_e+N_e)(1N_{e,𝐤})],$$
(37)
where $`p_{exc}Im|\lambda _0|^2/(\gamma _{e,s}(\mathrm{\Omega }_0+i(\gamma _{e,s}+\gamma _{h,s})))`$ and $`\gamma _{e,c}=_𝐪Im|g_{e,𝐪}|^2/(\gamma _{e,s}(\mathrm{\Omega }_q+i\gamma _{e,s}))`$. For low pumping intensity, $`N_{e,𝐤}`$ is linearly proportional to $`p_{exc}`$. In addition, we note that $`\gamma _{e,c}`$ is fairly small due to the phonon bottleneck effect of the QD. Similarly for holes in the wetting layer, we obtain
$$N_{h,𝐤}=p_{exc}(1N_{e,𝐤}N_{h,𝐤})\gamma _{h,c}[(1+n_q)(1n_h)$$
$$N_{h,𝐤}n_qn_h(1N_{h,𝐤})],$$
(38)
where $`\gamma _{h,c}=_𝐪Im|g_{h,𝐪}|^2/(\gamma _{h,s}(\mathrm{\Omega }_q+i\gamma _{h,s}))`$ is the hole capture rate. Next, we need to solve the equation of motion for $`n_{e,\sigma }`$ and $`n_{h,\sigma }`$. We have
$$\frac{d}{dt}n_{e,\sigma }=i\mathrm{\Gamma }_en_{e,\sigma }+\underset{𝐐}{}2Im\lambda ^{}ad_{e,\sigma }^{}d_{h,\sigma }^{}$$
$$+2Img_{e,𝐪}Ab_{e,𝐤}^{}d_{e,\sigma },$$
(39)
and
$$\frac{d}{dt}n_{h,\sigma }=i\mathrm{\Gamma }_{h,s}n_{h,\sigma }+\underset{𝐐}{}2Im\lambda ^{}d_{e,\sigma }^{}d_{h,\sigma }^{}$$
$$+2Img_{h,𝐪}Ab_{h,𝐤}^{}d_{h,\sigma },$$
(40)
where $`\mathrm{\Gamma }_e=\mathrm{\Gamma }_L+\mathrm{\Gamma }_R`$ and $`\mathrm{\Gamma }_{h,s}`$ denote the electron tunnelling rate and hole-impurity scattering rate, respectively. At low temperatures, the phonon-absorption process is negligible, since $`n_q0`$. Using Eqs. (35) and (20), we obtain the steady-state solution for $`n_{e,\sigma }`$
$$n_{e,\sigma }=\frac{\gamma _{e,c}N_{e,𝐤}(1N_e)}{\gamma _{e,c}N_{e,𝐤}+R_{eh}n_h+\mathrm{\Gamma }_e}.$$
(41)
Note that the electron-hole recombination rate $`R_{eh}`$ is given by
$$R_{eh}=\mathrm{\Omega }\frac{d\omega }{\pi }\rho (\omega )|Im𝒫(\omega )|,$$
(42)
where $`\mathrm{\Omega }=4\pi n_r^3\mu _r^2/(6c^3\mathrm{}^3ϵ_0)`$ and $`n_r`$ is the refractive index of the system. $`\rho (\omega )=\omega ^2`$ arises from the density of states of photons. Similarly, we obtain
$$n_{h,\sigma }=\frac{\gamma _{h,c}N_{h,𝐤}}{\gamma _{h,c}N_{h,𝐤}+R_{eh}(n_e+N_e)+\mathrm{\Gamma }_{h,s}}.$$
(43)
Because of the condition $`\mathrm{\Gamma }_e\gamma _{e,c}`$ and weak pumping intensity, we have $`N_{i,𝐤}p_{exc}`$, and $`n_e`$ is negligible. As for $`N_e`$, we can employ Eqs. (29) and (31) to obtain
$$N_{e,\sigma }=\frac{dϵ}{\pi }\frac{\mathrm{\Gamma }_Lf_L(ϵ)+\mathrm{\Gamma }_Rf_R(ϵ)}{\mathrm{\Gamma }_L+\mathrm{\Gamma }_R}ImG_{e,\sigma }^r(ϵ).$$
(44)
Figure Captions
Fig. 1: Schematic diagram of a single electron transistor under optical pumping. $`\mathrm{\Gamma }_L`$ and $`\mathrm{\Gamma }_R`$ denote, respectively, the tunnelling rate for electrons in the QD to the source and drain electrodes. $`R_{eh}`$ denotes the electro-hole recombination rate. The vertical wavy line indicates the optical pumping, and the horizontal wavy line indicates the single-photon emission.
Fig. 2: The lowest three energy levels of a quantum dot as functions of the QD size $`b`$ for (a) electrons and (b) holes.
Fig. 3: Intralevel Coulomb interactions $`U_e`$ and $`U_h`$ and interlevel Coulomb interaction $`U_{eh}`$ as functions of the QD base length $`b`$.
Fig. 4: Electron occupation number $`N_e`$ and tunnelling current as functions of gate voltage at zero temperature for various strengths of optical excitation. Current density is in units of $`J_0=2e\times meV/h`$.
Fig. 5: Tunnelling current as a function of gate voltage at zero temperature for various strengths of optical excitation.
Fig. 6: Electron occupation number $`N_e`$ and tunnelling current as functions of gate voltage for various temperatures at fixed optical excitation strength ($`I=0.9`$). Current density is in units of $`J_0=2e\times meV/h`$.
Fig. 7: Tunnelling current as a function of applied bias at temperature $`k_BT=2meV`$ for various strengths of optical excitation. Current density is in units of $`J_0=2e\times meV/h`$.
Fig. 8: Intensities of emission spectrum for various gate voltages at fixed optical excitation strength ($`I=0.9`$).
Fig. 9: Intensities of emission spectrum for various strengths of optical excitation at fixed gate voltage ($`V_g=125mV`$).
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# On the Unification of Gauge Symmetries in Theories with Dynamical Symmetry Breaking
## I Introduction
The standard model (SM) gauge group $`G_{SM}=\mathrm{SU}(3)_c\times \mathrm{SU}(2)_w\times \mathrm{U}(1)_Y`$ has provided a successful description of both strong and electroweak interactions. Although the standard model itself predicts zero neutrino masses, its fermion content can be augmented to incorporate neutrino masses and lepton mixing. One of the great triumphs of this model was its unification of weak and electromagnetic interactions. However, as has long been recognized, there are a number of properties that this model does not explain, including the quantization of the electric charges of elementary particles, the ratios of the values of the respective standard-model gauge couplings $`g_3`$, $`g_{2L}`$, and $`g_Y`$, and the interconnected manner in which quark and lepton contributions to gauge anomalies cancel each other (separately for each generation). These deficiencies motivated the effort to construct theories with higher unification of gauge symmetries. Almost all of the work toward this goal started with the standard model, including its Higgs mechanism, and subsequently, supersymmetric extensions of this model.
In this paper we shall investigate various approaches to the partial or complete unification of gauge symmetries from a different viewpoint, incorporating the standard-model gauge group but removing the Higgs mechanism of this model and replacing it with ingredients that can produce dynamical electroweak symmetry breaking (EWSB) and also dynamical breaking of higher gauge symmetries. The origin of electroweak symmetry breaking is one of the most important outstanding questions in particle physics, and dynamical EWSB remains an interesting alternate to the Higgs approach. We shall focus, in particular, on models in which electroweak symmetry breaking is due to the formation of a bilinear condensate of new fermions interacting via an asymptotically free, vectorial gauge interaction, generically denoted technicolor (TC) tc ; dsb , that becomes strong at a scale $`\mathrm{\Lambda }_{TC}300`$ GeV. To communicate the electroweak symmetry breaking to the quarks and leptons (which are technisinglets) and generate masses for these fermions, one adds to the technicolor theory additional gauge degrees of freedom that transform technifermions into standard-model fermions and vice versa etc ; etcrev . These are denoted extended technicolor (ETC) gauge bosons.
Here we shall consider several types of unification of gauge symmetries with dynamical symmetry breaking:
1. Models that involve sufficient unification to quantize electric charge without embedding all of the three factor groups of the standard model in a (semi)simple Lie group simple . Since $`Q=T_{3L}+Y/2`$ in the standard model, a sufficient condition for this quantization is that the weak hypercharge $`Y`$ be expressed as a linear combination of generators of nonabelian gauge groups.
2. Models that attempt to unify the three standard-model gauge interactions in a simple grand unified (GU) group $`G_{GU}`$,
$$G_{GU}G_{SM}$$
(1)
and then combine this in a direct product, $`G_{ETC}\times G_{GU}`$ with the ETC gauge group, $`G_{ETC}`$. A successful model of this type would explain charge quantization and the relative sizes of SM gauge couplings (but not the relative size of the SM and ETC gauge couplings).
3. Most ambitiously, models that attempt to unify the three SM gauge interactions together with technicolor, or a larger gauge symmetry described by a group $`G_{SC}G_{TC}`$, in a simple Lie group $`G`$,
$$GG_{SC}\times G_{GU}.$$
(2)
In the models that we consider of types (1) and (2) above, the technicolor group $`G_{TC}`$ is embedded in a larger, extended technicolor group, $`G_{ETC}`$: $`G_{TC}G_{ETC}`$. As indicated by this subgroup relation, in these types of models, the infinitesimal generators of the Lie algebra of $`G_{TC}`$ close upon themselves, as do the generators of the Lie algebra of $`G_{ETC}`$. Furthermore, in these models $`[G_{ETC},G_{SM}]=0`$, so that the ETC gauge bosons do not carry any SM quantum numbers. In contrast, in models of type (3), although $`[G_{TC},G_{SM}]=0`$, the commutators of the (linear combinations of) generators of the Lie algebra of $`G`$ that transform technicolor indices to SU(3)<sub>c</sub> and SU(2)<sub>L</sub> indices and vice versa (corresponding to the ETC gauge bosons) generate the full Lie algebra of $`G`$, so that there is no ETC group, per se, that forms a subgroup of $`G`$ that is smaller than $`G`$ itself. In these models these ETC gauge bosons generically do carry SM quantum numbers and the corresponding generators of $`G`$ do not commute with the generators of $`G_{SM}`$. As will be seen below, the origin of standard-model fermion generations is different in models of types (1) and (2), on the one hand, and models of type (3) on the other.
The present paper is organized as follows. In section II we review some relevant properties of technicolor and extended technicolor. In Section III we consider a partial unification model of type (1) in which charge is quantized and all symmetry breaking is dynamical, and we address the question of how one might try to unify this model further. Sections IV and V contain analyses and critical assessments of models of type (2) and (3). In Section VI we discuss the issue of how to break a unified gauge symmetry dynamically. Section V contains some concluding remarks and ideas for further work. Certain general formulas that are used throughout the paper are contained in an appendix.
## II Technicolor and Extended Technicolor Models
In this section we discuss some relevant properties of technicolor and extended technicolor theories. We take the technicolor group to be $`G_{TC}=\mathrm{SU}(N_{TC})`$. For models of types (1) and (2), the technifermions will comprise a standard-model family, i.e., they transform according to the following representations of $`G_{SM}`$:
$`Q_L^{ai}=\left({\displaystyle \genfrac{}{}{0pt}{}{U^{ai}}{D^{ai}}}\right)_L:(N_{TC},3,2)_{1/3,L};`$ (3)
(4)
$`U_R^{ai}:(N_{TC},3,1)_{4/3,R},D_R^{ai}:(N_{TC},3,1)_{2/3,R}`$ (5)
(6)
$`L_L^i=\left({\displaystyle \genfrac{}{}{0pt}{}{N^i}{E^i}}\right)_L:(N_{TC},1,2)_{1,L}`$ (7)
(8)
$`\{N_R^i\}:(N_{TC},1,1)_{0,R},E_R^i:(N_{TC},1,1)_{2,R}`$ (9)
(10)
(11)
where here the indices $`a`$ and $`i`$ are color and technicolor indices, respectively, and the numerical subscripts refer to weak hypercharge. Since the $`N_R^i`$ are SM-singlets and the group SU(2) is free of anomalies in gauged currents, it follows that if $`N_{TC}=2`$, then there can be more than just one $`N_R^i`$; we indicate this by the brackets and denote this number $`N_{N_R}`$. In models of type (3) we will encounter different types of technicolor fermion sectors.
For models of types (1) and (2), which have a well-defined ETC gauge group, we take this group to be
$$G_{ETC}=\mathrm{SU}(N_{ETC}).$$
(12)
For these models, a natural procedure in constructing the ETC theory is to gauge the generation index, assigning the first $`N_{gen.}`$ components of a fundamental representation of SU($`N_{ETC}`$) to be the standard-model fermions of these three generations, followed by $`N_{TC}`$ components which are the technifermions with the same standard-model quantum numbers. The fact that, a priori, the value of $`N_{gen.}`$ is arbitrary except for the requirement that the ETC theory be asymptotically free, distinguishes these types of models from models of type (3), where the origin and number of SM fermion generations are more highly constrained. Thus, models of types (1) and (2) can automatically accomodate the observed value of SM fermion generations, $`N_{gen.}=3`$, whereas, in contrast, this success is not guaranteed for a particular model of type (3). Given the way in which models of types (1) and (2) gauge the generational index, it follows that for these models
$`N_{ETC}`$ $`=`$ $`N_{gen.}+N_{TC}=3+N_{TC}.`$ (13)
The relation (13) and the requirement that $`N_{TC}2`$ for a nontrivial nonabelian SU($`N_{TC}`$) group (as required for asymptotic freedom) together imply that $`N_{ETC}5`$ for models of types (1) and (2). The minimal choice, $`N_{TC}=2`$ and hence $`N_{ETC}=5`$, has been used for a number of recent studies of ETC models at94 -kt . The choice $`N_{TC}=2`$ is motivated for a number of reasons; (a) with the one-family structure of eq. (11), amounting to $`N_{TF}=2(N_c+1)=8`$ vectorially coupled technifermions in the fundamental representation of SU(2)<sub>TC</sub>, it can yield an approximate infrared fixed point and associated slow running (“walking”) of the TC gauge coupling wtc from $`\mathrm{\Lambda }_{TC}`$ up to an ETC scale wtc , (b) it minimizes the technicolor contributions to the electroweak $`S`$ parameter scalc , and (c) it makes possible a mechanism to account for light neutrinos in an extended technicolor context nt ; lrs . Although the value $`N_{TC}=2`$ is thus favored, we often shall let $`N_{TC}`$ be arbitrary in the present paper (subject to the requirement of asymptotic freedom of the TC and ETC theories) in order to show the generality of certain results.
The condition $`[G_{ETC},G_{SM}]=0`$ in models of type (2) means that all components of a given representation of $`G_{SM}`$ transform according to the same representation of $`G_{ETC}`$ and all components of a given representation of $`G_{ETC}`$ transform according to the same representation of $`G_{SM}`$. However, this does not imply that all of the representations of $`G_{SM}`$ transform according to the same representation of $`G_{ETC}`$. For example, in Ref. kt we studied a class of ETC models in which $`[G_{ETC},G_{SM}]=0`$ and the left-handed and right-handed representations of the charge $`Q=1/3`$ quarks and techniquarks transform according to relatively conjugate representations of $`G_{ETC}`$ (which were the fundamental and conjugate fundamental representations), and similarly with the charged leptons and technileptons, while the charge $`Q=2/3`$ quarks and techniquarks of both chiralities transformed according to the same representations of $`G_{ETC}`$.
ETC models of this type (2) can be classified further according to whether (i) the ETC gauge interactions are vectorial on the SM quarks and charged leptons, or (ii) some ETC gauge interactions are chiral on these SM fermions. In both cases, the TC interaction must be vectorial; this is automatically satisfied by the SU(2)<sub>TC</sub> group since it has only (pseudo)real representations. (In this SU(2)<sub>TC</sub> case, the number of chiral doublets is $`15+N_{N_R}`$, and this must be even to avoid a global SU(2) anomaly, so $`N_{N_R}`$ must be odd.) These two options (i) and (ii) were labelled VSM and CSM in Ref. ckm , where the V and C referred to the corresponding vectorial and (relatively) conjugate ETC representations of the SM quarks and charged leptons, respectively. While both of these classes of models have promising features, neither is fully realistic. In Ref. ckm it was shown that constraints from neutral flavor-changing current processes were not as severe for VSM-type models as had been previously thought. However, these models require additional ingredients to produce mass splittings within each generation, such as $`m_t>>m_b,m_\tau `$ (without excessive contributions to the parameter $`\rho =m_W^2/(m_Z^2\mathrm{sin}^2\theta _W)`$). The CSM-type models in which charge $`1/3`$ quarks and leptons of opposite chiralities transform according to relatively conjugate representations of SU(5)<sub>ETC</sub> while the charge $`2/3`$ quarks have vectorial ETC couplings can produce these requisite intragenerational mass splittings and also some CKM mixing; however, these models do have problems with flavor-changing neutral current processes. Hence, we shall concentrate on VSM-type ETC models here. Moreover, additional ingredients are necessary to avoid overly light Nambu-Goldstone bosons.
As an illustration, for VSM models of type (2) the fermions with SM quantum numbers are assigned to representations of the group $`\mathrm{SU}(5)_{ETC}\times G_{SM}`$ which are obtained from those listed in eq. (11) by letting the index $`i`$ range over the full set of ETC indices, $`i=1,\mathrm{},N_{ETC}`$. Here and below, it will often be convenient to use the compact notation $`Q_L`$, $`L_L`$, $`u_R`$, $`d_R`$, and $`e_R`$ to denote these ETC multiplets, so that, for example, $`e_R`$, written out explicitly, is
$$e_R(e^1,e^2,e^3,e^4,e^5)_R(e,\mu ,\tau ,E^4,E^5)_R.$$
(14)
Since the fermion content of the SU(5)<sub>ETC</sub> theory is chosen so that it is asymptotically free b0cases , as the energy scale decreases from large values, the ETC coupling increases in strength. Eventually, this coupling becomes large enough to produce fermion condensates, and the ETC sector is constructed to be a chiral gauge theory, so that a bilinear fermion condensate generically self-breaks the ETC gauge symmetry. In order to obtain the desired sequential breaking of SU(5)<sub>ETC</sub> to SU(2)<sub>TC</sub>, Refs. at94 -kt incorporated an additional asymptotically free gauge interaction which becomes strongly coupled at roughly the same energy scale as the ETC interaction. This additional interaction was called “hypercolor” (HC) and the corresponding gauge group was chosen to be SU(2)<sub>HC</sub>. The symmetry breaking of SU(5)<sub>ETC</sub> occurs as a combination of self-breaking and couplings to the auxiliary strongly interacting group, SU(2)<sub>HC</sub>. The fermions involved in producing the ETC symmetry-breaking condensates are nonsinglets under the ETC group and are singlets under the standard-model group; they include both singlets and nonsinglets under the HC group. The determination of which condensation channels are dynamically favored is a difficult, nonperturbative problem involving strong coupling. The procedure makes use of the “most attractive channel” (MAC) criterion rds and tools such as approximate solutions of the Schwinger-Dyson equation for the relevant fermion propagator gap -tcvac (see appendix). As the energy scale decreases, the first breaking occurs at a scale denoted $`\mathrm{\Lambda }_1`$, where $`\mathrm{SU}(5)_{ETC}\mathrm{SU}(4)_{ETC}`$; here the first generation fermions split off from the rest in each ETC multiplet. Since this is the highest ETC symmetry-breaking scale, we shall label it more generally as $`\mathrm{\Lambda }_{ETC,max}`$. Similarly, one has the successive breakings $`\mathrm{SU}(4)_{ETC}\mathrm{SU}(3)_{ETC}`$ at $`\mathrm{\Lambda }_2`$, where the second-generation fermions split off, and $`\mathrm{SU}(3)\mathrm{SU}(2)_{TC}`$ at $`\mathrm{\Lambda }_3`$, where the third-generation fermions split off, leaving the exact residual technicolor gauge symmetry. This can account, at least approximately, for the observed fermion masses while satisfying other constraints such as those from flavor-changing neutral current processes, if one takes $`\mathrm{\Lambda }_110^3`$ TeV, $`\mathrm{\Lambda }_210^2`$ TeV, and $`\mathrm{\Lambda }_34`$ TeV, as in Refs. at94 -kt . The generational ETC scales are bounded above by the requirement that the resultant SM quark and lepton masses of the $`j`$’th generation, $`m_{f_j}\eta _j\mathrm{\Lambda }_{TC}^3/\mathrm{\Lambda }_j^2`$ be sufficiently large (where $`\eta _j`$ is an enhancement factor present in theories with walking technicolor wtc and can be of order $`\mathrm{\Lambda }_3/\mathrm{\Lambda }_{TC}`$.)
An important general feature of ETC theories, illustrated in the specific models mentioned above, is that the ETC symmetry-breaking scales are far below the conventional grand unification scale $`M_{GU}10^{16}`$ GeV. Recall that, before engaging in detailed calculations of gauge coupling evolution, one knows that it would be difficult to have a generic grand unification scale lower than $`10^{15}10^{16}`$ GeV without producing excessively rapid nucleon decay. Therefore, insofar as one studies the possibility of unifying standard-model and technicolor gauge symmetries, it is not technicolor itself, but rather the higher symmetry associated with extended technicolor, that enters into this unification. That is, one must take account of the fact that the effective theory at energy scales far below the grand unification scale is already invariant under a larger symmetry involving gauge degrees of freedom transforming technicolor indices to standard-model (color and electroweak) indices.
## III Example of Partial Unification of SM Gauge Symmetries and Attempt at Higher Unification
In this section we consider the model of type (1) from Ref. lrs ; nag , which successfully achieves the important goals of electric charge quantization via partial unification of standard-model gauge interactions with all symmetry breaking dynamical, and we investigate how one might try to unify it further to be a model of type (2). A sufficient condition for the quantization of electric charge $`Q`$ (or equivalently, the quantization of weak hypercharge $`Y`$, given the $`Q=T_{3L}+Y/2`$ relation) is that $`Q`$ be expressed as a linear combination of generators of nonabelian gauge groups. An early realization of this condition was provided by the Pati-Salam unification of color SU(3)<sub>c</sub> with U(1)<sub>B-L</sub> in SU(4)<sub>PS</sub>, where $`B`$ and $`L`$ denote baryon and lepton number ps . In this type of theory, the strong and electroweak gauge groups are enlarged to the group
$$G_{422}=\mathrm{SU}(4)_{PS}\times \mathrm{SU}(2)_L\times \mathrm{SU}(2)_R.$$
(15)
where $`\mathrm{SU}(2)_L\mathrm{SU}(2)_w`$ of the SM. The left- and right-handed SM fermions of each generation are assigned to the representations
$$\left(\begin{array}{cc}u^{ia}& \nu ^i\\ d^{ia}& e^i\end{array}\right)_\chi ,\chi =L,R,$$
(16)
transforming as (4,2,1) and (4,1,2), respectively, under the group $`G_{422}`$. The superscripts $`a`$ and $`i`$ in eq. (16) refer, as before, to color and generation. This fermion content thus requires the addition of three right-handed neutrinos to the standard model. The SU(4)<sub>PS</sub> gauge symmetry is vectorial, while the SU(2)<sub>L</sub> and SU(2)<sub>R</sub> symmetries are chiral. The electric charge operator is given by
$`Q`$ $`=`$ $`T_{3L}+T_{3R}+(1/2)(BL)`$ (17)
$`=`$ $`T_{3L}+T_{3R}+(2/3)^{1/2}T_{PS,15}`$ (19)
$`=`$ $`T_{3L}+T_{3R}+(1/6)\mathrm{diag}(1,1,1,3),`$ (21)
where $`T_{PS,15}=(2\sqrt{6})^1\mathrm{diag}(1,1,1,3)`$ is the third diagonal generator in the SU(4)<sub>PS</sub> Lie algebra. The (quantized) hypercharge generator is $`Y=T_{3R}+(2/3)^{1/2}T_{PS,15}`$. Since we will analyze the question of unification for the gauge couplings of the group $`G_{422}`$, we show their explicit normalization via the covariant derivative ,
$`D_\mu `$ $`=`$ $`_\mu ig_{PS}𝐓_{PS}𝐀_{PS,\mu }`$ (24)
$`ig_{2L}𝐓_L𝐀_{L,\mu }ig_{2R}𝐓_R𝐀_{R,\mu }.`$
The model of Refs. lrs ; nag uses the gauge group
$$\mathrm{SU}(5)_{ETC}\times \mathrm{SU}(2)_{HC}\times G_{422}$$
(25)
with the fermion representations
$$(5,1,4,2,1)_L,(5,1,4,1,2)_R.$$
(26)
(This model also contains fermions that are singlets under $`G_{422}`$.) In addition to the successful quantization of electric charge and partial unification of quarks with leptons (and techniquarks with technileptons), the SU(4)<sub>PS</sub> gauge interactions connecting quarks and leptons gives mass to the $`P^0`$ and $`P^3`$ Nambu-Goldstone bosons corresponding to the generators $`I_{2V}\times T_{PS,15}`$ and $`(T_3)_{2V}\times T_{PS,15}`$ binsik , where $`I`$ denotes the identity and the subscript $`2V`$ refers to vectorial isospin, SU(2)<sub>V</sub>.
In the model of Ref. lrs , as the energy decreases below a scale $`\mathrm{\Lambda }_{PS}\stackrel{>}{}\mathrm{\Lambda }_110^6`$ GeV, the $`G_{422}`$ gauge symmetry is broken to $`G_{SM}`$ by the formation of a bilinear fermion condensate. This value for $`\mathrm{\Lambda }_{PS}`$ satisfies experimental constraints such as those from upper limits on right-handed charged weak currents and on the branching ratio for the decays $`K_L\mu ^\pm e^{}`$. The matching relations for the (running) coupling constants at this scale $`\mathrm{\Lambda }_{PS}`$ are $`g_3=g_{PS}`$ and lrs
$$\frac{1}{g_Y^2}=\frac{1}{g_{2R}^2}+\frac{2}{3g_{PS}^2}.$$
(27)
Further symmetry breaking at lower scales is the same as in the model with gauge group (12). The relation between the electromagnetic coupling $`e`$ and the above gauge couplings is
$$\frac{1}{e^2}=\frac{1}{g_{2L}^2}+\frac{1}{g_Y^2}=\frac{1}{g_{2L}^2}+\frac{1}{g_{2R}^2}+\frac{2}{3g_{PS}^2}$$
(28)
from which one can calculate the weak mixing angle $`\mathrm{sin}^2\theta _W=e^2/g_{2L}^2`$ as
$$\mathrm{sin}^2\theta _W=\left[1+\frac{g_{2L}^2}{g_{2R}^2}+\frac{2g_{2L}^2}{3g_{PS}^2}\right]^1.$$
(29)
This partial gauge coupling unification is consistent with the precision determination of the three SM gauge couplings lrs . Evolving the SM couplings from $`m_Z`$ to $`\mathrm{\Lambda }_{PS}`$, one finds, at the latter scale, the values $`\alpha _3=0.064`$, $`\alpha _{2L}=0.032`$, and $`\alpha _{PS}=0.008`$ (where $`\alpha _j=g_j^2/(4\pi )`$), so that the matching equations can be satisfied with $`\alpha _{2R}(\mathrm{\Lambda }_{PS})0.013`$, i.e., $`g_{2R}/g_{2L}0.64`$ at this scale.
The results of Refs. lrs ; nag motivate one to investigate the possibility of embedding the three factor groups of $`G_{422}`$ in a simple group $`G_{GU}`$, which could form a direct product with $`G_{ETC}`$ (and possible other groups such as an auxiliary hypercolor group). This would promote this model of type (1) to a model of type (2). We observe that $`\mathrm{SU}(4)\mathrm{SO}(6)`$ and $`\mathrm{SU}(2)\times \mathrm{SU}(2)\mathrm{SO}(4)`$, so SO(10) contains, as a maximal subgroup, the direct product $`\mathrm{SO}(6)\times \mathrm{SO}(4)`$. Similarly, there is a natural embedding of the three-fold direct product $`G_{422}`$ as a maximal subgroup in SO(10). A necessary condition for this unification is that the three gauge couplings, $`g_{PS}`$, $`g_{2L}`$, and $`g_{2R}`$ could plausibly evolve, as the energy scale increases, to the single SO(10) coupling $`g`$. For an exploratory study of the feasibility of this, it will be sufficient to use one-loop renormalization group evolution equations. Some relevant general formulas are listed in the appendix. We need the leading coefficients in the beta functions for each factor group for the energy interval above $`\mathrm{\Lambda }_{PS}`$. These are, for SU(4)<sub>PS</sub>,
$`b_0^{(PS)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}[444(N_{gen.}+N_{TC})]`$ (30)
$`=`$ $`{\displaystyle \frac{1}{3}}(324N_{TC})`$ (32)
and, for SU(2)<sub>L</sub> and SU(2)<sub>R</sub>,
$`b_0^{(2L)}=b_0^{(2R)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}[224(N_{gen.}+N_{TC})]`$ (33)
$`=`$ $`{\displaystyle \frac{1}{3}}(104N_{TC})`$ (35)
where the formulas are given for general $`N_{TC}`$ to show the fact that our conclusions concerning unification hold for arbitrary values of this parameter. We find that, with the fermion content as specified above, the couplings $`g_{PS}`$, $`g_{2L}`$, and $`g_{2R}`$ do not unify at any higher energy scale. In particular, since $`\alpha _{2L}`$ and $`\alpha _{2R}`$ are unequal at $`\mathrm{\Lambda }_{PS}`$ and have the same beta functions, the respective $`\alpha _{2L}^1`$ and $`\alpha _{2L}^1`$ evolve as a function of $`\mathrm{ln}\mu `$ as two parallel lines, which precludes unification. This is still true if one augments the fermion content of the hypothetical SO(10) theory, since the representations of SO(10) treat the SU(2)<sub>L</sub> and SU(2)<sub>R</sub> subgroups symmetrically. Thus, we find that it appears to be difficult to increase the partial unification of the SM gauge symmetries in this model to a full unification of these symmetries in a direct product group containing $`\mathrm{SO}(10)\times \mathrm{SU}(5)_{ETC}`$. Nevertheless, the model of Refs. lrs ; nag does provide an example of partial unification of SM gauge symmetries explaining charge quantization in a fully dynamical framework.
## IV Prospects for Models with a $`G_{ETC}\times G_{GU}`$ Symmetry Group
### IV.1 Evolution of SM Gauge Couplings in an ETC Framework
In this section we assess the prospects for attempting to unify the three gauge groups of the standard model, SU(3)<sub>c</sub>, SU(2)<sub>w</sub>, and U(1)<sub>Y</sub>, in a simple grand unified group $`G_{GU}`$ which forms a direct product with the ETC group (and possibly other groups such as hypercolor) at a high scale, $`M_{GU}`$. In terms of the classification given at the beginning of the paper, these are models of type (2). Here and elsewhere in the paper the adjective “grand unified” is used with its historical meaning, as referring to the unification of the three SM gauge interactions only, not additional interactions such as (extended) technicolor. We assume that at $`M_{GU}`$, $`G_{GU}`$ breaks to the three-fold direct product group comprising $`G_{SM}`$, so that in the interval of energies extending downward from $`M_{GU}`$ to the highest ETC scale, $`\mathrm{\Lambda }_{ETC,1}10^6`$ GeV, the effective field theory is invariant under $`G_{ETC}\times G_{SM}`$ (times possible auxiliary groups such as hypercolor). As before, we take $`G_{ETC}=\mathrm{SU}(N_{ETC})`$ and, to show the generality of our results, we keep $`N_{TC}`$ arbitrary (subject to the requirement of asymptotic freedom for the ETC and TC group).
A prerequisite for this unification is the condition that the three SM gauge couplings unify at the hypothetical scale $`M_{GU}`$. The normalization of the abelian gauge coupling is determined by the embedding of the fermions with SM quantum numbers in the unified group, $`G_{GU}`$. We shall assume that $`G_{GU}`$ is either of the well-known unification groups SU(5) gg or SO(10) guts , with the usual assignments of SM fermions; in both cases, the U(1) coupling that unifies with $`g_3`$ and $`g_{2L}`$ is $`g_1=\sqrt{5/3}g_Y`$. We will take $`\mu =m_Z`$ as the starting point for the evolution of the SM couplings to higher scales. For our analysis, it will be adequate to use the one-loop approximations to the respective beta functions, as given in eqs. (204) and (205) of the appendix; these depend on the leading-order coefficients $`b_0^{(j)}`$ for each factor group $`G_j`$. It will also be sufficient to take the top quark to be dynamical at the electroweak scale, i.e., to include its contribution in the calculation of the beta functions. The values of the $`b_0^{(j)}`$ for the standard model with its Higgs boson are well known: $`b_0^{(3)}=(1/3)(332N_q)=7`$, $`b_0^{(2)}=(1/3)(22N_d1/2)=19/6`$, and $`b_0^{(1)}=(3/5)b_0^{(Y)}=41/10`$, where $`N_q=2N_{gen.}=6`$ denotes the number of active quarks and $`N_d=N_{gen.}(N_c+1)=12`$ denotes the number of SU(2)<sub>w</sub> doublets. It is also well known that, if one evolves these couplings individually without further new physics at intermediate scales, they do not unify at any one scale.
To calculate the evolution of the SM gauge couplings in the framework of an ETC theory, we first remove the SM Higgs and, for energies above $`\mu \mathrm{\Lambda }_{TC}300`$ GeV, where the technifermions are active, we add their contributions to the $`b_0^{(j)}`$. The deletion of the Higgs boson from the theory removes a term $`1/6`$ from $`b_0^{(2)}`$, which becomes $`b_0^{(2)}=10/3`$, and a term $`1/6`$ from $`b_0^{(Y)}`$, so that $`b_0^{(1)}`$ becomes $`b_0^{(1)}=4`$ (and leaves $`b_0^{(3)}`$ unchanged). A caveat is that, even if the SM gauge couplings are small at a given scale $`\mu `$, their evolution may still be significantly affected by nonperturbative, strong couplings of the SM-nonsinglet technifermions. Since we start our integration of the renormalization group equations at $`\mu =m_Z`$, which is comparable to, and, indeed, slightly less than, the technicolor scale, $`\mathrm{\Lambda }_{TC}300`$ GeV, these strong technifermion interactions produce some uncertainty in the evolution of the SM gauge couplings. This is a consequence of the fact that in the beta function calculations, one treats the technifermions as weakly interacting, but this is only a good approximation for $`\mu >>\mathrm{\Lambda }_110^6`$ GeV.
Generalizing $`N_q`$ to refer to both quarks and techniquarks, and letting $`N_d`$ denote the total number of SU(2)<sub>L</sub> doublets, we calculate
$`b_0^{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(11N_c2N_q)`$ (36)
$`=`$ $`{\displaystyle \frac{1}{3}}[334(N_{gen.}+N_{TC})]`$ (38)
$`=`$ $`7{\displaystyle \frac{4}{3}}N_{TC}`$ (40)
$`b_0^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(11N_wN_d)`$ (41)
$`=`$ $`{\displaystyle \frac{1}{3}}[22(N_c+1)(N_{gen.}+N_{TC})]`$ (43)
$`=`$ $`{\displaystyle \frac{10}{3}}{\displaystyle \frac{4}{3}}N_{TC}`$ (45)
and
$`b_0^{(1)}`$ $`=`$ $`{\displaystyle \frac{3}{5}}b_0^{(Y)}={\displaystyle \frac{4}{3}}(N_{gen.}+N_{TC})`$ (46)
$`=`$ $`4{\displaystyle \frac{4}{3}}N_{TC}.`$ (48)
As is evident in these results, the respective beta functions, and, in particular, the leading coefficients $`b_0^{(j)}`$, depend on $`N_{TC}`$ only through the combination $`N_{gen.}+N_{TC}=N_{ETC}`$. Consequently, the addition of the one family of technifermions to the fermion content of the standard model is equivalent to the addition of $`N_{TC}`$ additional generations of SM fermions. Now we recall that the addition of one or more (complete) generations of SM fermions leaves the differences $`(\mathrm{\Delta }b_0)_{ij}b_0^{(i)}b_0^{(j)}`$, $`ij=12,13,23`$ invariant gu . (Explicitly, $`(\mathrm{\Delta }b_0)_{32}=11/3`$, $`(\mathrm{\Delta }b_0)_{31}=11`$, and $`(\mathrm{\Delta }b_0)_{21}=22/3`$.) Hence, for a given set of values of $`\alpha _j`$, $`j=1,2,3`$ at $`\mu =m_Z`$, the scales $`\mu _{ij}`$ where $`\alpha _i=\alpha _j`$ are independent of $`N_{TC}`$, in the same way as these differences are independent of $`N_{gen.}`$. Therefore, just as there was no gauge coupling unification in the SM and the SM without a Higgs, so also, this remains true for the theory with one family of technifermions added. For reference, we note that the scales $`\mu _{ij}`$ at which pairwise equalities of couplings occur are roughly $`\mu _{23}10^{18}`$ GeV, $`\mu _{13}10^{14.5}`$ GeV, and $`\mu _{12}10^{12.8}`$ GeV.
### IV.2 An Implication of Unification Involving $`G_{ETC}\times G_{GU}`$
Even if one were able to achieve unification of the three standard-model gauge symmetries in a simple group $`G_{GU}`$ which commutes with the ETC group in a model of type (2), one would encounter another problem. The unification groups $`G_{GU}`$ that have been studied have the property that for at least one fermion $`f`$, both $`f_L`$ and $`f_L^c`$ are contained in a given representation of $`G_{GU}`$. For example, in the SU(5) model of Ref. gg , $`u_L`$ and $`u_L^c`$ are both assigned to the rank-2 antisymmetric tensor representation, the $`10_L`$. In the SO(10) model, all of the left-handed components and conjugates of the right-handed components of quarks and leptons are contained in the 16-dimensional spinor representation. The property that the full gauge group contains the direct product $`G_{ETC}\times G_{GU}`$ means that $`[G_{ETC},G_{GU}]=0`$. It follows that any $`f_L`$ and $`f_L^c`$ that belong to the same representation of $`G_{GU}`$ also transform according to the same representation $`_{ETC}`$ of $`G_{ETC}`$, or equivalently, $`f_L`$ and $`f_R`$ transform according to relatively conjugate representations $`_{ETC}`$ and $`\overline{}_{ETC}`$ lrs . Here we use $`f_L`$ and $`f_R`$ to refer to all of the fermions of these respective chiralities with the same SM quantum numbers, as illustrated for $`e_R`$ in eq. (14). This strongly suppresses the mass $`m_{f^j}`$ that is produced, for a given generation $`j`$. This is true for the $`Q=2/3`$ quarks ($`f=u`$) in the case $`G_{GU}=\mathrm{SU}(5)_{GU}`$ and for all of the SM fermions in the case $`G_{GU}=\mathrm{SO}(10)`$. This suppression would render it very difficult to obtain adequate quark and lepton masses, in particular, the top quark mass. (Recall that the CSM models of Ref. kt always used vectorial ETC representations for $`Q=2/3`$ quarks and techniquarks.) Moreover, there would be serious problems associated with excessive ETC contributions to flavor-changing neutral current processes, as was shown in Ref. kt for the CSM case where the $`Q=1/3`$ quarks and leptons of opposite chiralities transform according to relatively conjugate representations of $`G_{ETC}`$.
## V On the Unification of SM and TC Gauge Symmetries in a Simple Group
### V.1 General
Here we consider models of type (3), which attempt to achieve the unification of the three standard-model gauge symmetries with a group $`G_{SC}`$ which contains technicolor,
$$G_{SC}G_{TC},$$
(49)
and possibly also some generational symmetries, in a unified gauge symmetry described by the group, $`G`$, as specified in eq. (2). The physics is invariant under the symmetry group $`G`$ at energies above the unification scale, $`M_{GU}`$, and this symmetry breaks at $`M_{GU}`$. The subscript $`SC`$ indicates that the $`G_{SC}`$ gauge interaction becomes strongly coupled at a scale which is roughly comparable to conventional ETC scales, although it is small and perturbative at the high scale $`M_{GU}`$. These are the most ambitious of all of the three types of models considered here. They entail the unification of the three SM gauge couplings and the SC gauge coupling at $`M_{GU}`$. Thus, the absence of this gauge coupling unification would, by itself, be enough to exclude such models. However, the analysis of the evolution of the relevant gauge couplings is more complicated for these models because of the nonperturbative behavior and associated dynamical symmetry breaking that occurs at intermediate energy scales between $`m_Z`$ and the hypothesized $`M_{GU}`$; as a consequence of this, one cannot use the perturbative evolution equation (205) for all of the relevant couplings. We will discuss this further below.
As soon as one hypothesizes a technicolor gauge symmetry as a dynamical mechanism for electroweak symmetry breaking, it is natural to explore the idea of trying to unify technicolor with the three SM factor groups - color, weak isospin, and weak hypercharge - in a simple group. The motivations for this are similar to the motivations for the original grand unification program, including a unified description of the fermion representations and an explanation of the relative coupling strengths at lower energies, including, in particular, the property that the TC interaction becomes strongly coupled at a scale $`\mathrm{\Lambda }_{TC}300`$ GeV, which is essentially the electroweak scale, and which is about $`10^3`$ times larger than the scale $`\mathrm{\Lambda }_{QCD}0.2`$ GeV at which the color coupling becomes large. In this approach, one would, a priori, plausibly hope to explain the large ratio $`\mathrm{\Lambda }_{TC}/\mathrm{\Lambda }_{QCD}`$ and hence also $`m_Z/\mathrm{\Lambda }_{QCD}`$, as a consequence of moderate differences in the relevant beta functions, together with the property of slow, logarithmic running over the energy interval between $`M_{GU}`$ and the highest scale at which some couplings grow to be of order unity. Early studies on the possibility of grand unification of technicolor and standard-model symmetries were fs ; fr .
Perhaps the simplest notion of unification of technicolor with the three standard-model gauge interactions would be to have a simple gauge group $`G`$ that contains all four of these interactions as subgroups, $`GG_{TCSM}`$, where $`{}_{TCSM}{}^{}=G_{TC}\times G_{SM}`$, such that $`G`$ breaks to $`G_{TCSM}`$ at a high scale $`M_{GU}`$. This hypothetical theory would be constructed so that the technicolor beta function would be more negative than the SU(3)<sub>c</sub> beta function, $`\beta _{TC}<\beta _{SU(3)_c}<0`$, and hence, as the energy scale decreases, the technicolor gauge coupling would become sufficiently large to cause a technifermion condensate at a scale $`\mathrm{\Lambda }_{TC}`$ well above the scale $`\mathrm{\Lambda }_{QCD}`$ at which the SU(3)<sub>c</sub> coupling gets large and produces the $`\overline{q}q`$ condensate. However, this approach is excluded immediately by the fact that the gauge bosons in $`G`$ that transform technifermions into the technisinglet standard-model fermions and hence communicate the electroweak symmetry breaking to the latter and give them masses lie in the coset $`G/G_{TCSM}`$ and hence pick up masses of order $`M_{GU}`$. The effective ETC scale would thus be the grand unification scale, $`M_{GU}`$, resulting in standard-model fermion masses that are much too small; for example, for the illustrative value $`M_{GU}=10^{16}`$ GeV, these fermion masses would be of order $`\mathrm{\Lambda }_{TC}^3/M_{GU}^210^{25}`$ GeV. It should be noted that this early approach to the unification of technicolor and SM gauge symmetries led to the inference that $`N_{TC}`$ had to be greater than $`N_c=3`$. But since this attempt at unification is immediately ruled out by its failure to obtain fermion masses of adequate size, its requirement concerning $`N_{TC}`$ is only of historical interest, and, indeed, many recent TC models at94 -kt use $`N_{TC}=2`$ for the reasons that we have discussed in Section II bcond .
Here we consider a different approach to this goal of unification, in which the ETC gauge bosons have masses in the usual ETC range, and not all of the fermion generations arise from the representations of the unified group $`G`$ but instead, some arise from sequential symmetry breaking of a smaller subgroup of $`G`$ at ETC-type scales. Let us denote $`N_{gh}`$ and $`N_g\mathrm{}`$ as the numbers of standard-model fermion generations arising from these two sources, respectively, where the subscripts $`gh`$ and $`g\mathrm{}`$ refer to generations from the representation content of the high-scale symmetry group and from the lower-scale breaking. Together, these equal the observed number of SM fermion generations:
$$N_{gen.}=3=N_{gh}+N_g\mathrm{}.$$
(50)
Note that in this approach involving sequential symmetry breaking, one does not calculate the beta functions of the low-energy SC or TC sectors by enumerating the fermion content at the unification scale since some subset of these fermions would be involved in condensates formed at intermediate energy scales, hence would gain dynamical masses of order these scales, and would be integrated out before the energy decreases to the scale relevant for the evolution of SC or TC gauge couplings. It should also be remarked that at this stage the number $`N_g\mathrm{}`$ is only formal; that is, we set up a given model so that, a priori, it can have the possibility that a subgroup of $`G`$ such as $`G_{SC}`$ might break in such a manner as to peel off $`N_g\mathrm{}`$ SM fermion generations. In fact, we will show that, at least in the models that we study, it is very difficult to arrange that this desired symmetry breaking actually takes place.
The requirement that the ETC gauge bosons have masses of the necessary scales means that $`G`$ cannot break to the direct product group $`G_{TCSM}`$ at the unification scale $`M_{GU}`$, and also cannot break at this scale to the larger subgroup $`{}_{SCSM}{}^{}=G_{SC}\times G_{SM}`$. Instead, $`G`$ must break to a direct product group such that one or more of the factor groups that are residual symmetries between $`M_{GU}`$ and $`\mathrm{\Lambda }_{ETC,max}`$ contain gauge bosons that transform technifermions into technisinglet standard-model fermions, i.e. are ETC gauge bosons. As the energy scale decreases, this intermediate symmetry should break at $`\mathrm{\Lambda }_{ETC,max}`$ so that some of the ETC gauge bosons get masses of this order, and so forth for other lower sequential ETC scales.
To provide an explicit context for our analysis, let us consider unifying the SU($`N_{SC})`$ symmetry containing technicolor with the SM gauge symmetries by using a group $`G=\mathrm{SU}(N)`$ as in eq. (2), with
$$N=N_{SC}+N_c+N_w=N_{SC}+5.$$
(51)
Thus, $`GG_{SCSM}`$. Here we shall take $`G_{TC}=\mathrm{SU}(N_{TC})`$, $`G_{SC}=\mathrm{SU}(N_{SC})`$, and $`G_{GU}`$ to be the group SU(5)<sub>GU</sub> of Ref. gg . The fermion representations are determined by the structure of the fundamental representation, which we take to be
$$\psi _R=\left(\begin{array}{c}(N^c)^\tau \\ d^a\\ e^c\\ \nu ^c\end{array}\right)_R$$
(52)
where $`d`$, $`e`$, and $`\nu `$ are generic symbols for the fermions with these quantum numbers. Thus, the indices on $`\psi _R`$ are ordered so that the indices in the SC set, which we shall denote $`\tau `$, take on the values $`\tau =1,\mathrm{}N_{SC}`$ and then the remaining five indices are those of the $`5_R`$ of SU(5)<sub>GU</sub> gg . The components of $`N_R^c`$ transform according to the fundamental representation of SU($`N_{SC}`$), are singlets under SU(3)<sub>c</sub> and SU(2)<sub>w</sub>, and have zero weak hypercharge (hence also zero electric charge). Our choice to write these components as $`(N^c)_R^\tau `$ instead of $`N_R^\tau `$ is a convention. The quantum numbers of components of any representation of $`G`$ are determined by the structure of the fundamental representation (52). This structure is concordant with the direct product in eq. (2) and the corresponding commutativity property
$$[G_{SC},G_{GU}]=0$$
(53)
which, since $`G_{SC}G_{TC}`$, implies
$$[G_{TC},G_{GU}]=0.$$
(54)
These properties have important consequences for fermion masses. We recall the theorem from Ref. lrs discussed in Section IVB, that $`[G_{ETC},G_{GU}]=0`$ implies that for one or more fermions $`f`$, since $`f_L`$ and $`f_L^c`$ are both contained in the same representation of $`G_{GU}`$, $`f_L`$ and $`f_R`$ transform according to relatively conjugate representations of $`G_{ETC}`$. By the same argument, the commutativity property (53) implies that for one or more fermions $`f`$, since $`f_L`$ and $`f_L^c`$ are both contained in the same representation of $`G_{GU}`$, $`f_L`$ and $`f_R`$ transform according to relatively conjugate representations of $`G_{SC}`$ and hence also of $`G_{TC}`$. For the case $`G_{GU}=\mathrm{SU}(5)_{GU}`$ on which we focus here, this includes the charge 2/3 techniquarks. In the models of types (1) and (2) discussed in Section IVB, this would lead to the strong suppression of the masses of the TC-singlet SM fermions with these quantum numbers, i.e., the charge 2/3 quarks. Here, it will also lead to the strong suppression of certain SM fermion masses, but because the ETC vector bosons carry color and charge in the present type-(3) models, the fermions with suppressed masses will be leptons.
Corresponding to the subgroup decomposition $`GG_{SC}\times G_{GU}`$, the Lie algebra of $`G`$ contains subalgebras for $`G_{SC}`$ and $`G_{GU}`$. There are also (linear combinations of) generators of $`G`$ that transform the $`N_{SC}`$ indices $`\tau =1,\mathrm{}N_{SC}`$ to the last five indices in the fundamental representation, and vice versa. The gauge bosons corresponding to these generators include some of the ETC gauge bosons and have nontrivial SM quantum numbers. We label the basic transitions as
$$(N^c)_R^\tau d_R^a+V_a^\tau $$
(55)
$$(N^c)_R^\tau e_R^c+(U^{})^\tau $$
(56)
and
$$(N^c)_R^\tau \nu _R^c+(U^0)^\tau $$
(57)
where $`\tau =1,\mathrm{},N_{SC}`$, and the $`V_a^\tau `$ and $`\left(\genfrac{}{}{0pt}{}{U^0}{U^{}}\right)^\tau `$ are ETC gauge bosons. Under the group $`G_{SCSM}`$ these transform according to
$$V_a^\tau :(N_{SC},\overline{3},1)_{2/3},Q=1/3$$
(58)
and
$$\left(\genfrac{}{}{0pt}{}{U^0}{U^{}}\right)^\tau :(N_{SC},1,2)_{1/2}.$$
(59)
These are thus quite different from the ETC gauge bosons of models (1) and (2), which carry no SM quantum numbers. It is important to note that the (commutators of the) generators to which these ETC gauge bosons correspond do not close to yield a subalgebra smaller than the full Lie algebra of $`G`$, so that there is no ETC group $`G_{ETC}`$, as such. This is analogous to the fact that the (linear combinations of) generators that transform color indices to electroweak indices in the conventional SU(5) grand unified theory do not close to form an algebra smaller than the Lie algebra of SU(5) itself. This is one of the features that distinguish these models of type (3) from the models of types (1) and (2), which do have well-defined ETC gauge groups.
To delineate the remaining ETC gauge bosons, let us divide the SU($`N_{SC}`$) indices into (a) the ones for SU($`N_{TC}`$), say $`\tau =1,\mathrm{},N_{TC}`$, and (b) the ones for the coset $`G_{SC}/G_{TC}`$, $`\widehat{\tau }=N_{TC}+1,\mathrm{},N_{SC}`$. Consider the process $`(N^c)_R^\tau (N^c)_R^{\widehat{\tau }}+V_{\widehat{\tau }}^\tau `$; the $`V_{\widehat{\tau }}^\tau `$’s constitute the remaining ETC gauge bosons. These carry no SM quantum numbers and are similar in this regard to the ETC gauge bosons of models (1) and (2).
Given that $`G`$ cannot break completely to $`G_{TCSM}`$ or $`G_{SCSM}`$ at the unification scale $`M_{GU}`$ in a viable model, we next investigate which subgroup it could break to at this scale. The breaking pattern must be such as to satisfy the upper limits on the decays of protons and otherwise stable bound neutrons. The gauge bosons of $`G_{GU}`$ that contribute to these decays are the set $`\left(\genfrac{}{}{0pt}{}{X_a}{Y_a}\right)`$, (where $`a`$ is a color index) which transform as $`(\overline{3},2)_{5/6}`$ under $`\mathrm{SU}(3)_c\times \mathrm{SU}(2)_w\times \mathrm{U}(1)_Y`$ (whence $`Q_X=4/3`$, $`Q_Y=1/3`$), and their adjoints. These 12 gauge bosons span the coset $`\mathrm{SU}(5)_{GU}/G_{SM}`$ and must gain masses of order $`M_{GU}`$. A priori, the breaking at $`M_{GU}`$ could leave an invariant subgroup $`\mathrm{SU}(3)_c\times G_{SCW}`$, where $`G_{SCW}`$ is a group of transformations on the $`N_{SC}`$ indices of $`\mathrm{SU}(N_{SC})`$ and the two electroweak indices of SU(2)<sub>w</sub> which naturally takes the form $`\mathrm{SU}(N_{SC}+2)`$. However, it appears quite difficult to construct a model of this sort because of the conflicting requirements that the $`G_{SCW}`$ coupling get large, as is necessary for self-breaking (see further below) and that the SU(2)<sub>w</sub> coupling, which is supposed to have a small, perturbative value of $`\alpha _2=\alpha _{GU}`$ at the presumed unification scale, $`M_{GU}`$ and then evolve to its similarly small, perturbative value of $`\alpha _2(m_Z)=0.034`$ at the electroweak scale.
An alternative is that the breaking of $`G`$ at $`M_{GU}`$ would leave an invariant subgroup $`\mathrm{SU}(2)_w\times G_{SCC}`$, where $`G_{SCC}`$ is a group of transformations on the $`N_{SC}`$ indices of $`\mathrm{SU}(N_{SC})`$ and the three color indices of SU(3)<sub>c</sub>, which would naturally take the form $`\mathrm{SU}(N_{SCC})`$ with
$$N_{SCC}=N_{SC}+N_c=N_{SC}+3.$$
(60)
Thus,
$$\mathrm{SU}(N_{SCC})\mathrm{SU}(N_{SC})\times \mathrm{SU}(3)_c.$$
(61)
All representations of SU($`N_{SCC}`$) are determined by the fundamental representation, which follows directly from eq. (52); again, it is convenient to write this as a right-handed field,
$$\left(\begin{array}{c}(N^c)^\tau \\ d^a\end{array}\right)_R.$$
(62)
As is evident from this, the components of a representation of $`\mathrm{SU}(N_{SCC})`$ do not, in general, have the same weak hypercharge $`Y`$, so this representation does not have a well-defined value of $`Y`$. Considerations of the relative sizes of gauge couplings would favor this option over the one involving $`G_{SCW}`$ because the color SU(3)<sub>c</sub> coupling $`\alpha _3`$ increases substantially from the presumed unification value at $`M_{GU}`$ to $`\alpha _3(m_Z)=0.118`$ at the electroweak scale. To begin with, several possibilities could be envisioned for the symmetry breaking of $`\mathrm{SU}(N_{SCC})`$. One would be that, as the energy scale decreases from $`M_{GU}`$, $`\alpha _{SCC}`$ becomes sufficiently large at a scale $`\mathrm{\Lambda }_{ETC,max}`$ for the breaking $`\mathrm{SU}(N_{SCC})\mathrm{SU}(N_{SC})\times \mathrm{SU}_c`$ to occur, and then, if $`N_{SC}>N_{TC}`$, the further sequential breakings of SU($`N_{SC}`$), eventually yielding the exact symmetry group SU($`N_{TC})`$, peeling off the $`N_g\mathrm{}`$ SM fermion generations. A different scenario would be one in which, as the energy scale decreases below $`M_{GU}`$, SU($`N_{SCC}`$) breaks into smaller simple groups in a sequence of $`N_g\mathrm{}`$ steps, and then the residual smaller simple group finally splits into the direct product $`\mathrm{SU}(N_{TC})\times \mathrm{SU}(3)_c`$. A third possibility would be a combination of these two types of breakings, in which $`\mathrm{SU}(N_{SCC})`$ breaks to smaller simple groups in $`k<N_g\mathrm{}`$ stages, then splits to a two-fold direct product group one of whose factor groups is SU(3)<sub>c</sub>, and then the other factor group sequentially self-breaks $`N_g\mathrm{}k`$ times to the residual exact SU($`N_{TC}`$) group. In all of these scenarios, the $`V_a^\tau `$’s and $`V_\tau ^{}^\tau `$ would have masses in the usual ETC range, from $`\mathrm{\Lambda }_1`$ down to $`\mathrm{\Lambda }_3`$ while the $`(U^0)^\tau `$ and $`(U^{})^\tau `$ would have masses of order $`M_{GU}`$. (Hence, although $`(U^0)^\tau `$ and $`(U^{})^\tau `$ are formally ETC gauge bosons, they would play a negligible role in producing masses for SM fermions.)
However, this scenario with a $`G_{SCC}`$ gauge symmetry characterizing the effective theory between $`M_{GU}`$ and $`\mathrm{\Lambda }_{ETC,max}`$ also encounters a complication. Consider, for example, the version in which SU($`N_{SCC})`$ would break to $`\mathrm{SU}(N_{SC})\times \mathrm{SU}(3)_c`$ at roughly $`\mathrm{\Lambda }_{ETC,max}10^6`$ GeV. The matching conditions for gauge couplings imply that at this energy scale the SU(3)<sub>c</sub> and $`G_{SC}`$ couplings are equal, since they inherit the value that the SU($`N_{SCC})`$ coupling had just above $`\mathrm{\Lambda }_{ETC,max}`$. But, assuming that the unified theory with gauge group $`G`$ is self-contained, i.e., there is no direct product at the high scale $`M_{GU}`$ with an auxiliary group like hypercolor, the only mechanism for dynamical symmetry breaking below $`M_{GU}`$ is self-breaking. This implies that at the above energy scale of $`10^6`$ GeV, $`\alpha _{SCC}`$ should be $`O(0.1)O(1)`$, depending on the attractiveness of the relevant fermion condensation channel(s) as measured by the respective values of $`\mathrm{\Delta }C_2`$, defined in eq. (209) in the appendix. This is difficult to reconcile with the SU(3)<sub>c</sub> beta function, which has leading coefficient $`b_0^{(3)}=7`$ due to the SM fermions, and, in the energy interval above $`\mathrm{\Lambda }_{TC}`$, $`b_0^{(3)}=7(4/3)N_{TC}`$ for one family of massless technifermions. That is, at an energy scale of $`10^6`$ GeV, the value of $`\alpha _3`$ would already be larger than its value at the electroweak scale, which would imply that it would have to decrease, rather than increasing, as the scale $`\mu `$ decreases from $`10^6`$ GeV to $`\mu =m_Z`$, i.e., that the SU(3)<sub>c</sub> sector would have to be non-asymptotically free in this interval. Another problem is that in the models that we have constructed and studied, we find that the $`\mathrm{SU}(N_{SCC})`$ theory is unlikely to break in the necessary manner; to explain this, it is first necessary to describe the fermion representations, which we do below.
Let us proceed with the construction and critical evaluation of the prospects for this class of unified models. The following conditions are equivalent:
$$G_{SC}=G_{TC}N_{gh}=3,N_g\mathrm{}=0,$$
(63)
i.e., in this case, all of the fermion generations would arise from the fermion representations of $`G`$. The other formal possibility is that $`N_g\mathrm{}1`$ so that the coset $`G_{SC}/G_{TC}`$ is nontrivial, with $`G_{SC}`$ containing some gauged generational structure. For generality, it should be noted that although we use $`G_{GU}`$ to classify fermion representations, the breaking of $`G`$ may be such that the lower-energy effective field theory is not actually symmetric under a direct product group in which $`G_{GU}`$ occurs as a factor group. From the property (2), it follows that the ranks satisfy
$$r(G)r(G_{SC})+r(G_{GU})$$
(64)
where, with our assumption that $`G_{SC}=\mathrm{SU}(N_{SC})`$, we have $`r(G_{SC})=N_{SC}1`$. In addition to the subgroup decomposition (2), we will also use the subgroup decomposition
$$\mathrm{SU}(N)\mathrm{SU}(2)_w\times \mathrm{SU}(N_{SCC}).$$
(65)
We shall assume that at energies below the unification scale $`M_{GU}`$ all subsequent breaking of gauge symmetries is dynamical. Since $`M_{GU}`$ is not very far below the Planck scale, which certainly constitutes an upper limit to the possible validity of the theory, owing to the lack of inclusion of quantum gravity, it is not clear that one needs to assume that the initial breaking of $`G`$ is dynamical. We shall comment on this further below. The dynamical symmetry breaking at energies below $`M_{GU}`$ can be classified as being of two general types: (i) self-breaking (“tumbling”), in which an asymptotically free chiral gauge symmetry group has an associated coupling that becomes large enough to produce a fermion condensate that breaks the gauge symmetry, and (ii) induced breaking in which a gauge symmetry is weakly coupled, but is broken by the formation of condensates involving fermions that are nonsinglets under a strongly coupled group (which is the way that electroweak symmetry is broken by technifermion condensates); (iii) a combination of the two, as in the sequential breaking of the SU(5)<sub>ETC</sub> symmetry in Refs. at94 -kt . Since we only consider unification in a single, simple group $`G`$ here, we are led to focus on self-breaking below $`M_{GU}`$. Note that with our choice of the minimal GU group as SU(5) with rank 4, the inequality (64) becomes $`r(G)r(G_{SC})+4`$. Our choice $`G=\mathrm{SU}(N)`$ with $`N`$ given by eq. (51) satisfies this inequality as an equality; i.e., we are choosing the minimal $`G`$ for a given value of $`N_{SC}`$. In order to maintain the nonabelian structure of the TC group and hence the asymptotic freedom that leads to confinement and the formation of the EWSB bilinear fermion condensate, we require that $`N_{TC}2`$. This yields the inequality
$$N7$$
(66)
and $`r(G)6`$.
For any of the possible types of sequential breakings of $`G_{SCC}`$ and/or $`G_{SC}`$ described above that produce the $`N_g\mathrm{}`$ SM fermion generations, one has
$$N_g\mathrm{}=N_{SCC}(N_{TC}+N_c).$$
(67)
In particular, if $`G_{SCC}`$ first splits to $`G_{SC}\times \mathrm{SU}(3)_c`$ and $`G_{SC}`$ then sequentially breaks to produce these $`N_g\mathrm{}`$ generations, then
$$N_g\mathrm{}=N_{SC}N_{TC}.$$
(68)
The requirement that $`N_{TC}2`$ in order for the technicolor interactions to be asymptotically free, combined with eq. (68), implies
$$N_g\mathrm{}N_{SC}2.$$
(69)
Thus, for $`N_{SC}=2`$ all of the SM fermion generations must arise via $`N_{gh}`$.
We next specify the fermion representations of $`G=\mathrm{SU}(N)`$. Without loss of generality, we shall usually deal with left-handed fermions (or antifermions). In order to avoid fermion representations of SU(3)<sub>c</sub> and SU(2)<sub>w</sub> other than those experimentally observed, namely singlets and fundamental or conjugate fundamental representations, one restricts the fermions to lie in $`k`$-fold totally antisymmetrized products of the fundamental or conjugate fundamental representation of SU($`N`$) g79 ; we denote these as $`[k]_N`$ and $`[\overline{k}]_N=\overline{[k]}_N`$. The notational correspondence with Young tableaux is (suppressing the dependence on $`N`$), $`[1]\text{ }\text{ }\text{ }\text{ }\text{ }`$, $`[2]\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }`$, etc. Some elementary properties of the representation $`[k]_N`$ are listed in the appendix; these include its dimensionality and expressions for the Casimir invariants $`C_2([k]_N)`$ and $`T([k]_N)`$. A set of (left-handed) fermions $`\{f\}`$ transforming under $`G`$ is thus given by
$$\{f\}=\underset{k=1}{\overset{N1}{}}n_k[k]_N$$
(70)
where $`n_k`$ denotes the multiplicity (number of copies) of each representation $`[k]_N`$. We use a compact vector notation
$$𝐧(n_1,\mathrm{},n_{N1})_N.$$
(71)
If $`k=N\mathrm{}`$ is greater than the integral part of $`N/2`$, we shall work with $`[\overline{\mathrm{}}]_N`$ rather than $`[k]_N`$; these are equivalent with respect to SU($`N`$) (see eq. (196) in the appendix). An optional additional constraint would be to require that the numbers in the set of $`n_k`$, $`k=1,\mathrm{},N1`$ have no common factors greater than unity, i.e., the greatest common divisor $`GCD(\{n_k\})=1`$. This might be viewed as a kind of irreducibility condition. Although we will not impose this condition here, the two candidate models that we consider that have $`GCD(\{n_k\})2`$ are excluded anyway because the SCC theory is not asymptotically free. A fermion field corresponding to $`[k]_N`$ is denoted generically by $`\psi _L^{i_1\mathrm{}i_k}`$. It will sometimes be convenient to deal with the charge-conjugate right-handed field.
Before proceeding, it is appropriate to summarize the requirements on the choice of fermion representations:
1. The theory must contain a mechanism to break the unified $`G`$ gauge symmetry, eventually down to the symmetry group operative above the electroweak scale, $`\mathrm{SU}(N_{TC})\times G_{SM}`$. The breaking scales must be such as to obey upper bounds on the decay rate for protons and bound neutrons.
2. The contributions from various fermions to the total SU($`N`$) gauge anomaly must cancel each other, yielding zero gauge anomaly.
3. The resultant TC-singlet, SM-nonsinglet left-handed fermions must comprise a set of generations, i.e., must have the form $`N_{gen.}[(1,\overline{5})_L+(1,10)_L]`$, where the first number in parentheses signifies that these are TC-singlets and the second number denotes the dimension of the SU(5)<sub>GU</sub> representation.
4. For a fully realistic model, one requires $`N_{gen.}=3`$.
5. In order to account for neutrino masses, one needs to have TC-singlet, electroweak-singlet neutrinos to produce Majorana neutrino mass terms that can drive an appropriate seesaw nt . In the present context, these are also singlets under $`SU(5)_{GU}`$.
6. The model must contain ETC gauge bosons with masses in the general range from a few TeV to $`10^3`$ TeV so as to produce acceptable SM fermion masses. As explained above, a plausible way to satisfy this requirement is for $`G`$ to break to the subgroup (65) containing the factor group $`G_{SCC}`$ which contains both $`\mathrm{SU}(N_{TC})`$ and SU(3)<sub>c</sub> and is naturally SU($`N_{SCC}`$) with $`N_{SCC}`$ given by eq. (60). Thus, the effective field theory at energy scales between $`M_{GU}`$ and $`\mathrm{\Lambda }_{ETC,max}10^6`$ GeV is invariant under this direct product (65). The dynamics should be such that SU($`N_{SCC}`$) breaks at ETC scales, in one of the ways delineated above, eventually yielding the residal exact symmetry group $`\mathrm{SU}(2)_{TC}\times \mathrm{SU}(3)_c`$, with the requisite three SM fermion generations emerging. If at least some of the stages of this process involve self-breaking, then the SCC sector should be an asymptotically free chiral gauge theory.
7. If $`N_{SC}>N_{TC}`$, then there must be a mechanism to break $`\mathrm{SU}(N_{SC})`$ down in $`N_g\mathrm{}`$ stages to $`\mathrm{SU}(N_{TC})`$. Again, if this is to be a self-breaking, then the SC sector should be an asymptotically free chiral gauge theory so that the associated coupling will increase sufficiently as the energy scale decreases to produce the requisite condensate(s).
8. The color SU(3)<sub>c</sub> interaction must be asymptotically free in the energy interval at and below the electroweak scale, where the associated coupling has been measured.
9. The technicolor interaction must asymptotically free, so that the associated gauge coupling will increase sufficiently, as the energy scale decreases, to produce a technifermion condensate and break the electroweak symmetry; further, the technicolor symmetry must be vectorial so that the technifermions are confined and the technifermion condensate does not self-break $`G_{TC}`$.
10. When evolved down to the low energies, the respective SM gauge couplings must agree with their measured values.
Let us check that the first and sixth of these constraints can be simultaneously satisfied. This requires that one confirm that the masses of the (mass eigenstates corresponding to the interaction eigenstates) $`V_a^\tau `$ needed for their role as ETC vector bosons are consistent with the upper bounds on the decays of protons and bound neutrons. These decays are induced by the $`s`$-channel transitions (1) $`u^a+u^bX_c`$, (2) $`u^a+d^bY_c`$ (and, if the theory contained another vector boson, $`\mathrm{\Xi }_c`$ with $`Q_\mathrm{\Xi }=2/3`$, also (3) $`d^a+d^b\mathrm{\Xi }_c`$), where $`a,b,c`$ are (different) color indices. Corresponding transitions in the $`t`$ and $`u`$ channel also contribute. Among the $`G_{SC}`$-nonsinglet gauge bosons, the only ones that transform in the right way under color and electric charge to contribute to these decays are the $`V_a^\tau `$, with charge 1/3. Among these, the subset with SC indices $`\tau `$ in $`G_{TC}`$ cannot contribute to these decays, since the $`G_{TC}`$ technicolor symmetry is exact, but the quarks in a nucleon are technisinglets. If $`N_g\mathrm{}=0`$, then $`G_{SC}=G_{TC}`$, so all of the SC indices are in $`G_{TC}`$. If $`N_g\mathrm{}1`$, then there is also a subset of $`V_a^{\widehat{\tau }}`$ with indices $`\widehat{\tau }`$ in the coset $`G_{SC}/G_{TC}`$. The exact symmetries (color, electric charge, and technicolor) allow these to mix with the $`Y_a`$, via one-loop and higher-loop nondiagonal propagator corrections, so that the actual vector boson mass eigenstates would be, for $`N_g\mathrm{}=1`$,
$`V_{a,heavy}`$ $`=`$ $`\mathrm{cos}\omega Y_a+\mathrm{sin}\omega V_a^{\widehat{\tau }}`$ (72)
$`V_{a,light}`$ $`=`$ $`\mathrm{sin}\omega Y_a+\mathrm{cos}\omega V_a^{\widehat{\tau }}`$ (74)
and similarly in the case where $`N_g\mathrm{}2`$. Since this mixing would be forbidden at energy scales where $`G_{SC}`$ or $`G_{SCC}`$ is still an exact symmetry, and since the maximum of the relevant breaking scales is of order the highest ETC scale, $`\mathrm{\Lambda }_{ETC,max}10^6`$ GeV, it follows that
$$|\omega |\frac{\mathrm{\Lambda }_{ETC,max}}{M_{GU}}1.$$
(75)
Hence, the mixing would lead to a diagram for nucleon decay with a propagator for $`V_{a,light}`$, $`1/\mathrm{\Lambda }_{ETC,max}^2`$ multiplied by the mixing factors for each of the two vertices involving SM fermions, which have size $`\stackrel{<}{}(\mathrm{\Lambda }_{ETC,max}/M_{GU})^2`$. The product is then $`\stackrel{<}{}1/M_{GU}^2`$, the same as for the usual contribution from $`V_{a,heavy}`$. Hence (given that $`M_{GU}`$ is sufficiently large so that the usual contributions to nucleon decay are not excessive) this mixing does not significantly increase the rate of nucleon decay. This shows that these ETC gauge bosons can, indeed, have ETC-scale masses, as required to give SM fermions their masses.
We next consider the constraint that there be no anomaly in gauged currents of the unified theory invariant under the group $`G`$. For this purpose, we define a $`(N1)`$-dimensional vector of anomalies
$$𝐚=(A([1]_N),\mathrm{},A([N1]_N))$$
(76)
where $`A([k]_N)`$ is given in eq. (197) in the appendix. Then the constraint that there be no $`G`$ gauge anomaly is the condition
$$𝐧𝐚=0.$$
(77)
which is a diophantine equation for the components of the vector of multiplicities $`𝐧`$, subject to the constraint that the components $`n_k`$ are non-negative integers. Geometrically, if $`𝐚`$ and $`𝐧`$ were vectors in $`^{N1}`$, then the solution set of eq. (77) would be the $`(N2)`$-dimensional subspace of $`^{N1}`$ orthogonal to the vector $`𝐚`$; the situation here is more complicated because of the diophantine requirement that $`n_k`$ is a nonnegative integer. The actual solution is also subject to additional conditions, as we shall discuss shortly.
The most natural way to satisfy the third requirement, that the TC-singlet, SM-nonsinglet fermions should form a well-defined set of SM generations, is to impose this separately on the subset of these fermions that arise from the fermion representations of $`G`$ and on the complementary subset that arise from the sequential symmetry breaking of SU($`N_{SCC})`$. The fermion representations of $`G`$ transform according to $`(_{SC},_{GU})`$ with respect to the subgroup decomposition (2). In terms of these, the condition on the former subset yields the two conditions
$$N_{(1,\overline{5})}=N_{(1,10)}$$
(78)
and, for nonsinglet $`_{SU(5)_{GU}}`$,
$$N_{(1,_{SU(5)_{GU}})}=0\mathrm{if}_{SU(5)_{GU}}\overline{5}\mathrm{or}10,$$
(79)
i.e., the number of TC-singlet left-handed (anti)fermions from these representations transforming as $`\overline{5}`$ and $`10`$ of SU(5)<sub>GU</sub> must be equal and the theory must not contain any other TC-singlet, SM-nonsinglet fermions. Regarding the fermions in the complementary subset, we note that the breaking of SU($`N_{SCC}`$) can be viewed as the breaking of the SU($`N_{SC}`$) part of this group, since the SU(3)<sub>c</sub> part remains unbroken. For these, the requirements that we impose are the analogues of eqs. (78) and (79) with the multiplets considered to refer to $`(_{TC},_{GU})`$.
As noted earlier, in models of type (3), for a given choice of fermion representations of the unified group $`G`$, it is not guaranteed that the resultant (TC-singlet) SM fermions come in well-defined generations, and even for a choice which does satisfy this constraint, it is not guaranteed that the model can accomodate three such generations. In both respects, these models are different from models of type (1) and (2), where one can automatically satisfy both of these conditions.
We now incorporate the constraint that the standard-model fermions that arise from the fermion representations of $`G`$ comprise well-defined generations. Each of these SM fermion generations is equivalent to the set $`\{(1,\overline{5})+(1,10)\}`$ of representations of left-handed fermions, under the direct product $`\mathrm{SU}(N_{SC})\times \mathrm{SU}(5)_{GU}`$. A $`(1,\overline{5})`$ representation arises in two ways (i) from a $`[N_{SC}+4]_N=[N1]_N[\overline{1}]_N`$ representation, $`\psi _{i,L}`$, when $`i`$ takes values in SU(5)<sub>GU</sub>; and (ii) from a $`[4]_N`$ representation, $`\psi _L^{i_1i_2i_3i_4}`$, when all of the indices take on values in SU(5)<sub>GU</sub>. If one were to choose $`N_{SC}=0`$ and hence $`N=5`$, these sources would coincide; for the relevant case of a nonabelian TC group, for which $`N_{SC}N_{TC}2`$, they constitute two different sources. Hence, for $`N_{SC}2`$,
$$N_{(1,\overline{5})}=n_{_{N_{SC}+4}}+n_4$$
(80)
where, equivalently, $`n_{_{N_{SC}+4}}=n_{_{N1}}`$. A (1,10) representation also arises in two ways: (i) from a $`[2]_N`$ representation, $`\psi _L^{i_1i_2}`$, when both of the indices take on values in SU(5)<sub>GU</sub>, and (ii) from a $`[N_{SC}+2]_N`$ representation, when $`N_{SC}`$ of the indices take on values in SU($`N_{SC}`$), thereby producing a singlet under this group, and the remaining two indices take on values in SU(5)<sub>GU</sub>. (Since $`[N_{SC}+2]_N[\overline{N3}]_N`$, one can equivalently describe source (ii) as arising from $`\psi _{i_1i_2i_3,L}`$ when all of the three indices take on values in SU(5)<sub>GU</sub>.) Again, for $`N_{SC}=0`$ and hence $`N=5`$, these sources (i) and (ii) coincide; for the relevant nonabelian case $`N_{TC}2`$ and hence $`N_{SC}2`$, they are different, so that
$$N_{(1,10)}=n_2+n_{_{N_{SC}+2}}.$$
(81)
Thus, the requirement that the left-handed SC-singlet, SM-nonsinglet (anti)fermions comprise equal numbers of $`(1,\overline{5})`$ and (1,10)’s implies the condition
$$n_{_{N_{SC}+4}}+n_4=n_2+n_{_{N_{SC}+2}}$$
(82)
and the number of SM fermion generations $`N_{gh}`$ produced by the representations of $`G`$ is given by either side of this equation;
$$N_{gh}=n_2+n_{_{N_{SC}+2}}.$$
(83)
The remaining $`N_g\mathrm{}`$ generations of SM fermions arise via the breaking of $`G_{SCC}`$ and/or $`G_{SC}`$.
We next determine the implications of the constraint (79) excluding SC-singlet fermions that have unphysical nonsinglet SM transformation properties. These can be identified in terms of their SU(5)<sub>GU</sub> representations. Since the only possibilities for nonsinglet $`[k]_5`$ are $`[1]_5=5`$, $`[2]_5=10`$, $`[3]_5=\overline{10}`$, and $`[4]=\overline{5}`$, these unphysical representations are the $`(1,5)`$ and $`(1,\overline{10})`$. Now a $`[1]_N`$ representation, $`\psi _L^i`$, yields a (1,5) when the index $`i`$ takes on values in SU(5)<sub>GU</sub>. Further, a $`[N_{SC}+1]_N`$ representation, $`\psi _L^{i_1\mathrm{}i_{N_{SC}+1}}`$, also yields a (1,5) when $`N_{SC}`$ of the indices take on values in SU($`N_{SC}`$), thereby yielding an SC-singlet, and the one remaining index takes on values in SU(5)<sub>GU</sub>. The requirement that there be no (1,5)’s is therefore
$$n_1=0,n_{_{N_{SC}+1}}=0.$$
(84)
A $`[3]_N`$ representation, $`\psi _L^{i_1i_2i_3}`$, yields a $`(1,\overline{10})`$ if all of the three indices take values in SU(5)<sub>GU</sub>. In addition, a $`[N_{SC}+3]_N`$ representation, $`\psi _L^{i_1\mathrm{}i_{N_{SC}+3}}`$, yields a $`(1,\overline{10})`$ in the case when $`N_{SC}`$ of the indices take on values in SU($`N_{SC}`$) and the remaining three indices take on values in SU(5)<sub>GU</sub>. (Since $`[N_{SC}+3]_N[\overline{N2}]_N`$, the latter source is equivalent to $`\psi _{i_1i_2,L}`$ with both indices taking on values in SU(5)<sub>GU</sub>.) Hence, the requirement that there be no $`(1,\overline{10})`$ is
$$n_3=0,n_{_{N_{SC}+3}}=0.$$
(85)
The representations $`[1]_N`$ and $`[N1]_N[\overline{1}]_N`$, when decomposed with respect to the subgroup (65), will yield a term (2,1), i.e., a doublet under SU(2)<sub>w</sub> which is a singlet under SU($`N_{SCC}`$). This is a lepton doublet, such as $`\left(\genfrac{}{}{0pt}{}{\nu _e}{e}\right)_L`$. The fact that it is a singlet under SU($`N_{SCC}`$) means that neither of the component fermions couples directly to the ETC gauge bosons, and hence both have strongly suppressed masses. This is acceptable for the neutrino, but $`m_e`$ is only about a factor of 10 less than $`m_u`$, so this strong mass suppression may be problematic for the electron. In order to prevent this, one could require that $`n_1=0`$ and $`n_{N1}=0`$. We shall not do this here, but it should be borne in mind that models with nonzero values for $`n_1`$ and/or $`n_{N1}`$ will have this property.
In the present type-3 models, SC-singlet, SM-singlet fermions (1,1), which can be identified as electroweak-singlet neutrinos, arise, in general, from two sources: (i) $`[N_{SC}]_N`$, when all of the $`N_{SC}`$ indices take values in SU($`N_{SC}`$); and (ii) $`[5]_N`$, when all of the indices take values in SU(5)<sub>GU</sub>. In the special case $`N_{SC}=5`$, these each contribute. Hence,
$$N_{(1,1)}=n_{_{N_{SC}}}+n_5.$$
(86)
If $`G_{SC}`$ is the same as the TC group, then this is the full set of TC-singlet, electroweak-singlet neutrinos, so that the the right-hand side of eq. (86) should be nonzero. If $`N_{SC}>N_{TC}`$, then electroweak-singlet neutrinos can also arise from SC-nonsinglet representations when the SC group breaks to the TC group.
The constraint concerning the breaking of $`\mathrm{SU}(N_{SCC})`$ and the behavior of the SU($`N_{TC})`$ technicolor group that is operative below the lowest ETC breaking scale entails several parts. Since the TC theory emerges from the breaking of the SCC theory, and since at the unification scale the squared coupling $`\alpha _{SCC}=\alpha `$ is small, one wants the SCC theory to be asymptotically free in order for $`\alpha _{SCC}`$ to increase as the energy scale decreases, yielding, after breakings, a TC coupling that is sufficiently large to produce the eventual technifermion condensate. The asymptotic freedom of the SCC theory is required if, as assumed here, one or more of the sequential breakings of the SU($`N_{SCC}`$) theory are self-breakings. The constraint $`\beta _{TC}<\beta _{SU(3)_c}<0`$ in the original approach to the unification of TC and SM gauge symmetries does not appear here because the SU(3)<sub>c</sub> group is subsumed within the SU($`N_{SCC}`$) group in the relevant range of energies $`\mathrm{\Lambda }_{ETC,max}<\mu <M_{GU}`$. Note that, since $`N_{SCC}5`$, it follows that $`b_0^{(SCC)}b_0^{(2)}=N_{SCC}23`$, where here $`b_0^{(2)}`$ is the leading coefficient of the only other nonabelian subgroup of $`G`$, namely SU(2)<sub>w</sub> (cf. eq. (65)). Hence, provided that the SU($`N_{SCC}`$) theory is asymptotically free, its beta function is more negative than that of the SU(2)<sub>w</sub> sector, as should be the case to account for the observed values of the SM gauge couplings at the electroweak scale.
Before analyzing specific models, we mention some challenges that can be anticipated at the outset. First, for models with $`N_{gh}=3`$ so that $`N_{SC}=N_{TC}`$, there is a single effective ETC mass scale that governs the origin of the SM fermion masses, namely the scale at which the breaking
$$\mathrm{SU}(N_{SCC})G_{TC}\times \mathrm{SU}(3)_c$$
(87)
occurs. (There would, in general, also be a U(1) factor; here we concentrate on the nonabelian symmetries.) The only other scale that enters into the generation of the masses for these SM fermions is the technicolor scale. With only a single effective ETC scale to work with, one cannot satisfactorily reproduce the observed SM fermion mass hierarchy. In models with $`N_g\mathrm{}1`$, there is, a priori, the formal possibility of having enough ETC mass scales to produce the observed SM generational hierarchy, making use of the sequential breaking scales of $`G_{SCC}`$ and/or $`G_{SC}`$ that are supposed to yield the $`N_g\mathrm{}`$ additional generations.
However, when we actually examine these models based on a simple unification group $`G`$ (without auxiliary groups such as hypercolor), we find that the requisite sequential dynamical symmetry breaking of $`\mathrm{SU}(N_{SCC})`$ is unlikely to occur. The breaking of $`G_{SCC}`$ should eventually yield, after the sequential self-breaking, the residual exact nonabelian symmetry group on the right-hand side of eq. (87). Given that the $`\mathrm{SU}(N_{SCC})`$ theory is asymptotically free, the associated gauge coupling $`\alpha _{SCC}`$ will increase as the energy scale decreases from the unification scale, $`M_{GU}`$, and, when $`\alpha _{SCC}`$ is sufficiently large, the theory will form bilinear fermion condensate(s). If the $`\mathrm{SU}(N_{SCC})`$ gauge interaction is vectorlike, i.e., if (neglecting other gauge interactions) the nonsinglet fermion content of $`\mathrm{SU}(N_{SCC})`$ consists of the set of left-handed fermions $`\{_{}+\overline{}\}`$, then the most attractive channel for this condensation process is
$$\times \overline{}1,$$
(88)
i.e., it yields a condensate that is a singlet under $`\mathrm{SU}(N_{SCC})`$. Hence, the model would fail to break $`\mathrm{SU}(N_{SCC})`$ at all, let alone to the residual subgroup (87). With the same initial set of representations, one also has the channel
$$\times \overline{}Adj,$$
(89)
where here $`Adj`$ refers to the adjoint representation of SU($`N_{SCC}`$). This channel could lead to the desired breaking of SU($`N_{SCC}`$) in eq. (87). However, channel (88) is always more attractive than channel (89). From our studies of specific models, we find that in most cases where $`G_{SCC}`$ is asymptotically free, it is vectorlike, and hence, applying the MAC criterion, one would conclude that the necessary dynamical symmetry breaking would not take place. Even in a model with $`N_{SC}=5`$ and with $`𝐧`$ given in eq. (174), where SU($`N_{SCC}`$) is an asymptotically free chiral gauge theory, we find that it is unlikely to break in the desired manner.
A related problem with the dynamical symmetry breaking is that in many cases, not only does the condensation in the most attractive channel not break SU($`N_{SCC}`$), it breaks SU(2)<sub>w</sub> at a scale which is higher than the ETC scales where SU($`N_{SCC}`$) should break. A possible way to avoid undesired condensation channels of this sort could be to invoke a “generalized most attractive channel” (GMAC) criterion at94 ; nt ; ckm , which makes use of vacuum alignment and related energy minimization arguments to suggest that if the condensate formation can avoid breaking a certain symmetry, it will tcvac .
Yet another complication can occur in cases where $`N_{SC}`$ is odd so that $`N_{SCC}`$ is even, say $`N_{SCC}=2p`$. In these cases, there can occur a most attractive channel of the form
$$[p]_{2p}\times [p]_{2p}1$$
(90)
with
$$\mathrm{\Delta }C_2=2C_2([p]_{2p})=\frac{p(2p+1)}{2}=\frac{N_{SCC}(N_{SCC}+1)}{4}.$$
(91)
The associated condensate is
$$ϵ_{i_1\mathrm{}i_{2p}}\psi _L^{1_1\mathrm{}i_pT}C\psi _L^{1_{p+1}\mathrm{}i_{2p}}.$$
(92)
This condensate is symmetric (antisymmetric) under interchange of $`\psi _L^{1_1\mathrm{}i_p}`$ and $`\psi _L^{1_{p+1}\mathrm{}i_{2p}}`$ if $`p`$ is even (odd). Since the condensate (92) is invariant under $`G_{SCC}`$, it is, a fortiori, invariant under SU($`N_{SC}`$) and SU(3)<sub>c</sub>. The only way to construct the requisite SU(3)<sub>c</sub>-invariant contractions involves product(s) $`ϵ_{\mathrm{}mn}d^{\mathrm{}}d^md^n`$, each of which has weak hypercharge $`Y=2`$ (and electric charge $`1`$). Hence, this condensate violates weak hypercharge and electric charge. It may be noted that in a hypothetical world in which only the SU(2)<sub>w</sub> interaction were strongly coupled, the same kind of violation would presumably occur. Consider, say, the first two generations of lepton doublets, $`\psi _{g1,L}=\left(\genfrac{}{}{0pt}{}{\nu _e}{e}\right)`$ and $`\psi _{g2,L}=\left(\genfrac{}{}{0pt}{}{\nu _\mu }{\mu }\right)`$. These would form an SU(2)<sub>w</sub>-invariant condensate of type (92) with $`p=1`$, namely
$$ϵ_{jk}\psi _{1g,L}^{jT}C\psi _{2g,L}^k=2\nu _{eL}^TC\mu _Le_L^TC\nu _{\mu L}.$$
(93)
where $`j,k`$ are SU(2)<sub>w</sub> indices.
### V.2 $`N_{SC}=2`$, $`G=\mathrm{SU}(7)`$
We proceed to study a number of specific models of type (3) to explore their properties. We begin by considering the minimal nontrivial case, $`N_{SC}=2`$. Since this is the smallest value for a nonabelian group, it follows that $`G_{SC}=G_{TC}`$ and hence we denote $`G_{SCC}G_{TCC}`$; further, it follows that $`N_g\mathrm{}=0`$, so that an acceptable model would have to have $`N_{gh}=3`$. The vector $`𝐧`$ has the form $`𝐧=(n_1,\mathrm{},n_6)_7`$. From eqs. (84) and (85) we have $`n_1=n_3=n_5=0`$. Equation (82) yields $`n_6=n_2`$. The no-anomaly condition, eq. (77), reads $`3n_22n_4n_6=0`$; substituting $`n_6=n_4`$ in this equation gives the result $`n_2=n_4`$. Hence, $`𝐧=n_2(0,1,0,1,0,1)_7`$ Taking $`n_2=1`$ yields
$$𝐧=(0,1,0,1,0,1)_7.$$
(94)
More generally, for an SU($`N`$) group with $`N`$ odd, say $`N=2n+1`$, the chiral fermion content $`\{f\}=_{\mathrm{}=1}^n[2\mathrm{}]_N`$ is anomaly-free g79 . This property was used in Ref. fs for a study of the possible unification of TC and SM symmetries in SU(7) and SU(9). For the fermion set in eq. (94) the number of SM generations is given by eq. (83) as $`N_{gen.}=2n_2=2`$, so the requirement that $`N_{gen.}=3`$ cannot be satisfied, and this model is not acceptable. One could, nevertheless, consider it as a toy model. The simplest special case of this toy model has $`n_2=1`$, so that there would be two SM generations. Taking the next higher value, $`n_2=2`$ is not acceptable because it would yield the unphysical result of four SM generations. Applying eq. (86), we note that there is one electroweak-singlet neutrino.
With respect to the group (2), namely,
$$\mathrm{SU}(2)_{TC}\times \mathrm{SU}(5)_{GU},$$
(95)
the fermions have the following decompositions:
$$[2]_7=(1,1)+(2,5)+(1,10)$$
(96)
$$[4]_7[\overline{3}]_7=(1,\overline{5})+(2,\overline{10})+(1,10)$$
(97)
$$[6]_7[\overline{1}]_7=(2,1)+(1,\overline{5})$$
(98)
where here and below we use the equivalences $`2\overline{2}`$ for SU(2) and (196) for SU($`N`$).
The technifermions in this model are
$`U_{\tau L}^a,D_L^{\tau a},\left({\displaystyle \genfrac{}{}{0pt}{}{U^{\tau a}}{D^{\tau a}}}\right)_R`$ (99)
(100)
$`N_{\tau L},E_{\tau L},\left({\displaystyle \genfrac{}{}{0pt}{}{N_\tau }{E_\tau }}\right)_R`$ (101)
where, as before, $`\tau =1,2`$ is the technicolor index and $`a=1,2,3`$ is the color index. Note how, in accordance with our general discussion above, the left- and right-handed chiral components of the charge 2/3 techniquark, $`U_L`$ and $`U_R`$, transform according to relatively conjugate representations of SU(2)<sub>TC</sub>. In this case, since SU(2) has only (pseudo)real representations, these are equivalent, and the technicolor theory is a vectorial gauge theory.
Since both the subgroup (2) and the subgroup (65),
$$\mathrm{SU}(2)_w\times \mathrm{SU}(5)_{TCC},$$
(102)
are abstractly SU(2) $`\times `$ SU(5), the fermions have formally the same decompositions with respect to (102) as in eqs. (96)-(98), although the component fields are different. It should be noted that a fermion that is a singlet under color and technicolor, and hence is a lepton, can occur as a component of a TCC nonsinglet representation. For example, with respect to the subgroup (102), the $`[2]_7`$ yields a term
$$(1,10)=\left(\begin{array}{ccccc}0& \nu ^c& D^{11}& D^{12}& D^{13}\\ & 0& D^{21}& D^{22}& D^{23}\\ & & 0& u_3^c& u_2^c\\ & & & 0& u_1^c\\ & & & & 0\end{array}\right)_L$$
(103)
where the upper indices $`\tau a`$ on $`D^{\tau a}`$ refer to technicolor and color, and the entries left blank are equal to minus the transposed entries. The $`\nu _L^c`$ field illustrates the general point made above.
The SU(5)<sub>TCC</sub> interaction (neglecting other interactions) is vectorial, consisting of the left-handed fermion content $`2\{5+\overline{5}+10+\overline{10}\}`$. The SU(5)<sub>TCC</sub> gauge interaction is asymptotically free, with leading beta function coefficient $`b_0^{(TCC)}=13`$. Hence, as the energy scale decreases below $`M_{GU}`$, $`\alpha _{TCC}`$ increases to the point where the theory forms bilinear fermion condensates. Let us consider condensations of these fermions, with the representations classified according to the subgroup of eq. (102) as $`(_{SU(2)_w},_{SU(5)_{TCC}})`$. In this notation, the most attractive channel (MAC) is
$$(1,10)\times (2,\overline{10})(2,1).$$
(104)
If, indeed, this condensate formed, it would rule out this SU(7) model, since it would break SU(2)<sub>w</sub>, at much too high a scale; indeed, this SU(2)<sub>w</sub>-breaking scale would be greater than the ETC scales where SU(5)<sub>TCC</sub> should break, clearly an unphysical situation. The attractiveness of the channel, as measured by $`\mathrm{\Delta }C_2`$, is given by $`\mathrm{\Delta }C_2=36/5`$. If one uses as a rough guide the critical value of $`\alpha _{TCC}`$ given by the Schwinger-Dyson gap equation (eq. (208) in the appendix), one finds the value $`\alpha _{TCC}=5\pi /540.3`$. As an illustrative value, assume that $`\alpha _{TCC}=\alpha _{GU}0.04`$ at $`M_{GU}`$. Substituting the above critical value into eq. (210), one obtains the rough estimate that $`\alpha _{TCC}`$ increases to the critical value for the condensate (104) to form as the scale $`\mu `$ decreases through the value $`\mu _c(3\times 10^5)M_{GU}`$. For the hypothetical unification scale $`M_{GU}=10^{16}`$ GeV, this would mean that electroweak symmetry would be broken at $`10^{11}`$ GeV, clearly far too high a scale. In addition, the channel (104) would fail to break the SU(5)<sub>SCC</sub> group to $`\mathrm{SU}(2)_{TC}\times \mathrm{SU}(3)_c`$.
Other condensation channels which are, a priori possible, are listed below, together with their $`\mathrm{\Delta }C_2`$ values,
$$(1,10)\times (1,10)(1,\overline{5}),\mathrm{\Delta }C_2=24/5$$
(105)
$$(1,\overline{5})\times (2,5)(2,1),\mathrm{\Delta }C_2=24/5$$
(106)
$$(1,\overline{5})\times (1,10)(1,5),\mathrm{\Delta }C_2=18/5$$
(107)
$$[(2,5)\times (2,\overline{10})]_a(1,\overline{5}),\mathrm{\Delta }C_2=18/5$$
(108)
$$[(2,5)\times (2,\overline{10})]_s(3,\overline{5}),\mathrm{\Delta }C_2=18/5$$
(109)
where the subscripts $`a`$ and $`s`$ in eqs. (108) and (109) refer to antisymmetric and symmetric combinations of representations. None of these channels is acceptable in a viable model. Consider, for example, channel (105). The condensate for this channel is
$$ϵ_{ijk\mathrm{}n}\psi _L^{jkT}C\psi _L^\mathrm{}n$$
(110)
where the indices are in SU(5)<sub>TCC</sub> (with the ordering as in eq. (52)) and $`\psi _L^{jk}`$ is the fermion field transforming as (1,10) under $`\mathrm{SU}(2)_w\times \mathrm{SU}(5)_{TCC}`$. The free index must take on one of the two SU(2)<sub>TC</sub> values, $`i=1`$ or $`i=2`$, in order to avoid breaking SU(3)<sub>c</sub>; with no loss of generality, we may choose $`i=1`$. The condensate (110) is then proportional to
$$D_L^{21T}Cu_{1L}^c+D_L^{22T}Cu_{2L}^c+D_L^{23T}Cu_{3L}^c$$
(111)
where the indices on $`D^{\tau a}`$ are as in eq. (103). This condensate violates weak hypercharge and electric charge, and leaves only one unbroken technicolor index, so that technicolor becomes an abelian symmetry. Channel (106) breaks SU(2)<sub>w</sub> at too high a scale and fails to break SU(5)<sub>TCC</sub>. Channels (107) and (108) leave technicolor as an abelian symmetry. Channel (109) breaks SU(2)<sub>w</sub> in the wrong way and at too high a scale, and leaves technicolor as an abelian symmetry.
Thus, none of these condensation channels produces the desired symmetry breaking of SU(5)<sub>TCC</sub> to $`\mathrm{SU}(2)_{TC}\times \mathrm{SU}(3)_c`$. In the absence of an actual mechanism to produce this breaking, it is difficult to analyze the properties of the hypothetical resultant SU(2)<sub>TC</sub> theory. With the one SM family of technifermions listed in eq. (101), the SU(2)<sub>TC</sub> interaction would be asymptotically free, with leading beta function coefficient $`b_0^{(TC)}=2`$, but since this is smaller than the corresponding $`b_0^{(3)}`$ for SU(3)<sub>c</sub> (which is $`b_0^{(3)}=25/38.3`$ for this toy two-generation model), and $`\alpha _{TC}=\alpha _{SU(3)_c}`$ at the energy scale where SU(5)<sub>TCC</sub> splits to $`\mathrm{SU}(2)_{TC}\times \mathrm{SU}(3)_c`$, the color coupling would grow considerably faster than the technicolor coupling as the energy scale decreased, leading to the unphysical prediction that $`\mathrm{\Lambda }_{QCD}>\mathrm{\Lambda }_{TC}`$.
### V.3 $`N_{SC}=3`$, $`G=\mathrm{SU}(8)`$
Next, we consider $`N_{SC}=3`$, $`G=\mathrm{SU}(8)`$, so that $`𝐧=(n_1,\mathrm{},n_7)_8`$. In this case, a priori, one has the two options $`N_{TC}=N_{SC}=3`$ with $`N_g\mathrm{}=0`$, or $`N_{TC}=2`$ with $`N_g\mathrm{}=1`$. In general, eqs. (84) and (85) yield
$$n_1=n_3=n_4=n_6=0$$
(112)
and eq. (82) reads
$$N_{gh}=n_2+n_5=n_4+n_7.$$
(113)
The no-anomaly condition is
$$4n_25n_5n_7=0.$$
(114)
For a given value of $`N_{gh}=3N_g\mathrm{}`$, these are three nondegenerate linear equations for the three quantities $`n_2`$, $`n_5`$, and $`n_7`$. We display the formal solution, with the understanding that it is physical only for positive nonnegative integer values of the $`n_k`$:
$`n_2`$ $`=`$ $`{\displaystyle \frac{2N_{gh}}{3}}`$ (115)
$`n_5`$ $`=`$ $`{\displaystyle \frac{N_{gh}}{3}}`$ (117)
$`n_7`$ $`=`$ $`N_{gh}.`$ (119)
In order for $`n_2`$ and $`n_5`$ to be nonnegative integers, $`N_{gh}=0`$ mod 3; the value $`N_{gh}=0`$ is not permitted because this would require $`N_g\mathrm{}=3`$, but $`N_g\mathrm{}N_{SC}(N_{TC})_{min}=N_{SC}2=1`$. Hence, the only possibility is
$$N_{gh}=3,N_g\mathrm{}=0,$$
(120)
whence
$$n_2=2,n_5=1,n_7=3,$$
(121)
so that
$$𝐧=(0,2,0,0,1,0,3)_8,$$
(122)
i.e., the fermion content of the model is
$$\{f\}=2[2]_8+[5]_8+3[7]_82[2]_8+[\overline{3}]_8+3[\overline{1}]_8.$$
(123)
Further, for this case,
$$G_{SC}=G_{TC}=\mathrm{SU}(3)_{TC}$$
(124)
and
$$G_{SCC}=\mathrm{SU}(6)_{SCC}=\mathrm{SU}(6)_{TCC},$$
(125)
where, since $`G_{SC}`$ is just the technicolor group, we have indicated this explicitly in the subscript. Note that, by eq. (86),
$$N_{(1,1)}=1.$$
(126)
In Table 1 we list properties of this model and others that we have studied.
Let us analyze this SU(8) model further. With respect to the subgroup given by (2), viz.,
$$\mathrm{SU}(8)\mathrm{SU}(3)_{TC}\times \mathrm{SU}(5)_{GU},$$
(127)
we have the decomposition
$$[2]_8=(\overline{3},1)+(3,5)+(1,10)$$
(128)
$$[5]_8[\overline{3}]_8=(1,1)+(3,\overline{5})+(\overline{3},\overline{10})+(1,10)$$
(129)
$$[7]_8[\overline{1}]_8=(\overline{3},1)+(1,\overline{5}).$$
(130)
Since both technicolor and color are described by SU(3) subgroups of SU(8), the theory is formally symmetric under the interchange of technicolor indices $`i=1,2,3`$ and color indices, $`i=4,5,6`$, and eqs. (128)-(130) also describe the decomposition of the fermion representations with respect to the subgroup $`\mathrm{SU}(3)_c\times \mathrm{SU}(5)`$, where this SU(5) involves technicolor and electroweak indices. The nonsinglet fermion content under color or technicolor consists of 15 copies of $`\{3+\overline{3}\}`$. Evidently, both color and technicolor are vectorial gauge symmetries.
With respect to the subgroup given by (65), namely
$$\mathrm{SU}(8)\mathrm{SU}(2)_w\times \mathrm{SU}(6)_{TCC}$$
(131)
the fermion representations have the decompositions
$$[2]_8=(1,1)+(2,[1]_6)+(1,[2]_6)$$
(132)
$$[5]_8[\overline{3}]_8=(1,[\overline{1}]_6)+(2,[\overline{2}]_6)+(1,[\overline{3}]_6)$$
(133)
$$[7]_8[\overline{1}]_8=(2,1)+(1,[\overline{1}]_6).$$
(134)
Here and below, it is convenient to use the $`[k]_N`$ notation for larger Lie groups; the corresponding dimensionalities for representations of SU(6)<sub>TCC</sub> are $`[2]_6=15`$ and $`[\overline{3}]_6=\overline{20}`$. Thus, the SU(6)<sub>TCC</sub> theory is vectorial, with nonsinglet fermion content (neglecting other interactions) consisting of the set of (left-handed) fermions
$$4([1]_6+[\overline{1}]_6)+2([2]_6+[\overline{2}]_6+[\overline{3}]_6)$$
(135)
(Here, with respect to SU(6), $`[3]_6[\overline{3}]_6`$.) Let us assume that SU(8) breaks in such a manner as to yield an effective theory at lower energies that has a $`\mathrm{SU}(2)_w\times \mathrm{SU}(6)_{TCC}`$ symmetry (ignoring an abelian factor). Since the TCC theory is asymptotically free, with leading beta function coefficient $`b_0^{(TCC)}=12`$, the TCC coupling increases as the energy scale decreases. Because the SU(6)<sub>TCC</sub> theory is vectorial, when the energy scale decreases sufficiently that $`\alpha _{TCC}O(1)`$, the TCC interaction will naturally form SU(6)<sub>TCC</sub>-invariant fermion condensates rather than breaking to $`\mathrm{SU}(3)_c\times \mathrm{SU}(3)_{TC}`$, as is necessary in order to separate color and the TC interaction. The most attractive channel, written in terms of representations of $`\mathrm{SU}(2)_w\times \mathrm{SU}(6)_{TCC}`$, together with its $`\mathrm{\Delta }C_2`$ value, is
$$(1,[\overline{3}]_6)\times (1,[\overline{3}]_6)(1,1),\mathrm{\Delta }C_2=\frac{21}{2}.$$
(136)
This condensation channel is of the form of (the conjugate of) (90) with $`p=3`$, $`N_{SC}=6`$. By our general argument given above, the associated condensate violates weak hypercharge and electric charge. If it did occur, this, by itself, would rule out the present SU(8) model. The rough estimate for the corresponding critical coupling is $`\alpha _{TCC,c}2\pi /630.1`$ from the Schwinger-Dyson equation. Substituting this into eq. (210) with $`\alpha _{TCC}=\alpha `$ at $`\mu =M_{GU}`$ and using the illustrative value for the unified coupling $`\alpha =0.04`$, we find that $`\alpha _{TCC}`$ would be large enough for this condensate to form at $`\mu _c(3\times 10^5)M_{GU}`$, i.e., about $`3\times 10^{11}`$ GeV for the hypothetical unification scale $`M_{GU}=10^{16}`$.
Other possible condensation channels have smaller values of the attractiveness measure $`\mathrm{\Delta }C_2`$; they include the following, in order of descending $`\mathrm{\Delta }C_2`$:
$$(1,[2]_6)\times (2,[\overline{2}]_6)(2,1),\mathrm{\Delta }C_2=\frac{28}{3}9.3$$
(137)
$$(1,[2]_6)\times (1,[\overline{3}]_6)(1,[\overline{1}]_6),\mathrm{\Delta }C_2=7$$
(138)
$$(2,[\overline{2}]_6)\times (1,[\overline{3}]_6)(2,[\overline{5}]_6)(2,[1]_6),\mathrm{\Delta }C_2=7$$
(139)
$$(1,[\overline{1}]_6)\times (1,[2]_6)(1,[1]_6),\mathrm{\Delta }C_2=\frac{35}{6}5.8$$
(140)
$$(2,[1]_6)\times (1,[\overline{1}]_6)(2,1),\mathrm{\Delta }C_2=\frac{35}{6}.$$
(141)
Channels (137) and (141) fail to break SU(6)<sub>TCC</sub> and break SU(2)<sub>w</sub> at a scale higher than the ETC scale where SU(6)<sub>TCC</sub> should break. Channels (138) and (140) break SU(6)<sub>TCC</sub> to SU(5)<sub>TCC</sub> rather than $`\mathrm{SU}(3)_c\times \mathrm{SU}(3)_{TC}`$. Channel (139) breaks SU(2)<sub>w</sub> at too high a scale and breaks SU(6)<sub>TCC</sub> to SU(5)<sub>TCC</sub>.
It may be noted that even if there were some way to produce a breaking of SU(6)<sub>TCC</sub> that yielded a lower-energy theory invariant under $`\mathrm{SU}(3)_{TC}\times \mathrm{SU}(3)_c`$, the interchange symmetry between of SU(3)<sub>TC</sub> and SU(3)<sub>c</sub> would imply that the respective technicolor and color gauge couplings would evolve in the same way as the energy decreases below the scale at which this condensate occurred. Both of these sectors are asymptotically free, and condensates would form, but the model would still not be realistic, since the scale $`\mathrm{\Lambda }_{TC}`$ would be the same as $`\mathrm{\Lambda }_{QCD}`$.
### V.4 $`N_{SC}=4`$, $`G=\mathrm{SU}(9)`$
Here we consider $`N_{SC}=4`$, so that $`G=\mathrm{SU}(9)`$ and $`𝐧=(n_1,\mathrm{},n_8)_9`$. A priori, one has three possibilities regarding the origins of the SM fermion generations: (i) $`N_{gh}=3`$, or equivalently, $`N_g\mathrm{}=0`$, whence $`N_{TC}=4`$; (ii) $`N_{gh}=2`$ so that $`N_g\mathrm{}=1`$, which would be associated with a breaking of $`\mathrm{SU}(4)_{SC}`$ to $`\mathrm{SU}(3)_{TC}`$; (iii) $`N_{gh}=1`$, or equivalently, $`N_g\mathrm{}=2`$, so that one SM generation would arise initially from the representations of $`G`$, and the other two would arise via the sequential breaking $`\mathrm{SU}(4)_{SC}\mathrm{SU}(3)_{SC}`$ and then $`\mathrm{SU}(3)_{SC}\mathrm{SU}(2)_{TC}`$. In general, equations (84) and (85) yield
$$n_1=n_3=n_5=n_7=0$$
(142)
and eq. (82) is
$$N_{gh}=n_2+n_6=n_4+n_8.$$
(143)
The condition of zero gauge anomaly (77) is
$$5(n_2+n_4)9n_6n_8=0.$$
(144)
For a given value of $`N_{gh}=3N_g\mathrm{}`$, these are three nondegenerate linear equations for the four quantities $`n_2`$, $`n_4`$, $`n_6`$, and $`n_8`$. Taking $`n_2`$, say, as the independent variable, initially free to take on values $`n_2=0,1,\mathrm{},N_{gh}`$, we find the formal solution of these equations to be
$`n_6`$ $`=`$ $`N_{gh}n_2`$ (145)
$`n_4`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(5N_{gh}7n_2\right)`$ (147)
$`n_8`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(7n_22N_{gh}\right).`$ (149)
Consider first the a priori possible value $`N_{gh}=3`$, whence $`n_4=5(7/3)n_2`$ and $`n_8=(7/3)n_22`$. In order for $`n_2`$ and $`n_8`$ to be nonnegative integers, $`n_2=0`$ mod 3. Now $`n_2`$ cannot be zero, because this would make $`n_8`$ negative. But $`n_2`$ also cannot be equal to 3, because this would make $`n_4`$ negative. Hence, the value $`N_{gh}=3`$ is not allowed.
Consider next the value
$$N_{gh}=2$$
(150)
whence $`N_g\mathrm{}=1`$ and
$$G_{SC}=\mathrm{SU}(4)_{SC},G_{TC}=\mathrm{SU}(3)_{TC}.$$
(151)
Equations (149) yield $`n_4=(1/3)(107n_2)`$ and $`n_8=(1/3)(7n_24)`$. From its formal range for this case, $`n_2=0,1,2`$, the only allowed value is $`n_2=1`$ which gives $`n_4=n_8=n_6=1`$, so that
$$𝐧=(0,1,0,1,0,1,0,1)_9$$
(152)
i.e., the fermion content is given by $`\{f\}=[2]_9+[4]_9+[6]_9+[8]_9`$. For this case $`N_{(1,1)}=1`$.
Finally, we examine the minimal possible value, $`N_{gh}=1`$, corresponding to the maximal possible value of $`N_g\mathrm{}`$, namely $`N_g\mathrm{}=N_{SC}(N_{TC})_{min}=42=2`$. Here eqs. (149) read $`n_4=(1/3)(57n_2)`$ and $`n_8=(1/3)(7n_22)`$, where, a priori, $`n_2`$ can take values in the set $`\{0,1\}`$ . Evidently, neither of these values would make $`n_4`$ and $`n_8`$ nonnegative integers and hence neither is allowed.
Let us return to the case with $`N_{gh}=2`$. An initial comment is that the corresponding value $`N_{TC}=3`$ is disfavored, relative to $`N_{TC}=2`$, since it leads to larger technicolor contributions to precision electroweak quantities and could reduce the likelihood of walking behavior for the technicolor theory. Notwithstanding this concern, let us investigate this case. The decomposition of the fermion representations with respect to the subgroup (2), which for this case reads
$$\mathrm{SU}(4)_{SC}\times \mathrm{SU}(5)_{GU},$$
(153)
is listed below:
$$[2]_9=(6,1)+(4,5)+(1,10)$$
(154)
$$[4]_9=(1,1)+(\overline{4},5)+(6,10)+(4,\overline{10})+(1,\overline{5})$$
(155)
$$[6]_9[\overline{3}]_9=(4,1)+(\overline{6},\overline{5})+(\overline{4},\overline{10})+(1,10)$$
(156)
$$[8]_9[\overline{1}]_9=(\overline{4},1)+(1,\overline{5}).$$
(157)
(Note that $`6\overline{6}`$, i.e., $`[2]_4[\overline{2}]_4`$, in SU(4)<sub>SC</sub>.) This decomposition shows that, neglecting the $`G_{GU}`$ couplings relative to those of SU(4)<sub>SC</sub>, the latter interaction is vectorial, involving nonsinglet fermions comprising a set of 16 copies of $`\{4+\overline{4}+6\}`$ left-handed fermions.
As before, we consider the implications of a scenario in which the unified group, here SU(9), breaks to yield a theory at lower energy scales that is invariant under the gauge symmetry (ignoring an abelian factor) of the form (65), which for the present model is
$$\mathrm{SU}(2)_w\times \mathrm{SU}(7)_{SCC}.$$
(158)
With respect to this direct product symmetry group, the fermions have the decomposition
$$[2]_9=(1,1)+(2,[1]_7)+(1,[2]_7)$$
(159)
$$[4]_9=(1,[2]_7)+(2,[3]_7)+(1,[\overline{3}]_7)$$
(160)
$$[6]_9[\overline{3}]_9=(1,[\overline{1}]_7)+(2,[\overline{2}]_7)+(1,[\overline{3}]_7)$$
(161)
$$[8]_9[\overline{1}]_9=(2,1)+(1,[\overline{1}]_7).$$
(162)
Here, the dimensionalities include $`\mathrm{dim}([2]_7)=21`$ and $`\mathrm{dim}([3]_7)=35`$. Thus, the SU(7)<sub>SCC</sub> theory is vectorial, with nonsinglet fermion content (neglecting other interactions) given by
$$2\{[1]_7+[\overline{1}]_7+[2]_7+[\overline{2}]_7+[3]_7+[\overline{3}]_7\}.$$
(163)
The SU(7)<sub>SCC</sub> gauge interaction is asymptotically free, with leading beta function coefficient $`b_0^{(SCC)}=13/3`$. However, as in the models that we studied above, because of the vectorlike nature of the SCC gauge symmetry, when the energy decreases sufficiently so that $`\alpha _{SCC}`$ grows large enough to produce fermion condensates, these will preferentially be in the channels $`\times \overline{}1`$, where $``$ refers to the SU(7)<sub>SCC</sub> representation. Thus, these preserve the SU(7)<sub>SCC</sub> invariance rather than breaking it down to a direct product which includes $`\mathrm{SU}(3)_c\times \mathrm{SU}(4)_{SC}`$, as is necessary to separate color from the strongly coupled SU(4)<sub>SC</sub> group. With respect to the $`\mathrm{SU}(2)_w\times \mathrm{SU}(7)_{SCC}`$ subgroup, the most attractive channel, with its measure of attractiveness $`\mathrm{\Delta }C_2`$, is
$$(2,[3]_7)\times (1,[\overline{3}]_7)(2,1),\mathrm{\Delta }C_2=\frac{96}{7}13.7.$$
(164)
In addition to its failure to break the SU(7)<sub>SCC</sub> symmetry, this MAC breaks SU(2)<sub>w</sub> at a scale that would be higher than the ETC scales where the SU(7)<sub>SCC</sub> should break. The same problems characterize a channel with a somewhat smaller value of $`\mathrm{\Delta }C_2`$, namely,
$$(1,[2]_7)\times (2,[\overline{2}]_7)(2,1),\mathrm{\Delta }C_2=\frac{80}{7}11.4.$$
(165)
As before, one can examine other possible condensation channels with still smaller values of $`\mathrm{\Delta }C_2`$, which include
$$(1,[2]_7)\times (1,[\overline{3}]_7)(1,[\overline{1}]_7),\mathrm{\Delta }C_2=\frac{74}{7}10.6$$
(166)
$$[(2,[2]_7)\times (2,[\overline{3}]_7)]_a(1,[\overline{1}]_7),\mathrm{\Delta }C_2=\frac{74}{7}$$
(167)
and
$$[(2,[2]_7)\times (2,[\overline{3}]_7)]_s(3,[\overline{1}]_7),\mathrm{\Delta }C_2=\frac{74}{7}.$$
(168)
Channel (168) is unacceptable because it breaks SU(2)<sub>w</sub> in the wrong way and at too high a scale. Channels (166) and (167) are not forbidden but would only break SU(7)<sub>SCC</sub> to SU(6)<sub>SCC</sub>, thereby necessitating a further breaking to $`\mathrm{SU}(3)_c\times \mathrm{SU}(3)_{TC}`$; moreover, because of their subdominant $`\mathrm{\Delta }C_2`$ values, it is difficult to make a convincing argument that they would predominate. Another conceivable breaking pattern is $`\mathrm{SU}(7)_{SCC}\mathrm{SU}(3)_c\times \mathrm{Sp}(4)`$ fs , but it is not clear what dynamical fermion condensation channel could produce this breaking.
### V.5 $`N_{SC}=5`$, $`G=\mathrm{SU}(10)`$
We proceed to examine the case where $`N_{SC}=5`$, so that $`G=\mathrm{SU}(10)`$ and $`𝐧=(n_1,\mathrm{},n_9)_{10}`$. A priori, one has four possibilities for the manner in which the SM fermion generation arise, as specified by $`(N_{gh},N_g\mathrm{},N_{TC})`$, namely (i) (3,0,5), (ii) (2,1,4), (iii) (1,2,3), and (iv) (0,3,2). The minimization of technicolor contributions to electroweak corrections favor the last of these options. The conditions (84) and (85) forbidding $`5_L`$ and $`\overline{10}_L`$ yield
$$n_1=n_3=n_6=n_8=0$$
(169)
and eq. (82) is
$$N_{gh}=n_2+n_7=n_4+n_9.$$
(170)
The condition of zero gauge anomaly (77) is
$$6n_2+14(n_4n_7)n_9=0.$$
(171)
For a given value of $`N_{gh}=3N_g\mathrm{}`$, these are three nondegenerate linear equations for the five quantities $`n_2`$, $`n_4`$, $`n_5`$, $`n_7`$, and $`n_9`$. Taking $`n_5`$ and $`n_7`$, say, as the two independent variables, with $`n_7`$ constrained by eq. (170) to take on values in the set $`\{0,1,\mathrm{},N_{gh}\}`$, we find the formal solution of these equations to be
$$n_4=\frac{1}{3}\left(4n_7N_{gh}\right)$$
(172)
$$n_9=\frac{4}{3}\left(N_{gh}n_7\right).$$
(173)
From eq. (173) and the requirement that $`n_9`$ be a nonnegative integer, it follows that $`N_{gh}n_7=0`$ mod 3. This could be satisfied for $`N_{gh}=3`$ and $`n_7=0`$, but this choice is excluded because it would produce a negative value for $`n_4`$. The other choice is $`N_{gh}=n_7`$, which gives $`n_9=0`$ and $`n_4=n_7`$. Substituting $`n_9=0`$ in eq. (170) yields $`n_4=N_{gh}`$ so that $`n_7=N_{gh}`$ also; substituting the latter in eq. (170) then gives $`n_2=0`$. These conditions leave $`n_5`$ free, subject to the additional requirement that SU(5)<sub>SCC</sub> be asymptotically free. Thus, we have
$$𝐧=(0,0,0,N_{gh},n_5,0,N_{gh},0,0)_{10}.$$
(174)
The number of SC-singlet, SU(5)<sub>GU</sub>-singlet fermions is $`N_{(1,1)}=2n_5`$.
Let us consider first the value $`n_5=0`$, so that
$$𝐧=(0,0,0,N_{gh},0,0,N_{gh},0,0)_{10}.$$
(175)
This allows one to choose $`(N_{gh},N_g\mathrm{})=(3,0)`$, (2,1), or (1,2). If one were to apply the irreducibility condition that $`GCD(\{n_k\})=1`$, it would imply that $`N_{gh}=1`$ in eq. (175), which requires $`N_g\mathrm{}=2`$, i.e., the TC group is SU(3)<sub>TC</sub> and the SC group should undergo two sequential breakings, $`\mathrm{SU}(5)_{SC}\mathrm{SU}(4)_{SC}`$, followed by $`\mathrm{SU}(4)_{SC}\mathrm{SU}(2)_{TC}`$. Although we will not impose the irreducibility here, we will exclude all of the reducible solutions because they lead to excessively many fermions, which render the SU(8)<sub>SCC</sub> theory non asymptotically free.
The fermions in each type of representation have the decomposition, with respect to the subgroup (2) for this case,
$$\mathrm{SU}(5)_{SC}\times \mathrm{SU}(5)_{GU},$$
(176)
of
$$[4]_{10}=(\overline{5},1)+(\overline{10},5)+(10,10)+(5,\overline{10})+(1,\overline{5})$$
(177)
$$[7]_{10}=[\overline{3}]_{10}=(10,1)+(\overline{10},\overline{5})+(\overline{5},\overline{10})+(1,10).$$
(178)
Thus, the SU(5)<sub>SC</sub> sector forms a chiral gauge theory, with left-handed (nonsinglet) fermion content consisting of
$$\{f\}=N_{gh}[10(5+\overline{10})+11(\overline{5}+10)].$$
(179)
With respect to the subgroup (65),
$$\mathrm{SU}(2)_w\times \mathrm{SU}(8)_{SCC},$$
(180)
the fermions in each representation have the decomposition
$$[4]_{10}=(1,[2]_8)+(2,[3]_8)+(1,[4]_8)$$
(181)
$$[7]_{10}[\overline{3}]_{10}=(1,[\overline{1}]_8)+(2,[\overline{2}]_8)+(1,[\overline{3}]_8).$$
(182)
Some dimensionalities of relevant SU(8) representations are $`\mathrm{dim}([2]_8)=28`$, $`\mathrm{dim}([3]_8)=56`$, and $`\mathrm{dim}([4]_8)=70`$; note that $`[4]_8=[\overline{4}]_8`$. Thus, assuming that the SU(10) unified theory breaks in such a way as to yield, for a range of lower energies, an effective theory with SU(8)<sub>SCC</sub> gauge symmetry, this SCC sector is a chiral gauge theory. The SU(8)<sub>SCC</sub> gauge interaction has leading beta function coefficient $`b_0^{(SCC)}=4(2221N_{gh})/3`$, so it is asymptotically free only for the choice $`N_{gh}=1`$, for which $`b_0^{(SCC)}=4/3`$. As the energy scale decreases and the coupling $`\alpha _{SCC}`$ becomes sufficiently large, this theory will thus form bilinear fermion condensates. The most attractive channel is
$$(1,[4]_8)\times (1,[4]_8)(1,1),\mathrm{\Delta }C_2=18.$$
(183)
This is of the form of channel (90) with $`p=4`$ and thus violates weak hypercharge and electric charge.
Other channels with smaller values of $`\mathrm{\Delta }C_2`$ include
$$(2,[3]_8)\times (1,[\overline{3}]_8)(2,1),\mathrm{\Delta }C_2=\frac{135}{8}16.9$$
(184)
$$(1,[2]_8)\times (2,[\overline{2}]_8)(2,1),\mathrm{\Delta }C_2=\frac{63}{4}=15.75$$
(185)
$$(1,[\overline{3}]_8)\times (1,[4]_8)(1,[1]_8),\mathrm{\Delta }C_2=\frac{27}{2}=13.5$$
(186)
$$(1,[2]_8)\times (1,[\overline{3}]_8)(1,[\overline{1}]_8),\mathrm{\Delta }C_2=\frac{45}{4}=11.25.$$
(187)
Channels (184) and (185) break SU(2)<sub>w</sub> at too high a scale and fail to break SU(8)<sub>SCC</sub>. Channels (186) and (187) are allowed by symmetry considerations. However, since their $`\mathrm{\Delta }C_2`$ values are smaller than those of the leading channels, one cannot make a persuasive case that they would occur. Moreover, because of the relatively small value of $`b^{(SCC)}`$, estimates based on eq. (210) indicate that for a hypothetical unified coupling $`\alpha _{GU}0.04`$ at $`M_{GU}10^{16}`$ GeV, these condensates would form at much too small a scale for a viable model.
We consider next the choice $`N_{gh}=1`$, $`n_5=1`$, so that
$$𝐧=(0,0,0,1,1,0,1,0,0)_{10}.$$
(188)
For this choice the decomposition of the fermions with respect to the subgroup (176) is
$$[4]_{10}=(\overline{5},1)+(\overline{10},5)+(10,10)+(5,\overline{10})+(1,\overline{5})$$
(189)
$$[5]_{10}=2(1,1)+(\overline{5},5)+(5,\overline{5})+(\overline{10},10)+(10,\overline{10})$$
(190)
$$[7]_{10}[\overline{3}]_{10}=(10,1)+(\overline{10},\overline{5})+(\overline{5},\overline{10})+(1,10).$$
(191)
The SU(5)<sub>SC</sub> theory is a chiral gauge theory, consisting of the nonsinglet fermion content
$$\{f\}=16(\overline{5})+15(5)+21(10)+20(\overline{10}).$$
(192)
With respect to the subgroup (180) the fermions in the $`[4]_{10}`$ and $`[7]_{10}`$ have the decompositions given in eqs. (181) and (182), and
$$[5]_{10}=(1,[3]_8)+(2,[4]_8)+(1,[\overline{3}]_8).$$
(193)
Although the $`[5]_{10}`$ is self-conjugate, the $`[4]_{10}`$ and $`[7]_{10}`$ make this SU(8)<sub>SCC</sub> theory chiral. However, it is non-asymptotically free, with leading beta function coefficient $`b_0^{(SCC)}=22`$.
Finally, we consider the choice $`N_{gh}=0`$, $`n_5=1`$, so that
$$𝐧=(0,0,0,0,1,0,0,0,0)_{10}.$$
(194)
For this choice the decomposition of the fermions in the $`[5]_{10}`$ with respect to the subgroup (176) is given by eq. (190) and with respect to the subgroup (180) by eq. (193). Assuming that the breaking of SU(10) is such as to yield an SU(8)<sub>SCC</sub>-invariant theory for a range of lower energies, its coupling does grow as the energy decreases, as governed by the leading beta function coefficient $`b_0^{(SCC)}=6`$. However, as is evident from eq. (193), this theory is vectorial, so that when the SU(8)<sub>SCC</sub> coupling grows sufficiently large to produce a fermion condensate, this condensate will preferentially preserve the SU(8)<sub>SCC</sub> symmetry rather than breaking it, as is necessary, to $`\mathrm{SU}(5)_{TC}\times \mathrm{SU}(3)_c`$ (and thence, sequentially, breaking SU(5)<sub>TC</sub> to SU(2)<sub>TC</sub>). We have investigated higher values of $`N_{SC}`$ and thus $`N`$ but have found that they exhibit problems similar to those of the models above.
### V.6 Assessment
Thus we we find several general problems with the unification approach embodied in models of type 3. One, pertaining to the mechanism for breaking the unified gauge group $`G`$ in a weak-coupling framework, also applies to models of type (1) and (2) and will be discussed in the next section. A second problem is that, even if one could arrange some mechanism to break the $`G`$ symmetry in the desired manner to yield a lower-energy theory invariant, presumably, under an SCC gauge symmetry combining the strongly coupled group $`G_{SC}`$ with the color group, it appears very difficult to get this SCC symmetry to break in the requisite manner. This is especially true when, as is often the case, the SU($`N_{SCC}`$) gauge interaction is vectorial. Even when it is chiral (and asymptotically free), as in the model with $`N_{SC}=5`$ and $`𝐧`$ given by eq. (175) with $`N_{gh}=1`$, the most attractive condensation channels do not lead to the requisite breaking.
Since the determination of the resultant SC and TC sectors depends on having a viable SCC breaking pattern, this prevents one from proceeding very far with the analysis of these lower-energy effective field theories in the context of these models. However, we note that in cases where $`N_g\mathrm{}1`$, it could also be challenging to get the SC symmetry to break sequentially down to the resultant exact TC symmetry. Because of the previous problems, one cannot obtain very definite predictions for fermion masses. A general concern pertains to reproducing the observed mass hierarchy of the three SM fermion generations. In order to do this, one tends to need three different ETC mass scales, essentially the $`\mathrm{\Lambda }_j`$, $`j=1,2,3`$ discussed in Section II. In the present approach, one has $`N_g\mathrm{}+1`$ ETC-type mass scales. This is illustrated by the scenario in which the SU($`N_{SCC}`$) symmetry first breaks to $`\mathrm{SU}(N_{SC})\times \mathrm{SU}(3)_c`$ and then $`\mathrm{SU}(N_{SC})`$ breaks sequentially at $`N_g\mathrm{}`$ lower scales, finally yielding the residual exact SU($`N_{TC}`$) symmetry. Consequently, unless $`N_g\mathrm{}2`$, one does not have enough mass scales to account for the SM fermion mass hierarchy. This problem reaches its most acute form when $`N_g\mathrm{}=0`$, $`N_{gh}=3`$.
Because of the presence of intermediate scales between $`m_Z`$ and $`M_{GU}`$ with nonperturbative behavior and the feature that the ETC gauge bosons involved in symmetry breakings at these scales carry SM quantum numbers, the calculation of the evolution of the SM gauge couplings is more complicated in models of type (3) than in models of types (1) or (2). However, exploratory analysis of plausible evolution of gauge couplings for the various scenarios that we have examined indicate that satisfying the constraint of gauge coupling unification is still a very restrictive requirement.
## VI Dynamical Breaking of Unified Gauge Symmetries
In addition to other issues that we have addressed concerning prospects for unification of gauge symmetries in a dynamical context, there is another one which is quite general. For the present discussion let us assume that one has a model that does achieve gauge coupling unification. The resultant value of the unified gauge coupling at $`M_{GU}`$ is generically expected to be small. But if one is trying to construct a theory in which all gauge symmetry breaking is dynamical, this would normally require there to be some strongly coupled gauge interaction at the relevant scale. For example, the dynamical breaking of the electroweak symmetry in a technicolor theory requires that technicolor be an asymptotically free gauge interaction that becomes strongly coupled at the electroweak scale. Although the ETC interaction is strongly coupled at the scale $`\mathrm{\Lambda }_1=\mathrm{\Lambda }_{ETC,max}10^6`$ GeV, in typical models of type (1) and (2) the ETC coupling evolves to relatively small values as the energy scale ascends to the region of $`M_{GU}`$, so one could not use ETC interactions to break $`G_{GU}`$ at $`M_{GU}`$ in these models. Moreover, in specific models such as those of Refs. at94 -kt , the ETC (and HC) condensates involve SM-singlet fields which, in the present context, would naturally be $`G_{GU}`$-singlet fields, so their condensates would not break $`G_{GU}`$.
One way to break a hypothetical symmetry $`G_{GU}`$ unifying SM gauge interactions at a high scale $`M_{GU}`$ would be to expand the theory to include an additional gauge interaction, associated with a group $`G_a`$, that is strongly coupled at this scale, together with fermions that transform as nonsinglets under both $`G_{GU}`$ and $`G_a`$. The great disparity between the coupling strengths of the $`G_{GU}`$ and $`G_{ETC}`$ gauge bosons, on the one hand, and the assumed $`G_a`$ gauge bosons, on the other, is a striking property of such a model. Another approach would be to envision a nonperturbative unification of gauge symmetries, but by its very nature, this is difficult to study reliably using tools such as perturbative evolution of SM couplings nonpert .
For completeness, one should note that a major objection to the use of Higgs fields for symmetry breaking at lower scales is the instability of a Higgs sector to large radiative corrections, necessitating fine tuning to keep the Higgs bosons light compared with an ultraviolet cutoff. But from considerations of nucleon stability alone, not to mention gauge coupling evolution, one knows that in the (four-dimensional) theories considered here, $`M_{GU}`$ is generically quite high, not too far below the Planck scale where the theories are certainly incomplete, since they do not include quantum gravity. Hence, this objection would not be as strong for the breaking of $`G_{GU}`$ as for symmetry breakings that occur at substantially lower scales.
## VII Conclusions
In this paper we have analyzed approaches to the partial or complete unification of gauge symmetries in theories with dynamical symmetry breaking. We considered three main types of models with progressively greater degrees of unification, including those that (1) involve sufficient unification to quantize electric charge, (2) attempt to unify the three standard-model gauge interactions in a simple group that forms a direct product with an extended technicolor group, and, (3) attempt to unify the standard-model gauge interactions with (extended) technicolor in a simple group. The model of Refs. lrs ; nag is a successful example of theories of type (1). Models of type (3) provide an interesting contrast to those of types (1) and (2) in the different way in which standard-fermion generations are produced and in the property of having ETC gauge bosons that carry standard-model quantum numbers. We have pointed out a number of challenges that one faces in trying to construct viable models of types (2) and (3). There are certainly further avenues for research in this area. For example, one could investigate direct product groups involving auxiliary hypercolor-type groups. Another idea would be to study ways to unify top-color tc2 with standard-model gauge symmetries. In conclusion, it is possible that electroweak symmetry breaking is dynamical, involving a new strongly coupled gauge symmetry, technicolor, embedded in an extended technicolor theory to give masses to standard-model fermions. There is a strong motivation to understand how the associated symmetries can be unified with the color and weak isospin and hypercharge gauge symmetries. We hope that the results of the present paper will be of use in the further study of this unification program.
We thank T. Appelquist for helpful discussions, in particular for the collaborative work on the partial unification model in Refs. lrs ; nag . This present research was partially supported by the grant NSF-PHY-00-98527.
## VIII Appendix
We gather in this appendix some standard formulas that are used in the calculations reported in the text.
### VIII.1 Some Group-Theoretic Properties of Representations $`[k]_N`$ of SU($`N`$)
We denote the completely antisymmetric $`k`$-fold products of the fundamental and conjugate fundamental representation of SU($`N`$) as $`[k]_N`$ and $`[\overline{k}]_N`$, respectively. These can be displayed as tensors with $`k`$ upper indices, $`\psi ^{i_1\mathrm{}i_k}`$ and $`k`$ lower indices, $`\psi _{i_1\mathrm{}i_k}`$, and have the (same) dimension
$$\mathrm{dim}([k]_N)=\left(\genfrac{}{}{0pt}{}{N}{k}\right).$$
(195)
These representations satisfy the equivalence property
$$[Nk]_N[\overline{k}]_N$$
(196)
under SU($`N`$), as follows by contraction with the totally antisymmetric tensor density $`ϵ_{i_1\mathrm{}i_N}`$. The fact that these representations have the same dimension is evident from the identity $`\left(\genfrac{}{}{0pt}{}{N}{k}\right)=\left(\genfrac{}{}{0pt}{}{N}{Nk}\right)`$. Further, if $`N`$ is even, say $`N=2k`$, then $`[N/2]_N`$ is self-conjugate with respect to SU($`N`$).
The contribution of the (left-handed) fermions in the representation $`[k]_N`$ to the gauge anomaly for SU($`N`$) is g79 ; anomaly
$$A([k]_N)=\frac{(N2k)(N3)!}{(Nk1)!(k1)!},$$
(197)
where the normalization is such that contribution of the fundamental representation is 1. The property $`A([k]_N)=A([\overline{k}]_N)=A([Nk]_N)`$ for $`1kN1`$ is evident in eq. (197).
The quadratic Casimir invariant $`C_2()`$ for the representation $``$ is defined by
$$\underset{a=1}{\overset{o(G)}{}}\underset{j=1}{\overset{\mathrm{dim}()}{}}𝒟_{}(T_a)_j^i𝒟_{}(T_a)_k^j=C_2()\delta _k^i,$$
(198)
where $`𝒟_{}(T_a)`$ is the $``$-representation of the generator $`T_a`$ and $`o(G)`$ is the order of the group $`G`$.
The contribution of a fermion loop, for fermions of representation $``$ of SU($`N`$), to the beta function coefficient $`b_0^{(j)}`$, involves the invariant $`T()`$ defined by
$$\underset{i,j=1}{\overset{\mathrm{dim}()}{}}(𝒟_{}(T_a))_j^i(𝒟_{}(T_b))_i^j=T()\delta _{ab}.$$
(199)
These invariants satisfy the elementary relation
$$C_2()\mathrm{dim}()=T()o(G).$$
(200)
For SU($`N`$),
$$C_2([k]_N)=\frac{k(N+1)(Nk)}{2N}$$
(201)
and
$$T([k]_N)=\frac{1}{2}\left(\genfrac{}{}{0pt}{}{N2}{k1}\right).$$
(202)
As the value of $`N`$ (in eq. (51)) increases, the number of (left-handed) fermions in a generic set $`𝐧`$ tends to increase rapidly. For example, for even $`N_{SC}`$ and hence odd $`N=N_{SC}+52m+1`$, the set $`𝐧`$ with $`n_2\mathrm{}=1`$ for $`\mathrm{}=1,..m`$ and $`n_{2\mathrm{}+1}=0`$ for $`\mathrm{}=0,..,m1`$, has a total number of fermions given by
$$\underset{\mathrm{}=1}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{2m+1}{2\mathrm{}}\right)=2^{2m}1=2^{N1}1.$$
(203)
Evidently, this grows exponentially rapidly with $`N_{SC}`$ and $`N`$, which quickly renders the SCC theory non-asymptotically free. Note that for this choice of $`𝐧`$, since $`n_2=n_{_{N_{SC}+2}}`$, it follows that $`N_{gh}=2`$.
### VIII.2 Formulas for the Evolution of Gauge Couplings
We consider a factor group $`G_j`$ with gauge coupling $`g_j`$ and denote $`\alpha _j=g_j^2/(4\pi )`$. The evolution of the gauge couplings as a function of the momentum scale $`\mu `$ is given by the renormalization group equation
$$\beta _j=\frac{d\alpha _j}{dt}=\frac{\alpha _j^2}{2\pi }\left(b_0^{(j)}+\frac{b_1^{(j)}}{2\pi }\alpha _j+O(\alpha _j^2)\right),$$
(204)
where $`t=\mathrm{ln}\mu `$, and the first two terms $`b_0^{(j)}`$ and $`b_1^{(j)}`$ are scheme-independent. The beta function with perturbatively calculated coefficients is appropriate to describe the running of the respective couplings in the energy ranges where the respective gauge fields are dynamical (i.e., above corresponding scales at which $`G_j`$ is broken) and where the couplings $`\alpha _j`$ are not too large. For our analyses of the perturbative evolution of gauge couplings, it will be sufficient to keep only the $`b_0^{(j)}`$ term; the well-known solution of eq. (204) is then given by
$$\alpha _j^1(t_2)=\alpha _j^1(t_1)+\frac{b_0^{(j)}}{2\pi }(t_2t_1).$$
(205)
If an effective field theory involves the direct product of two gauge groups $`G_j`$ and $`G_k`$ for energy scales between $`\mu _{\mathrm{}}`$ and a larger scale $`\mu _{jk}`$ where the associated couplings $`\alpha _j`$ and $`\alpha _k`$ are equal, one has, to this order,
$$\mathrm{ln}\left(\frac{\mu _{jk}}{\mu _{\mathrm{}}}\right)=\frac{2\pi [\alpha _j^1(\mu _{\mathrm{}})\alpha _k^1(\mu _{\mathrm{}})]}{b_0^{(k)}b_0^{(j)}}.$$
(206)
The three SM gauge couplings are accurately determined at $`\mu =m_Z`$, with the results pdg ; ew $`\alpha _3(m_Z)0.118`$, $`\alpha _{em}(m_Z)^1128`$, and $`(\mathrm{sin}^2\theta _W)_{\overline{MS}}(m_Z)0.2312`$. The SU(2)<sub>L</sub> and U(1)<sub>Y</sub> couplings $`gg_{2L}`$ and $`g^{}g_Y`$ are given by $`e=g\mathrm{sin}\theta _W=g^{}\mathrm{cos}\theta _W`$ and have the values (quoted to sufficient accuracy for our present purposes) $`\alpha _2(m_Z)=0.033`$ and $`\alpha _Y(m_Z)=0.010`$. The evolution of these couplings to scales $`\mu >m_Z`$ depends on the type of gauge symmetry unification that one is considering.
### VIII.3 Fermion Condensation
Consider massless fermions that transform according to some representations $``$ of a nonabelian gauge group $`G`$. In the approximation of a single-gauge boson exchange , the critical coupling for condensation in the channel
$$_1\times _2R_{cond.}$$
(207)
is given by gap
$$\alpha _c=\frac{2\pi }{3\mathrm{\Delta }C_2}$$
(208)
where
$$\mathrm{\Delta }C_2=C_2(_1)+C_2(_2)C_2(_{cond.}).$$
(209)
Because $`\alpha O(1)`$ where fermion condensation occurs, the one-gauge boson approximation is only a rough guide to the actual critical value of $`\alpha `$. Corrections to this have been estimated in Ref. gap . In addition to gauge boson exchange diagrams, nonperturbative processes involving instantons are also important instanton .
For condensation due to the asymptotically free gauge interaction with gauge group $`G_j`$ and associated squared coupling $`\alpha _j`$ obeying a renormalization group equation with leading beta function coefficient $`b_0^{(j)}`$, the solution in eq. (205) together with the condition (208), yield, for the mass scale at which the condensation takes place, the rough estimate
$$\mu _{c,j}M_{GU}\mathrm{exp}\left[\frac{2\pi }{b_0^{(j)}}\left(\alpha _j(M_{GU})^1\frac{3\mathrm{\Delta }C_2}{2\pi }\right)\right]$$
(210)
where $`\mathrm{\Delta }C_2`$ is the value appropriate for this channel, as given by eq. (209).
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# Derivation of Newton’s law of gravitation based on a fluidic continuum model of vacuum and a sink flow model of particles
## I Introduction
The Newton’s law of gravitation can be written as
$$\text{F}_{21}=G\frac{m_1m_2}{r^2}\widehat{\text{r}}_{21},$$
(1)
where $`m_1`$ and $`m_2`$ are the masses of two particles, $`r`$ is the distance between the two particles, $`G`$ is the gravitational constant, $`\text{F}_{12}`$ is the force exerted on the particle with mass $`m_2`$ by the particle with mass $`m_1`$, $`\widehat{\text{r}}_{21}`$ denotes the unit vector directed outward along the line from the particle with mass $`m_1`$ to the particle with mass $`m_2`$.
The main purpose of this paper is to derive the Newton’s law of gravitation by means of fluid mechanics based on sink flow model of particles.
The motive of this paper is to seek a mechanism of gravitational phenomena. The reasons why new models of gravity are interesting may be summarized as follows.
Firstly, there exists some astronomical phenomena that could not be interpreted by the present theories of gravitation, for instance, the Titius-Bode law Nietro (1972). New theories of gravity may view these problems from new angles.
Secondly, whether the gravitational constant $`G`$ depends on time and space is still unknown Gilvarry and Muller (1972); Tomaschitz (2000); Gaztanaga et al. (2001); Copi et al. (2004); Khoury and Weltman (2004); Nagata et al. (2004); Biesiada and Malec (2004). It is known that the gravitational constant $`G`$ is a constant in the Newton’s theory of gravitation and in theory of general relativity.
Thirdly, the mechanism of the action-at-a-distance gravitation remains an unsolved problem in physics for more than 300 years Whittaker (1953); Hesse (1961); Jaakkola (1996). Although theory of general relativity is a field theory of gravityFock (1959), the concept of field is different from that of continuum mechanics Truesdell (1966); Fung (1977); Eringen (1980); Landau and Lifshitz (1987) because of the absence of a continuum in theory of general relativity. Thus, theory of general relativity can only be regarded as a phenomenological theory of gravity.
Fourthly, we do not have a satisfactory quantum theory of gravity presently Carlip (2001); Amelino-Camelia (2002); Ahluwalia (2002); Mitrofanov (2003); Christian (2005). One of the challenges in theoretical physics is to reconcile quantum theory and theory of general relativity Carlip (2001); Collins et al. (2004). New theories of gravity may open new ways to solve this problem.
Fifthly, one of the puzzles in physics is the problem of dark matter and dark energy Ellis (2003); Linder (2004); Bernardeau (2003); Bergstrom et al. (2005); Beacom et al. (2005); Akerib et al. (2006); Feng et al. (2006); P. Fayet (2006); Carena et al. (2006). New theories of gravity may provide new methods to attack this problemLinder (2004); Bernardeau (2003).
Finally, we do not have a successful unified field theory presently. Great progress has been made towards an unification of the four fundamental interactions in the universe in the 20th century. However, gravitation is still not unified successfully. New theories of gravity may shed some light on this puzzle.
To conclude, it seems that new considerations on gravitation is needed. It is worthy keeping an open mind with respect to all the theories of gravity before the above problems been solved.
Now let us briefly review the long history of mechanical interpretations of gravitational phenomena. Many philosophers and scientists, such as LaoziLaozi (1995), Thales, Anaximenes, believed that everything in the universe is made of a kind of fundamental substanceWhittaker (1953). Descartes was the first to bring the concept of aether into science by suggesting that it has mechanical propertiesWhittaker (1953). Since the Newton’s law of gravitation was published in 1687Newton (1962), this action-at-a-distance theory was criticized by the French CartesianWhittaker (1953). Newton admitted that his law did not touch on the mechanism of gravitationCohen (1980). He tried to obtain a derivation of his law based on Descartes’ scientific research program Newton (1962). Newton himself even suggested an explanation of gravity based on the action of an aetherial medium pervading the spaceNewton (1931); Cohen (1980). Euler attempted to explain gravity based on some hypotheses of a fluidic aetherWhittaker (1953).
In a remarkable paper published in 1905, Einstein abandoned the concept of aetherEinstein (1905). However, Einstein’s assertion did not cease the explorations of aether Whittaker (1953); Vigier (1980); Barut (1988); Oldershaw (1989a, b); Carvalho and et al. (2003); Arminjon (2002); Carvalho and Oliveira (2003); Jacobson (2004); Davies (2005); Levin and Wen (2006). Einstein changed his view later and introduced his new concept of etherEinstein (1983); Kostro (2000). I regret to admit that it is impossible for me to mention all the works related to this field in history. Adolphe Martin and Roy KeysMartin (2005a, b); Martin and Keys (1994) proposed a gas model of vacuum to explain the physical phenomena such as electromagnetism, gravitation, quantum mechanics and the structure of elementary particles.
Inspired by the aforementioned thoughts and othersLagally (1922); Landweber and Yih (1956); Yih (1969); Faber (1995); Currie (2003), we show that the Newton’s law of gravitation is derived based on the assumption that all the particles are made of singularities of a kind of ideal fluid.
During the preparation of the manuscript, I noticed that John C. Taylor had proposed an idea that the inverse-square law of gravitation may be explained based on the concept of source or sink Taylor (2001).
## II Forces acting on sources and sinks in ideal fluids
The purpose of this section is to calculate the forces between sources and sinks in inviscid incompressible fluids which is called ideal fluids usually.
Suppose the velocity field u of an ideal fluid is irrotational, then we have Lamb (1932); Kochin et al. (1964); Yih (1969); Wu (1982); Landau and Lifshitz (1987); Faber (1995); Currie (2003),
$$\text{u}=\varphi ,$$
(2)
where $`\varphi `$ is the velocity potential, $`=\text{i}\frac{}{x}+\text{j}\frac{}{y}+\text{k}\frac{}{z}`$ is the Hamilton operator.
It is known that the equation of mass conservation of an ideal fluid becomes Laplace’s equation Lamb (1932); Kochin et al. (1964); Yih (1969); Wu (1982); Faber (1995); Currie (2003),
$$^2\varphi =0,$$
(3)
where $`\varphi `$ is velocity potential, $`^2=\frac{^2}{x^2}+\frac{^2}{y^2}+\frac{^2}{z^2}`$ is the Laplace operator.
Using spherical coordinates$`(r,\theta ,\phi )`$, a general form of solution of Laplace’s equation (3) can be obtained by separation of variables asCurrie (2003)
$$\varphi (r,\theta )=\underset{l=0}{\overset{\mathrm{}}{}}\left(A_lr^l+\frac{B_l}{r^{l+1}}\right)P_l(\mathrm{cos}\theta ),$$
(4)
where $`A_l`$ and $`B_l`$ are arbitrary constants, $`P_l(x)`$ are Legendre’s function of the first kind which is defined as
$$P_l(x)=\frac{1}{2^ll!}\frac{\mathrm{d}^l}{\mathrm{d}x^l}(x^21)^l.$$
(5)
If there exists a velocity field which is continuous and finite at all points of the space, with the exception of individual isolated points, then these isolated points are called singularities usually.
###### Definition 1
Suppose there exists a singularity at point $`P_0=(x_0,y_0,z_0)`$. If the velocity field of the singularity at point $`P=(x,y,z)`$ is
$$\text{u}(x,y,z,t)=\frac{Q}{4\pi r^2}\widehat{\text{r}},$$
(6)
where $`r=\sqrt{(xx_0)^2+(yy_0)^2+(zz_0)^2}`$, $`\widehat{\text{r}}`$ denotes the unit vector directed outward along the line from the singularity to the point $`P=(x,y,z)`$, then we call this singularity a source if $`Q>0`$ or a sink if $`Q<0`$. $`Q`$ is called the strength of the source or sink.
Suppose a static point source with strength $`Q`$ locates at the origin $`(0,0,0)`$. In order to calculate the volume leaving the source per unit time, we may enclose the source with an arbitrary spherical surface $`S`$ with radius $`a`$. A calculation shows that
$$_S\text{u}\text{n}dS=_S\frac{Q}{4\pi a^2}\widehat{\text{r}}\text{n}dS=Q,$$
(7)
where n denotes the unit vector directed outward along the line from the origin of the coordinates to the field point$`(x,y,z)`$. Equation (7) shows that the strength $`Q`$ of a source or sink evaluates the volume of the fluid leaving or entering a control surface per unit time.
From (4), we see that the velocity potential $`\varphi (r,\theta )`$ of a source or sink is a solution of Laplace’s equation (3).
###### Theorem 2
Suppose (1) there exists an ideal fluid (2) the ideal fluid is irrotational and barotropic, (3) the density $`\rho `$ is homogeneous, that is $`\rho /x=\rho /y=\rho /z=\rho /t=\mathrm{\hspace{0.17em}0},`$ (4) there are no external body forces exerted on the fluid, (5)the fluid is unbounded and the velocity of the fluid at the infinity is approaching to zero. Suppose a source or sink is stationary and is immersed in the ideal fluid. Then, there is a force
$$\text{F}_Q=\rho Q\text{u}_0$$
(8)
exerted on the source by the fluid, where $`\rho `$ is the density of the fluid, $`Q`$ is the strength of the source or the sink, $`\text{u}_0`$ is the velocity of the fluid at the location of the source induced by all means other than the source itself.
Proof. Only the proof of the case of a source is needed. Let us select the coordinates that is attached to the static fluid at the infinity.
We set the origin of the coordinates at the location of the source. We surround the source by two arbitrary spherical surfaces $`S_\epsilon `$ and $`S`$ centered at the origin of the coordinates. The radii of the surfaces $`S_\epsilon `$ and $`S`$ are $`r_\epsilon `$ and $`r`$ respectively. The radius $`r_\epsilon `$ of the surfaces $`S_\epsilon `$ is smaller than $`r`$ and is arbitrarily small. The outward unit normal to the surface $`S`$ is denoted by n.
Let $`\tau (t)`$ denotes the mass system of fluid enclosed in the volume between the surface $`S_\epsilon `$ and $`S`$ at time $`t`$. Let $`\text{F}_Q`$ denotes the hydrodynamic force exerted on the source by the mass system $`\tau `$, then a reaction $`\text{F}_Q`$ of this force $`\text{F}_Q`$ must act on the the surfaces $`S_\epsilon `$ enclosing the source. Let $`\text{F}_S`$ denotes the hydrodynamic force exerted on the mass system $`\tau `$ due to the pressure distribution on the surface $`S`$.
As an application of the Newton’s second law of motion to the mass system $`\tau `$, we have
$$\frac{\mathrm{D}\text{K}}{\mathrm{D}t}=\text{F}_Q+\text{F}_S,$$
(9)
where K denotes momentum of the mass system $`\tau `$, $`\mathrm{D}/\mathrm{D}t`$ represents the material derivative in the lagrangian system Lamb (1932); Kochin et al. (1964); Yih (1969); Wu (1982); Landau and Lifshitz (1987); Faber (1995); Currie (2003).
The expressions of the momentum K and the force $`\text{F}_S`$ are
$$\text{K}=_\tau \rho \text{u}dV,\text{F}_S=_S(p)\text{n}dS,$$
(10)
where the first integral is volume integral, the second integral is surface integral, n denotes the unit vector directed outward along the line from the origin of the coordinates to the field point$`(x,y,z)`$.
Since the velocity field is irrotational, we have the following relation
$$\text{u}=\varphi ,$$
(11)
where $`\varphi `$ is the velocity potential.
According to Ostrogradsky–Gauss theorem (see, for instance, Kochin et al. (1964); Yih (1969); Wu (1982); Faber (1995); Currie (2003)), we have
$$_\tau \rho \text{u}dV=_\tau \rho \varphi dV=_S\rho \varphi \text{n}dS.$$
(12)
Note that the mass system $`\tau `$ does not include the singularity at the origin. According to Reynolds’ transport theorem Kochin et al. (1964); Yih (1969); Wu (1982); Faber (1995); Currie (2003), we have
$$\frac{\mathrm{D}}{\mathrm{D}t}_\tau \rho \text{u}𝑑V=\frac{}{t}_V\rho \text{u}𝑑V+_S\rho \text{u}(\text{u}\text{n})𝑑S,$$
(13)
where $`V`$ is the volume fixed in space which coincide with the mass system $`\tau (t)`$ at time $`t`$, that is $`V=\tau (t)`$.
Then, using (13) , (10) and (12), we have
$$\frac{\mathrm{D}\text{K}}{\mathrm{D}t}=_S\rho \frac{\varphi }{t}\text{n}dS+_S\rho \text{u}(\text{u}\text{n})dS.$$
(14)
According to Lagrange–Cauchy integral Kochin et al. (1964); Yih (1969); Wu (1982); Faber (1995); Currie (2003), we have
$$\frac{\varphi }{t}+\frac{(\varphi )^2}{2}+\frac{p}{\rho }=f(t),$$
(15)
where $`f(t)`$ is an arbitrary function of time $`t`$.
According to the assumption that the velocity u of the fluid at the infinity is approaching to zero, we have
$$\text{u}0,r\mathrm{}.$$
(16)
Noticing (4), $`\varphi (t)`$ must be of the following form
$$\varphi (r,\theta ,t)=\underset{l=0}{\overset{\mathrm{}}{}}\frac{B_l(t)}{r^{l+1}}P_l(\mathrm{cos}\theta ),$$
(17)
where $`B_l(t),l0`$ are functions of time $`t`$. Thus, we have the following estimations at the infinity of the velocity field
$$\varphi =O\left(\frac{1}{r}\right),\frac{\varphi }{t}=O\left(\frac{1}{r}\right),r\mathrm{},$$
(18)
where $`\phi (x)=O(\psi (x)),xa`$ stands for $`\overline{lim}_{xa}\phi (x)/`$$`\psi (x)=k,(0k<+\mathrm{}).`$
Applying (15) at the infinity and using (18), we have
$$\frac{\varphi }{t}0,pp_{\mathrm{}},r\mathrm{},$$
(19)
where $`p_{\mathrm{}}`$ is a constant.
Putting (19) and (16) into (15), we have
$$f(t)=\frac{p_{\mathrm{}}}{\rho }.$$
(20)
Putting (20) into (15), we have
$$p=p_{\mathrm{}}\rho \frac{\varphi }{t}\frac{\rho (\text{u}\text{u})}{2}.$$
(21)
Using (10) and (21), we have
$$\text{F}_S=_S\rho \frac{\varphi }{t}\text{n}dS+_S\frac{\rho (\text{u}\text{u})\text{n}}{2}dS.$$
(22)
Using (9), (14), (22), we have
$$\text{F}_Q=_S[\frac{1}{2}\rho (\text{u}\text{u})\text{n}\rho \text{u}(\text{u}\text{n})]dS.$$
(23)
Now let us calculate this velocity u in order to obtain $`\text{F}_Q`$. Since the velocity field induced by the source $`Q`$ is (6), then according to the superposition principle of velocity field of ideal fluids, the velocity on the surface $`S`$ is
$$\text{u}=\frac{Q}{4\pi r^2}\text{n}+\text{u}_0,$$
(24)
where n denotes the unit vector directed outward along the line from the origin of the coordinates to the field point$`(x,y,z)`$. Using (23) and (24), we have
$`\text{F}_Q`$ $`=`$ $`\rho {\displaystyle }{\displaystyle _S}[{\displaystyle \frac{Q^2}{32\pi ^2r^4}}\text{n}+{\displaystyle \frac{1}{2}}(\text{u}_0\text{u}_0)\text{n}`$ (25)
$`{\displaystyle \frac{Q}{4\pi r^2}}\text{u}_0(\text{u}_0\text{n})\text{u}_0]dS.`$
Since the radius $`r`$ can be arbitrarily small, the velocity $`\text{u}_0`$ can be treated as a constant in the integral of (25). Thus, (25) turns out to be
$$\text{F}_Q=\rho _S\frac{Q}{4\pi r^2}\text{u}_0dS.$$
(26)
Since again $`\text{u}_0`$ can be treated as a constant, (26) turns out to be (8). This completes the proof. $`\mathrm{}`$
Remark. Lagally Lagally (1922), Landweber and Yih Landweber and Yih (1956); Yih (1969), Faber Faber (1995) and Currie Currie (2003) obtained the same result of Theorem 2 for the special case where the velocity field is steady.
Theorem 2 only considers the situation that the sources or sinks are at rest. Now let us consider the case that the sources or sinks are moving in the fluid.
###### Theorem 3
Suppose the presuppositions (1), (2), (3), (4) and (5) in Theorem 2 are valid and a source or a sink is moving in the fluid with a velocity $`\text{v}_s`$, then there is a force
$$\text{F}_Q=\rho Q(\text{u}_f\text{v}_s)$$
(27)
is exerted on the source by the fluid, where $`\rho `$ is the density of the fluid, $`Q`$ is the strength of the source or the sink, $`\text{u}_f`$ is the velocity of the fluid at the location of the source induced by all means other than the source itself.
Proof. The velocity of the fluid relative to the source at the location of the source is $`\text{u}_f\text{v}_s`$. Let us select the coordinates that is attached to the source and set the origin of the coordinates at the location of the source. Then (27) can be arrived following the same procedures in the proof of Theorem 2. $`\mathrm{}`$
Applying Theorem 3 to the situation that a source or sink is exposed to the velocity field of another source or sink,we have:
###### Corollary 4
Suppose the presuppositions (1), (2), (3), (4) and (5) in Theorem 2 are valid and a source or a sink with strength $`Q_2`$ is exposed to the velocity field of another source or sink with strength $`Q_1`$, then the force $`\text{F}_{21}`$ exerted on the singularity with strength $`Q_2`$ by the velocity field of the singularity with strength $`Q_1`$ is
$$\text{F}_{21}=\rho Q_2\frac{Q_1}{4\pi r^2}\widehat{\text{r}}_{21}+\rho Q_2\text{v}_2,$$
(28)
where $`\widehat{\text{r}}_{21}`$ denotes the unit vector directed outward along the line from the singularity with strength $`Q_1`$ to the singularity with strength $`Q_2`$, $`r`$ is the distance between the two singularities, $`\text{v}_2`$ is the velocity of the source with strength $`Q_2`$.
## III Derivation of inverse-square-law of gravitation
Since quantum theory shows that vacuum is not empty and has physical effects, e.g., the Casimir effectLamoreaux (2005); Intravaia and Lambrecht (2005); Guo and Zhao (2004); Davies (2005), it is valuable to probe vacuum by introducing the following hypotheses:
###### Assumption 5
Suppose the universe is filled by an ideal fluid named $`\mathrm{\Omega }(0)`$ substratum; the ideal fluid fulfil the conditions (2), (3), (4), (5) in Theorem 2.
This fluid may be named $`\mathrm{\Omega }(0)`$ substratum in order to distinguish with Cartesian aether.
The idea that all microscopic particles are source or sink flows in a fluidic substratum is not new. For instance, in order to compare fluid motions with electric fields, J. C. Maxwell introduced an analogy between source or sink flows and electric charges (Whittaker (1951), p243). H. Poincar$`\stackrel{´}{e}`$ also speculates that matters are holes in fluidic aether (Poincare (1997), p171). A. Einstein and L. Infeld said (Einstein and Infeld (1938), p256-257):”Matter is where the concentration of energy is great, field where the concentration of energy is small. $`\mathrm{}`$ What impresses our senses as matter is really a great concentration of energy into a comparatively small space. We could regard matter as the regions in space where the field is extremely strong.” It seems that they are suggesting that particles are some kinds of singularities of field.
Following these researchers, we adopt the idea that particles may be looked as singularities in fields. Noticing (28), it is nature to introduce the following:
###### Assumption 6
All the microscopic particles were made up of a kind of elementary sinks of $`\mathrm{\Omega }(0)`$ substratum. These elementary sinks were created simultaneously. The initial masses and the strengths of the elementary sinks are the same.
We may call these elementary sinks as monads. Suppose a particle with mass $`m`$ is composed of $`N`$ monads. Then, according to Assumption 6, we have:
$`\text{ }m_0(t)=m_0(0)+\rho q_0t,`$ (29)
$`\text{ }Q=Nq_0,m(t)=Nm_0(t)={\displaystyle \frac{Q}{q_0}}m_0(t),`$ (30)
$`\text{ }{\displaystyle \frac{dm_0}{dt}}=\rho q_0,{\displaystyle \frac{dm}{dt}}=\rho Q,`$ (31)
where $`m_0(t)`$ is the mass of monad at time $`t`$, $`q_0(q_0>0)`$ is the strength of a monad, $`m(t)`$ is the mass of a particle at time $`t`$, $`Q`$ is the strength of the particle, $`N`$ is the number of monads that make up the particle, $`\rho `$ is the density of the $`\mathrm{\Omega }(0)`$ substratum, $`t0`$.
From (31), we see that the mass $`m_0`$ of a monad is increasing since $`q_0`$ evaluates the volume of the $`\mathrm{\Omega }(0)`$ substratum fluid entering the monad per unit time. From (31), we also see that the mass of a monad or a particle is increasing linearly.
Based on Assumption 5 and Assumption 6, the motion of a particle is determined by:
###### Theorem 7
The equation of motion of a particle is
$$m(t)\frac{\mathrm{d}\text{v}}{\mathrm{d}t}=\frac{\rho q_0}{m_0(t)}m(t)\text{u}\frac{\rho q_0}{m_0(t)}m(t)\text{v}+\text{F},$$
(32)
where $`m_0(t)`$ is the mass of monad at time $`t`$, $`q_0`$ is thestrength of a monad, $`m(t)`$ is the mass of a particle at time $`t`$, v is the velocity of the particle, u is the velocity of the $`\mathrm{\Omega }(0)`$ substratum at the location of the particle induced by all means other than the particle itself, F denotes other forces.
Proof. Applying the Newton’s second law and Theorem 3 to this particle, we have
$$m\frac{d\text{v}}{dt}=\rho Q(\text{u}\text{v})+\text{F}.$$
(33)
Putting (30) into (33), we get (32). $`\mathrm{}`$
Formula (32) shows that there exists a universal damping force
$$\text{F}_d=\frac{\rho q_0}{m_0}m\text{v}$$
(34)
exerted on each particle.
Now let us consider a system consists of two particles. Based on Assumption 6, applying Theorem 7 to this system, we have:
###### Corollary 8
Suppose there is a system consists of two particles and there are no other forces exerted on the particles, then the equations of motion of this system are
$`m_1{\displaystyle \frac{d\text{v}_1}{dt}}={\displaystyle \frac{\rho q_0}{m_0}}m_1\text{v}_1{\displaystyle \frac{\rho q_0^2}{4\pi m_0^2}}{\displaystyle \frac{m_1m_2}{r^2}}\widehat{\text{r}}_{12}`$ (35)
$`m_2{\displaystyle \frac{d\text{v}_2}{dt}}={\displaystyle \frac{\rho q_0}{m_0}}m_2\text{v}_2{\displaystyle \frac{\rho q_0^2}{4\pi m_0^2}}{\displaystyle \frac{m_1m_2}{r^2}}\widehat{\text{r}}_{21},`$ (36)
where $`m_{i=1,2}`$ is the mass of the particles, $`\text{v}_{i=1,2}`$ is the velocity of the particles, $`m_0`$ is the mass of a monad, $`q_0`$ is the strength of a monad, $`\rho `$ is the density of the $`\mathrm{\Omega }(0)`$ substratum, $`\widehat{\text{r}}_{12}`$ denotes the unit vector directed outward along the line from the particle with mass $`m_2(t)`$ to the particle with mass $`m_1(t)`$, $`\widehat{\text{r}}_{21}`$ denotes the unit vector directed outward along the line from the particle with mass $`m_1(t)`$ to the particle with mass $`m_2(t)`$.
Ignoring the damping forces in (36), we have:
###### Corollary 9
Suppose (1) $`\text{v}_{i=1,2}\text{u}_{i=1,2}`$, where $`\text{v}_i`$ is the velocity of the particle with mass $`m_i`$, $`\text{u}_i`$ is the velocity of the $`\mathrm{\Omega }(0)`$ substratum at the location of the particle with mass $`m_i`$ induced by the other particle, (2) there are no other forces exerted on the particles, then the force $`\text{F}_{21}(t)`$ exerted on the particle with mass $`m_2(t)`$ by the velocity field of $`\mathrm{\Omega }(0)`$ substratum induced by the particle with mass $`m_1(t)`$ is
$$\text{F}_{21}(t)=G(t)\frac{m_1(t)m_2(t)}{r^2}\widehat{\text{r}}_{21},$$
(37)
where
$$G(t)=\frac{\rho q_0^2}{4\pi m_0^2(t)},$$
(38)
$`\widehat{\text{r}}_{21}`$ denotes the unit vector directed outward along the line from the particle with mass $`m_1(t)`$ to the particle with mass $`m_2(t)`$, $`r`$ is the distance between the two particles.
Corollary 9 is similar to Newton’s inverse-square-law of gravitation (1) except for two differences. The first difference is that $`m_{i=1,2}`$ are constants in the Newton’s law (1) while in Corollary 9 they are functions of time $`t`$. The second difference is that $`G`$ is a constant in the Newton’s theory while here $`G`$ is a function of time $`t`$.
Let us now introduce an assumption that $`G`$ and the masses of particles are changing so slowly relative to the time scale of human beings that they can be treated as constants approximately. Thus, the Newton’s law of gravitation may be considered as a result of Corollary 9 based on this assumption.
## IV Superposition principle of gravitational field
The definition of gravitational field g of a particle with mass $`m`$ is
$$\text{g}=\frac{\text{F}}{m_{test}},$$
(39)
where $`m_{test}`$ is the mass of a test point mass, F is the gravitational force exerted on the test point mass by the gravitational field of the particle with mass $`m`$.
Based on (39), Theorem 7 and Corollary 9, we have
$$\text{g}=\frac{\rho q_0}{m_0}\text{u},$$
(40)
where $`\rho `$ is the density of the $`\mathrm{\Omega }(0)`$ substratum, $`m_0`$ is the mass of a monad, $`q_0`$ is the strength of a monad, u is the velocity of the $`\mathrm{\Omega }(0)`$ substratum at the location of the test point mass induced by the particle mass $`m`$.
From (40), we see that the superposition principle of gravitational field is deduced from the superposition theorem of the velocity field of ideal fluids.
## V Equations of gravitational field of continuously distributed particles
The definition of the volume density of continuously distributed sink is
$$\rho _s=\underset{\mathrm{}V0}{lim}\frac{\mathrm{}Q}{\mathrm{}V},$$
(41)
where $`\mathrm{}V`$ is a small volume, $`\mathrm{}Q`$ is the sum of strengthes of all the sinks in the volume $`\mathrm{}V`$.
Now let us to derive the continuity equation of the taothe $`\mathrm{\Omega }(0)`$ substratum from the principle of mass conservation. Consider an arbitrary volume $`V`$ bounded by a closed surface $`S`$ and fixed in space. Suppose there are some point sources or sinks continuously distributed in the volume $`V`$. The total mass in volume $`V`$ is
$$M=_V\rho 𝑑V,$$
(42)
where $`\rho `$ is the density of the $`\mathrm{\Omega }(0)`$ substratum.
The rate of increase of the total mass in volume $`V`$ is
$$\frac{M}{t}=\frac{}{t}_V\rho 𝑑V.$$
(43)
The rate of mass outflow through the surface $`S`$ is
$$_S\rho (\text{u}\text{n})𝑑S,$$
(44)
where u is the velocity field of the $`\mathrm{\Omega }(0)`$ substratum.
The rate of mass created inside the volume $`V`$ is
$$_V\rho \rho _s𝑑V.$$
(45)
Now according to the principle of mass conservation, and making use of (43), (44) and (45), we have
$$\frac{}{t}_V\rho 𝑑V=_V\rho \rho _s𝑑V_S\rho (\text{u}\text{n})𝑑S$$
(46)
According to Ostrogradsky–Gauss theorem (refer to, for instance, Kochin et al. (1964); Yih (1969); Wu (1982); Faber (1995); Currie (2003)), we have
$$_S\rho (\text{u}\text{n})𝑑S=_V(\rho \text{u})𝑑V,$$
(47)
Using (47), (46) becomes
$$\frac{}{t}_V\rho 𝑑V=_V\rho \rho _s𝑑V_V(\rho \text{u})𝑑V$$
(48)
Since the volume $`V`$ is arbitrary, from (48) we have
$$\frac{\rho }{t}+(\rho \text{u})=\rho \rho _s.$$
(49)
Since the $`\mathrm{\Omega }(0)`$ substratum is homogeneous, i.e., $`\rho /x=\rho /y=\rho /z=\rho /t=0,`$ (49) becomes
$$\text{u}=\rho _s.$$
(50)
Thus, from (40) and (50), we have
$$\text{G}=\frac{\rho q_0\rho _s}{m_0}.$$
(51)
Let $`\rho _m`$ denotes the volume mass density of continuously distributed particles. According to Assumption 6, and using (30)and (41), we have
$$\rho _m=\frac{m_0\rho _s}{q_0}.$$
(52)
Thus, noticing that the velocity field u of the $`\mathrm{\Omega }(0)`$ substratum is irrotational and $`G=\rho q_0^2/(4\pi m_0^2)`$, the equations of gravitational field of continuously distributed particles can be summarized as
$$\{\begin{array}{cc}\times \text{G}=\text{0},\hfill & \\ \text{G}=4\pi G\rho _m.\hfill & \end{array}$$
(53)
## VI The equivalence of inertial mass and gravitational mass
According to Assumption 6 and Corollary 9, we have
$$m_{inertial}=m_{gravitational},$$
(54)
where $`m_{inertial}`$ is the inertial mass of a particle, $`m_{gravitational}`$ is the gravitational mass of the particle.
## VII Time dependence of gravitational constant $`G`$ and mass
The time dependence of the gravitational constant $`G`$ can be seen from (38).
The time dependence of the gravitational mass can be seen from (31).
## VIII Possible space dependence of gravitational constant $`G`$
If the density $`\rho `$ of the $`\mathrm{\Omega }(0)`$ substratum varies from place to place, i.e., $`\rho =\rho (\text{r})`$, then the space dependence of the gravitational constant $`G`$ can be seen from (38).
## IX Discussion
Although the new formula of gravitation (37) is similar to Newton’s inverse-square-law of gravitation (1), there exists the following five differences between this theory and Newton’s theory.
1. The gravitational masses are constants in Newton’s law, while in (37) they are functions of time $`t`$.
2. The gravitational constant $`G`$ is a constant in Newton’s theory, while in (38) $`G`$ depends on time $`t`$.
3. In this theory, the parameter $`G`$ depends on the density $`\rho _{\mathrm{\Omega }(0)}`$ of the $`\mathrm{\Omega }(0)`$ substratum. If $`\rho _{\mathrm{\Omega }(0)}`$ varies from place to place, then the space dependence of the gravitational constant $`G`$ can be seen from (38).
4. In Newton’s theory, the gravity is action-at-a-distance Whittaker (1953). In this theory, the gravity is transmitted by the $`\mathrm{\Omega }(0)`$ substratum.
5. Newton’s law of gravitation is an assumption. In this theory, (37) is derived by methods of classical fluid mechanics based on some assumptions.
## X Conclusion
We suppose that the universe may be filled with a kind of fluid which may be called the $`\mathrm{\Omega }(0)`$ substratum. Thus, the inverse-square law of gravitation is derived by methods of hydrodynamics based on a sink flow model of particles. There are two features of this theory of gravitation. The first feature is that the gravitational interactions are transmitted by a kind of fluidic medium. The second feature is the time dependence of gravitational constant and gravitational mass. Newton’s law of gravitation is derived if we introduce an assumption that $`G`$ and the masses of particles are changing so slowly that they can be treated as constants. As a byproduct, it is shown that there exists a universal damping force exerted on each particle.
## Acknowledgments
I am grateful to Robert L. Oldershaw for informing me his researches in the field of Self-Similar Cosmological Paradigm (SSCP) and discrete fractal cosmological models during the preparation of the manuscript. I wish to express my thanks to Dr. Roy Keys for providing me three interesting articles Martin (2005a, b); Martin and Keys (1994).
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# EFFECTS OF MAGNETIC FLUX AND OF ELECTRON MOMENTUM ON THE TRANSMISSION AMPLITUDE IN THE AHARONOV-BOHM INTERFEROMETER 1footnote 11footnote 1This work forms a part of the unpublished project report of the first author submitted for partial fulfillment of the degree of Master of Science in the department of Physics of Sri Sathya Sai Institute of Higher Learning (Deemed University), Vidya Giri, Prasanthi Nilayam- 516 134, A.P., India.
## I Introduction
The Aharonov-Bohm interferometer (ABI) in a two-terminal configuration has come to stay as a reliable tool in the study of mesoscopic/ nanoscopic systems. Figure1 presents a schematic diagram of the ABI.
The phase of the electron wave function has a topological contribution arising due to the magnetic flux confined within the ring, threading perpendicular to the ring plane. It also has a dynamic contribution arising due to the electron momentum. In a simplistic view $`t`$ is often assumed to be of the form
$$t=t_1+t_2e^{i\varphi }$$
where $`\varphi `$ is the Aharonov-Bohm (AB) phase and $`t_1`$, $`t_2`$ are the TA’s through the two arms of the ring. Whereas this view is appropriate to the original free-space double-slit setup, it is not so to the two-leaded ABI. More complicated dependence on $`\varphi `$ as well as the dynamical phases $`\chi _1`$ and $`\chi _2`$ contributed by the two arms dictate the behavior of the TAbutt-im-azbel ; cahay . It is the purpose of the present paper to present our analysis of the detailed effects of these phases on the TA. Chen et.al.chen\_shi report the effect of $`\varphi `$ on anti-resonances of the transmittance $`T(=t^2`$ in the ABI. They correlate the nature of constant-$`T`$ contour plots in the complex electron-energy plane to the variation of $`T`$ with the electron energy. Kim et.al.kim\_cho\_kim\_ryu study the effects of broken time-reversal symmetry on the transmission zeros of the ABI. They propose a classification scheme of the phase trajectories of $`t`$ in the complex $`t`$-plane and study the effect of the variation of the $`\chi `$’s. Taniguchi and Buttikertanibutt , in a pioneering work, analyze the Friedel phases and the TA phases for a one-dimensional resonant tunnel double barrier, a wire with a side branch and also for the combination of the above two structures.
In this work we consider the transmission properties of the two-leaded AB ring (see Fig.1). We have repeated the calculation of Xiaxia for the ABI for obtaining the TA. The model and the assumptions that we make are same as the ones made in his workxia . The wave functions on the various segments are assumed to be of the following form:
$`\psi _1(x_1)`$ $`=`$ $`e^{ikx_1}+ae^{ikx_1},`$
$`\psi _2(x_2)`$ $`=`$ $`c_1e^{ik_1x_2}+d_1e^{ik_2x_2},`$
$`\psi _3(x_3)`$ $`=`$ $`c_2e^{ik_2x_3}+d_2e^{ik_1x_3},`$
$`\psi _4(x_4)`$ $`=`$ $`te^{ikx_4}.`$ (1)
The subscripts on the $`\psi `$’s correspond to the labels used on the segments (see Fig.1). The lengths of the segments marked $`2`$ and $`3`$ are $`l_2`$ and $`l_3`$ and the total circumference of the ring is $`l=l_2+l_3`$. Here $`x_1`$, $`x_2`$ and $`x_3`$ have common origin at A and $`x_4`$ has its origin at B. The boundary conditions used at the junctions are
$`\psi _1(0)=\psi _2(0)=\psi _3(0),`$
$`\psi _2(l_2)=\psi _3(l_3)=\psi _4(0),`$
$`\psi _1^{}(0)=\psi _2^{}(0)+\psi _3^{}(0),`$
$`\psi _2^{}(l_2)+\psi _3^{}(l_3)=\psi _4^{}(0).`$
(2)
Here $`k_1=k+\eta `$ and $`k_2=k\eta `$ and $`\eta =(2\pi /l)(\mathrm{\Phi }/\mathrm{\Phi }_0)=(2\pi /l)\theta `$, $`\mathrm{\Phi }`$ denotes the magnetic flux threading the ring and $`\mathrm{\Phi }_0=hc/e`$ is the flux quantum. It is convenient to use a dimensionless wave number $`q`$ defined by $`k=(2\pi /l)q`$ and $`\theta `$, the dimensionless flux. We have used the software MAPLE to obtain analytic expression for the TA.
We report on the characteristic features of the complex $`t`$-plane plots when the magnetic flux is made to vary over one or several flux periods. We refer to such plots as type I $`t`$-plot. Next we examine similar plots when the electron momentum is varied. We refer to such plots as type II plots.
## II Type I $`t`$-plots
We have analyzed the complex $`t`$-plane plots for a symmetric (i.e., $`l_2=l_3`$) as well as for an asymmetric (i.e., $`l_2l_3`$) ring. We report our findings under separate headings for these two cases. Note that our usage of the phrases “symmetric” and “asymmetric” ABI is different from the conventional usage.
### II.1 The symmetric ring interferometer
For a symmetric ring ABI, for arbitrary values of $`q`$, the typical behavior of $`T`$ as $`\theta `$ is varied is shown in the curves marked $`1`$ and $`2`$ in Fig.2.
Note the presence of minima of $`T`$ at integer $`q`$ values. In contrast, sharp transmission ones with characteristic appearance of resonances arise for values of $`q`$ in close proximity of integer $`q`$’s corresponding to discrete energy eigenstates. These resonances were missed in the calculation of Xiaxia , but their existence has been anticipated by Cahay et.al.cahay . Note that the perfect ring (i.e., a ring with no attached leads) supports eigenstates of $`q`$ for integer $`q`$’s corresponding to energy eigenvalues $`E_n=n^2h^2/2\mu l^2`$ (corresponding to $`q=n`$), where $`\mu `$ is the effective mass of the electron.
The typical behavior of $`T`$ versus $`\theta `$ (for arbitrary $`q`$) exhibiting transmission zeros for half odd integer values of $`\theta `$ and local minima for integer values of $`\theta `$ is shown in the thin line curve of Fig.2. However, for special values of $`q`$ in the proximity of integer $`q`$ values, sharp transmission ones’ appear (thick line in Fig.2).
The flux periodicity of $`1`$ reflected in the pattern of Fig.2 may give the impression that the phenomenon of electron transmission in the ABI is also flux periodic with the same period $`1`$. However, this expectation is not borne out by the following results relating to the TA. For $`q=n+0.01`$ ($`n`$ being an integer), i.e., for $`q`$-values very close to integer values $`n`$ corresponding to the transmission ones, the trajectory in the complex $`t`$-plane is shown in Fig.3.
It starts from a point $`A`$ for $`\theta =0`$, covers the path $`ABO`$ as $`\theta `$ is made to vary from $`0`$ to $`1/2`$ and thereafter follows the path $`OCD`$ as $`\theta `$ varies from $`1/2`$ to $`1`$, reaching the point $`D`$. If $`\theta `$ is varied further from $`1`$ to $`2`$, the path $`DCOBA`$ is followed, thus retracing the path to generate a periodic trajectory. Thus two flux periods are necessary to complete one periodic cycle in the $`t`$-plane. The above description holds also for $`q=n0.01`$; only the shape of the trajectory appears reflected about the real axis. We note that the above two values of $`q`$ corresponds to energy values–one slightly above and the other slightly below the discrete eigenvalues of the electron energy (in the perfect ring).
For $`q`$ values covering the range $`nqn+\frac{1}{2}`$, the shape of the trajectory remains unaltered; only the curve gets rotated progressively in the anticlockwise (or clockwise, depending on the direction of the flux) direction to attain the maximum rotation by an angle $`\pi /2`$. As $`q`$ is increased further to cover the range $`n+\frac{1}{2}qn+1`$, the trajectory has a “flipping point” when it’s shape gets completely reflected about the Im($`t`$) axis as $`q`$ crosses the value $`n+\frac{1}{2}`$, but keeps rotating in the same direction by a further amount of $`\pi /2`$. The final shape is not back to its initial shape. Further variation of $`q`$ from $`n+1`$ to $`n+2`$ is required to generate one complete periodic cycle. This behavior is depicted in Fig.4.
### II.2 The asymmetric ring interferometer
In order to fix the ideas, we first consider the specific example of a ring with $`l_1:l_2=2:3`$. The cases corresponding to other rational fraction ratios are similar. The first interesting observation is that the trajectories form closed curves only upon varying $`\theta `$ over five ($`=2+3`$, the sum of the numbers in the ratio $`l_1:l_2`$) flux periods. A typical plot is presented in Fig.5 for $`q=1/4`$ (i.e., $`k=\pi /2l`$) when the flux is varied in the range $`0\theta 1`$. Note that the shape of this curve is similar to the curve in Fig.3, but the curve does not pass through the origin. In the case of the symmetric ABI, further variation beyond the above range only results in the repetition of this trajectory.
Hut in the present case, further variation of flux leads to evolution along different directions and when the variation covers five flux-periods (i.e., , $`0\theta 5`$, a closed flower-like pattern results (see Fig.7). Further variation of flux leads to reparation of this pattern. This behavior is indicative of an over-all periodicity of five in the underlying phenomenon. Note that for this value of $`q`$, the $`T`$ versus $`\theta `$ plot does not reveal this $`5`$-fold periodicity (see Fig.6).
Some plots for a few more typical $`q`$ values are shown in Figs.7.
Note that the passage of the trajectory through the origin signifies the presence of a transmission zero.
When the value of the ratio $`l_2:l_3`$ is not a rational fraction, the pattern generated is not a closed curve. Figure8 presents an example of such a pattern for $`l_2=1/\sqrt{5}`$ and $`l_3=11/\sqrt{5}`$. In this case the trajectory never closes and as $`\theta `$ is varied over longer and longer intervals, the pattern generated turns out progressively denser–the behavior reflected in the figure.
Clearly, the phenomenon is not flux-periodic. We now report on the behavior of the complex $`t`$-plots when, in stead of the flux, the momentum $`q`$ is made to vary.
## III Type II $`t`$-plots
The complex plane plot of $`t`$ as the electron momentum is varied have been discussed in the literature earlierkim\_cho\_kim\_ryu for the case of an ABI with an embedded quantum dot and also for an embedded double-barrier. We are not aware of such a calculation for a bare ABI.
For the symmetric ring ABI, we obtain the typical behaviors for different cases corresponding to fixed typical values of the flux. Figure 9 displays these behaviors.
Note that when the flux is zero (case (a)), the closed curve is an ellipse and it does not pass through the origin indicating the absence of transmission zeros. However, the presence of even a tiny flux causes transmission zeros to appear; thus the plots (cases b,c and d) exhibit curves passing through the origin. This qualitative difference between the two cases is remarkable. We note in passing that increase of flux to higher magnitudes does not alter the appearance shown in Fig.9b.
The case of an asymmetric ABI with the ratio $`l_2:l_3`$ being a rational fraction is shown in Figs.10.
The variation of $`q`$ over larger and larger intervals does not generate additional features. In fact, the plot for $`0q50`$ is no different from that presented in Fig.10b. This implies a kind of (rather complicated) $`q`$-periodicity in the phenomenon.
When the ratio $`l_2:l_3`$ is irrational, the type II $`t`$-plots, just as the type I $`t`$-plots, do not exhibit closed curves. Figure11 illustrates this point.
Note that as the variation of $`q`$ spreads over larger intervals, the pattern generated becomes denser. With these results we end with our conclusions.
## IV conclusion
Our results establish the utility of the type I $`t`$-plots in the characterization of the ABI. The information obtained via both types of plots is very rich and perhaps need to be pursued further. The fact that the flux periodicity nature of the transmission coefficient $`T`$ does not imply corresponding periodicity in the underlying phenomenon is clearly brought out in the present calculations. It may be of interest to study the changes arising in the plots when an impurity is embedded in one of the arms of the interferometer. We have obtained a host of new results by way of constant transmission contour plots and also investigated the effect of electric field acting on one of the arms of the interferometer. These results will be reported later.
## Acknowledgements
MVAK wishes to thank Prof. K. Venkataramaniah for his encouragement. DS wishes to thank the Director, Institute of Physics for his kind hospitality.
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# Number counts in homogeneous and inhomogeneous dark energy models
## 1 Introduction
In the present general picture of cosmology, converging evidences suggest that the matter density parameter is low and that the largest fraction of the energy density of the universe has an unknown nature leading to an accelerating phase. These indications come primarily from supernovae Ia data (e.g. Riess *et al.* (1998); Perlmutter *et al.* (1999); Riess *et al.* (2004)) and are corroborated by cosmic microwave background radiation (e.g. Spergel *et al.* (2003)) and large scale structure observations (e.g.Cole *et al.* (2005)), although there are different interpretations for the data (e.g. Blanchard *et al.* (2003); Shanks (2004)). A cosmological constant can explain the acceleration of the universe, however, the disagreement by $`120`$ orders of magnitudes with predictions from theoretical particle physics has shown the need for resorting to an alternative explanation. This is why several theoretical models were recently proposed to explain this dark energy in the universe. This new component can be identified with a slowly varying, self-interacting, neutral scalar “quintessence” field (Wetterich (1988); Ratra and Peebles (1988); Ferreira and Joyce (1997); Zlatev *et al.* (1999)) which can be minimally coupled, non-minimally coupled (e.g. Uzan (1999); Amendola (2000); Baccigalupi *et al.* (2000)), a phantom (e.g. Caldwell (1999)), a tachyon (e.g. Bagla *et al.* (2003)) or of purely kinetic nature (e.g. Armendariz-Picon *et al.* (2000)), known as K-essence. An alternative to dark energy is a fluid with a Chaplygin gas type of equation of state (e.g. Kamenshchik *et al.* (2001)). See Sahni (2004) and references therein for a review on some of these, and other, models.
Regardless of its nature, dark energy as a dominant component, plays a role in the structure formation and thus is likely to modify the number of formed structures. The evolution of linear perturbations in a scalar field like quintessence and the effects on structure formation have already been investigated theoretically (e.g. Ferreira and Joyce (1997); Perrotta and Baccigalupi (2002); Amendola (2003)). The effects on the abundance of collapsed structures and its evolution with redshift were also widely explored and suggested as a tool to constrain the dark energy’s nature and evolution (e.g. Haiman *et al.* (2001); Weller *et al.* (2002); Weinberg and Kamionkowski (2002); Battye and Weller (2003); Wang *et al.* (2004); Mohr (2004); Horellou and Berge (2005)).
Recently, numerical simulations including a dark energy component were performed by several groups to complement the analytical computations and to study the effects of dark energy at the structure level (e.g. shape of the dark matter halo, mass function) (Linder and Jenkins (2003); Lokas *et al.* (2003); Klypin *et al.* (2003); Kuhlen *et al.* (2004)). In such studies, the scalar field associated with dark energy is assumed not to have density fluctuations on scales of galaxy clusters or below. If dark energy influences the perturbations on small scales as proposed for example by Arbey *et al.* (2001), Bean and Magueijo (2002), Padmanabhan and Choudhury (2002) or Bagla *et al.* (2003), the collapse of structures as well as their properties will be affected.
Mota and van de Bruck (2004) have shown that the properties of collapsed halos (density contrast, virial radius) depend strongly on the shape of the potential, the initial conditions, the time evolution of the equation of state and on the behaviour of the scalar field in non-linear regions. This is what we will refer to as the inhomogeneity of the scalar field. More recently, Nunes and Mota (2004) have investigated how inhomogeneous quintessence models have a specific signature even in the linear regime of structure formation. They have shown that the time of collapse is affected by the inhomogeneity of dark energy and they have computed the resulting effect on the linearly extrapolated density threshold $`\delta _c`$. Moreover, they examined the evolution of matter overdensity as a function of time varying equation of state in homogeneous and inhomogeneous assumptions. Maor and Lahav (2005) have generalized the formalism to allow for a smooth transition between the homogeneous and inhomogeneous scenarios. They have concluded that, if only matter virializes, the final size of the system is fundamentally distinct from the one reached if the full system virializes.
In the present study, we extend the work of Nunes and Mota (2004) to investigate how the quintessence field affects the abundance of collapsed halos when the field follows the background evolution (homogeneous) and more specifically when it collapses with the dark matter (inhomogeneous). We compare the two assumptions for models with constant equation of state and more general cases of time-varying equation of state. To compute the structure abundances and their evolution with redshift, we use the canonical Press and Schechter (1974) formalism. Its theoretical expression allows us to account for the effects of inhomogeneous quintessence field through $`\delta _c`$, the growth factor as well as the volume element. Finally, we focus on the effects of the different models for structures with masses ranging between $`10^{13}`$ and $`10^{16}h^1M_{}`$.
This paper is organised as follows. In Section 2 we introduce the fundamental equations that describe the evolution of the quintessence field and the collapse of structure in the homogeneous and inhomogeneous hypothesis. In Section 3 we describe the method used to compute the number density of collapsed objects (mass function) in both of these scenarios, including the case where the equation of state of dark energy is allowed to vary with time. We give results and discuss the effects of the normalisation of the mass function on the predicted number counts in Section 4, and present concluding remarks in Section 5.
## 2 Theoretical background
In a spatially flat Friedmann-Robertson-Walker Universe the cosmic dynamics is determined by a background pressureless fluid (dark and visible matter), radiation and dark energy. The governing equations of motion are
$`\dot{H}`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}\left(\rho _B+p_B+\rho _{\mathrm{de}}+p_{\mathrm{de}}\right),`$ (1)
$`\dot{\rho }_B`$ $`=`$ $`3H(\rho _B+p_B),`$ (2)
$`\dot{\rho }_{\mathrm{de}}`$ $`=`$ $`3H(\rho _{\mathrm{de}}+p_{\mathrm{de}}),`$ (3)
with
$`H^2={\displaystyle \frac{\kappa ^2}{3}}\left(\rho _B+\rho _{\mathrm{de}}\right).`$ (4)
Here $`H=\dot{a}/a`$ is the expansion rate of the Universe, $`a(t)`$ is the scale factor, $`\kappa ^2=8\pi G`$ and $`\rho _B`$ and $`p_B`$ are the energy density and pressure of the background fluid, respectively. In this work the background is taken to be dominated by non-relativistic matter, hence, $`\rho _B=\rho _\mathrm{m}a^3`$. If dark energy is a perfect fluid its energy density and pressure are related by the equation of state $`\rho _{\mathrm{de}}=w\rho _{\mathrm{de}}`$ and $`\rho _{\mathrm{de}}=\mathrm{\Omega }_{\mathrm{de}}\rho _0/a^{3(w+1)}`$. Alternatively, dark energy can be described by a dynamical evolving scalar field rolling down its potential $`V(\varphi )`$. In this case, its energy density and pressure are defined as, $`\rho _\varphi =\dot{\varphi }^2/2+V`$ and $`p_\varphi =\dot{\varphi }^2/2V`$, respectively where the relation $`w=p_\varphi /\rho _\varphi `$, still holds. The equation of motion for the scalar field is,
$$\ddot{\varphi }=3H\dot{\varphi }\frac{dV}{d\varphi }.$$
(5)
We explore models of dark energy that have a significant contribution at high redshift unlike models such as the inverse power law (Zlatev *et al.* (1999)) or SUGRA models (Brax and Martin (1999)). We compare the results of such models with the cosmological constant model with $`w=1`$ and we assume $`\mathrm{\Omega }_{\mathrm{de}}=0.7`$ today for all models considered in the present study. Moreover, we choose the present value of the equation of state, $`w_0`$, and its running, $`dw/dz(z=0)`$, to be within the current observational limits (Riess *et al.* (2004)).
Figure 1 shows the evolution of the dark energy density and equation of state with redshift for the models we consider in this paper. We focus on dark energy models for which $`w=0.8`$, $`w=1.2`$ (phantom energy, Caldwell (1999)) and two cases where the dark energy results from a slowly evolving scalar field in a potential with two exponential terms (2EXP) (Barreiro *et al.* (2000))
$$V(\varphi )=V_0\left(e^{\alpha \kappa \varphi }+e^{\beta \kappa \varphi }\right).$$
(6)
We choose as in Nunes & Mota (2004) the pairs $`(\alpha ,\beta )=(6.2,0.1)`$ (2EXP<sub>1</sub>) and $`(\alpha ,\beta )=(20.1,0.5)`$ (2EXP<sub>2</sub>). Both provide an equation of state at present $`w_0=0.95`$. The equation of state for 2EXP<sub>1</sub> approaches zero faster than for 2EXP<sub>2</sub> (see bottom panel in Fig. 1). The two models differ also by their contributions at high redshift. As seen from top panel of Fig. 1, 2EXP<sub>1</sub> provides a contribution of dark energy that is non negligible at high redshifts ($`\mathrm{\Omega }_\varphi =0.1`$ at $`z=5`$ for 2EXP<sub>1</sub> whereas $`\mathrm{\Omega }_\varphi =0.02`$ for 2EXP<sub>2</sub> at same redshift).
In this work we use the spherical collapse model to describe the gravitational collapse of an overdense region of radius $`r`$ and density contrast $`\delta `$ such that $`1+\delta =\rho _{\mathrm{m}}^{}{}_{c}{}^{}/\rho _\mathrm{m}=(a/r)^3`$, where $`\rho _{\mathrm{m}}^{}{}_{c}{}^{}`$ and $`\rho _\mathrm{m}`$ are the energy densities of pressureless matter in the cluster and in the background, respectively. We have considered the possibility that dark energy also clusters, i.e. $`\rho _{\mathrm{de}}^{}{}_{c}{}^{}\rho _{\mathrm{de}}`$. Following Mota and van de Bruck (2004) and Nunes and Mota (2004), we study two extreme limits for the evolution of dark energy in the overdensity region. First, we assume that dark energy is “homogeneous”, i.e. the value of $`\rho _{\mathrm{de}}`$ inside the overdensity is the same as in the background. Second, dark energy is “inhomogeneous” and collapses with dark matter. In general terms, the evolution of dark energy inside a cluster can be written as in Mota and van de Bruck (2004)
$$\dot{\rho }_{\mathrm{de}_c}=3\frac{\dot{r}}{r}(\rho _{\mathrm{de}}^{}{}_{c}{}^{}+p_{\mathrm{de}}^{}{}_{c}{}^{})+\mathrm{\Gamma }_{\mathrm{de}},$$
(7)
where $`\mathrm{\Gamma }_{\mathrm{de}}`$ represents the dark energy loss inside the cluster and the ratio $`\dot{r}/r`$ is related to the Hubble ratio and the evolution of the density contrast through
$$\frac{\dot{r}}{r}=\frac{\dot{a}}{a}\frac{1}{3}\frac{\dot{\delta }}{1+\delta }.$$
(8)
If dark energy is homogeneous, then
$$\mathrm{\Gamma }_{\mathrm{de}}=3\left(\frac{\dot{a}}{a}\frac{\dot{r}}{r}\right)(\rho _{\mathrm{de}}^{}{}_{c}{}^{}+p_{\mathrm{de}}^{}{}_{c}{}^{}).$$
(9)
However if dark energy is inhomogeneous and collapses with dark matter, $`\mathrm{\Gamma }_{\mathrm{de}}=0`$ and the equation of motion of a scalar field inside the cluster is
$$\ddot{\varphi }_c=3\frac{\dot{r}}{r}\dot{\varphi _c}\frac{dV(\varphi _c)}{d\varphi _c}.$$
(10)
Finally, the evolution of the linear density contrast $`\delta _L`$ is determined by
$`\ddot{\delta }_L=2H\dot{\delta }_L+{\displaystyle \frac{\kappa ^2}{2}}\left[\rho _m\delta _L+(1+3w)\delta _{\mathrm{de}}\rho _{\mathrm{de}}+3\rho _{\mathrm{de}}\delta w\right],`$
where we have defined the linear density contrast in dark energy as $`\delta _{\mathrm{de}}=\delta \rho _{\mathrm{de}}/\rho _{\mathrm{de}}`$ and $`\delta w`$ is the linear perturbation in the equation of state.
## 3 Mass function
In hierarchical models, cosmic structures form from the gravitational amplification of small initial density perturbations. The time evolution of the structure abundances is determined mainly by the rate at which the perturbations grow until they reach the collapse, or virialization.
An analytical computation, proposed by Press and Schechter (1974), gives the comoving number density of collapsed dark matter halos of mass $`M`$ in the interval $`dM`$ at a given redshift of collapse, $`z`$ by
$`{\displaystyle \frac{dn}{dM}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\rho _{\mathrm{m0}}}{M}}{\displaystyle \frac{\delta _c(z)}{\sigma (M,z)}}{\displaystyle \frac{d\mathrm{ln}\sigma (M,z)}{dM}}\mathrm{exp}\left[{\displaystyle \frac{\delta _c(z)^2}{2\sigma (M,z)^2}}\right],`$
where $`\rho _{\mathrm{m0}}`$ is the present matter mean density of the universe and $`\delta _c(z)`$ is the linearly extrapolated density threshold above which structures collapse, i.e. $`\delta _c(z)=\delta _L(z=z_{\mathrm{col}})`$. In an Einstein-de Sitter Universe, an overdensity region collapses with a linear contrast $`\delta _c=1.686`$ (see e.g. Padmanabhan (1993)). This canonical value was used in the seminal paper of Press and Schechter (1974). The linearly extrapolated density threshold has recently been re-computed (Nunes & Mota 2004) in a more general case accounting for homogeneous and inhomogeneous dark energy. They show significant variations of $`\delta _c(z)`$ with redshift and with dark energy models considered. In the present study, we compute halo abundances using the results on $`\delta _c(z)`$ obtained by Nunes & Mota (2004), to which we refer the reader for further details.
The quantity $`\sigma (M,z)=g(z)\sigma _M`$ is the linear theory rms density fluctuation in spheres of radius $`R`$ containing the mass $`M`$ and $`g(z)=\delta _L(z)/\delta _L(z=0)`$ is the linear growth factor. The smoothing scale $`R`$ is often specified by the mass within the volume defined by the window function at the present time (e.g. Peebles 1980). In our analysis the variance of the smoothed overdensity containing a mass $`M`$ is given by
$$\sigma _M=\sigma _8\left(\frac{M}{M_8}\right)^{\gamma /3},$$
(13)
where $`M_8=6\times 10^{14}\mathrm{\Omega }_\mathrm{m}h^1M_{}`$, the mass inside a sphere of radius $`R_8=8h^1\mathrm{Mpc}`$, and $`\sigma _8`$ is the variance of the overdensity field smoothed on a scale of size $`R_8`$. The index $`\gamma `$ is a function of the mass scale and the shape parameter, $`\mathrm{\Gamma }`$, of the matter power spectrum (Viana and Liddle (1996))
$$\gamma =(0.3\mathrm{\Gamma }+0.2)\left[2.92+\frac{1}{3}\mathrm{log}\left(\frac{M}{M_8}\right)\right].$$
(14)
In our study we use $`\mathrm{\Gamma }=0.167`$ (Spergel et al. 2003).
For a fixed $`\sigma _8`$ (power spectrum normalization) the predicted number density of dark matter halos given by the above formula is uniquely affected by the dark energy models through the ratio $`\delta _c(z)/g(z)`$. The underlying assumption of this approach is that the transfer function used in the computation of $`\sigma _M`$ is that of a cosmological constant model. This is a good approximation at cluster scales for homogeneous dark energy models (Ma *et al.* (1999)), which remains to be theoretically investigated in the inhomogeneous hypothesis.
For a fixed local halo abundance, the important quantity to consider is thus $`\xi =\delta _c(z)/\sigma _8g(z)`$. We have verified that apart from $`w=1.2`$ (phantom energy) all homogeneous models have a ratio $`\delta _c/\sigma _8g`$ below that of the $`\mathrm{\Lambda }`$-model. This means that, at higher redshifts, models $`w=0.8`$, 2EXP<sub>1</sub> and 2EXP<sub>2</sub> are expected to give larger halo densities (whereas the phantom energy model is expected to give lower abundances) when compared to the cosmological constant model. For inhomogeneous dark energy we find that all models, except the $`w=0.8`$ model, have larger $`\delta _c/\sigma _8g`$ than the cosmological constant model, and therefore we expect at high redshift lower hallo densities for this models (higher abundances in the case of $`w=0.8`$) when compared to the $`\mathrm{\Lambda }`$-model.
To obtain the same halo abundance at $`z=0`$ we scaled $`\sigma _8`$ according to $`\sigma _8=\delta _c(z=0)\sigma _{8,\mathrm{\Lambda }}/\delta _{c,\mathrm{\Lambda }}(z=0)`$, where the index ‘<sub>Λ</sub>’ represents the cosmological constant model, and we fix $`\sigma _{8,\mathrm{\Lambda }}=0.9`$. Table 1 lists the values of $`\sigma _8`$ obtained in this way. Note the much larger dispersion of $`\sigma _8`$ between models in the inhomogeneous case (third column) as compared to homogeneous dark energy (second column) where models require practically the same $`\sigma _8`$ to reproduce the present-day halo abundance of the $`\mathrm{\Lambda }`$-model.
It is also interesting to see that the 2EXP<sub>1</sub> model requires the lowest and the $`w=1.2`$ model the highest $`\sigma _8`$ normalizations to reproduce the same local abundance of halos. This means that structures form early, at a slow rate, if the universe is dominated by phantom dark energy. They form late, at a faster rate, in the 2EXP<sub>1</sub> model.
In Fig. 2 we plot the redshift evolution of the mass function of objects with mass $`10^{14}h^1M_{}`$ for both homogeneous (top panel) and inhomogeneous (bottom panel) dark energy using Eq. (LABEL:eq:mf). As discussed in the previous paragraphs the $`w=0.8`$, 2EXP<sub>1</sub> and 2EXP<sub>2</sub> models give larger halo abundances than the $`w=1`$ model, whereas $`w=1.2`$ gives the lowest densities. We find this same qualitative behaviour within the mass range of interest for this paper, $`10^{13}10^{16}h^1M_{}`$, for both homogeneous and inhomogeneous dark energy models. The counts were normalized so that at $`z=0`$ all models give the same halo abundance as the cosmological constant ($`w=1`$) model gives for $`\sigma _8=0.9`$. Note that the overdensity contrast at collapse, $`\delta _c`$, is different for different dark energy models (see Fig. 2 in Nunes & Mota 2004) and therefore the number density of halos at $`z=0`$ is different if we assume the same $`\sigma _8`$ for all models. The embedded panels in Fig. 2 show this situation (zoomed near $`z=0`$), where $`\sigma _8`$ was set equal to 0.9 for all models. Larger differences are found in the case of inhomogeneous dark energy, because this is where $`\delta _c(z=0)`$ presents larger deviations from its value in a $`\mathrm{\Lambda }`$model (see Nunes & Mota 2004).
In what follows we assume the halo number densities of models normalized to the present-day halo abundance of the $`\mathrm{\Lambda }`$-model. In Section 4.3 we discuss the effects of this assumption on our results.
## 4 Predicted number counts
In this paper we investigate the modifications caused by a dark energy component on the number of dark matter halos. We test for one model with a cosmological constant ($`w=1`$) and for four quintessence models (defined in Section 2). The computations are done in the case where dark energy is homogeneous and in the case where it may cluster, i.e. inhomogeneous dark energy. Moreover, we choose to explore the effects on the integrated number of dark matter halos in mass bins \[$`M_{\mathrm{inf}},M_{\mathrm{sup}}`$\] illustrating different classes of cosmological structures, namely $`10^{13}10^{14}`$, $`10^{14}10^{15}`$, and $`10^{15}10^{16}`$ in units of $`h^1M_{}`$.
We study the effect of dark energy on the number of dark matter halos by computing two quantities. The first is the all sky number of halos per unit of redshift, in the mass bin
$$𝒩^{\mathrm{bin}}\frac{dN}{dz}=_{4\pi }𝑑\mathrm{\Omega }_{M_{\mathrm{inf}}}^{M_{\mathrm{sup}}}\frac{dn}{dM}\frac{dV}{dzd\mathrm{\Omega }}𝑑M,$$
(15)
where the volume element is given by $`dV/dzd\mathrm{\Omega }=r^2(z)/H(z)`$, with $`r(z)=_0^zH^1(x)𝑑x`$. The redshift evolution of $`dV/dzd\mathrm{\Omega }`$ for different models of dark energy is depicted in Fig. 3. Note that the volume element does not depend on the hypothesis of dark energy clustering, it is thus the same for both homogeneous and inhomogeneous dark energy models. Note also that the phantom model is the only one having a volume element larger than the volume element in the cosmological constant model.
The second quantity we compute is the all sky integrated number counts above a given mass threshold, $`M_{\mathrm{inf}}`$, and up to redshift $`z`$:
$`N(z,M>M_{\mathrm{inf}})={\displaystyle _{4\pi }}𝑑\mathrm{\Omega }{\displaystyle _{M_{\mathrm{inf}}}^{\mathrm{}}}{\displaystyle _0^z}{\displaystyle \frac{dn}{dM}}{\displaystyle \frac{dV}{dz^{}d\mathrm{\Omega }}}𝑑M𝑑z^{}.`$
Our knowledge of both these quantities for galaxy clusters will improve enormously with upcoming (underway or planned) cluster surveys operating at different wavebands. These include the Planck Surveyor satellite and South Pole Telescope (SPT) (Ruhl *et al.* (2004)) Sunyaev-Zel’dovich surveys, the XMM-Newton serendipitous X-ray cluster survey (XCS) (Romer *et al.* (2001)) and the recently proposed deep multiband optical Dark Energy Survey (DES) designed to probe almost the same sky region of SPT.
### 4.1 Number counts in mass bins
Figures 4, 5 and 6 show the number counts, $`𝒩^{\mathrm{bin}}=dN/dz`$, obtained from Eq. (15) (top panels) together with the difference of counts of the various dark energy models to the $`\mathrm{\Lambda }`$-model, $`\mathrm{\Delta }𝒩^{\mathrm{bin}}=𝒩^{\mathrm{bin}}𝒩_\mathrm{\Lambda }^{\mathrm{bin}}`$, (bottom panels). First let us concentrate on the homogeneous dark energy case, i.e. on the left panels of these figures. Below redshift unity the volume element has the most important role in the integral of Eq. (15), as it increases by orders of magnitude. Above this redshift the volume element does not vary much and it is the mass function that decreases by orders of magnitude. It decreases faster for models of larger $`\delta _c/\sigma _8g`$. Therefore, for our particular models, we expect to see at low redshifts, in decreasing order of number of counts per redshift, the following sequence of models: $`w=1.2`$, $`w=1`$, $`2\mathrm{E}\mathrm{X}\mathrm{P}_2`$, $`2\mathrm{E}\mathrm{X}\mathrm{P}_1`$ and $`w=0.8`$ and the reverse order at high redshifts. In practice this is only true if the maximum of counts occurs at a redshift around or greater than unity, i.e after the volume element does not vary much. We can verify that this requirement is only satisfied in the lowest mass bin $`10^{13}<M/(h^1M_{})<10^{14}`$, hence the low redshift order might differ from the one stated in the two largest mass bins. Nonetheless, the high redshift description is accurate for all mass bins. More generally, the differences between counts can be understood in terms of the relative contribution to $`𝒩^{\mathrm{bin}}`$ of the integrand in Eq. (15). The counts thus depend not only on the volume element and on the growth factor but also on the small differences in $`\sigma _8`$. At high redshift, models with smaller $`\delta _c/\sigma _8g`$ ratios imply larger halo abundances and this effect dominates that of the volume element. At low redshifts the differences between counts from one model to another depend also on the masses of interest. As a matter of fact, the number of low mass structures is mostly sensitive to the volume element whereas for massive structures the number counts become sensitive to the normalisation $`\sigma _8`$. The left lower panels of the figures illustrate quantitatively the relative evolution of the number counts. For $`w=1.2`$, the figures depict the expected excesses and deficits of counts compared to the cosmological constant model below and above redshift $`z_\mathrm{t}`$ (defined as $`𝒩^{\mathrm{bin}}(z_\mathrm{t})=𝒩_\mathrm{\Lambda }^{\mathrm{bin}}(z_\mathrm{t})`$). Conversely, we have deficits and excesses below and above $`z_\mathrm{t}`$, respectively, for all the other models. It is also worth pointing out that larger mass bins imply a larger ratio $`\delta _c/\sigma `$ (see Eq. (13)), which makes the mass function to dominate at lower redshifts in Eq. (15) and consequently $`z_\mathrm{t}`$ to move to smaller values as depicted in the figures.
Now, we turn to the comparison between the inhomogeneous and homogeneous hypotheses (right versus left panels in Figs. 4, 5 and 6). Figure 7 is helpful to understand how the differences arise. Indeed, when compared to the homogeneous case, we verify that the quantity $`\delta _c/\sigma _8g`$ is smaller in the inhomogeneous case for $`w=0.8`$ and larger for all the other models. Given this, it is clear that for $`w=0.8`$ we see, as expected, smaller deficits and larger excesses of counts comparatively to the homogeneous case and for all the other models the inverse trend. The effects of the inhomogeneous hypothesis are more evident for high mass bins where count deficits and excesses show larger differences between homogeneous and inhomogeneous dark energy. For example, in Fig. 6 we see that the $`w=0.8`$ model has $`2`$ times larger count excess in the inhomogeneous case.
### 4.2 Integrated number counts
The integrated number of collapsed structures above a given mass (Eq. (4)) is an important observable quantity. In this section we present results for the integrated number counts of structures with masses above $`M_{\mathrm{inf}}=10^{13}h^1M_{}`$ and $`10^{14}h^1M_{}`$. These are displayed in the upper panels of Figs. 8 and 9, respectively (here we omit displaying results for $`M>M_{\mathrm{inf}}=10^{15}h^1M_{}`$ because, as it will be clear below, integrated counts in this mass range can be directly estimated from the curves in Fig. 6). We also calculate the difference in integrated number counts with respect to the cosmological constant model shown in the lower panels of the same figures. As in the previous section, we compare the predicted numbers in the homogeneous (left panels of figures) and inhomogeneous hypotheses (right panels) for the cosmological models we consider in this study.
For each mass bin, the differences between integrated counts result from the combination of effects (that act on different mass ranges) discussed in the previous section. In our work we have obviously not performed the integration in the mass range in Eq. (4) all the way up to infinity but only to $`M_{\mathrm{sup}}=10^{16}h^1M_{}`$. In the cases $`M>M_{\mathrm{inf}}=10^{13}h^1M_{}`$ and $`M>M_{\mathrm{inf}}=10^{14}h^1M_{}`$, integrated number counts reflect mainly the behaviour of curves in the lowest and middle mass bin, respectively. This is because $`N(z,M>M_{\mathrm{inf}})`$ is dominated by the contribution of the lower bound of the mass integration range. Hence, it is legitimate to concentrate on the lowest mass when a qualitative description of the integrated number counts is concerned.
We have seen in the Figs. 4, 5 and 6 that in the homogeneous case, the model with $`w=1.2`$ presents excesses with respect to the cosmological constant at low redshift and deficits at high redshift. All the other models give the opposite behaviour as their $`\delta _c/\sigma _8g`$ is always lower than that of the $`w=1`$ model. Therefore, for $`w=1.2`$ we expect the integrated number counts to be larger than for the cosmological constant model until the redshift $`z_\mathrm{t}`$. The remaining models must show the opposite behaviour. The difference between the integrated number counts of a model with respect to the cosmological constant must decrease with redshift for $`z>z_\mathrm{t}`$. Eventually above a redshift $`z_{\mathrm{flat}}`$, the integrated number counts should become constant with redshift because, as we have noted before, the number counts (Eq. (4)) decrease exponentially with redshift hence contributions above $`z_{\mathrm{flat}}`$ become negligible. Note that $`z_{\mathrm{flat}}`$ becomes progressively smaller for structures with larger mass limits ($`M_{\mathrm{inf}}`$) because a smaller number of these objects form at higher redshifts. This simply reflects the hierarchical nature of structure formation in models described by Eqs. (LABEL:eq:mf)-(14). Models give quite different integrated count differences ($`NN_\mathrm{\Lambda }`$) depending on the maximum redshift of integration. For homogeneous dark energy, maximum deviations from the $`\mathrm{\Lambda }`$-model are generally obtained near $`z_\mathrm{t}`$ of the dominant class of structures, whereas for inhomogeneous dark energy maximum deviations generally occur at much higher redshifts, $`zz_{\mathrm{flat}}`$.
In the inhomogeneous scenario, the interpretation of the redshift dependence of the integrated number counts is similar to that of the homogeneous case. As we have seen, the inhomogeneous case yields higher excesses and lower deficits for the $`w=0.8`$ model, hence we expect at high redshift, a positive difference in integrated counts compared to the cosmological constant model. This is quite visible in the right lower panels of Figs. 8 and 9. Through a similar argument, we expect at high redshift, a negative difference of integrated number counts with respect to the cosmological constant for the $`w=1.2`$ model. Moreover, because there are hardly any excesses in the 2EXP models, we must expect the difference of integrated number counts to flatten out at $`z_\mathrm{t}`$ at large negative values for these two models.
### 4.3 Normalisation
So far in this work we have normalised the various dark energy models such that they reproduce the abundance of dark matter halos at redshift zero (i.e. the local abundance) of a $`\mathrm{\Lambda }`$-model with $`\sigma _8=0.9`$. However, observations that are being used to quantify the number of observed structures trace baryonic (emitting) gas and not dark matter halos directly. For systems where non-gravitational physics is important, i.e. the less massive structures, this may lead to an important miss-match between dark matter halos and emitting structure counts. Therefore, we dedicate the rest of this section to discuss the implications on our results of dropping the constraint of normalizing models to the same abundance at $`z=0`$. Instead, we consider that all models have the same $`\sigma _8=0.9`$ normalization, as given by present day observations (e.g. Spergel et al. 2003).
We illustrate the effects of inhomogeneous dark energy on one single mass bin. Figure 10 shows number counts in the mass bin $`10^{14}<M/(h^1M_{})<10^{15}`$ and Fig. 11 the integrated number counts above mass $`M>10^{14}h^1M_{}`$. Both figures consider inhomogeneous dark energy models. In the left panels of these figures we are assuming that models are normalized to the same halo abundance at $`z=0`$ (i.e. same curves as those in the right panels of Figs. 5 and 9), whereas right panels assume that all models have a fixed normalisation, $`\sigma _8=0.9`$. Models in the homogeneous hypothesis have practically the same $`\sigma _8`$ when they are normalized to the halo abundance at $`z=0`$. Therefore we do not expect to observe much difference in the halo abundances from one dark energy model to another. The situation is quite different when dark energy is inhomogeneous. The comparison between panels in Fig. 10 indicates that fixing $`\sigma _8`$ in all models causes much larger departures from the $`\mathrm{\Lambda }`$-model than in the case where models are normalized to reproduce the same local halo abundances. At the maximum of $`𝒩^{\mathrm{bin}}`$, the differences between dark energy models come from the fact that for a fixed $`\sigma _8`$, the mass function reflects the variations of $`\delta _c/g`$ with $`z`$, which effects dominate those of the volume element in Eq. (15) for this mass bin. It is interesting to note that the structure of the curves can change dramatically depending weather we fix the local abundance or $`\sigma _8`$.
As this work was being completed Manera and Mota (2005) made public an analysis similar to the one presented here for the particular scenario of coupled quintessence. In their work, they have fixed $`\sigma _8`$ for the two models under study rather than the local halo abundance. We have verified that indeed, also in the coupled quintessence model, the departures from a cosmological constant model are larger if one takes the former approach.
## 5 Conclusions
Number counts of collapsed structures are commonly proposed as a tool to probe dark energy models. In the simplest case of a constant equation of state, galaxy cluster number counts may constrain the dark energy, provided we have a good knowledge of the cluster physics and their redshifts (see for example Majumdar and Mohr (2004); Wang *et al.* (2004)). In the present study we have investigated, using the Press-Shechter mass function, two complications to the general picture: First, we have revisited the modifications introduced by a varying equation of state or by including a phantom component; second we have explored the effect of an inhomogeneous dark energy which collapses with the dark matter.
More specifically, we have considered a scalar potential built out of two exponential terms for which two sets of parameters (models 2EXP<sub>1</sub> and 2EXP<sub>2</sub>) were explored. Although their equations of state at present are almost indistinguishable, they undergo quite different evolutions at higher redshifts and generally give different results. In our analysis we have also included a phantom dark energy model ($`w=1.2`$), compatible with present observations. An interesting feature of phantom dark energy is that it induces opposite departures from the $`\mathrm{\Lambda }`$-model as compared with the other models considered in this paper. That is, we expect an excess of sources in a phantom energy model when other models predict a deficit.
Recently Mota and van de Bruck (2004) proposed that dark energy may cluster in forming structures (inhomogeneous dark energy). They have investigated the growth and collapse of cosmological structures under the inhomogeneous hypothesis. Going a step forward, we have investigated in the present study the implications of this hypothesis on the number density of collapsed objects. We have found that inhomogeneous dark energy generally enhances departures from the $`\mathrm{\Lambda }`$-model. This includes the models with a time varying equation of state, which can present several times larger departures (from the $`\mathrm{\Lambda }`$-model) as compared to the homogeneous case. Yet our results indicate that the inhomogeneous dark energy hypothesis causes maximum deviations no larger than $`15`$% in mass bins with comfortably large numbers of collapsed halos. Another interesting feature is that maximum departures from the $`\mathrm{\Lambda }`$-model are generally obtained at higher redshift for inhomogeneous dark energy than for the homogeneous case, which generally show maximum departures near the maximum of $`𝒩^{\mathrm{bin}}(z)`$. This may be a helpful feature to test for the inhomogeneous hypothesis. Larger departures from the $`\mathrm{\Lambda }`$-model are also stronger for the more massive structures, but these are quite rare objects, which makes it difficult to statistically distinguish between models. Our analysis reveals that the inhomogeneous dark energy hypothesis has the greatest impact on the 2EXP1 model.
In this work, we have assumed that the matter transfer function remains unchanged at cluster scales. We have further assumed that models are normalized to reproduce the same abundance of dark mater halos at redshift zero. In the homogeneous hypothesis all models have practically the same $`\sigma _8`$ and there are not much differences in the halo abundances from one dark energy model to another. When dark energy is inhomogeneous, $`\sigma _8`$ differ by a few percent and the departures from the $`\mathrm{\Lambda }`$-model are much larger. It is, however, worth noting that the gas physics which rules the observed quantities adds a degree of degeneracy. We have evaluated the effects of alternatively fixing $`\sigma _8`$ to a specific value regardless of the dark energy model. In this case we verified that the departures from a $`\mathrm{\Lambda }`$ cosmology are further enhanced.
Our results show that constraining dark energy models from structure counts is complicated when models have time varying equation of state. It becomes an even more complicated task when the possibility of inhomogeneous dark energy is taken into account. Therefore in order to constrain dark energy models, we need to explore as many observable quantities as possible. Our results suggest that besides redshift distribution of structures, considering structures in mass ranges significantly increases the number of observables. Indeed, each theoretical model provides specific predictions for the redshift evolution of number counts and integrated count differences in different mass bins. The comparison of such quantities with observations can be used for testing models against observational data. It may further allow to distinguish between homogeneous and inhomogeneous dark energy models. However, this requires good knowledge of the gas physics, redshifts of observed structures and, more precisely, a good understanding of the selection function of the observations.
###### Acknowledgements.
We thank David Mota and Morgan LeDelliou for invaluable discussions. NJN is supported by the Department of Energy under contract DE-FG02-94ER40823 at the University of Minnesota. AdS acknowledges support by CMBnet EU TMR network and Fundação para a Ciência e Tecnologia under contract SFHR/BPD/20583/2004, Portugal. This work was partly supported by CNES. We would like to thank the referee for his/her comments.
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# ISO’s Contribution to the Study of Clusters of Galaxies
## 1 Introduction
### 1.1 ISO looks deep
The most strongly star forming galaxies are heavily dust obscured, and estimates of their star formation rates (SFR) made at visual wavelengths often fall one or two orders of magnitude below their true values. Frequently, very actively star forming galaxies occur in associations or groups. At the same time it has been believed that in the dense environments of galaxy clusters the interactions of galaxies with each other, with the cluster tidal field, and with the intra-cluster medium (via ram pressure) strip galaxies of their reserves of gas, and eventually suppress star formation.
Much has already been learned about the evolution of the cosmic star formation rate in field galaxies from observations with the European Space Agency’s Infrared Space Observatory (ISO) satellite (Kessler et al. 1996). Thanks to deep surveys in the mid-infrared (MIR) and far-infrared (FIR) (e.g. Elbaz et al. 1999, 2002; Serjeant et al. 2000, 2004; Gruppioni et al. 2002; Lari et al. 2001; Metcalfe et al. 2003; Sato et al. 2003; Rodighiero et al. 2003; Kawara et al. 2004; Rowan-Robinson et al. 2004; Fadda et al. 2004, and several others) conducted, respectively, with ISOCAM<sup>1</sup><sup>1</sup>1Throughout this paper the ISOCAM filters having reference wavelengths 6.75 and 14.3 $`\mu `$m will respectively be referred to as the 7 $`\mu `$m and 15 $`\mu `$m filters. Both bandpasses are referred to as mid-infrared (MIR). (Cesarsky et al. 1996) and ISOPHOT (Lemke et al. 1996), we now know that the comoving density of IR-bright galaxies has a very rapid evolution from $`z0`$ to $`z1`$. This evolution has been interpreted as the result of an increased rate of galaxy-galaxy interactions, coupled with an increase in the gas content of the galaxies (Franceschini et al. 2001). Dust obscuration plays an important role in concealing star formation in field galaxies. Its role in relation to the star formation activity of cluster galaxies is less known. Published ISO results to–date addressing the fields of galaxy clusters (Duc et al. 2002, 2004; Coia et al. 2004a & b; Biviano et al. 2004; Metcalfe et al. 2003; Fadda et al. 2000; Altieri et al. 1999; Soucail et al. 1999; Barvainis et al. 1999; Quillen et al. 1999; Lémonon et al. 1998; Pierre et al. 1996.) concern a dozen clusters at $`z<0.6`$, and point to an important role for dust, evolving with redshift, also in cluster galaxies. The fraction of IR-bright, star-forming cluster galaxies changes significantly from cluster to cluster, without a straightforward correlation with either the redshift, or the main cluster properties (such as the cluster mass and luminosity).
The presence of dust in cluster galaxies may have hampered our efforts to fully understand the evolution of galaxies in clusters. Optical-band observations have provided a few, but very fundamental results, which cannot yet however be combined in a unique, well constrained scenario of cluster galaxy evolution.
### 1.2 Some properties of cluster galaxies
Perhaps the most fundamental phenomenology that all scenarios of cluster galaxy evolution (e.g. Dressler 2004) must explain is the so called morphology density relation (hereafter MDR; Dressler 1980), whereby early-type galaxies, i.e. ellipticals and S0s, dominate rich clusters, while late-type galaxies, i.e. spirals and irregulars, are more common in the field. The fact that early-type galaxies seem to reside in the cluster centres since $`z\mathrm{\hspace{0.17em}1}`$–2, while the colour-magnitude relation (CMR, Visvanathan & Sandage 1977) remains very tight even at $`z\mathrm{\hspace{0.17em}1}`$ suggests that the MDR is established at the formation of galaxy clusters and that the early-type galaxies defining the CMR are uniformly old and passively evolving since their formation redshift, $`z_f>2`$ (e.g. Ellis et al. 1997). Similar conclusions are obtained by analysing the fundamental plane (FP; Dressler et al. 1987), relating basic properties of early-type galaxies (their effective radius, internal velocity dispersion, and effective surface brightness). The FP, like the CMR, still holds for early-type galaxies in $`z\mathrm{\hspace{0.17em}1}`$ clusters, and its scatter is similar to that seen in nearby clusters (van Dokkum & Stanford 2003).
However, these conclusions could only be true on average. Independent analyses suggest that at least part of the cluster galaxies have undergone significant evolution over the last 3–8 Gyr. First and foremost is the observational evidence for an increasing fraction of blue cluster galaxies with redshift, the so called ‘Butcher-Oemler’ (BO) effect (Butcher & Oemler 1978, 1984; Margoniner et al. 2001). Approximately 80% of galaxies in the cores of nearby clusters are ellipticals or S0s, i.e. red galaxies (Dressler 1980), but the fraction of blue galaxies increases with redshift. These blue galaxies are typically disk systems with ongoing star formation (Lavery & Henry 1988), with spectra characterized by strong Balmer lines in absorption (typically, EW(H$`\delta `$) $`>3`$ Å) and no emission lines, and have been named ‘E+A’ (or also ‘k+a’) galaxies (Dressler & Gunn 1983). Modelling of their spectra indicates that star formation stopped typically between 0.05 and 1.5 Gyr before the epoch of observation, in some cases after a starburst event (Poggianti et al. 1999, 2001).
Similarly, the fraction of spirals increases with redshift (Dressler et al. 1997; Fasano et al. 2000; van Dokkum et al. 2001), at the expense of the fraction of S0’s. Maybe, the colour and the morphological evolution of the cluster galaxy population are two aspects of the same phenomenon. Field spirals are being accreted by clusters (Tully & Shaya 1984; Biviano & Katgert 2004), and the accretion rate was higher in the past (Ellingson et al 2001). It is therefore tempting to identify the blue galaxies responsible for the BO-effect in distant clusters with the recently accreted field spirals.
Possibly, these evolutionary trends could be reconciled with the passive evolution inferred from the CMR and FP studies by taking into account the so-called ‘progenitor bias’ (van Dokkum et al. 2000), namely the fact that only the most evolved among cluster early-type galaxies are indeed selected in studies of the colour-magnitude and fundamental plane relations.
### 1.3 A hostile environment
There is no shortage of plausible physical mechanisms that could drive the evolution of a galaxy in a hostile cluster environment. Among these, the most popular today are ram pressure, collisions, and starvation from tidal stripping (Dressler 2004), all in principle capable of depleting a spiral of its gas reservoir, thereby making it redder and more similar to a local S0.
Ram pressure from the dense intra-cluster medium can sweep cold gas out of the galaxy stellar disk (Gunn & Gott 1972) and induce star formation via compression of the gas that remains bound to the galaxy. Collisions or close encounters between galaxies generate tidal forces that tend to funnel gas towards the galaxy centre (Barnes & Hernquist 1991) eventually fueling a starburst that ejects gas from the galaxy. The cumulative effect of many minor collisions (named ‘harassment’, Moore et al. 1996), can lead to the total disruption of low surface brightness galaxies (Martin 1999). The collision of a group with a cluster can also trigger starbursts in cluster galaxies, as a consequence of the rapidly varying tidal field (Bekki 2001). Finally, the so called ‘starvation’ mechanism (Larson et al. 1980) affects the properties of a galaxy by simply cutting off its gaseous halo reservoir. This can occur because of tidal stripping, a mechanism effective in galaxy-galaxy encounters, but also when galaxies pass through the deep gravitational potential well of their cluster.
The common outcome of all these processes is galaxy gas depletion, ultimately leading to a decrease of the star formation activity for lack of fuel, and, hence, to a reddening of the galaxy stellar population. However, some of these processes induce a starburst phase before the gas depletion, and some do not.
To date, it remains unclear which physical process dominates in the cluster environment. Useful constraints can be obtained by finding where the properties of cluster galaxies change with respect to the field, since the different processes become effective at different galaxy or gas densities. Recently it has been found (Kodama et al. 2001; Gómez et al. 2003; Lewis et al. 2002) that a major change in the star-formation properties of cluster galaxies occurs in the outskirts of clusters (at $`\mathrm{\hspace{0.17em}1.5}`$ cluster virial radii). These results would seem to exclude ram-pressure stripping as a major factor in cluster galaxy evolution, since the density of the intra-cluster medium is too low in the cluster outskirts.
Recent observations of high-$`z`$ clusters seem to have complicated, rather than simplified, the issue of cluster galaxy evolution. A surprisingly high fraction of red merger systems has been found in these distant clusters (van Dokkum et al. 2001). The red colours of these merging galaxies and the lack of emission lines in their spectra suggest that their stellar populations were formed well before the merger events, but the occurrence of relatively recent starburst events in these galaxies is instead suggested by detailed analyses of their spectra (Rosati 2004).
### 1.4 An uncluttered view
A better understanding of the evolutionary processes affecting cluster galaxies can come from observations in the infrared. Dust, if present, is capable of obscuring most of a galaxy’s stellar radiation, making the observed galaxy red and dim at optical wavelengths, and affecting optical estimates of the galaxy star formation activity. The effects of dust can be particularly severe if the galaxy is undergoing a starburst (Silva et al. 1998), so that we might be missing a substantial part of the evolutionary history of cluster galaxies by observing them at optical wavelengths. Since the dust-reprocessed stellar radiation is re-emitted at IR wavelengths, the IR luminosity is a much more reliable indicator of a galaxy’s star formation activity (Elbaz et al. 2002).
The plan of this review is as follows: in Section 2 the development of knowledge of the infrared properties of galaxy clusters from early IRAS observations, through the “ISO-era” to the present is described. Section 2.1 considers the accumulation of data on the Virgo cluster, while Section 2.2 addresses other nearby clusters, e.g. the Fornax, Hydra, Coma and Hercules clusters. Section 2.3 discusses the significant progress that has been possible with ISO in the study of cluster galaxy properties out to moderate redshifts ($`<`$ 0.6) and attempts to draw some comparison among the still rather heterogeneous sample of cluster observations. Section 2.4 reviews the status of attempts to directly observe diffuse intra-cluster dust in the infrared. Finally, Section 3 summarises the current status of the field and remarks on the important opportunity represented by the Spitzer Observatory to decisively extend the field.
## 2 Cluster galaxies in the infrared
### 2.1 The Virgo galaxy cluster
Being the most nearby relatively massive galaxy cluster, the Virgo cluster has been studied extensively at all wavelengths. Already in 1983, Scoville et al. obtained 10 $`\mu `$m data with IRTF for 53 Virgo spiral galaxies and concluded that star formation is occurring in the nuclei of virtually all spiral galaxies independent of spectral type, with an average star formation rate (SFR) of 0.1 M/yr. No correlation was found between 10 $`\mu `$m emission and gross properties (barred or normal, early or late spiral morphology, total optical luminosity), nor with location in the cluster. Using far- to mid-infrared flux correlations, they inferred a far-infrared (FIR) luminosity $`2\times \mathrm{\hspace{0.17em}10}^{10}L_{}`$ for the brightest 10 $`\mu `$m source in the Virgo cluster, NGC4388.
Leggett et al. (1987) made optical identifications of 145 IRAS galaxies in the 113 square degree field centred on the Virgo cluster and concluded that they were mostly seeing the spirals, and that the infrared properties of the Virgo cluster galaxies are indistinguishable from those of field galaxies at similar redshift. Virgo spirals were indistinguishable from field disc galaxies with normal SFRs. The typical infrared luminosity for the sample was $`L_{IR}\mathrm{\hspace{0.17em}10}^9L_{}`$, and the most luminous confirmed cluster sources were found to have $`L_{IR}=`$ a few $`\times \mathrm{\hspace{0.17em}10}^{10}L_{}`$. They concluded that the cluster environment has no effect on galaxy IR properties, even for very HI deficient galaxies. Their conclusion was however criticized by Doyon & Joseph (1989) who, using a sample of 102 Virgo spirals detected at 60 and 100 $`\mu `$m, were able to show that HI-deficient galaxies have lower IR fluxes, lower star formation activity, and cooler IR colour temperatures than those with normal HI content.
If cluster environment affects the interstellar medium (ISM) of cluster galaxies, one might expect the FIR, radio and FIR-radio correlation to be affected. Niklas et al. (1995) looked at the FIR-radio correlation in Virgo galaxies and suggested that, in this respect, most of the Virgo galaxies behave like normal field galaxies. Only a few early type spirals with strong central sources as well as very disturbed galaxies show a high radio excess. These are in the inner part of the cluster, where galaxies have disturbed HI-distributions, and truncated HI-disks, as shown by 21 cm observations (Cayatte et al. 1990). Measurements of molecular gas emission in radio continuum (Kenney & Young 1986, 1989) show a much less pronounced stripping effect.
In the ISO era, Tuffs et al. (2002) and Popescu et al. (2002a, b) used ISOPHOT at 60, 100 and 170 $`\mu `$m, to study a luminosity- and volume-limited sample of 63 S0a or later-type galaxies in the core and periphery of the Virgo cluster. They reached sensitivities 10 times better than IRAS in the two shorter-wavelength bandpasses, and the confusion limit at 170 $`\mu `$m. These programmes sought to extend knowledge of FIR SEDs to lower limits covering a complete sample of normal<sup>2</sup><sup>2</sup>2A ‘Normal’ galaxy is understood to mean a galaxy not dominated by an active nucleus, with SFR sustainable for a substantial fraction of a Hubble time. late-type galaxies over a range of morphological type and star formation activity.
A significant cold dust component (with a temperature of around 18 K) was found in all morphological classes of late-type galaxies, from early giant spirals to irregular galaxies and Blue Compact Dwarfs (BCDs), and which could not have been recognized by IRAS. These results required a revision of the masses and temperatures of dust in galaxies. On average, dust masses are raised by factors of 6 to 13 with respect to IRAS results. The FIR/radio correlation is confirmed for the warm FIR dust, and is found for the cold dust. The predominance of the very cold dust component (down to less than 10 K) in BCDs was remarkable, with a few 10s of percent of the UV/Optical component appearing in the cold dust emission. The cold emission might be due to collisional dust heating in dust swept up in the intergalactic medium (proto-galactic cloud) by a galactic wind from the BCD or, alternatively, to photon heating of the dust particles in an optically thick disk, indicative of a massive gas/dust accretion phase that makes BCDs sporadically bright optical/UV sources when viewed out of the disk equatorial plane.
On average, 30% of the stellar light of spirals in Virgo is re-radiated by dust, with a strong dependence on morphological type, ranging from 15% for early spirals to 50% for some late spirals, and even more for some BCDs (Popescu & Tuffs 2002). Fig. 1, taken from Popescu & Tuffs (2002), shows the dependency of the dust contribution on the Hubble type in the Virgo cluster, and it illustrates a sequence of increasing FIR-to-total bolometric output running from normal to gas-rich-dwarf galaxies.
Leech et al. (1999) studied a sample of 19 Virgo cluster spiral galaxies with ISO’s Long Wavelength Spectrometer (LWS) (Clegg et al. 1996) obtaining spectra around the \[CII\] 157.7 $`\mu `$m fine-structure line for 14 of them. \[CII\] line radiation provides the most important gas cooling mechanism in normal, i.e. non-starburst, late-type galaxies, balancing photoelectric heating from grains. In field galaxies \[CII\] is typically 10<sup>-3</sup> to 10<sup>-2</sup> of total galaxy FIR. The sample, drawn from both the cluster core and the cluster periphery (and being a sub-sample of the Tuffs et al. 2002 sample), spanned the S0a–Sc morphological range to probe any difference in the \[CII\] emission between different types or between core and periphery galaxies.
A good correlation was found between the strength of the \[CII\] line and the FIR flux, as measured by IRAS. Moreover, the \[CII\]-to-$`K^{}`$-band flux ratio shows a two order of magnitude difference between early-type galaxies and late spiral types. Galaxies with large \[CII\]/FIR ratios tend to have later Hubble types. No apparent relationship was found between \[CII\] strength and galaxy position in the cluster, nor between \[CII\] and HI-mass surface density. Any influence of the Virgo cluster environment on the \[CII\] emission was found to be small compared with the strong dependence of the line emission on basic measurables such as morphology or bulk mass of the stellar component, as measured by the near-IR ($`K^{}`$-band) luminosity.
In a series of papers, Boselli et al. (1997a,b, 1998, 2003a,b, 2004) explored the properties of a large sample of ISO-detected spiral and irregular galaxies in the Virgo cluster, in measurements made with ISOCAM at 7 and 15 $`\mu `$m. 71 objects were in the cluster periphery ($``$ 4 degrees from M87) and 28 in the core ($``$ 2 degrees from M87) in order to allow study of the effects of the environment on dust emission and evolution. A further 24 Virgo cluster galaxies were serendipitously observed, bringing the full Virgo sample up to 123 galaxies. S0 and elliptical galaxies observed by chance in the ISO fields were used to establish the stellar contribution to the IR emission. Objects had luminosities in the range $`10^{7.4}L_{FIR}\mathrm{\hspace{0.17em}10}^{10.1}L_{}`$.
Thirty four of the Virgo objects were fully resolved by ISOCAM, and MIR images could be presented for these, along with radial light and colour profiles and other morphological and structural information.
Boselli et al. showed that the MIR emission of optically-selected, normal early-type galaxies is dominated by the Rayleigh-Jeans tail of the cold stellar component, while that of late-type galaxies is dominated by the thermal emission from dust, but partly contaminated by stellar emission, especially at 7 $`\mu `$m in early-type spirals (Sa). While the MIR emission (per unit mass) of the spirals and later systems are comparable, the average 7 $`\mu `$m and 15 $`\mu `$m to $`K^{}`$ flux ratio of E, S0 and S0/a galaxies is significantly lower than that of spirals. At 15 $`\mu `$m, where the difference is clearer, spirals have on average a MIR emission per unit mass higher by more than one order of magnitude than E-S0/a. BCDs have, on average, a MIR emission per unit mass comparable to that of spirals.
The IR emission carriers and their behaviour are consistent with quantitative expectations based on MIR studies of the ISM in our Galaxy. In spiral and irregular galaxies the 7 $`\mu `$m emission is almost entirely due to the UIB<sup>3</sup><sup>3</sup>3The Unidentified Infrared Bands (UIB) dominate the 5 to 12 micron MIR spectrum of a wide range of celestial sources, and are usually assumed to arise from polycyclic aromatic hydrocarbon molecules (PAHs). carriers. The MIR fluxes are proportional to the SFR when it is not too large, but fall off in the presence of a high SFR suggesting that the UIBs are destroyed by the UV field. However, in galaxies with a high SFR, there is an additional diffuse contribution (i.e. not well correlated with HII regions or H<sub>α</sub> sources) to the 15 $`\mu `$m flux from very small, three-dimensional, grains. As a consequence, MIR dust emission is not an optimal tracer of star formation in normal, late-type galaxies, and MIR luminosities are better correlated with FIR luminosities than with more direct tracers of the young stellar population such as the H$`\alpha `$ and the UV luminosity. This conclusion, valid for nearby normal, late-type galaxies, may not apply to luminous starburst galaxies (Luminous Infrared Galaxies, or LIRGs), such as those detected in the ISO deep surveys (Förster-Schreiber et al. 2004). The MIR emission traces well the FIR and bolometric emission (Boselli et al. 1998; Elbaz et al. 2002).
### 2.2 The Coma, Fornax, Hydra, Hercules & other nearby ($`z<\mathrm{\hspace{0.17em}0.1}`$) clusters
Located at much larger distance than Virgo, but with a much larger mass, Coma has often been a favourite observational target, also in the IR.
Wang et al. (1991) made optical identifications of a total of 231 IRAS point sources in the regions of the Fornax, Hydra and Coma clusters, and identified respectively 13, 29 and 26 cluster galaxies. They concluded that the cluster environment has no detectable influence on galaxy infrared properties. Bicay & Giovanelli (1987) came to the same conclusion, based on a sample of 200 FIR-emitting galaxies in seven nearby clusters (including Coma). However, they remarked that LIRGs were much less common in the clusters than in the field (see also Section 2.3).
The most distant cluster studied with IRAS was A2151 (aka Hercules, at $`z=0.036`$). Out of 41 sources detected at 60$`\mu `$m in a 1.6$`\times `$2.5 sq.deg. field, Young et al. (1984) correlated 24 with late-type spiral galaxies of the cluster remarking the total absence of IR emission from E and S0 galaxies. Odenwald (1986) found CO emission in 3 of the 9 most optically luminous galaxies in the spiral-rich Ursa Major I(S) galaxy group, seeking evidence for galaxy interactions in a cluster lacking an intra-cluster medium, but found little evidence that processes unrelated to the galaxies themselves had influenced their histories. Studying a sample of 200 galaxies in seven nearby clusters<sup>4</sup><sup>4</sup>4A262, Cancer, A1367, A1656 (Coma), A2147, A2151 (Hercules), and Pegasus, Bicay & Giovanelli (1987) pointed out the absence of luminous IR galaxies (LIRGs, i.e. galaxies with $`L_{FIR}>\mathrm{\hspace{0.17em}10}^{11}L_{}`$). The sample consisted almost entirely of IR normal galaxies ($`L_{FIR}<\mathrm{\hspace{0.17em}10}^{10}L_{}`$) in contrast to the rather high percentage of LIRGs (20%) detected by IRAS in the field. Moreover, the lack of a strong correlation between galaxy HI content and IR emission led them to conclude that SFR is not enhanced by interaction with the ICM, and might even be quenched by it. On the contrary, the suppressed FIR-to-radio ratio of spiral galaxies found in rich clusters with respect to poor clusters (Andersen & Owen, 1995) seemed to suggest that ram pressure enhances the radio emission in rich clusters while galaxy-galaxy interactions play a more important role in poor clusters where velocity dispersion, and so encounter velocities, are smaller.
Quillen et al. (1999) observed 7 E+A galaxies plus one emission-line galaxy at 12 $`\mu `$m with ISOCAM. They found that E+A galaxies have mid- to near-IR flux ratios typical of early-type quiescent galaxies, while the emission-line galaxy had enhanced 12 $`\mu `$m emission relative to the near-IR. Galaxies with ongoing star formation have a different velocity distribution in the cluster from galaxies with stopped SF, suggesting that the ongoing infall of field spirals into the cluster potential may first trigger and then quench star formation.
Further observations of Coma cluster galaxies in the MIR came from Boselli et al. (1998) and Contursi et al. (2001). They also observed the cluster A1367, located in the Coma supercluster, detecting, in total, 18 spiral/irregular galaxies in the MIR and FIR with ISO.
Confirming results found in Virgo galaxies, these authors concluded that most IR-detected Coma galaxies display diffuse MIR emission unrelated to their H$`\alpha `$ emission. The aromatic carriers are not only excited by UV photons, but also by visible photons from the general ISM. When the UV radiation field is too intense, it can even destroy the aromatic carriers, and overall the MIR emission is dominated by photo-dissociation regions rather than HII-like regions. A cold dust component was detected in all galaxies, at temperature of $`22`$ K, more extended than the warm dust. Only a very weak trend was found between the total dust mass and the gas content of the galaxies, even if some galaxies are very HI-deficient, and there was no detection of any relation between the MIR/FIR properties and the environment.
All these results seemed to suggest very little (if any) dependence of the IR properties of galaxies on the environment. If anything, IR emission was thought to be quenched in the cluster environment. However, only nearby clusters had been studied, in which most of the galaxies are early type with little gas (and dust). With the launch of ISO it became possible for the first time to study galaxy clusters out to redshifts where a significant change in the composition of the cluster galaxy population had already been (Butcher & Oemler 1984) or was soon to be (Dressler et al. 1999) observed in the optical.
### 2.3 Galaxy clusters at intermediate redshift
ISO’s mid-infrared camera, ISOCAM, with its vastly improved sensitivity and spatial resolution with respect to IRAS, has successfully observed several galaxy clusters out to redshifts at which significant evolution might be expected to occur. At the same time, while studies of galaxies in nearby clusters are frequently done by targeting galaxies selected at other wavelengths, distant clusters can be completely surveyed due to their smaller angular size, producing an unbiased sample of infrared-emitting galaxies. Published ISO observations to date for clusters at redshifts above 0.1 (Coia et al. 2004a, b; Biviano et al. 2004; Metcalfe et al. 2003; Duc et al. 2002, 2004; Fadda et al. 2000; Barvainis et al. 1999; Altieri et al. 1999; Lémonon et al. 1998 and Pierre et al. 1996) address seven clusters (see Table. 1) spanning the redshift range $`0.17<z<\mathrm{\hspace{0.17em}0.6}`$ and yield MIR data for around 110 cluster galaxies, slightly over 40 of these seen at 15 $`\mu `$m, and the rest only in the 7 $`\mu `$m bandpass.
Almost half of the cluster galaxies detected at 15 $`\mu `$m prove to be LIRGs. (At the higher redshifts of the above sample only LIRGs fall above the sensitivity limit of the observations.) About 60% of these cluster galaxies were detected in observations originally intended to study distant field galaxies via the gravitational lensing amplification of the foreground clusters (Metcalfe et al. 2003; Barvainis et al. 1999; Altieri et al. 1999), and which, being generally very deep spatially-oversampled measurements, were able to provide insights about the lensing cluster galaxy populations (Biviano et al. 2004; Coia et al. 2004a and b). Fig. 2, taken from Coia et al. (2004b) is a V-band image of the $`z=0.39`$ galaxy cluster CL0024+1564 overlaid with contours of an ISO 15 $`\mu `$m map. The capacity of ISOCAM to detect and assign MIR counterparts unambiguously to numerous galaxies in the field is evident.
The first published ISOCAM observations of a distant cluster were those of the $`z=0.193`$ cluster A1732 by Pierre et al. (1996), which they observed at 7 and 15 $`\mu `$m over an $`8\times \mathrm{\hspace{0.17em}8}`$ arcmin<sup>2</sup> field. They found some evidence for a deficiency of spirals and star forming galaxies in the cluster, identifying only four cluster sources (at 7 $`\mu `$m, no cluster sources were detected at 15 $`\mu `$m) and 10 MIR galaxies in total in the field, most of them judged to be foreground. Nevertheless, these were the faintest MIR extragalactic sources reported up to that point and underlined the need for ultra-deep observations to detect cluster members at 15 $`\mu `$m.
Lémonon et al. (1998) reported evidence for an active star-forming region in a cooling flow (later ‘cool-core’) from 7 and 15 $`\mu `$m observations of the inner square arcminute of the well known lensing cluster A2390 ($`z=0.23`$), with an attendant SFR of as much as 80 $`M_{}yr^1`$ in the central cD. But this was later found to be compatible with non-thermal emission from a jet associated with the cD (Edge et al. 1999). The estimated cluster mass deposition rates in cooling flows have since been lowered by one or two orders of magnitude (Böhringer et al. 2002).
A1689 ($`z=0.181`$) was the first distant cluster for which detailed ISO observations were reported. Fadda et al. (2000) detected numerous cluster members (30 at 7 $`\mu `$m and 16 at 15 $`\mu `$m) within 0.5 Mpc of the cluster centre, and they found a correlation between the B-15 $`\mu `$m colour and cluster-centric distance of the galaxies. The 15 $`\mu `$m galaxies are blue outliers with respect to the colour/magnitude relation for the cluster and become brighter going from the center to the outer parts of the cluster. Coupled with the systematic excess of the distribution of the B-15 $`\mu `$m colours with respect to nearby clusters (Virgo and Coma), this suggested the existence of an IR analogue of the Butcher-Oemler effect in A1689 (see Fig. 3). A follow-up optical study of these infrared galaxies (Duc et al. 2002) showed that the morphology of the 15 $`\mu `$m sources in A1689 is generally spiral-like, with disturbances reminiscent of tidal interactions. No LIRGs were found in A1689. The highest total IR luminosity found for a cluster galaxy was $`6.2\times 10^{10}L_{}`$, corresponding to a SFR of $`11M_{}`$ per year. The median SFR for A1689, derived from 15 $`\mu `$m measurements, was $`2M_{}`$ per year, while the median found from \[OII\] (optical) measurements was only $`0.2M_{}`$ per year.
This paper revealed the importance of IR observations in the study of star-formation in clusters. About one third of the 15 $`\mu `$m sources show no sign of star formation in their optical spectra. Moreover, comparing the star-formation estimates from IR and \[OII\] (see Fig. 4), Duc et al. (2002) deduced that at least 90% of the star formation activity taking place in A1689 is obscured by dust.
The ISO gravitational lensing survey programme (Metcalfe et al. 2003), and related work, led to deep observations of the core of several distant clusters. Large numbers of cluster galaxies were detected in the fields of A2218, A2390 and CL0024+1654. So far, the cases of A2218 and CL0024+1654 have been treated in dedicated papers.
In the analysis of the galaxies in the field of A2218, a rich cluster at $`z=0.175`$, Biviano et al. (2004) found 9 cluster members at 15$`\mu `$m inside a radius of 0.4 Mpc. In contrast to the case of A1689, which is at almost identical redshift ($`z=0.181`$) and for which the median SFR is about 2$`M_{}yr^1`$ and median $`L_{IR}`$ is around $`10^{10}L_{}`$ , only one of the A2218 MIR galaxies is a blue Butcher-Oemler galaxy. The MIR luminosity of A2218 galaxies is moderate and the inferred star formation rate is typically less than 1$`M_{}yr^1`$ with a median $`L_{IR}`$ of only $`6\times \mathrm{\hspace{0.17em}10}^8L_{}`$. The absence of a MIR BO effect in A2218 might be a consequence of the small area observed, about 20.5 square arcminutes ($`r<0.4Mpc`$), and yet the area studied for A1689 was not much larger, at about 36 square arcminutes. Coia et al. (2004b) suggest that the difference between these two clusters may be traced to their having different dynamical status.
The same sort of comparison can be drawn between the clusters CL0024+1654 ($`z=0.39`$) and A370 ($`z=0.37`$). These two clusters at similar redshift and mapped in almost identical ways, exhibit very different numbers of LIRGs (Table 2), probably, according to Coia et al., because an ongoing cluster merger gives rise to enhanced star-forming activity in CL0024+1654.
As can be appreciated from Fig. 5 taken from Coia et al. (2004b), the median infrared luminosity of ISO-detected cluster galaxies in CL0024+1654 is around $`1\times \mathrm{\hspace{0.17em}10}^{11}L_{}`$. Star formation rates derived from the 15 $`\mu `$m data range from 8 to 77 $`M_{}yr^1`$, with median(mean) value of 18(30) $`M_{}yr^1`$.
Because of the different sky areas mapped with ISO for several of the clusters discussed here, and the different sensitivities achieved, it is not straightforward to compare the results for different clusters. A useful approach is to compare the number of LIRGs detected in each cluster, since these IR-bright galaxies would be expected to be seen for any of the observations considered. Such a comparison, taken from Coia et al. 2004b, is presented in Table 2. For each cluster an “expected” number of LIRGs is derived to test the hypothesis that the cluster is similar to the LIRG-rich cluster CL0024+1654<sup>5</sup><sup>5</sup>5In fact, to avoid throwing away several sources close to the LIRG flux threshold of $`1\times \mathrm{\hspace{0.17em}10}^{11}L_{}`$ and thereby degrading the statistics of the comparison, the flux threshold for the comparison of detected luminous sources was set to $`9\times \mathrm{\hspace{0.17em}10}^{10}L_{}`$. The LIRG count in CL0024+1654 is multiplied by the ratios of (a) virial mass per unit area of the cluster to that of CL0024+1654, the square of the respective distances to the cluster and to CL0024+1654, and the observed solid-angle for the cluster to that of CL0024+1654. The resulting column of the table can then be compared with the column listing the actual observed number of LIRGs for each cluster.
The two most distant clusters observed with ISO are CL0024+1654 ($`z=0.39`$, Coia et al. 2004) and J1888.16CL ($`z=0.56`$, Duc et al. 2004). These are among the deepest ISOCAM observations and could detect several cluster members. (For CL0024+1654, 13 out of 35 sources found at 15 $`\mu `$m are spectroscopically confirmed to be cluster sources. For J1888.16CL, 6 out of 44 sources found at 15 $`\mu `$m are so confirmed.) A common feature of these two clusters is the high star formation rate inferred from their MIR luminosities. These two observations were also the most extended cluster maps performed in terms of absolute cluster area covered at the cluster, $`2.3\times \mathrm{\hspace{0.17em}2.3}`$ and $`1.3\times \mathrm{\hspace{0.17em}6}`$ square Mpc for CL0024+1654 and J1888.16CL, respectively. This fact, and evolutionary effects detectable around $`z\mathrm{\hspace{0.17em}0.5}`$ (see Dressler et al. 1999), may explain the high IR luminosity of the 15 $`\mu `$m sources found in these clusters.
In the case of J1888.16CL, Duc et al. (2004) estimate star-formation rates ranging between 20 and 120 $`M_{}`$ per year. At least six galaxies belong to the cluster and have IR luminosities above $`1.3\times \mathrm{\hspace{0.17em}10}^{11}L_{}`$. In CL0024+1654, Coia et al. (2004b) report ten sources brighter than $`9\times \mathrm{\hspace{0.17em}10}^{10}L_{}`$. The star formation rates inferred from the MIR flux are one to two orders of magnitudes greater than those based on the \[OII\] flux (though in this case the comparison was only possible for the three sources for which \[OII\] data was available.) This is compatible with the result in A1689 (Fig.4) and implies similar dust extinction characteristics.
Interestingly, the galaxies emitting at 15$`\mu `$m appear to have a spatial distribution and a velocity dispersion slightly different from the other cluster galaxies. Galaxies in CL0024+1654 are detected preferentially at larger radii, with the velocity dispersion of 15$`\mu `$m sources being greater than that of the galaxies in the cluster. In J1888.16CL, Duc et al. (2004) estimate that to explain the number of sources detected on the basis of infall of galaxies from the field an infall rate of about 100 massive galaxies per 100 Myr is required, which seems unrealistic. Numerical simulations and X–ray observations show however that accretion onto clusters from the field is not a spherically symmetric process, but occurs along filaments or via mergers with other groups and clusters. One therefore cannot exclude the possibility that the LIRGs observed in these distant clusters belonged to such a recently accreted structure. An alternative possibility is that the collision with an accreted group of galaxies stimulated star formation in the galaxies of the group as a consequence of a rapidly varying tidal field (Bekki 1999). This could be the case for CL0024+1654 and A1689, clusters which show evidence of accreting groups of galaxies in their multi-modal velocity distributions. CL0024+1654 is in the process of interacting with a smaller cluster.
### 2.4 A diffuse intra-cluster dust component ?
The hot intra-cluster material contains metals, and so is not entirely primordial. Might not the stars which produced the metals also have deposited dust in the ICM? The first to note that emission from intra-cluster material might be observable were Yahil & Ostriker (1973), based on a galactic dust-to-gas ratio and the observed intra-cluster gas. Ostriker & Silk (1973) and Silk & Burke (1974) developed expressions for the lifetime of dust in a hot intra-cluster medium. Pustilnik (1975), drawing upon contemporaneous reports of optical absorption in clusters, attributed it to dust and estimated that cluster emission at 100 $`\mu `$m would be in the $`10^3`$ to $`10^4`$ Jy range for 6 nearby clusters. Voshchinnikov & Khersonskii (1984) also attributed claimed reddenning of galaxies in distant clusters to dust absorption, and estimated that the total FIR emission from the Coma or Perseus clusters should be $`10^5`$ to $`10^6`$ Jy (tens of Jy/arcminute<sup>2</sup>) in the 50 to 100 $`\mu `$m range. They estimated the sputtering lifetime of intra-cluster dust grains to be up to $`10^8`$ years. Hu et al. (1985), noting that intra-cluster dust must be short-lived, predicted FIR dust emission of a few Jy per square degree, close to the IRAS limit, for a sample of X-ray luminous clusters.
IRAS measurements failed to bear out even the most modest of the above predictions. Kelly & Rieke (1990) co-added IRAS scans across 71 clusters with $`0.3z0.92`$ to arrive at an average 60 $`\mu `$m value for cluster emission of 26$`\pm `$5 mJy per cluster, and 46$`\pm `$22 mJy at 100 $`\mu `$m. Dwek et al. (1990) refined models of intra-cluster dust and its interactions and calculated an upper limit of 0.2 MJy/sr for dust-emission from the Coma cluster, consistent with IRAS observations. Then total cluster emission would not be more than a few Jy at the peak wavelength (around 100 $`\mu `$m). They concluded that dust in the cluster centre could not explain the visual extinction, nor could cluster galaxies or their halos. Dust in the outskirts could, if it were un-depleted. But they saw no mechanism for the production of such dust. Wise et al. (1993) analysed 56 clusters at 60 and 100 $`\mu `$m from clusters with a range of X-ray emission, and some without cDs. For the only two clusters (A262 and A2670) showing a far-infrared excess lacking an immediate explanation (in terms of point sources or cirrus) they concluded that the result was likely to be due to discrete sources in the clusters. Averaged over the sample as a whole there was evidence of excess FIR at the 2-$`\sigma `$ level. No large FIR excesses associated with cooling flows were found. Bregman et al. (1990) looked for evidence of star formation in 27 cD galaxies. In half of their sample of X-ray-bright clusters they found IR, X-ray and blue luminosities to be comparable, consistent with dust grains heated by the X-ray emitting gas, thereby suggesting that dust cooling can compare with thermal bremsstrahlung as a cooling mechanism for the intra-cluster gas.
Cox et al. (1995) studied a much larger sample of 158 Abell clusters, again at 60 and 100 $`\mu `$m, and after making a more rigorous correction for spurious sources due to galactic cirrus, they concluded that only about 10% of cD galaxies in rich clusters have significant FIR emission, but with luminosities ten times greater than the X-ray luminosities produced in the cores of clusters, a condition which they therefore regarded as transient for any individual cluster. If the FIR emission comes from dust heated by the intra-cluster thermal electrons, significant dust sputtering is expected on timescales of several 10<sup>8</sup> years. Dust must then be replenished to account for continuous IR emission, presumably through the mechanisms discussed for stripping material from cluster galaxies (see Section 1).
By the launch of ISO one could imagine a “life-cycle” of dust in a cluster, tracing the flow of material - gas and dust - out of infalling galaxies, the destruction of dust in the high-temperature intra-cluster medium, and its possible eventual re-deposition, through the mechanism of cooling flows, into the cluster-dominant galaxies. But only upper-limits or occasionally, and with marginal significance, global or average cluster FIR emission, could constrain scenarios for the role of dust in the physics of the clusters as a whole.
Stickel et al. (1998, 2002) used ISOPHOT to observe extended FIR emission of six Abell clusters. Strip scanning measurements were performed at 120 $`\mu `$m and 180 $`\mu `$m. The raw profiles of the I$`_{120\mu m}`$/I$`_{180\mu m}`$ surface brightness ratio including zodiacal light show a bump towards Abell 1656 (Coma), dips towards Abell 262 and Abell 2670, and are without clear structure towards Abell 400, Abell 496, and Abell 4038. After subtraction of the zodiacal light and allowance for cirrus emission, only the bump towards Abell 1656 (Coma) is still present (Fig. 6). This excess of $`\mathrm{\hspace{0.17em}0.2}`$ MJy/sr seen at 120 $`\mu `$m towards Abell 1656 (Coma) is interpreted as thermal emission from intra-cluster dust distributed in the hot X-ray emitting Coma intra-cluster medium. The integrated excess flux within the central region of 10 to 15 diameter is $``$ 2.8 Jy. Since the dust temperature is poorly constrained only a rough estimate of the associated dust mass of $`M_D\mathrm{\hspace{0.17em}10}^7M_{}`$ can be derived. The associated visual extinction is negligible (A<sub>V</sub>$``$ 0.1 mag) and much smaller than claimed from optical observations. No evidence is found for intra-cluster dust in the other five clusters observed.
Quillen et al. (1999) suggested integrated emission from the cluster galaxies as the most likely source for the detected signal at the centre of Coma. Stickel et al. (2002) replied that if this was indeed the case, the same signal should have been detected in all clusters observed.
The absence of any signature for intra-cluster dust in five clusters and the rather low inferred dust mass in Abell 1656 indicates that intra-cluster dust is probably not responsible for the excess X-ray absorption reported in cooling flow clusters (White et al. 1991). These observations thus represent a further unsuccessful attempt to detect the presumed final stage of the cooling flow material. This agrees with numerous previous studies at other wavelengths, while spectroscopic observations with ESA’s XMM-Newton X-ray observatory have shown that intra-cluster gas does not cool beyond $`1`$ keV (see, e.g., Molendi & Pizzolato 2001), so dust deposition in cooling flows is not actually expected. Not surprisingly, then, further attempts by Hansen et al. (1999, 2000a, b) also failed to detect dust associated with the cooling flows presumed to exist at the centres of most galaxy clusters.
## 3 Conclusions
Building upon the legacy of IRAS, ISO could consolidate a number of important fundamenatal conclusions about the properties of galaxies in the Virgo, Coma and other nearby galaxy clusters. The correlation between galaxy Hubble type and mid- and far-infrared properties was firmly established. A major new cold dust component was identified by ISOPHOT observations, which could not have been found by IRAS. Moderate resolution FIR spectroscopy was possible for Virgo cluster galaxies using LWS, establishing a correlation between the strength of the \[CII\] line and the FIR flux and a two order of magnitude difference in \[CII\] to near-IR ratio between early type galaxies and late (spiral) types. Mid-infrared emission (5 to 18 $`\mu `$m) was found to correlate with star formation, but to trace it less faithfully when star formation rates become high enough for UV photons to disrupt the infrared-emitting materials. However, the MIR emission traces well the FIR and bolometric emission. In general the properties of galaxies in nearby clusters ($`z<\mathrm{\hspace{0.17em}0.1}`$) were found to exhibit little dependence on the cluster environment.
ISO, for the first time, could extend mid-infrared observations to clusters beyond z = 0.1, and in so doing has detected over 100 galaxies in clusters in the redshift range 0.17 to almost 0.6. Although the collected observations on seven such clusters were rather heterogeneous, a number of important trends are found in the ISO results. There is a clear tendency for clusters at higher redshift to exhibit higher average rates of star formation in their galaxies, and numerous LIRGs have been observed in such clusters. There is tentative evidence for an association between the infrared luminosities found and galaxy infall from the field, but substantial evidence to link high levels of star-formation in cluster galaxies to the dynamical status of a cluster, and to interactions with other (sub-)clusters. It seems clear that star formation rates deduced for cluster galaxies from optical tracers often fall one to two orders of magnitude below rates derived from MIR emission levels, so that star formation in cluster galaxies may have been seriously underestimated in some cases, in the past.
ISO found little evidence for widespread infrared emission from dust in the intra-cluster medium.
Clearly, what ISO has not been able to supply has been systematic large area mapping of a substantial sample of galaxy clusters out to high redshift. Observing times to achieve such coverage with ISO would have been prohibitive, given the multi-position and heavily-overlapped rasters that would have been required. If galaxy evolution within clusters is to be further explored and related to galaxy evolution in the field, clusters of different masses and dynamical status must be studied systematically in the MIR and FIR out to well over 1 virial radius in order to understand how and why the IR properties of galaxies change from cluster to cluster, and from cluster to field. Large area coverage is important since it is known that significant modifications of the galaxy properties already occur in the outskirts of galaxy clusters (Gómez et al. 2003; Kodama et al. 2001; Lewis et al. 2002).
Several programmes underway or planned with the Spitzer Space Observatory should thoroughly address these challenges.
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# Trace Formulas in Connection with Scattering Theory for Quasi-Periodic Background
## 1. Introduction
Scattering theory for Jacobi operators $`H`$ with periodic (respectively more general) background has attracted considerable interest recently. In Volberg and Yuditskii have treated the case where $`H`$ has a homogeneous spectrum and is of Szegö class exhaustively. In Egorova and we have established direct and inverse scattering theory for Jacobi operators which are short range perturbations of quasi-periodic finite-gap operators. For further information and references we refer to these articles and .
In the case of constant background it is well-known that the transmission coefficient is the perturbation determinant in the sense of Krein , see e.g., or . The purpose of the present paper is to establish this result for the case of quasi-periodic finite-gap background, thereby establishing the connection with Krein’s spectral shift theory.
Moreover, scattering theory for Jacobi operators is not only interesting in its own right, it also constitutes the main ingredient of the inverse scattering transform for the Toda hierarchy (see, e.g., , , , or ). Since the transmission coefficient is invariant when our Jacobi operator evolves in time with respect to some equation of the Toda hierarchy, the corresponding trace formulas provide the conserved quantities for the Toda hierarchy in this setting.
## 2. Notation
Let
(2.1)
$$H_qf(n)=a_q(n)f(n+1)+a_q(n1)f(n1)+b_q(n)f(n)$$
be a quasi-periodic Jacobi operator in $`\mathrm{}^2()`$ associated with the Riemann surface of the function
(2.2)
$$R_{2g+2}^{1/2}(z),R_{2g+2}(z)=\underset{j=0}{\overset{2g+1}{}}(zE_j),E_0<E_1<\mathrm{}<E_{2g+1},$$
$`g`$. The spectrum of $`H_q`$ is purely absolutely continuous and consists of $`g+1`$ bands
(2.3)
$$\sigma (H_q)=\underset{j=0}{\overset{g}{}}[E_{2j},E_{2j+1}].$$
For every $`z`$ the Baker-Akhiezer functions $`\psi _{q,\pm }(z,n)`$ are two (weak) solutions of $`H_q\psi =z\psi `$, which are linearly independent away from the band-edges $`\{E_j\}_{j=0}^{2g+1}`$, since their Wronskian is given by
(2.4)
$$W_q(\psi _{q,}(z),\psi _{q,+}(z))=\frac{R_{2g+2}^{1/2}(z)}{_{j=1}^g(z\mu _j)}.$$
Here $`\mu _j`$ are the Dirichlet eigenvalues at base point $`n_0=0`$. We recall that $`\psi _{q,\pm }(z,n)`$ have the form
$$\psi _{q,\pm }(z,n)=\theta _{q,\pm }(z,n)w(z)^{\pm n},$$
where $`\theta _{q,\pm }(z,n)`$ is quasi-periodic with respect to $`n`$ and $`w(z)`$ is the quasi-momentum. In particular, $`|w(z)|<1`$ for $`z\backslash \sigma (H_q)`$ and $`|w(z)|=1`$ for $`z\sigma (H_q)`$.
We assume that the reader is familiar with this class of operators and refer to and for further information.
## 3. Asymptotics of Jost solutions
After we have these preparations out of our way, we come to the study of short-range perturbations $`H`$ of $`H_q`$ associated with sequences $`a`$, $`b`$ satisfying $`a(n)a_q(n)`$ and $`b(n)b_q(n)`$ as $`|n|\mathrm{}`$. More precisely, we will make the following assumption throughout this paper:
Let $`H`$ be a perturbation of $`H_q`$ such that
(3.1)
$$\underset{n}{}\left(|a(n)a_q(n)|+|b(n)b_q(n)|\right)<\mathrm{},$$
that is, $`HH_q`$ is trace class.
We first establish existence of Jost solutions, that is, solutions of the perturbed operator which asymptotically look like the Baker-Akhiezer solutions.
###### Theorem 3.1.
Assume (3.1). Then there exist (weak) solutions $`\psi _\pm (z,.)`$, $`z\backslash \{E_j\}_{j=0}^{2g+1}`$, of $`H\psi =z\psi `$ satisfying
(3.2)
$$\underset{n\pm \mathrm{}}{lim}w(z)^n\left(\psi _\pm (z,n)\psi _{q,\pm }(z,n)\right)=0,$$
where $`\psi _{q,\pm }(z,.)`$ are the Baker-Akhiezer functions. Moreover, $`\psi _\pm (z,.)`$ are continuous (resp. holomorphic) with respect to $`z`$ whenever $`\psi _{q,\pm }(z,.)`$ are and have the following asymptotic behavior
(3.3)
$$\psi _\pm (z,n)=\frac{z^n}{A_\pm (n)}(\begin{array}{c}n1\\ ^{}\\ j=0\end{array}a_q(j))^{\pm 1}\left(1+\left(B_\pm (n)\pm \begin{array}{c}n\\ ^{}\\ j=1\end{array}b_q(j{\scriptscriptstyle \genfrac{}{}{0pt}{}{0}{1}})\right)\frac{1}{z}+O(\frac{1}{z^2})\right),$$
where
$`A_+(n)`$ $`=`$ $`{\displaystyle \underset{j=n}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a(j)}{a_q(j)}},B_+(n)={\displaystyle \underset{m=n+1}{\overset{\mathrm{}}{}}}(b_q(m)b(m)),`$
(3.4) $`A_{}(n)`$ $`=`$ $`{\displaystyle \underset{j=\mathrm{}}{\overset{n1}{}}}{\displaystyle \frac{a(j)}{a_q(j)}},B_{}(n)={\displaystyle \underset{m=\mathrm{}}{\overset{n1}{}}}(b_q(m)b(m)).`$
###### Proof.
The proof can be done as in the periodic case (see e.g., , , or , Section 7.5). There a stronger decay assumption (i.e., first moments summable) is made, which is however only needed at the band edges $`\{E_j\}_{j=0}^{2g+1}`$. ∎
For later use we note the following immediate consequence
###### Corollary 3.2.
Under the assumptions of the previous theorem we have
(3.5)
$$\underset{n\pm \mathrm{}}{lim}w(z)^n\left(\psi _\pm ^{}(z,n)n\frac{w^{}(z)}{w(z)}\psi _\pm (z,n)\psi _{q,\pm }^{}(z,n)\pm n\frac{w^{}(z)}{w(z)}\psi _{q,\pm }(z,n)\right)=0,$$
where the prime denotes differentiation with respect to $`z`$.
###### Proof.
Just differentiate (3.2) with respect to $`z`$, which is permissible by uniform convergence on compact subsets of $`\backslash \{E_j\}_{j=0}^{2g+1}`$. ∎
We remark that if we require our perturbation to satisfy the usual short range assumption as in (i.e., first moments summable), then we even have $`w(z)^n(\psi _\pm ^{}(z,n)\psi _{q,\pm }^{}(z,n))0`$.
From Theorem 3.1 we obtain a complete characterization of the spectrum of $`H`$.
###### Theorem 3.3.
Assume (3.1). Then we have $`\sigma _{ess}(H)=\sigma (H_q)`$, the point spectrum of $`H`$ is confined to $`\overline{\backslash \sigma (H_q)}`$. Furthermore, the essential spectrum of $`H`$ is purely absolutely continuous except for possible eigenvalues at the band edges.
###### Proof.
An immediate consequence of the fact that $`HH_q`$ is trace class and boundedness of the Jost solutions inside the essential spectrum. ∎
Our next result concerns the asymptotics of the Jost solutions at the other side.
###### Lemma 3.4.
Assume (3.1). Then the Jost solutions $`\psi _\pm (z,.)`$, $`z\backslash \sigma (H)`$, satisfy
(3.6)
$$\underset{n\mathrm{}}{lim}|w(z)^n(\psi _\pm (z,n)\alpha (z)\psi _{q,\pm }(z,n))|=0,$$
where
(3.7)
$$\alpha (z)=\frac{W(\psi _{}(z),\psi _+(z))}{W_q(\psi _{q,}(\lambda ),\psi _{q,+}(z))}=\frac{_{j=1}^g(z\mu _j)}{R_{2g+2}^{1/2}(z)}W(\psi _{}(z),\psi _+(z)).$$
###### Proof.
Since $`HH_q`$ is trace class, we have for the difference of the Green’s functions
$$\underset{n\pm \mathrm{}}{lim}G(z,n,n)G_q(z,n,n)=\underset{n\pm \mathrm{}}{lim}\delta _n,((Hz)^1(H_qz)^1)\delta _n=0$$
and hence
$$\underset{n\mathrm{}}{lim}\psi _{q,}(z,n)(\psi _+(z,n)\alpha (z)\psi _{q,+}(z,n))=0,$$
which is the claimed result. ∎
Note that $`\alpha (z)`$ is just the inverse of the transmission coefficient (see, e.g., or , Section 7.5). It is holomorphic in $`\backslash \sigma (H_q)`$ with simple zeros at the discrete eigenvalues of $`H`$ and has the following asymptotic behavior
(3.8)
$$\alpha (z)=\frac{1}{A}(1+\frac{B}{z}+O(z^2)),A=A_{}(0)A_+(0),B=B_{}(1)+B_+(0).$$
## 4. Connections with Krein’s spectral shift theory and Trace formulas
To establish the connection with Krein’s spectral shift theory we next show:
###### Lemma 4.1.
We have
(4.1)
$$\frac{d}{dz}\alpha (z)=\alpha (z)\underset{n}{}\left(G(z,n,n)G_q(z,n,n)\right),z\backslash \sigma (H),$$
where $`G(z,m,n)`$ and $`G_q(z,m,n)`$ are the Green’s functions of $`H`$ and $`H_q`$, respectively.
###### Proof.
Green’s formula (, eq. (2.29)) implies
(4.2)
$$W_n(\psi _+(z),\psi _{}^{}(z))W_{m1}(\psi _+(z),\psi _{}^{}(z))=\underset{j=m}{\overset{n}{}}\psi _+(z,j)\psi _{}(z,j),$$
hence the derivative of the Wronskian can be written as
$`{\displaystyle \frac{d}{dz}}W(\psi _{}(z),\psi _+(z))=W_n(\psi _{}^{}(z),\psi _+(z))+W_n(\psi _{}(z),\psi _+^{}(z))`$
$`=`$ $`W_m(\psi _{}^{}(z),\psi _+(z))+W_n(\psi _{}(z),\psi _+^{}(z)){\displaystyle \underset{j=m+1}{\overset{n}{}}}\psi _+(z,j)\psi _{}(z,j).`$
Using Corollary 3.2 and Lemma 3.4 we have
$`W_m(\psi _{}^{}(z),\psi _+(z))`$ $`=`$ $`W_m(\psi _{}^{}+m{\displaystyle \frac{w^{}}{w}}\psi _{},\psi _+)`$
$`{\displaystyle \frac{w^{}}{w}}\left(mW(\psi _{},\psi _+)a(m)\psi _{}(m+1)\psi _+(m)\right)`$
$``$ $`\alpha W_{q,m}(\psi _{q,}^{}+m{\displaystyle \frac{w^{}}{w}}\psi _{q,},\psi _{q,+})`$
$`\alpha {\displaystyle \frac{w^{}}{w}}\left(mW_q(\psi _{q,},\psi _{q,+})a_q(m)\psi _{q,}(m+1)\psi _{q,+}(m)\right)`$
$`=`$ $`\alpha (z)W_m(\psi _{q,}^{}(z),\psi _{q,+}(z))`$
as $`m\mathrm{}`$. Similarly we obtain
$`W_n(\psi _{}(z),\psi _+^{}(z))`$ $``$ $`\alpha (z)W_n(\psi _{q,}(z),\psi _{q,+}^{}(z))`$
as $`n\mathrm{}`$ and again using (4.2) we have
$$W_m(\psi _{q,}^{}(z),\psi _{q,+}(z))=W_n(\psi _{q,}^{}(z),\psi _{q,+}(z))+\underset{j=m+1}{\overset{n}{}}\psi _{q,+}(z,j)\psi _{q,}(z,j).$$
Collecting terms we arrive at
$`W^{}(\psi _{}(z),\psi _+(z))`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(\psi _+(z,j)\psi _{}(z,j)\alpha (z)\psi _{q,+}(z,j)\psi _{q,}(z,j)\right)`$
$`+\alpha (z)W_q^{}(\psi _{q,}(z)\psi _{q,+}(z)).`$
Now we compute
$`{\displaystyle \frac{d}{dz}}\alpha (z)`$ $`=`$ $`{\displaystyle \frac{d}{dz}}\left({\displaystyle \frac{W}{W_q}}\right)=\left({\displaystyle \frac{1}{W_q}}\right)^{}W+{\displaystyle \frac{1}{W_q}}W^{}`$
$`=`$ $`{\displaystyle \frac{W_q^{}}{W_q^2}}W+{\displaystyle \frac{1}{W_q}}\left({\displaystyle \underset{j}{}}\left(\psi _+\psi _{}\alpha \psi _{q,+}\psi _{q,}\right)+\alpha W_q^{}\right)`$
$`=`$ $`{\displaystyle \frac{1}{W_q}}{\displaystyle \underset{j}{}}\left(\psi _+\psi _{}\alpha \psi _{q,+}\psi _{q,}\right),`$
which finishes the proof. ∎
As an immediate consequence, we can identify $`\alpha (z)`$ as Krein’s perturbation determinant () of the pair $`H`$, $`H_q`$.
###### Theorem 4.2.
The function $`A\alpha (z)`$ is Krein’s perturbation determinant:
(4.3)
$$\alpha (z)=\frac{1}{A}det\left(1\mathrm{l}+(H(t)H_q(t))(H_q(t)z)^1\right),A=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\frac{a(j)}{a_q(j)}.$$
By , Theorem 1, $`\alpha (z)`$ has the following representation
(4.4)
$$\alpha (z)=\frac{1}{A}\mathrm{exp}\left(_{}\frac{\xi _\alpha (\lambda )d\lambda }{\lambda z}\right),$$
where
(4.5)
$$\xi _\alpha (\lambda )=\frac{1}{\pi }\underset{ϵ0}{lim}\mathrm{arg}\alpha (\lambda +iϵ)$$
is the spectral shift function.
Hence
(4.6)
$$\tau _j=\mathrm{tr}(H^j(H_q)^j)=j_{}\lambda ^{j1}\xi _\alpha (\lambda )𝑑\lambda ,$$
where $`\tau _j/j`$ are the expansion coefficients of $`\mathrm{ln}\alpha (z)`$ around $`z=\mathrm{}`$:
$$\mathrm{ln}\alpha (z)=\mathrm{ln}A\underset{j=1}{\overset{\mathrm{}}{}}\frac{\tau _j}{jz^j}.$$
They are related to the expansion $`\alpha _j`$ coefficients of
$$\alpha (z)=\frac{1}{A}\underset{j=0}{\overset{\mathrm{}}{}}\frac{\alpha _j}{z^j},\alpha _0=1,$$
via
(4.7)
$$\tau _1=\alpha _1,\tau _j=j\alpha _j\underset{k=1}{\overset{j1}{}}\alpha _{jk}\tau _k.$$
## 5. Conserved quantities of the Toda hierarchy
Finally we turn to solutions of the Toda hierarchy $`\mathrm{TL}_r`$ (see, e.g., , , , or ). Let $`(a_q(t),b_q(t))`$ be a quasi-periodic finite-gap solution of some equation in the Toda hierarchy, $`\mathrm{TL}_r(a_q(t),b_q(t))=0`$, and let $`(a(t),b(t))`$ be another solution, $`\mathrm{TL}_r(a(t),b(t))=0`$, such that (3.1) holds for all $`t`$.
Since the transmission coefficient $`T(z,t)=T(z,0)T(z)`$ is conserved (see – formally this follows from unitary invariance of the determinant), so is $`\alpha (z)=T(z)^1`$.
###### Theorem 5.1.
The quantities
(5.1)
$$A=\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\frac{a(j,t)}{a_q(j,t)}$$
and $`\tau _j=\mathrm{tr}(H^j(t)H_q(t)^j)`$, that is,
$`\tau _1`$ $`=`$ $`{\displaystyle \underset{n}{}}b(n,t)b_q(n,t)`$
$`\tau _2`$ $`=`$ $`{\displaystyle \underset{n}{}}2(a(n,t)^2a_q(n,t)^2)+(b(n,t)^2b_q(n,t)^2)`$
$`\mathrm{}`$
are conserved quantities for the Toda hierarchy.
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# Stellar evolution with rotation and magnetic fields:
## 1 Introduction
Stellar rotation influences all the outputs of stellar evolution and nucleosynthesis and several grids of models have been made for massive stars (Langer et al. langeragb (1999); Heger et al. hegerI (2000); Meynet & Maeder MMV (2000)). The effects of rotation are even more important at metallicities $`Z`$ lower than solar (Maeder & Meynet MMVII (2001)). However, we do not know the role of the magnetic field in stellar evolution and it is the purpose of this series of works to approach this matter, which is considered to be critical one (cf. Roxburgh Rox03 (2003)). Spruit (Spruit02 (2002)) proposed a dynamo mechanism operating in stellar radiative layers in differential rotation. This dynamo is based on the conjectured Tayler instability, which is apparently the first one to occur in a radiative zone (Tayler Tay73 (1973); Pitts & Tayler Pitts86 (1986)). For now, there is no empirical or observational proof of the existence of this instability. According to Pitts & Tayler (Pitts86 (1986)) and Spruit (Spruit02 (2002)), even a very weak horizontal magnetic field is subject to this instability, which then creates a vertical field component, which is wound up by differential rotation. As a result, the field lines become progressively closer and denser and thus a strong horizontal field is created at the energy expense of differential rotation.
In a first paper (Maeder & Meynet Magn1 (2003), paper I), we have shown that in a rotating star a magnetic field can be created during MS evolution by the Spruit dynamo. We have examined the timescale for the field creation, its amplitude and the related diffusion coefficients. The clear result is that magnetic field and its effects are quite important. In the second paper (Maeder & Meynet Magn2 (2004), paper II), a generalisation of the equations of the dynamo was developed. The solutions fully agree with Spruit’s solution in the two limiting cases considered (Spruit Spruit02 (2002)), i.e. “Case 0” when the $`\mu `$–gradient dominates and “Case 1” when the $`T`$–gradient dominates with large non–adiabatic effects. Our more general solution encompasses all cases of $`\mu `$– and $`T`$–gradients, as well as all cases from the fully adiabatic to non–adiabatic solutions.
Paper II suggested that there is a complex feedback between the magnetic instability, which generates the field, and the thermal instability which drives the meridional circulation (Maeder & Meynet Magn2 (2004)). However, it was beyond the scope of paper II to make numerical models of the interaction between circulation currents and dynamo. In this paper we account for this feedback, the main steps of which are the following:
* Differential rotation creates the magnetic field.
* Magnetic field tends to suppress differential rotation.
* The absence of differential rotation strongly enhances meridional circulation,
* Meridional circulation increases differential rotation.
* Differential rotation creates the magnetic field (the loop is closed).
Thus, a delicate balance between the thermal and magnetic effects occurs during evolution. The balance of these effects influences the transport of chemical elements to the stellar surface and the internal transport of angular momentum, which determines the evolution of the surface rotation.
In Sect. 2, we collect in a consistent way the basic equations of the dynamo. In Sect. 3, we give the basic expressions for the transport coefficients. In Sect. 4, we calculate numerical models for the interaction of the dynamo and circulation on the internal stellar structure. In Sect. 5, the effects at the stellar surface and in particular the evolution of abundances are analysed. Sect. 6 gives the conclusions.
## 2 The general equations for the dynamo
The present set of equations for the dynamo based on Tayler–Spruit instability has three important advantages with respect to the system of equations by Spruit (Spruit02 (2002)).
1. The equations by Spruit apply to the non–adiabatic case. As most of the stellar interior is adiabatic, this situation is not appropriate. The present equations fully treat the problem whether it is adiabatic or non–adiabatic.
2. Either the $`\mu `$–gradient was considered to dominate or it was absent. Such a situation is not adequate when mild mixing effects, which produce small $`\mu `$–gradients, are studied. In addition, as always in all problems of mixing, the $`\mu `$–gradients need to be very carefully treated, otherwise the results are incorrect.
3. There is no need to interpolate between the two asymptotic solutions called “Case 0” or “Case 1”. The interpolation was made by Spruit (Spruit02 (2002)) as a function of the parameter $`q`$, which defines the amount of differential rotation (see definition given in Eq. 8). This is hazardous, as it expresses the interpolation of effects concerning the $`\mu `$–gradient and the degree of non–adiabaticity as a function of differential rotation. This is not physically justified and it introduces an additional spurious dependence on the parameter $`q`$ into the problem.
Let us briefly summarize the consistent system of equations. The energy density $`u_\mathrm{B}`$ of a magnetic field of intensity $`B`$ per volume unity is
$`u_\mathrm{B}={\displaystyle \frac{B^2}{8\pi }}={\displaystyle \frac{1}{2}}\rho r^2\omega _\mathrm{A}^2\mathrm{with}\omega _\mathrm{A}={\displaystyle \frac{B}{(4\pi \rho )^{\frac{1}{2}}r}},`$ (1)
where $`\omega _\mathrm{A}`$ is the Alfvén frequency in a spherical geometry, $`r`$ is the radius and $`\rho `$ the density. In stable radiative layers, there is in principle no particular motion. However, if due to the magnetic field or rotation, some unstable displacements of vertical amplitude $`l/2`$ occur around an average stable position, the restoring buoyancy force produces vertical oscillations around the equilibrium position with a frequency equal to the Brunt–Va̋isa̋la̋ frequency $`N`$. In a medium with both thermal and magnetic diffusivity $`K`$ and $`\eta `$, the Brunt–Väisälä frequency becomes (Maeder & Meynet Magn2 (2004))
$$N^2=\frac{\eta /K}{\eta /K+2}N_\mathrm{T}^2+N_\mu ^2,$$
(2)
with the radiative diffusivity $`K=\frac{4acT^3}{3\kappa \rho ^2C_\mathrm{p}}`$ and with the following partial frequencies
$$N_\mu ^2=\frac{g\phi }{H_\mathrm{P}}_\mu \mathrm{and}N_\mathrm{T}^2=\frac{g\delta }{H_\mathrm{P}}(_{\mathrm{ad}}).$$
(3)
The thermodynamic coefficients $`\delta `$ and $`\phi `$ are $`\delta =\left(\frac{\mathrm{ln}\rho }{\mathrm{ln}T}\right)_{\mathrm{P},\mu }`$ and $`\phi =\left(\frac{\mathrm{ln}\rho }{\mathrm{ln}\mu }\right)_{\mathrm{T},\mathrm{P}}`$. $`H_\mathrm{P}`$ is the pressure scale height. The restoring oscillations will have an average density of kinetic energy
$$u_\mathrm{N}f_\mathrm{N}\rho l^2N^2,$$
(4)
where $`f_\mathrm{N}`$ is a geometrical factor of the order of unity. If the magnetic field produces some instability with a vertical component, one must have $`u_\mathrm{B}>u_\mathrm{N}`$. Otherwise, the restoring force of gravity which acts at the dynamical timescale would immediately counteract the magnetic instability. From this inequality, one obtains $`l^2<\frac{1}{2f_\mathrm{N}}r^2\frac{\omega _\mathrm{A}^2}{N^2}`$. If $`f_\mathrm{N}=\frac{1}{2}`$, we have the condition for the vertical amplitude of the instability (Spruit Spruit02 (2002); Eq. 6),
$$l<r\frac{\omega _\mathrm{A}}{N},$$
(5)
where $`r`$ is the radius. This means that there is a maximum size of the vertical length $`l`$ of a magnetic instability. In order not to be quickly damped by magnetic diffusivity, the vertical length scale of the instability must satisfy
$$l^2>\frac{\eta }{\sigma _\mathrm{B}}=\frac{\eta \mathrm{\Omega }}{\omega _\mathrm{A}^2},$$
(6)
where $`\mathrm{\Omega }`$ is the angular velocity and $`\sigma _\mathrm{B}`$ is the characteristic frequency of the magnetic field. In a rotating star, this frequency is $`\sigma _\mathrm{B}=(\omega _\mathrm{A}^2/\mathrm{\Omega })`$ due to the Coriolis force (Spruit Spruit02 (2002); see also Pitts & Tayler Pitts86 (1986)). The combination of the limits given by Eqs. (5) and (6) gives for the case of marginal stability,
$$\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^4=\frac{N^2}{\mathrm{\Omega }^2}\frac{\eta }{r^2\mathrm{\Omega }}.$$
(7)
The equality of the amplification time of Tayler instability $`\tau _\mathrm{a}=N/(\omega _\mathrm{A}\mathrm{\Omega }q)`$ with the inverse of the characteristic frequency $`\sigma _\mathrm{B}`$ of the magnetic field leads to the equation (Spruit Spruit02 (2002))
$`{\displaystyle \frac{\omega _\mathrm{A}}{\mathrm{\Omega }}}=q{\displaystyle \frac{\mathrm{\Omega }}{N}}\mathrm{with}q={\displaystyle \frac{\mathrm{ln}\mathrm{\Omega }}{\mathrm{ln}r}}.`$ (8)
With account of the Brunt–Väisälä (Eq. 2), one has
$$\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^2=\frac{\mathrm{\Omega }^2q^2}{N_\mathrm{T}^2\frac{\eta /K}{\eta /K+\mathrm{\hspace{0.33em}2}}+N_\mu ^2}.$$
(9)
By eliminating the expression of $`N^2`$ between Eqs. (7) and (9), we obtain an expression for the magnetic diffusivity,
$$\eta =\frac{r^2\mathrm{\Omega }}{q^2}\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^6.$$
(10)
Eqs. (7) and (9) form a coupled system relating the two unknown quantities $`\eta `$ and $`\omega _\mathrm{A}`$. Instead, one may also consider for example the system formed by Eqs. (9) and (10). Formally, if one accounts for the complete expressions of the thermal gradient $``$, the system of equations would be of degree 10 in the unknown quantity $`x=\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^2`$ (paper II). The fact that the ratio $`\eta /K`$ is very small allows us to bring these coupled equations to a system of degree 4,
$$\frac{r^2\mathrm{\Omega }}{q^2K}\left(N_\mathrm{T}^2+N_\mu ^2\right)x^4\frac{r^2\mathrm{\Omega }^3}{K}x^3+2N_\mu ^2x2\mathrm{\Omega }^2q^2=0.$$
(11)
The solution of this equation, which is easily obtained numerically, provides the Alfvén frequency and by Eq. (10) the thermal diffusivity. As shown in paper II, the various peculiar cases studied by Spruit (Spruit02 (2002)) are all contained in the more general solution given here.
## 3 Coefficients for the transport of angular momentum and chemical elements
The condition that the ratio $`\omega _\mathrm{A}/\mathrm{\Omega }`$ given by (Eq. 8) is equal to or larger than the minimum value defined by Eq. (7) leads to a condition on the minimum shear for the dynamo to work (cf. Spruit Spruit02 (2002))
$$q>\left(\frac{N}{\mathrm{\Omega }}\right)^{7/4}\left(\frac{\eta }{r^2N}\right)^{1/4}.$$
(12)
There, $`N`$ is given by Eq. (2) above and $`\eta `$ by Eq. (10). If this condition is not fulfilled, there is no stationary situation and we consider that there is no magnetic field present. In practice, we will see that this situation occurs in the outer stellar envelope. There is a second condition. We need to check that $`\mathrm{\Omega }>\omega _\mathrm{A}`$. If this last condition is not realized, the present system of equations does not apply and we should consider the case of very slow rotation. This case is treated in the Appendix. In the present models, the value of $`\mathrm{\Omega }`$ considered is large enough so that we did not have to apply the solution for the case of very slow rotation.
The azimuthal component of the magnetic field is much stronger that the radial one in the Tayler–Spruit dynamo. We have for these components (Spruit Spruit02 (2002))
$$B_\phi =(4\pi \rho )^{\frac{1}{2}}r\omega _\mathrm{A}\mathrm{and}B_\mathrm{r}=B_\phi (l_\mathrm{r}/r),$$
(13)
where $`\omega _\mathrm{A}`$ is the solution of the general equation (11) and $`l_\mathrm{r}`$ is given by $`l_\mathrm{r}=r\frac{\omega _\mathrm{A}}{N}`$, which is obtained by assuming the marginal stability in Eq. (5).
Turning towards the transport of angular momentum by magnetic field, we first write the azimuthal stress by volume unity due to the magnetic field
$`S={\displaystyle \frac{1}{4\pi }}B_\mathrm{r}B_\phi ={\displaystyle \frac{1}{4\pi }}\left({\displaystyle \frac{l_\mathrm{r}}{r}}\right)B_\phi ^2=\rho r^2\left({\displaystyle \frac{\omega _\mathrm{A}^3}{N}}\right).`$ (14)
Then, the viscosity $`\nu `$ for the vertical transport of angular momentum can be expressed in terms of $`S`$ (Spruit Spruit02 (2002)),
$$\nu =\frac{S}{\rho q\mathrm{\Omega }}=\frac{\mathrm{\Omega }r^2}{q}\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^3\left(\frac{\mathrm{\Omega }}{N}\right).$$
(15)
This is the general expression of $`\nu `$ with $`\omega _\mathrm{A}`$ given by the solution of Eq. (11) and with $`N`$ by Eq. (2). We have the full set of expressions necessary to obtain the Alfvén frequency $`\omega _\mathrm{A}`$ and the magnetic diffusivity $`\eta `$. The parameter $`\eta `$ also expresses the vertical transport of the chemical elements, while the viscosity $`\nu `$ determines the vertical transport of the angular momentum by the magnetic field.
In paper II, we also checked that the rate of magnetic energy production $`W_\mathrm{B}=\frac{B^2}{8\pi }\frac{\omega _\mathrm{A}^2}{\mathrm{\Omega }}=\frac{1}{2}\rho r^2\frac{\omega _\mathrm{A}^4}{\mathrm{\Omega }}`$ per unit of time and volume is equal to the rate $`W_\nu =\frac{1}{2}\rho \nu \mathrm{\Omega }^2q^2`$ of the dissipation of rotational energy by the magnetic viscosity $`\nu `$ as given above.
## 4 Numerical applications
### 4.1 The models
We consider here models of 15 M, with a standard composition of $`X=0.705`$ and $`Z=0.02`$. The physics of the models, opacities, nuclear reactions, mass loss rates, structural rotational effects, shear mixing, meridional circulation, etc. are the same as in recent models of the Geneva group (Meynet & Maeder MMXI (2005)). We calculate 3 sets of models: the first is without rotation, the second with rotation and the third one with both rotation and magnetic field. The initial rotation velocity is 300 km s<sup>-1</sup>, which leads to an average velocity on the MS phase of about 240 km s<sup>-1</sup> when no magnetic field is present (see Fig. 3). In paper II, we did not include the effects of the meridional currents, in order to independently study the effects of the magnetic field. However, we showed that meridional circulation and the magnetic field may significantly interact, with possible consequences for the transport mechanisms. Thus, to examine this interplay, we also account here for meridional circulation currents. For the moment, we do not account for magnetic coupling by stellar winds, since there is no external convective zone in massive stars on the Main Sequence.
The choice of the time step is imposed by the most rapid process taking place. We checked that changing the time steps did not change the results. Here the fastest process is the transport of the angular momentum by the magnetic diffusivity. The diffusion coefficients $`\nu `$ for the angular momentum by the magnetic field reach in some cases values of $`10^{13}`$ cm<sup>2</sup> g<sup>-1</sup>, while the average value is 1–2 orders of magnitude lower. The cases of strong coefficients lead to diffusion timescales $`\tau (\mathrm{\Delta }R)^2/\nu `$ a few $`10^2`$ yr, where $`D`$ is the diffusion coefficient. For an appropriate treatment, we need to adopt very small time steps. In practice, we take time steps of the order of 20 years. This implies about $`610^5`$ models to cover in a detailed way the exact interplay of the effects of the magnetic field and of the meridional circulation during the MS phase !
Clearly, the model properties do not change significantly over such a time scale as shown for example by Fig. 2. Thus, in future, faster processes of calculations may be devised. However, by imposing very short time steps on the two processes, we set the study in the linear regime for both instabilities and we may thus proceed to an addition of their own different effects.
### 4.2 Evolution of the internal and superficial rotation
In a rotating star, the internal profile of angular velocity $`\mathrm{\Omega }(M_\mathrm{r})`$ changes with time during evolution, due to various effects (Meynet & Maeder MMV (2000)): central contraction and envelope expansion, transport of angular momentum by circulation and mass loss at the stellar surface. One assumes here that the mass lost by stellar winds just embarks its own angular momentum. The case of anisotropic stellar winds has been studied by Maeder (MIX (2002)), who has shown that such effects are important only for very massive stars ($`M60`$ M) in fast rotation. Here, the loss of angular momentum at the stellar surface only has a limited importance for the evolution of rotation. Fig. 1 shows the evolution of the internal profile of rotation. Differential rotation builds up during the MS phase, reaching about a factor of two between $`\mathrm{\Omega }`$ at the surface and in the convective core, near the end of the MS when the central H–content is $`X_\mathrm{c}=0.10`$. Then, fast central contraction accelerates the core rotation.
The situation is quite different when magnetic fields as described above are accounted for. As shown by Fig. 2, the angular velocity $`\mathrm{\Omega }`$ is almost constant throughout the stellar interior. It is not exactly constant, otherwise $`q`$ would be zero and the magnetic field would not be sustained anymore. Only at the very end of the MS phase does the fast central contraction bring about a small significant difference.
Fig. 3 shows the evolution of the surface rotation velocities as a function of age in the models with and without magnetic field. We notice the higher velocity during the whole MS phase of the model with magnetic field, which has $`v300`$ km s<sup>-1</sup>, compared to the model without a field, where $`v250`$ km s<sup>-1</sup>. There are 3 effects intervening:
1. The main part of the difference already occurs on the ZAMS. With rotation only, the model adjusts its rotation very quickly to an equilibrium profile determined by meridional circulation (Meynet & Maeder MMV (2000)) and this establishes an $`\mathrm{\Omega }`$–gradient which reduces the surface velocity. On the contrary, with a magnetic field, the model can keep a constant $`\mathrm{\Omega }`$ and there is no initial decrease.
2. The other two smaller effects occur during MS evolution. The magnetic coupling forces the surface to co–rotate with the convective core, as contraction makes it spin–up. With rotation only, the coupling (due to meridional circulation) is much weaker.
3. In the magnetic model, the surface enrichment in helium is slightly higher, this keeps the radius smaller. Thus, as the radius expansion is smaller, the decrease of surface rotation is also smaller.
The question of the differential rotation in the horizontal direction has also to be examined. In current rotating models without a magnetic field, it is assumed that the horizontal turbulence is strong enough to suppress the horizontal differential rotation, so that the hypothesis of shellular rotation with $`\mathrm{\Omega }=\mathrm{\Omega }(r)`$ applies (Zahn Zahn92 (1992)). In the presence of a magnetic field, the horizontal turbulence is likely reduced. The numerical values for the field $`B_\phi `$ found below is of the order of a few $`10^4`$ G (cf. Fig. 4). The horizontal coupling ensured by the magnetic field (cf. Fig. 4) is expressed by a coefficient $`D_{\mathrm{B}_\mathrm{h}}`$ (Maeder & Meynet Magn1 (2003))
$$D_{\mathrm{B}_\mathrm{h}}=r^2(\omega _\mathrm{A}^2/\mathrm{\Omega })=\frac{B_\phi ^2}{4\pi \rho \mathrm{\Omega }}\mathrm{for}\mathrm{\Omega }\omega _\mathrm{A}.$$
(16)
The value of $`D_{\mathrm{B}_\mathrm{h}}`$ is of the order of $`10^{11}`$ cm$`{}_{}{}^{2}/`$s in the present 15 $`M_{}`$ model. Thus, the magnetic field is large enough to ensure the horizontal coupling, so that the assumption of shellular rotation is still valid.
### 4.3 Diffusion coefficients for the transports of the angular momentum and chemical elements
We need to determine in which regions of the star the magnetic field is present, its intensity and the run of the various diffusion coefficients. As an example, we examine the 15 M model at the beginning of the evolution, when the central H–content is $`X_\mathrm{c}=0.60`$. In this model as in all MS models, $`qq_{\mathrm{min}}`$ in a region which starts just above the core and extends over most of the envelope. In the model with $`X_\mathrm{c}=0.60`$, the magnetic field is present from $`M_\mathrm{r}=6.7`$ M to 13.4 M, as illustrated in Fig. 4. The average field is about $`210^4`$ G and the corresponding value of $`\omega _\mathrm{A}/\mathrm{\Omega }`$ is about $`510^4`$. Above 13.4 M, condition (12) is no longer satisfied. The region where magnetic field is present slightly increases during MS evolution as the convective core recedes. Near the end of the MS phase, when $`X_\mathrm{c}=0.05`$, the magnetic field is present from 4.2 M up to 13.8 M.
The runs of the various diffusion coefficients are illustrated in Fig. 5. The largest diffusion coefficient is $`\nu `$ which acts for the vertical transport of angular momentum, the large value of $`\nu `$ imposes the nearly constant $`\mathrm{\Omega }`$ in the interior. The value of $`\nu `$ is about 6 orders of magnitude larger than the diffusion coefficient $`\eta `$ for the transport of the chemical elements, so that the surface enrichments in products of the CNO cycle due to this magnetic diffusion coefficient only are rather limited. The coefficient $`D_{\mathrm{eff}}`$ which applies to the transport of chemical elements by meridional circulation is much larger than $`D_{\mathrm{shear}}`$, while in rotating stars without magnetic field, it is generally the opposite. This is due to the fact that the velocity of meridional circulation is much larger when $`\mathrm{\Omega }`$ is almost constant throughout the stellar interior (about $`10^2`$ cm s<sup>-1</sup>), while $`D_{\mathrm{shear}}`$ is much smaller. Indeed $`D_{\mathrm{shear}}`$ is about 4 orders of magnitude smaller than in rotating stars without a magnetic field, as a consequence of the near solid body rotation in magnetic models.
Fig. 5 also shows that the value $`\eta /K`$ is always very small, which justifies the simplifications made in deriving Eq. (11). Apart from the slight extension of the zone covered by the magnetic field as mentioned above, the orders of magnitude and relative ratios of the diffusion coefficients remain about the same during the whole MS evolution.
An important result of Fig. 5 is that the transport by meridional circulation (expressed by $`D_{\mathrm{eff}}`$) is in general 2–3 orders of magnitude larger than the transport of the elements by the magnetic field (expressed by $`\eta `$). This clearly shows that the direct effect of the magnetic instability is of little importance with respect to thermal instability, which drives meridional circulation, for the transport of chemical elements. In addition, meridional circulation is important for the transport of the elements from the external edge of the magnetic zone to the stellar surface. However, for the coupling of $`\mathrm{\Omega }`$, the magnetic field is much more efficient than meridional circulation. One can define a velocity $`U_\mathrm{B}`$ for the vertical transport of angular momentum by the magnetic field (cf. paper I)
$$U_\mathrm{B}=5\frac{\nu }{r}\frac{\mathrm{ln}\mathrm{\Omega }}{\mathrm{ln}r}.$$
(17)
This quantity may be compared to the vertical component of the velocity $`U(r)`$ of meridional circulation. In general $`U_\mathrm{B}`$ is much larger than $`U(r)`$, which confirms the dominant role of the magnetic field for the transport of angular momentum.
### 4.4 Internal distribution of hydrogen
Figs. 6 and 7 show the internal distribution of hydrogen in the cases without and with magnetic field. Firstly, rotation with a magnetic field leads to the formation of a larger core at the end of the MS phase compared to models with rotation only (models with rotation themselves have larger cores compared to models without rotation, cf. also Fig. 8).
In models with rotation only, there is a significant erosion of the $`\mu `$–gradient at the immediate edge of the core, developing during the second half of the MS phase. This erosion directly at the edge of the core is absent in models with a magnetic field. The reason is Eq. (10), which gives, if the $`\mu `$–term dominates in the Brunt–Va̋isa̋la̋ frequency ,
$$\eta _0=r^2\mathrm{\Omega }q^4\left(\frac{\mathrm{\Omega }}{N_\mu }\right)^6.$$
(18)
As noted in paper II, the mixing of chemical elements decreases strongly if the $`\mu `$–gradient grows and this effect limits the chemical mixing of elements by the Tayler–Spruit dynamo in the regions just above the convective core.
Finally, in models with a magnetic field the mixing in the envelope is much greater than in the case without a field. This effect appears during the entire MS phase and the final helium surface content reaches $`Y_\mathrm{s}=0.31`$. This effect is not due to the shear, which is essentially absent in magnetic models. It is also not due to the magnetic diffusivity, because it is rather small and in the very external regions there is no magnetic field. However, the transport of chemical elements by the circulation (an effect described by $`D_{\mathrm{eff}}`$) is larger and is generally the main effect in the envelope. It produces a slight transport of the elements, which enriches the stellar surface in elements of the CNO cycle (cf. Fig. 10). The models of paper II with a magnetic field, but without meridional circulation predict no N excesses and this was a difficulty. However, magnetic models, when the meridional circulation is included, lead to a significant N–enrichment, which compares with observations (cf. Sect. 5.2).
## 5 Evolution of observable parameters
The magnetic field created by the dynamo does not reach the stellar surface. However, there are consequences concerning the tracks in the HR diagram and especially the enrichment of surface abundances, which may lead to observable consequences.
### 5.1 Evolution of the core, lifetimes and tracks in the HR diagram
The evolution of the tracks in the HR diagram is determined by the evolution of the core mass fraction. Fig. 8 shows the evolution of the total stellar mass in the magnetic case. When $`X_\mathrm{c}=0.30`$, the actual mass is 14.875 M. Mass loss increases near the end of the MS phase and the mass is 14.389 M when $`X_\mathrm{c}=0.000`$. As is well known, rotational effects slightly increase the core mass fraction near the end of the MS phase and produces longer MS lifetimes. The account of the magnetic field further enlarges the core mass fraction as well as the MS lifetimes as illustrated in Fig. 8.
The effects on the tracks in the HR diagram are shown in Fig. 9. In addition to the structural changes, rotation also produces a distortion of the stellar surface and an increase of the average radius (estimated at $`P_2(\mathrm{cos}\vartheta )=0`$, where $`\vartheta `$ is the colatitude). This makes a redwards shift of the tracks of rotating models in the HR diagram, especially visible near the ZAMS, since no other effects intervene there. For the model with rotation only, we see the upwards shift due to the slightly larger core and a slightly bluer track resulting from a tiny surface He–enrichment. The model with rotation and magnetic field reaches a higher luminosity at the end of the MS due to the larger core. It is also bluer due to the larger surface enrichment in helium, which lowers the opacity and decreases the stellar radius. It is doubtful that the above effects are sufficient to infer the presence or absence of magnetic fields.
### 5.2 Evolution of the abundances of helium and CNO elements at the stellar surface
The evolution of the surface abundances in helium and CNO elements results from the internal profiles shown in Figs. 6 and 7. Fig. 10 shows for the 3 models considered the evolution of the helium abundance $`Y_\mathrm{s}`$ in mass fraction at the surface, the $`N/H`$, $`N/O`$ and $`N/C`$ ratios of numbers of atoms. We see the same trends in the four panels. There is no enrichment in absence of rotation. With rotation only, there are moderate enrichments, by a factor of 2 for $`N/H`$, 2.5 for $`N/C`$ and 1.8 in $`N/O`$. With both rotation and magnetic field, the surface enrichments are much larger, particularly for the helium abundance which reaches $`Y_\mathrm{s}=0.31`$ at the end of the MS phase. The increases in $`N/H`$ reach a factor of 5, 11 for $`N/C`$ and 6 for $`N/O`$.
Observationally, there are many estimates of N enrichments for OB stars. A recent review of the subject has been made by Herrero & Dufton (HerrDuft04 (2004)), who show clear evidence of rotationally–induced mixing in OB stars, also with the result that some fast rotating stars do not show N–enrichments. Four O9 stars were studied by Villamariz et al. (Villamariz02 (2002)); three low rotators have an excess of $`[N/O]<\mathrm{\hspace{0.17em}0.2}`$ dex, one fast rotating O9III star with $`v\mathrm{sin}i=430`$ km s<sup>-1</sup> has an excess of $`[N/O]=\mathrm{\hspace{0.17em}0.7}`$ dex. The two most massive stars with fast rotation in the association Cep 2 have an excess of $`[N/O]<\mathrm{\hspace{0.17em}0.3}`$ dex ( Daflon et al. Daflon01 (2001)). Two stars in Sher 25 have excesses in $`[N/O]`$ by a factor of 3 to 4 (Smartt et al. Smartt02 (2002)). Observations of the $`N/H`$ ratios by Venn and Przybilla (VennPr03 (2003)) for galactic A–F supergiants show an average excess of a factor of 3, with extreme values up to a factor of 8. The orders of magnitude of the predicted and observed enrichments are similar. However, the situation is still uncertain due to the relative lack of accurate observational data for MS stars.
## 6 Conclusions
The main result is that a magnetic field imposes nearly solid body rotation and this favours higher rotational velocities during MS evolution compared to cases where the magnetic field is not accounted for. Due to the nearly constant $`\mathrm{\Omega }`$ in the stellar interior, the transport of chemical elements by shear mixing is negligible. The transport of elements by Tayler–instability is also very limited. An interesting feature of the model is that meridional circulation is strongly enhanced by the flat $`\mathrm{\Omega }`$–curve and this is the main effect influencing the transport of the chemical elements in the present models.
There remain however some doubts as to whether the Tayler–Spruit dynamo is really active in stellar interiors. Up to now, magnetohydrodynamic models expressing the growth and evolution of the magnetic field in rotating stars have not yet confirmed the existence and efficiency of this particular instability in stellar interiors (Mathis Mathis05 (2005)). Even in the observationally and theoretically much better studied case of the Sun, the exact location, origin and evolution of the solar dynamo are still not fully understood.
The larger size of the predicted surface enrichments is well (or even better) supported by observations. However, this remains uncertain in view of the small number of accurate observations. In this respect, it might be crucial to also have such comparisons for stars in the Magellanic Clouds where the observed enrichments are much larger than in the Galaxy. The answer may come from asteroseismology. There are no p–modes expected, however g–modes may be present and yield some information on the internal $`\mathrm{\Omega }`$–distribution. At present, this seems to be the most compelling possible test.
## Appendix A Equations for the case of low rotation
The dynamo properties were established (Spruit Spruit02 (2002) and Sect. 2) for the case of fast rotation, with the condition that the rotation rate $`\mathrm{\Omega }`$ is larger than the Alfvén frequency $`\omega _\mathrm{A}`$, i.e. $`\mathrm{\Omega }\omega _\mathrm{A}`$. In the fast rotating case, the growth rate of the magnetic field is reduced by a factor $`\omega _\mathrm{A}/\mathrm{\Omega }`$ as firstly suggested by Pitts & Tayler (Pitts86 (1986)). The above equations of the dynamo have been derived under this condition. However, if we want to study very slowly rotating stars, we need also to consider the case where
$$\mathrm{\Omega }\omega _\mathrm{A}$$
(19)
This may be interesting also in central regions of radiative stars. There, as a result of the small value of $`r`$, the Alfvén frequency $`\omega _\mathrm{A}`$ becomes large so that the regime corresponding to Eq. (19) may apply instead of the usual one (cf. Sect. 2).
The energy of the Tayler instability (Tayler Tay73 (1973)) must be large enough to overcome the restoring force of buoyancy and this implies that the Alfvén frequency must be larger than some limit depending on $`N`$. The vertical extent $`l_\mathrm{r}`$ of the magnetic instability is also given by Eq. (5) above. In order that the perturbation is not too quickly damped by the diffusion of the magnetic field, one must have $`l_\mathrm{r}^2>\frac{\eta }{\sigma _\mathrm{B}}`$. At low rotation, in the absence of significant Coriolis force, the frequency $`\sigma _\mathrm{B}=\omega _\mathrm{A}`$. The combination of these two limits in the low rotation case leads to the condition $`\omega _\mathrm{A}^3>N^2\frac{\eta }{r^2}`$ and for the marginal situation corresponding to equality, we have
$$\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^3=\frac{N^2}{\mathrm{\Omega }^2}\frac{\eta }{r^2\mathrm{\Omega }}.$$
(20)
This equation relates the magnetic diffusivity $`\eta `$ and the Alfvén frequency $`\omega _\mathrm{A}`$ instead of Eq. (7).
The field amplitude may be fixed by the equality of the amplification time $`\tau _\mathrm{a}`$ of the Tayler instability and of the timescale $`\sigma _\mathrm{B}^1`$ of the field (Spruit Spruit02 (2002)). One has
$$\tau _\mathrm{a}=N/(\omega _\mathrm{A}\mathrm{\Omega }q)$$
(21)
Instead of the expression in (8), we get for the low rotation case with $`\sigma _\mathrm{B}=\omega _\mathrm{A}`$
$$\frac{N}{\mathrm{\Omega }}=q.$$
(22)
$`\omega _\mathrm{A}`$ has disappeared from the equation. How should we interpret this ? This can be done with the help of Eq.(2), which gives
$$\frac{N_\mathrm{T}^2}{\mathrm{\Omega }^2}\frac{\eta /K}{\eta /K+\mathrm{\hspace{0.33em}2}}+\frac{N_\mu ^2}{\mathrm{\Omega }^2}=q^2.$$
(23)
The two equations (20) and (23) form our basic system of equations with the two unknown quantities $`\eta `$ and $`\omega _\mathrm{A}`$. In order to have a diffusivity $`\eta `$ positively defined, Eq. (23) imposes that the differential rotation parameter $`q`$ must be larger than some minimum value $`q_{\mathrm{min}}`$,
$$q_{\mathrm{min}}^2=\frac{N_\mu ^2}{\mathrm{\Omega }^2}.$$
(24)
If this is the case, then according to Eq. (20) the same is true for $`\omega _\mathrm{A}`$ and thus a magnetic field is effectively created by the slow dynamo process.
Let us now estimate the transport coefficients. Eq. (23) immediately leads to the following expression for the magnetic diffusivity $`\eta `$,
$$\eta =\frac{2K\left(q^2\mathrm{\Omega }^2N_\mu ^2\right)}{N_\mathrm{T}^2+N_\mu ^2q^2\mathrm{\Omega }^2}.$$
(25)
In all cases, the ratio $`\eta /K`$ is very small, typically of the order of $`10^5`$ (cf. Fig. 5). Thus, we obtain the corresponding expression for the diffusivity,
$$\eta =\frac{2K}{N_\mathrm{T}^2}\left(q^2\mathrm{\Omega }^2N_\mu ^2\right).$$
(26)
This equation provides the magnetic diffusivity $`\eta `$, if $`|q|`$ is larger than $`|q_{\mathrm{min}}|`$. The coefficient $`\eta `$ determines the transport of the chemical elements by the magnetic instability.
Now, with Eq. (20), we obtain the Alfvén frequency $`\omega _\mathrm{A}`$ at each location $`r`$ in the star,
$$\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^3=\frac{2K}{N_\mathrm{T}^2}\frac{\left(q^2\mathrm{\Omega }^2N_\mu ^2\right)}{r^2\mathrm{\Omega }}q^2.$$
(27)
This confirms that for a star rotating with angular velocity $`\mathrm{\Omega }`$ and having a certain $`\mu `$–gradient, a magnetic field is created only if the actual value of the differential rotation parameter $`|q|`$ is larger than $`|q_{\mathrm{min}}|`$.
The azimuthal stress $`S`$ due to the magnetic field generated by Tayler instability becomes
$`S={\displaystyle \frac{1}{4\pi }}B_\mathrm{r}B_\phi ={\displaystyle \frac{1}{4\pi }}\left({\displaystyle \frac{l_\mathrm{r}}{r}}\right)B_\phi ^2=\rho r^2\left({\displaystyle \frac{\omega _\mathrm{A}^3}{N}}\right).`$ (28)
The viscosity $`\nu `$ is expressed in terms of $`S`$ and we get
$$\nu =\frac{S}{\rho q\mathrm{\Omega }}=\frac{\mathrm{\Omega }r^2}{q}\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^3\left(\frac{\mathrm{\Omega }}{N}\right).$$
(29)
With Eq. (27) for the Alfvén frequency, we get finally
$$\nu =\frac{\mathrm{\Omega }r^2}{q^2}\left(\frac{\omega _\mathrm{A}}{\mathrm{\Omega }}\right)^3=\frac{2K}{N_\mathrm{T}^2}\left(q^2\mathrm{\Omega }^2N_\mu ^2\right),$$
(30)
which determines the transport of angular momentum. The expression of $`\nu `$ is the same as the expression of $`\eta `$ (Eq. 26). This means that in the case of low rotation, the Tayler magnetic instability transports in a similar way the chemical elements and the angular momentum, while at high rotation the magnetic coupling is stronger by a factor $`\left(\frac{\mathrm{\Omega }}{\omega _\mathrm{A}}\right)^2`$ than the transport of the elements, according to the general expressions by Maeder & Meynet (Magn2 (2004)). This factor comes from the expression of $`\sigma _\mathrm{B}`$, and as in the present case the ratio $`\left(\frac{\mathrm{\Omega }}{\omega _\mathrm{A}}\right)`$ is absent in $`\sigma _\mathrm{B}`$, we get equality of $`\nu `$ and $`\eta `$.
Condition (22) must also be satisfied in order to have angular momentum transport. The equations (25) or (26) and (30) provide the transport coeffcients at each stellar layer as a function of the local quantities, such as $`\mathrm{\Omega }`$, $`q`$, $`N_\mathrm{T}^2`$, $`N_\mu ^2`$, etc… available in the stellar models. The Alfvén frequency and the magnetic field intensity are obtained by Eqs. (27) and (1).
The rate of magnetic energy production $`W_\mathrm{B}`$ per unit of time and volume must be equal to the rate $`W_\nu `$ of the dissipation of rotational energy by the magnetic viscosity $`\nu `$ as given above. This check of consistency was verified for the case of high rotation. One assumes here for simplification that all the energy dissipated is converted to magnetic energy. One has
$$W_\nu =\frac{1}{2}\rho \nu \mathrm{\Omega }^2q^2,$$
(31)
which with Eq. (30) gives the dissipation rate,
$$W_\nu =\rho \mathrm{\Omega }^2q^2\frac{K}{N_\mathrm{T}^2}\left(q^2\mathrm{\Omega }^2N_\mu ^2\right).$$
(32)
The magnetic energy per unit volume is $`\frac{B^2}{8\pi }`$, it is produced in a characteristic time given by $`\sigma _\mathrm{B}^1=\omega _\mathrm{A}^1`$. Thus, one has
$$W_\mathrm{B}=\frac{B^2}{8\pi }\omega _\mathrm{A}=\frac{1}{2}\rho r^2\omega _\mathrm{A}^3=\rho q^2\mathrm{\Omega }^2\frac{K}{N_\mathrm{T}^2}\left(q^2\mathrm{\Omega }^2N_\mu ^2\right),$$
(33)
where we have used the Eq. (27) for the Alfvén frequency. Thus, the expressions for $`W_\nu `$ and $`W_\mathrm{B}`$ are the same. We see that no magnetic energy is produced if condition (24) is not realized. This shows the consistency of the field expression for $`B_\phi `$ and of the transport coefficient $`\nu `$.
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# 1 INTRODUCTION
## 1 INTRODUCTION
The success of the $`B`$ Factories has significantly improved our knowledge of the $`CP`$ violation in the quark sector. In particular, the angle $`\beta `$ of the Unitarity Triangle has been measured to a $`5\%`$ accuracy from time-dependent $`CP`$ asymmetries in $`bc\overline{c}s`$ decays . On the other hand, experimental determination of the other two angles and of the lengths of the two sides (with the third side normalized to unit length) have yet to achieve comparable precision. One of the two sides, the one opposite to the angle $`\beta `$, is of particular interest. The uncertainty of this side is dominated by the smallest element $`|V_{ub}|`$, which is known to about 11% precision. Improved determination of $`|V_{ub}|`$ therefore translates directly to a more stringent test of the Standard Model.
Charmless semileptonic decays of the $`B`$ mesons provide the best probe for $`|V_{ub}|`$. Measurements can be done either exclusively or inclusively, i.e., with or without specifying the hadronic final state. Since both approaches suffer from significant theoretical uncertainties, it is important to pursue both types of measurements and test their consistency.
The exclusive $`BX_u\mathrm{}\nu `$ decay rates are related to $`|V_{ub}|`$ through form factors. In the simplest case of $`B\pi \mathrm{}\nu `$, the differential decay rate (assuming massless leptons) is given by
$$\frac{d\mathrm{\Gamma }(B^0\pi ^{}\mathrm{}^+\nu )}{dq^2}=2\frac{d\mathrm{\Gamma }(B^+\pi ^0\mathrm{}^+\nu )}{dq^2}=\frac{G_F^2|V_{ub}|^2}{24\pi ^3}|f_+(q^2)|^2p_\pi ^3,$$
(1)
where $`q^2`$ is the invariant-mass squared of the lepton-neutrino system, and $`p_\pi `$ is the momentum of the pion in the rest frame of the $`B`$ meson. The form factor (FF) $`f_+(q^2)`$ is calculated with a variety of theoretical models. In this paper, we consider recent calculations by Ball and Zwicky based on light-cone sum rules (LCSR) and by the HPQCD and FNAL Collaborations based on unquenched lattice QCD (LQCD). The LCSR and LQCD calculations provide the form factor with reliable uncertainties only in limited ranges of $`q^2`$. It is therefore necessary to extrapolate the calculated form factor using empirical functions, or to measure the partial decay rates $`\mathrm{\Delta }\mathrm{\Gamma }(B\pi \mathrm{}\nu )`$ with appropriate cuts on $`q^2`$, typically chosen as $`q^2<16\mathrm{GeV}^2`$ and $`q^2>16\mathrm{GeV}^2`$ for use with the LCSR and LQCD calculations, respectively.
Measurements of the partial branching fractions $`\mathrm{\Delta }(B\pi \mathrm{}\nu )`$ have been reported by CLEO , Belle , and by BABAR . The CLEO and BABAR measurements reconstructed $`B\pi \mathrm{}\nu `$ events by inferring the neutrino momentum from the missing momentum; the Belle measurement used $`B`$ mesons recoiling against another $`B`$ meson reconstructed in semileptonic decays. BABAR has also reported a measurement of the total branching fraction $`(B\pi \mathrm{}\nu )`$ using the recoil of fully-reconstructed hadronic $`B`$ decays .
In this paper, we report preliminary results from a study of the $`B^0\pi ^{}\mathrm{}^+\nu `$ decay, using an event sample tagged by $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ decays.<sup>4</sup><sup>4</sup>4Charge-conjugate modes are implied throughout this paper. A similar study of the $`B^+\pi ^0\mathrm{}^+\nu `$ decay is reported in a separate paper .
## 2 THE BABAR DETECTOR AND DATA SAMPLES
This measurement uses the $`e^+e^{}`$ colliding-beam data collected with the BABAR detector at the PEP-II storage ring. Charged particles are measured by a combination of five-layer silicon microstrip tracker and a 40-layer central drift chamber, both operating in a 1.5 T magnetic field. A detector of internally reflected Cherenkov light provides charged kaon identification. A CsI(Tl) electromagnetic calorimeter provides photon detection and, combined with the tracking detectors, electron identification. The instrumented flux return of the magnet identifies muons by their penetration through the iron absorber.
The data sample analyzed contains 232 million $`e^+e^{}B\overline{B}`$ events, where $`B\overline{B}`$ stands for $`B^+B^{}`$ or $`B^0\overline{B}^0`$. It corresponds to an integrated luminosity of 211$`\text{ fb}^1`$ on the $`\mathrm{{\rm Y}}(4S)`$ resonance. In addition, a smaller sample (22$`\text{ fb}^1`$) of off-resonance data recorded at approximately $`40\mathrm{MeV}`$ below the resonance is used for background subtraction and validation purposes.
We also use several samples of simulated $`e^+e^{}B\overline{B}`$ events for evaluating the signal and background efficiencies. Charmless semileptonic decays $`BX_u\mathrm{}\nu `$ are simulated as a mixture of exclusive channels ($`X_u=\pi `$, $`\eta `$, $`\eta ^{}`$, $`\rho `$, and $`\omega `$) based on the ISGW2 model and non-resonant $`BX_u\mathrm{}\nu `$ decays with hadronic masses above $`2m_\pi `$. For the signal channels, we give weights to the simulated events in order to reproduce the $`q^2`$ distribution predicted by the recent LCSR and unquenched LQCD calculations as well as the ISGW2 model.
## 3 ANALYSIS METHOD
The analysis method we use for event selection and signal yield extraction has been developed blind, i.e., without using the signal sample in the data. The procedure described in the following sections has been chosen and optimized using MC simulation to obtain the largest expected statistical significance of the partial signal yields in the three $`q^2`$ bins defined in Section 3.2.
The outline of the analysis is as follows: We look for combinations of a $`D^+`$ or $`D^+`$ meson and a lepton ($`e^{}`$ or $`\mu ^{}`$) that are kinematically consistent with $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ decays. For each such $`B`$ candidate, we define the recoil side as the tracks and calorimeter clusters that are not associated with the candidate. We search in the recoil side for a signature of a $`B^0\pi ^{}\mathrm{}^+\nu `$ decay. We take advantage of the simple kinematics of the $`B^0\pi ^{}\mathrm{}^+\nu `$ process to define discriminating variables, and extract the signal yield from their distributions in three bins of $`q^2`$. Finally we calculate the total and the partial branching fractions using the signal efficiencies predicted by a Monte Carlo (MC) simulation. We correct for the data-MC efficiency differences using a control sample in which both $`B`$ mesons decay to tagging modes.
### 3.1 Event Selection
We search for candidate $`B\overline{B}`$ events in which one $`B`$ meson decayed as $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$, where the $`D^+`$ meson is reconstructed in the $`D^+D^0\pi ^+`$ and $`D^+\pi ^0`$ channels. The $`D`$ mesons are reconstructed in the $`D^0K^{}\pi ^+`$, $`K^{}\pi ^+\pi ^{}\pi ^+`$, $`K^{}\pi ^+\pi ^0`$, $`K_S^0\pi ^+\pi ^{}`$, and $`D^+K^{}\pi ^+\pi ^+`$ channels. The widths of the signal regions around the nominal $`D`$-meson masses are between $`\pm 15\mathrm{MeV}`$ and $`\pm 30\mathrm{MeV}`$, which correspond approximately to $`\pm 3\sigma `$ of the mass resolution. We also define sideband regions, evenly split below and above the signal regions, which are used to subtract the combinatorial background. The sideband regions are chosen to be 1.5 times wider than the signal regions. The difference between $`D^+`$ and $`D`$ masses must be within $`\pm 3\mathrm{MeV}`$ of the nominal values.
The $`D^{()+}`$ candidates are combined with an identified electron or muon to form a $`D^{()+}\mathrm{}^{}\overline{\nu }`$ candidate. The lepton must have a center-of-mass momentum<sup>5</sup><sup>5</sup>5Variables denoted with a star ($`x^{}`$) are measured in the $`\mathrm{{\rm Y}}(4S)`$ rest frame; others are in the laboratory frame $`p_{\mathrm{}}^{}>0.8\mathrm{GeV}`$. If the $`D`$ meson was reconstructed with a charged kaon, the kaon charge and the lepton charge must have the same sign.
For each $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ candidate, we remove from the event the tracks and the neutral clusters that make up the $`D^{()+}\mathrm{}^{}\overline{\nu }`$ tag. We then search for a $`B^0\pi ^{}\mathrm{}^+\nu `$ candidate in the remaining part of the event. We require an identified lepton with $`p_{\mathrm{}}^{}>0.8\mathrm{GeV}`$, accompanied by an oppositely charged track that is not identified as either a lepton or a kaon. To allow for the $`B^0`$-$`\overline{B}^0`$ mixing, we do not require the two leptons in a candidate event to be oppositely charged.
The signal events we try to identify contain two neutrinos, one from each $`B`$ decay. Four-momentum conservation and the invariant masses of the $`B`$ mesons and neutrinos provide just the sufficient number of constraints (8) to determine the event kinematics. Referring to the $`D^{()+}\mathrm{}^{}`$ system as the “$`Y`$” system, we first calculate the cosine of $`\theta _{BY}`$, the angle between $`𝐩_B^{}`$ and $`𝐩_Y^{}`$, as
$$\mathrm{cos}\theta _{BY}=\frac{2E_B^{}E_Y^{}m_B^2m_Y^2}{2p_B^{}p_Y^{}}.$$
(2)
The energy $`E_B^{}`$ and momentum $`p_B^{}`$ of the $`B`$ meson are known from the beam energies and the $`B^0`$ mass. Equation (2) assumes that a $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ decay has been correctly reconstructed, and the only undetected particle in the final state is the neutrino. If that is the case, $`\mathrm{cos}\theta _{BY}`$ should be between $`1`$ and $`+1`$ within experimental resolution. If the tag has been incorrectly reconstructed, Equation (2) does not give a cosine of a physical angle, and $`\mathrm{cos}\theta _{BY}`$ is distributed more broadly.
Analogously, we can calculate the cosine of $`\theta _{B\pi \mathrm{}}`$, the angle between $`𝐩_B^{}`$ and $`𝐩_\pi \mathrm{}^{}`$, as
$$\mathrm{cos}\theta _{B\pi \mathrm{}}=\frac{2E_B^{}E_\pi \mathrm{}^{}m_B^2m_\pi \mathrm{}^2}{2p_B^{}p_\pi \mathrm{}^{}}.$$
(3)
This variable, again, should be between $`1`$ and $`+1`$ for the signal events, and distributed broadly for the background.
The momenta of the two $`B`$ mesons must be back-to-back in the center-of-mass frame. Given $`\mathrm{cos}\theta _{BY}`$, $`\mathrm{cos}\theta _{B\pi \mathrm{}}`$, and the directions of $`𝐩_Y^{}`$ and $`𝐩_\pi \mathrm{}^{}`$, we can determine the direction of the $`B`$ momenta up to a two-fold ambiguity. Denoting $`\varphi _B`$ to be the angle between $`𝐩_B^{}`$ and the plane defined by $`𝐩_Y^{}`$ and $`𝐩_\pi \mathrm{}^{}`$, we find
$$\mathrm{cos}^2\varphi _B=\frac{\mathrm{cos}^2\theta _{BY}+\mathrm{cos}^2\theta _{B\pi \mathrm{}}+2\mathrm{cos}\theta _{BY}\mathrm{cos}\theta _{B\pi \mathrm{}}\mathrm{cos}\gamma }{\mathrm{sin}^2\gamma },$$
(4)
where $`\gamma `$ is the angle between $`𝐩_Y^{}`$ and $`𝐩_\pi \mathrm{}^{}`$, as shown in Figure 1.
It should be noted that $`\mathrm{cos}^2\varphi _B`$ satisfies
$$\mathrm{cos}^2\varphi _B\mathrm{cos}^2\theta _{BY}\text{and}\mathrm{cos}^2\varphi _B\mathrm{cos}^2\theta _{B\pi \mathrm{}}$$
(5)
by construction. Equation (4) assumes that $`\mathrm{cos}\theta _{BY}`$ and $`\mathrm{cos}\theta _{B\pi \mathrm{}}`$ are correct, i.e., both the tag-side and the signal-side of the event have been correctly reconstructed with only one neutrino missing on each side. In that case, the variable $`\mathrm{cos}^2\varphi _B`$ must be between 0 and 1. If any part of the reconstruction is incorrect, $`\mathrm{cos}^2\varphi _B`$ is no longer a cosine squared of any physical angle, and its distribution spreads beyond $`+1`$.
We use $`\mathrm{cos}^2\varphi _B`$ as the principal discriminating variable for this analysis. Since $`\mathrm{cos}^2\varphi _B`$ uses the information of the complete event and is strongly correlated with $`\mathrm{cos}\theta _{BY}`$ and $`\mathrm{cos}\theta _{B\pi \mathrm{}}`$, we apply only very loose cuts, $`|\mathrm{cos}\theta _{BY}|<5`$ and $`|\mathrm{cos}\theta _{B\pi \mathrm{}}|<5`$, in the event selection.
If a signal event has been fully reconstructed, no other particles should be present. In reality, such an event often contains extra photons, some of which come from $`\pi ^0`$s and/or photons from decays of $`D^{}`$ or heavier charmed mesons. Although we use the $`D^+D^+\pi ^0`$ decay in the tags, the efficiency for reconstructing the soft $`\pi ^0`$ is low. We do not use the $`D^+D^+\gamma `$ decay due to a poor signal-to-background ratio. We identify the photons that may have come from $`D^{}D\pi ^0`$ and $`D\gamma `$ decays by combining them with the $`D`$ meson candidate; if the combination satisfies $`m_{D\gamma }m_D<150\mathrm{MeV}`$ and $`\mathrm{cos}\theta _{BY^{}}<1.1`$, where $`Y^{}`$ stands for the $`D\gamma \mathrm{}`$ system, the photon is considered as a part of the $`D^{()}\mathrm{}\nu `$ system and is removed from the recoil system. Another source of extra photons is Bremsstrahlung in the detector material. If either lepton is an electron, we identify and remove the Bremsstrahlung photons based on their proximity to the direction of the electron track.
At this point, we require that the event contains no charged tracks besides the $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ and $`B^0\pi ^{}\mathrm{}^+\nu `$ candidates. We further require that there be no residual photons with laboratory-frame energy $`E_\gamma >100\mathrm{MeV}`$.
A few additional cuts are applied to improve the signal-to-background ratio by rejecting specific types of known background events. In order to suppress non-$`B\overline{B}`$ background, we require that the ratio $`R_2`$ of the second and zeroth Fox-Wolfram moments , computed using all charged tracks and neutral clusters in the event, be smaller than 0.5. The invariant mass of the two leptons, if they are oppositely charged, must not satisfy $`2.95<m_{e^+e^{}}<3.15\mathrm{GeV}`$ or $`3.05<m_{\mu ^+\mu ^{}}<3.15\mathrm{GeV}`$ in order to suppress the background due to $`J/\psi \mathrm{}^+\mathrm{}^{}`$ decays. We also calculate the invariant mass of the $`\pi ^+\mathrm{}^{}`$ system assuming that the pion was a lepton of the same species as the identified lepton, and require this mass to be outside 3.06–3.12$`\mathrm{GeV}`$.
After all the selection cuts, a small fraction of events contain more than one candidate; simulated signal events that pass the selection contain 1.14 candidates on average. When there are multiple candidates in an event, we select the candidate with the smallest value of $`|\mathrm{cos}\theta _{BY}|`$. The candidates in the $`D`$-mass sidebands are included in the selection procedure.
### 3.2 Signal Yields
We find in the on-resonance data $`966\pm 31`$ and $`725\pm 34`$ candidate events in the $`D`$-mass signal region before and after the sideband subtraction, respectively. Table 1 summarizes the numbers of candidate events for each decay channel of $`D^{()+}`$.
We extract the signal from the $`\mathrm{cos}^2\varphi _B`$ distribution in three $`q^2`$ bins, namely, $`q^2<8\mathrm{GeV}^2`$, $`8<q^2<16\mathrm{GeV}^2`$, and $`q^2>16\mathrm{GeV}^2`$, as shown in Figure 2. For each event passing the selection, we calculate $`q^2`$ by
$$q^2=(p_{\mathrm{}}+p_\nu )^2(\stackrel{~}{p}_Bp_\pi )^2,$$
(6)
where $`\stackrel{~}{p}_B`$ is the approximate $`B`$ four-momentum defined, in the center-of-mass frame, as
$$\stackrel{~}{p}_B^{}=(\stackrel{~}{E}_B^{},\stackrel{~}{𝐩}_B^{})(\frac{m_{\mathrm{{\rm Y}}(4S)}}{2},0,0,0).$$
(7)
In other words, the center-of-mass motion of the $`B`$ meson is ignored. Since the $`B`$ momentum in the center-of-mass frame is small, the impact of this approximation is small, as will be shown in Section 3.3. It is worth noting that the experimental input to Equation (6) is the pion momentum and not the lepton momentum. As a result, electron energy loss due to unrecovered Bremsstrahlung has no impact on the $`q^2`$ resolution.
The backgrounds in this measurement are handled in three groups. First, the combinatoric background for the $`D`$ mesons is subtracted using the $`D`$-mass sidebands. The remaining backgrounds are mostly $`B\overline{B}`$ events, and are separated from the signal using the $`\mathrm{cos}^2\varphi _B`$ distribution. The possible contribution from the non-$`B\overline{B}`$ background is estimated, based on the $`\mathrm{cos}^2\varphi _B`$ distribution of the off-resonance data events that pass the event selection, to be smaller than 1.0 event of the total yield, and is included in the systematic error.
The raw signal yield is extracted in each $`q^2`$ bin by a simple binned $`\chi ^2`$ fit of the $`\mathrm{cos}^2\varphi _B`$ distribution of the on-resonance data to the weighted sum of the signal and $`B\overline{B}`$ background distributions from the MC simulation. The sources of the $`B\overline{B}`$ background can be $`BX_u\mathrm{}\nu `$ decays, $`BX_c\mathrm{}\nu `$ decays, and other (hadronic) $`B`$ decays. The $`\mathrm{cos}^2\varphi _B`$ distribution for the signal events peaks between 0 and 1, while that of the background is broad with a gradual fall off toward large values of $`\mathrm{cos}^2\varphi _B`$. The fall-off is faster for smaller $`q^2`$, as it can be seen in Figure 2. Since the $`\mathrm{cos}^2\varphi _B`$ distributions for the $`B\overline{B}`$ background from various sources are quite similar, we fix their relative abundances in the fit, and later vary them within their systematic uncertainties. The fit therefore has two free parameters: the normalization of the signal and the normalization of the $`B\overline{B}`$ background.
To maximize the statistical sensitivity, the first bin of the $`\mathrm{cos}^2\varphi _B`$ histogram should be at least as narrow as the signal peak. At the same time, each bin should contain sufficient entries to prevent the $`\chi ^2`$ fit from becoming biased. We use variable bin sizes to satisfy the two requirements. The bin boundaries are chosen to be $`\mathrm{cos}^2\varphi _B=0`$, 1, 2, 4, 7, and 12.
Table 2 summarizes the signal yield obtained from the fit of the on-resonance data, the $`\chi ^2`$ values (for 3 degrees of freedom) and the corresponding probabilities.
The combined $`\chi ^2`$ probability of all three $`q^2`$ bins is 0.033. The poor $`\chi ^2`$ value of the last $`q^2`$ bin is caused by the large number of events in $`1<\mathrm{cos}^2\varphi _B<2`$. Moving the first bin boundary from 1.0 to 1.5 improves the $`\chi ^2`$ and increases the yield significantly. We can find no reason for this behavior other than statistical fluctuation. We therefore retain the result obtained with the original binning, which was chosen before the data were unblinded.
### 3.3 Signal Efficiencies
In order to derive the partial branching fractions from the signal yields, we need the signal efficiency in each $`q^2`$ bin. Since the measured value of $`q^2`$ from Equation (6) in principle differs from the true value, we define the efficiency $`\epsilon _{ij}`$ as the probability, averaged over electrons and muons, of a $`B^0\pi ^{}\mathrm{}^+\nu `$ event whose true $`q^2`$ value belongs to the $`j`$-th bin to be found in the $`i`$-th measured-$`q^2`$ bin. With this definition, the signal yield obtained in the $`i`$-th reconstructed-$`q^2`$ bin is expressed as
$$N_i=2\underset{j}{}\epsilon _{ij}\mathrm{\Delta }_jN_B,$$
(8)
where the factor of two comes from using both electrons and muons, $`\mathrm{\Delta }_j`$ is the partial branching fraction in the $`j`$-th $`q^2`$ bin, and $`N_B`$ is the number of $`B^0`$ mesons in the data sample. Using the $`\mathrm{{\rm Y}}(4S)B^0\overline{B}^0`$ branching fraction $`f_{00}=0.488\pm 0.013`$ , the number $`N_B`$ equals $`2f_{00}N_{B\overline{B}}`$, where $`N_{B\overline{B}}`$ is the number of $`B\overline{B}`$ events in the data sample, and the factor of 2 comes from having two $`B`$ mesons in each event.
We use the Monte Carlo simulation to estimate $`\epsilon _{ij}`$ and correct for the known data-MC differences as
$$N_i=2\frac{\epsilon ^{\mathrm{data}}}{\epsilon ^{\mathrm{MC}}}\underset{j}{}\epsilon _{ij}^{\mathrm{MC}}\mathrm{\Delta }_jN_B.$$
(9)
Here we assumed that a single data-MC efficiency ratio $`\epsilon ^{\mathrm{data}}/\epsilon ^{\mathrm{MC}}`$ can be applied to all $`q^2`$ bins. This is reasonable so long as the ratio is close to unity, which is in fact the case as it will be shown below.
The efficiency determined from the MC simulation is
$$\epsilon _{ij}^{\mathrm{MC}}=\left(\begin{array}{ccc}1.142\pm 0.065& 0.050\pm 0.011& 0.004\pm 0.004\\ 0.074\pm 0.015& 1.232\pm 0.071& 0.035\pm 0.015\\ 0.007\pm 0.008& 0.063\pm 0.016& 1.350\pm 0.097\end{array}\right)\times 10^3.$$
(10)
The errors are due to Monte Carlo statistics. The selection efficiency averaged over the three $`q^2`$ bins is $`1.32\times 10^3`$. In addition to the experimental resolution, final-state radiation (FSR) changes the $`q^2`$ distribution. All MC samples used in this analysis are generated with PHOTOS to simulate FSR. The bin-to-bin migration due to FSR is predicted by the MC simulation to be less than 1.2%, and the full size of this effect is included in the systematic error.
We evaluate the data-MC difference of the $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ selection efficiencies using the double-tag events, in which both $`B`$ mesons decay to $`D^{()\pm }\mathrm{}\nu `$. Properties of the $`D^{()}\mathrm{}\nu `$ tags such as the composition of the $`D^{()}`$ decay channels are similar for the tagged-signal and double-tag events. The number of double-tag events is proportional to the square of the tagging efficiency after subtracting the small contribution from background.
The selection criteria for the double-tag events follow the main analysis as closely as possible. In each event, we look for two $`D^{()}\mathrm{}\nu `$ tags that do not share any particles. We remove all particles that are used in the two tags and require that there be no charged tracks and no neutral clusters remaining in the event.
After subtracting the $`D`$-mass sidebands, we find 1073.4 double-tag events in the on-resonance data. Figure 3 shows the $`\mathrm{cos}^2\varphi _B`$ distribution of the selected events.
The ‘signal’ in this case consists of the $`B^0\overline{B}^0`$ events in which the two $`B`$ mesons decay into $`D^{()}\mathrm{}\nu `$ and are correctly reconstructed as two tags. A small fraction of $`B^0\overline{B}^0`$ events with two $`D^{()}\mathrm{}\nu `$ decays are incorrectly tagged, i.e., wrong combinations of particles are selected as the two tags. Other sources of background events include $`B\tau \mathrm{}`$ cascade decays and lepton misidentification.
We extract from the study of the double-tag events the efficiency correction factor
$$\frac{\epsilon ^{\mathrm{data}}}{\epsilon ^{\mathrm{MC}}}=1.000\pm 0.047.$$
The stated error includes both the statistical and systematic uncertainties. For the latter, we considered a) the difference in the results when the selection criteria are relaxed to allow presence of neutral clusters, b) the residual background after the $`D`$-mass sideband subtraction, due to possible non-linearities in the backgrounds vs. $`D`$ mass, and c) the uncertainties in the exclusive $`BX_c\mathrm{}\nu `$ branching fractions. The relaxed criteria used in a) increase the double-tag yield and the background by approximately 50% and 100%, respectively.
## 4 SYSTEMATIC UNCERTAINTIES
The significant sources of systematic uncertainties and their impact on the measured total and partial branching fractions are summarized in Table 3.
For the $`B\pi \mathrm{}\nu `$ form factor, we use the Ball-Zwicky calculation for our central values, and consider the differences between that and the two LQCD calculations as the systematic uncertainties. The form factor affects the branching fractions through the $`q^2`$ dependence of the signal efficiency. As a result, only the shape and not the normalization of the form factor is relevant at this stage, while the normalization becomes important in the determination of $`|V_{ub}|`$ as discussed in Section 5. The $`B\rho \mathrm{}\nu `$ decays are significant sources of background at large $`q^2`$. We vary the branching fractions as $`(B^0\rho ^{}\mathrm{}^+\nu )=(2.69_{0.77}^{+0.74})\times 10^4`$ and $`(B^+\rho ^0\mathrm{}^+\nu )=(1.45_{0.41}^{+0.40})\times 10^4`$, based on the measurements and isospin symmetry. We also compare results using the $`B\rho \mathrm{}\nu `$ form factors calculated by Ball and Zwicky , Melikhov and Stech , and UKQCD .
The branching fractions for the $`BX_c\mathrm{}\nu `$ and $`BX_u\mathrm{}\nu `$ decays also significantly affect the backgrounds. We use the latest measurements for the branching fractions. Where appropriate, we combine the $`B^0`$ and $`B^+`$ branching fractions assuming isospin symmetry. The shape function parameters used for the simulation of the non-resonant $`BX_u\mathrm{}\nu `$ decay are varied according to Ref. .
The efficiency for the $`D^{()}\mathrm{}\nu `$ tagging has been discussed in Section 3.3. The systematic uncertainties in the track reconstruction and lepton identification efficiencies have been derived from studies of independent control samples. We vary the amount of migration between the $`q^2`$ bins due to resolution, given by the off-diagonal components in Equation (10), by $`\pm 50\%`$. We assign an additional $`\pm 1.2\%`$ error on the partial branching fractions for the $`q^2`$-bin migration due to final-state radiation, which is simulated using PHOTOS .
The largest source of systematic error is the shape of the $`\mathrm{cos}^2\varphi _B`$ distribution for the $`B\overline{B}`$ background events. We studied it using several control samples that are depleted of the $`\pi \mathrm{}\nu `$ signal, obtained, for example, by requiring one extra charged track or neutral cluster remaining in the event. The $`\mathrm{cos}^2\varphi _B`$ distributions of the control samples agree between the data and the MC simulation within the available statistics. We assign the systematic errors based on the statistical uncertainties of these tests.
We used the off-resonance data sample to set an upper limit on the residual non-$`B\overline{B}`$ background. MC simulation was used to determine the size of the combinatoric background that remains after the $`D`$-mass sideband subtraction, due to the non-linearities in the backgrounds vs. reconstructed $`D`$ mass.
The number of $`B\overline{B}`$ events in the on-resonance data sample is known to $`\pm 1.1\%`$. We use $`f_{00}=0.488\pm 0.013`$ as the $`\mathrm{{\rm Y}}(4S)B^0\overline{B}^0`$ branching fraction. Limited statistics of the MC samples affects the measurement primarily through the estimation of the signal efficiency $`\epsilon _{ij}^{\mathrm{MC}}`$ given in Equation (10).
In addition to the studies discussed above, we perform a large number of crosschecks and tests of cut-value dependences. We investigate all cases in which the variations exceed the expected statistical fluctuations, and find no indications of systematic problems.
## 5 RESULTS
From the signal yields and the efficiencies evaluated in Section 3, we extract the following preliminary result for the total branching fraction:
$$(B^0\pi ^{}\mathrm{}^+\nu )=(1.03\pm 0.25_{\mathrm{stat}.}\pm 0.13_{\mathrm{syst}.})\times 10^4.$$
The partial branching fractions and their statistical errors are given in Table 4.
Note that the errors in the partial branching fractions are negatively correlated because of the small migration across $`q^2`$ bins, so adding them in quadrature does not give the error in the total branching fraction. We extracted the branching fractions for each of four signal form-factor calculations (Ball-Zwicky, HPQCD, FNAL, ISGW2); the differences are small.
Figure 4 compares the measured partial branching fractions with the $`q^2`$ dependence predicted by Ball-Zwicky, HPQCD, FNAL, and ISGW2 calculations.
The calculations are normalized to the measured total branching fraction. Although the measured $`\mathrm{\Delta }`$ values depend slightly on the form-factor calculation, the differences are too small ($`<2.5\%`$) to be noticeable on this plot.
Table 5 summarizes the measurements of $`(B\pi \mathrm{}\nu )`$ by the BABAR collaboration.
Assuming isospin symmetry and the ratio of $`B`$ lifetimes $`\tau _{B^+}/\tau _{B^0}=1.081\pm 0.015`$ , the measurements agree with each other with $`\chi ^2=10.3`$ for 5 degrees of freedom, which corresponds to a one-sided probability of 7%.
From the measurement of the partial branching fractions $`\mathrm{\Delta }`$ described in this paper, we extract $`|V_{ub}|`$ using
$$|V_{ub}|=\sqrt{\frac{\mathrm{\Delta }}{\mathrm{\Delta }\zeta \tau _{B^0}}},$$
(11)
where $`\tau _{B^0}=1.536\pm 0.014\mathrm{ps}`$ is the $`B^0`$ lifetime, and $`\mathrm{\Delta }\zeta `$ is defined as
$$\mathrm{\Delta }\zeta =\frac{G_F^2}{24\pi ^3}_{q_{\mathrm{min}}^2}^{q_{\mathrm{max}}^2}|f_+(q^2)|^2p_\pi ^3𝑑q^2.$$
(12)
To minimize the theoretical error on $`|V_{ub}|`$, the range of $`q^2`$ should correspond to the region in which the form-factor calculation is most reliable: $`q^2<16\mathrm{GeV}^2`$ for LCSR and $`q^2>16\mathrm{GeV}^2`$ for LQCD. In order to extract $`|V_{ub}|`$ from the total branching fraction $``$, the form factor must be extrapolated to the full range of $`q^2`$. This is done in Refs. using empirical functions, and additional uncertainties are assigned for the extrapolation.
Table 6 summarizes the values of $`|V_{ub}|`$ extracted from the measured partial and total branching fractions.
The last errors on $`|V_{ub}|`$ come from the uncertainties in $`\mathrm{\Delta }\zeta `$, which in turn come from the uncertainties of the normalization of the form-factor calculations. The precision of the results obtained using the LQCD calculations in $`q^2>16\mathrm{GeV}^2`$ are limited by the large statistical error of the measured partial branching fraction. Using the total branching fraction reduces experimental errors on $`|V_{ub}|`$ at the cost of increased theoretical uncertainties due to the extrapolation of the form factor.
Instead of averaging results based on different theoretical calculations, we report the value of $`|V_{ub}|`$ obtained from the total branching fraction based on one of the LQCD calculations ,
$$|V_{ub}|=(3.3\pm 0.4_{\mathrm{stat}.}\pm 0.2_{\mathrm{syst}.}{}_{0.4}{}^{+0.8}{}_{\mathrm{FF}}{}^{})\times 10^3.$$
as a representative result. This result lies between the results based on the other two calculations, and carries the most conservative theoretical uncertainty.
## 6 SUMMARY
Using event samples tagged by $`\overline{B}{}_{}{}^{0}D^{()+}\mathrm{}^{}\overline{\nu }`$ decays, we obtain the exclusive branching fraction $`(B^0\pi ^{}\mathrm{}^+\nu )`$. The preliminary result for the total branching fraction is
$$(B^0\pi ^{}\mathrm{}^+\nu )=(1.03\pm 0.25_{\mathrm{stat}.}\pm 0.13_{\mathrm{syst}.})\times 10^4.$$
We also obtain the partial branching fractions in three bins of $`q^2`$ of the lepton-neutrino pair. The preliminary results are
$$\mathrm{\Delta }(B^0\pi ^{}\mathrm{}^+\nu )=\{\begin{array}{cc}(0.48\pm 0.17_{\mathrm{stat}.}{}_{0.05}{}^{+0.06}{}_{\mathrm{syst}.}{}^{})\times 10^4\hfill & q^2<8\mathrm{GeV}^2,\hfill \\ (0.34\pm 0.16_{\mathrm{stat}.}{}_{0.07}{}^{+0.06}{}_{\mathrm{syst}.}{}^{})\times 10^4\hfill & 8<q^2<16\mathrm{GeV}^2,\hfill \\ (0.21\pm 0.14_{\mathrm{stat}.}{}_{0.06}{}^{+0.05}{}_{\mathrm{syst}.}{}^{})\times 10^4\hfill & q^2>16\mathrm{GeV}^2.\hfill \end{array}$$
Using the measured total branching fraction and a form-factor calculation based on lattice QCD , we extract
$$|V_{ub}|=(3.3\pm 0.4_{\mathrm{stat}.}\pm 0.2_{\mathrm{syst}.}{}_{0.4}{}^{+0.8}{}_{\mathrm{FF}}{}^{})\times 10^3.$$
where the last error is due to normalization of the form-factor. We also use other recent calculations of the form factor and find values of $`|V_{ub}|`$, given in Table 6, that are consistent within the experimental and theoretical uncertainties.
## 7 ACKNOWLEDGMENTS
We are grateful for the extraordinary contributions of our PEP-II colleagues in achieving the excellent luminosity and machine conditions that have made this work possible. The success of this project also relies critically on the expertise and dedication of the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and the kind hospitality extended to them. This work is supported by the US Department of Energy and National Science Foundation, the Natural Sciences and Engineering Research Council (Canada), Institute of High Energy Physics (China), the Commissariat à l’Energie Atomique and Institut National de Physique Nucléaire et de Physique des Particules (France), the Bundesministerium für Bildung und Forschung and Deutsche Forschungsgemeinschaft (Germany), the Istituto Nazionale di Fisica Nucleare (Italy), the Foundation for Fundamental Research on Matter (The Netherlands), the Research Council of Norway, the Ministry of Science and Technology of the Russian Federation, and the Particle Physics and Astronomy Research Council (United Kingdom). Individuals have received support from CONACyT (Mexico), the A. P. Sloan Foundation, the Research Corporation, and the Alexander von Humboldt Foundation.
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# Two 2MASS-Selected Young Stellar Clusters: Photometry, Spectroscopy, and the IMF
## 1. INTRODUCTION
Despite their intrinsic rarity and short lifetimes, massive stars are extremely important in the evolution of galaxies. They play an important role in determining the course of the formation of less massive stars, though the nature of this role is still uncertain, and their stellar winds and eventual supernovae shape the interstellar medium. They produce most of the heavy elements in the universe, as well as much of the UV radiation in galaxies. Their rarity, combined with the effects of large Galactic extinctions, often results in the availability of more comprehensive studies of massive stars in external galaxies, where the entire stellar population can be observed at once, than within our own where massive stars must be studied individually and the census of massive stars is still very incomplete. High optical extinction within the galactic plane ($`A_V20`$) has limited optical studies of massive stars to relatively nearby regions (R$`{}_{\mathrm{solar}}{}^{}3.0`$ kpc, Massey, 2003). Even within that radius, optically selected catalogs of O stars have been found to be incomplete, especially in star-forming regions and young clusters (e.g. Hanson & Conti, 1995). This incompleteness necessitates the use of infrared, radio and X-ray observations, particularly in the inner regions of the Galaxy and in star formation regions. The near-infrared (NIR, 1-5 µm) is an especially useful regime for the study of massive stars; the stellar atmosphere is still observed directly, but since for example $`A_K0.11A_V`$, we can observe these stars in regions where dust, either along the line of sight or local to the star-forming region, makes them inaccessible at optical wavelengths. The discovery and characterization of stellar clusters observable only in the infrared can significantly enhance our understanding of obscured Galactic regions which harbor embedded massive stars or massive protostars.
Recent studies indicate that clusters may account for 70-90% of star formation and that embedded clusters (those still partially or fully enshrouded in their natal molecular cloud) may exceed the number of more traditional open clusters by a factor of $``$20 (Elmegreen et al., 2000; Lada & Lada, 2003). In the last decade, advancements in NIR observational capabilities resulted in the discovery and classification of some of the most massive young stellar clusters in the Galaxy, each containing dozens of O and WR stars (e.g. Nagata et al., 1995; Cotera et al., 1996; Figer, Morris, & McLean, 1996). Recent studies (Figer et al., 1999) have suggested that within these clusters, the initial mass function (IMF) does not follow the canonical Salpeter form with a slope $`\mathrm{\Gamma }=1.35`$, but instead is more heavily weighted toward massive stars; mass segregation has been proposed as a solution (Stolte et al., 2002). In the last several years a number of studies of well-known star formation regions have also been carried out in the NIR, (e.g Okumura et al., 2000; Blum, Damineli, & Conti, 2001; Conti & Blum, 2002; Figuerêdo et al., 2002). These studies have in most cases found an IMF consistent with the Salpeter value, and have uncovered candidate massive YSOs. In addition, within the past ten years, massive YSOs within molecular clouds have been studied in the NIR, (e.g. Chakraborty et al., 2000; Ishii et al., 2001) and in young stellar clusters (e.g. Hanson, Hayworth, & Conti, 1997). Massive YSOs, however, remain significantly less studied and are poorly understood in comparison with their lower-mass counterparts; many more must be identified and studied before we can adequately address how the formation of massive stars differs from that of low-mass stars.
The final release of the Two Micron All Sky Survey (2MASS) has fostered studies which can probe the entire Galaxy for previously unknown stellar clusters. Initial attempts were made which searched for stellar density enhancements, (e.g. Dutra & Bica, 2000, 2001; Dutra et al., 2003), but the identification of previously unknown clusters has met with limited success. For example, Dutra & Bica (2000) identified 52 candidate clusters, which subsequent observations (Dutra et al., 2003) indicated were in fact 10 confirmed clusters, 3 “probable” clusters, and 11 “dissolving cluster candidates”; the remainder were not clusters. Our observations of at least one of the Dutra et al. (2003) “confirmed clusters”, however, indicates that the “cluster” is most likely a region of low extinction rather than a true cluster (Cotera & Leistra, 2005). We have performed an independent search of the 2MASS archive, using color criteria in addition to stellar density enhancements. We have searched in the vicinity of regions identified as likely sites of star formation based on radio and IRAS far-infrared flux ratios, and are currently conducting a search of the entire 2MASS Point Source Catalog. We search the Point Source Catalog for regions of higher stellar density than the background (determined locally within a 5′radius) which are redder in $`HK`$ than the local field. This selects for embedded clusters, with the color criteria helping to eliminate chance superpositions and regions of low extinction. In contrast, Dutra & Bica (2000,2001) use only stellar density to select clusters. Our method has been relatively successful to date; correctly selecting 7 clusters out of 9 potential targets, including 4 candidates toward the inner Galaxy. We present NIR imaging and spectroscopy of the two confirmed clusters in the inner Galaxy in this paper, and discuss the two unconfirmed targets in detail in Cotera & Leistra (2005). The cluster near G305.3+0.2 was independently discovered by Dutra et al. (2003b). The additional 5 outer-galaxy targets are described in Paper II.
NIR imaging and spectroscopy of both young stellar clusters and nascent stellar clusters enables us to expand the study of the IMF in objects where there has been little to no stellar evolution off the main sequence or cluster evaporation, and where the cluster age can be constrained to within $`2`$ Myr. Spectral typing of the most massive stars in the cluster allows their masses to be determined relatively precisely, and when combined with photometry it facilitates a reliable determination of the masses of stars throughout the entire cluster (Massey, Johnson, & DeGioia-Eastwood 1995; Massey 2002), allowing the initial mass function of the cluster to be determined more accurately than photometry alone would permit. In this paper we present the results of NIR observations of two clusters found toward the inner Galaxy, which we designate by the Galactic coordinates of their centers, G353.4-0.36 (17:30:28 -34:41:36 J2000) and G305.3+0.2 (13:11:39.6 -62:33:13 J2000). In Paper II we will present the results of similar observations of five clusters in the outer Galaxy.
In §2 we present the observations and data reduction, in §3 we present the spectra and classifications of the spectroscopically observed cluster members as well as the color-magnitude diagrams, and in §4 we describe the luminosity function and the initial mass function.
## 2. OBSERVATIONS & DATA REDUCTION
We observed candidate young stellar clusters with the facility instrument IRIS2 on the 3.9m Anglo-Australian Telescope (AAT) on July 12-15, 2003. IRIS2 is an imaging spectrometer which uses a 1024x1024 Rockwell HAWAII-1 HgCdTe array with a platescale of 0$`\stackrel{}{\mathrm{.}}`$45/pixel, resulting in a 7$`\stackrel{}{\mathrm{.}}`$7$`\times `$7$`\stackrel{}{\mathrm{.}}`$7 field of view. Images were obtained in $`J`$ (1.25 µm), $`H`$ (1.63 µm), and $`K_s`$ (2.14 µm) filters. $`R2300`$ spectra of selected stars in each cluster candidate were obtained in $`K`$ for all candidates.
We selected a total of four cluster candidates in the southern hemisphere using the 2MASS Point Source Catalog based on color and density criteria. Two of the candidates observed appear to be regions of low extinction and are discussed elsewhere (Cotera & Leistra, 2005). The two confirmed clusters are near radio H II regions designated G305.3+00.2 and G353.4-0.4. We present three-color composites of the 8′$`\times `$8′ images of the G305.3+00.2 and G353.4-0.36 clusters in Figures 1 and 2 respectively.
G305.3+00.2 is an H II region which has been previously observed using radio recombination lines (Wilson & Mezger, 1970), C I emission in the submillimeter (Huang et al., 1999), and in the mid-infrared (MIR) by the Midcourse Space Experiment (MSX). The kinematic distance of $`3.5\pm 1.1`$ kpc obtained for this H II region (Wilson & Mezger, 1970) agrees well with the distance of 3.3 kpc for masers several arcminutes away (Caswell et al., 1995), suggesting they may be part of a single star-formation complex. A distance of 4 kpc is adopted as an upper limit to the radio kinematic distance by Clark & Porter (2004) in a study of the star clusters Danks 1 and 2 in this region. The situation is more complex for the G353.4-0.36 cluster, which is in a region known to be a site of massive star formation. There are numerous radio sources located within 1′ of the NIR cluster, which we discuss in detail in §3.2.2.
All photometric observations were done in excellent seeing conditions: 0$`\stackrel{}{\mathrm{.}}`$7-0$`\stackrel{}{\mathrm{.}}`$9. The images were reduced and combined automatically at the telescope using the ORAC-DR pipeline. ORAC-DR is a generic data reduction pipeline created at the Joint Astronomy Centre in Hawaii, originally for use with various UKIRT and JCMT instruments. Subsequent reprocessing did not noticeably improve the images, therefore the pipeline processed data has been used throughout. Source detection, PSF fitting, and photometry was carried out using IRAF-DAOPHOT, and is discussed in detail in §3.2.
All spectra were obtained with a 1″$`\times `$ 7$`\stackrel{}{\mathrm{.}}`$7 slit. The long-slit format combined with the high stellar density within the FOV resulted in the simultaneous observation of multiple stars. Total integration times ranged from 10 minutes to 30 minutes, and were chosen to provide adequate S/N for NIR spectral classification as described in Hanson, Conti, & Rieke (1996). After the data was flat-fielded, grism curvature was removed using the FIGARO<sup>1</sup><sup>1</sup>1FIGARO is part of the Starlink software package available at http://star-www.rl.ac.uk/ tasks cdist and sdist. Wavelength calibration was performed using the the OH<sup>-</sup> night sky lines and the FIGARO task arc. The uncertainty in the wavelength calibration fit was determined to be 2.18 Å. The FIGARO task irflux was used both to flux-calibrate the spectra and remove the telluric absorption using the G2V standards HD157017 and HD115496. Both of the standards had intrinsic Br $`\gamma `$ in absorption, with equivalent widths of 5.7 Å for HD157017 and of 5.6 Å for HD115496; in each case, the absorption line was removed by fitting a line to the continuum in the region of the line in the standard star spectrum prior to flux calibration. The individual spectra were obtained by extracting apertures 4-5 pixels wide from the full spectral array, then performing background subtraction using apertures of the same width on either side of the source, separated by 2 pixels (0$`\stackrel{}{\mathrm{.}}`$9). We also extracted off-source spectra in each cluster to characterize any nebular emission.
## 3. Analysis
### 3.1. Spectroscopy
The development of NIR spectral atlases of nearby massive stars of known spectral type (Hanson et al. 1996; Morris & Serabyn 1996; Blum et al. 1997), provides a valuable classification scheme for stars too heavily obscured by dust to permit optical spectroscopy. In the $`K`$ band, in addition to the Br$`\gamma `$ (2.165 µm) line, massive O stars have helium (He I 2.058 µm, He I 2.112 µm, He II 2.189 µm), carbon (C IV 2.078 µm), and nitrogen (N III 2.116 µm) lines in their spectra which allow for the determination of the spectral type to within a subtype if there is adequate ($`70`$) line signal to noise. Table 6 of Hanson et al. (1996) indicates that in many cases the mere presence of these lines in emission or absorption (without considering equivalent width) is sufficient to determine spectral type to within two subtypes for O stars. The situation is more complicated for B stars, which have fewer features in this part of the spectrum; however, they are still classifiable using only $`K`$-band spectra.
We obtained $`K`$-band spectra of five stars in the G305.3+0.2 cluster field and three stars in the G353.4-0.36 cluster. In order to reduce the level of foreground contamination, we imposed a color cut of $`HK>0.5`$ based on the 2MASS magnitudes and selected the brightest stars meeting this requirement. Despite this cutoff, two of the five stars observed in the G305.3+0.2 cluster proved to be foreground contaminants with sufficient line-of-sight extinction to push them over our threshold. The cluster sequence was much narrower and more well-separated from the foreground in the G353.4-0.36 cluster, and no obvious foreground contaminants were present in our spectroscopic sample. The G353.4-0.36 cluster was sufficiently red ($`HK_{cluster}1.3`$), that the time required to obtain a useful signal-to-noise in $`H`$-band spectra would have been prohibitively large, so only $`K`$-band were obtained.
#### 3.1.1 G305.3+0.2 Cluster
We present spectra for the three cluster members, which we label A1–A3, in Figure 3.
In Figure 4 we present a 106″$`\times `$ 120″ image of the cluster and label the positions of sources A1–A3.
The measured magnitudes (see §3.2) and observed spectral lines for A1–A3 are presented in Table 1. The other two stars for which we obtained high S/N spectra have late-type spectra, as indicated by strong CO absorption at 2.29 and 2.32 µm, suggesting they are either foreground objects or YSOs. The lack of nebular emission in the cluster and the presence of weak (nearly the same as in the G2V spectral standard) Br $`\gamma `$ absorption in one of the spectra suggest that these are foreground objects rather than YSOs. In addition, the $`K`$ magnitudes of these objects ($`K=9.43`$ and $`K=10.114`$) make them too bright to be low-mass YSOs at the cluster distance, and the presence of main-sequence O and B stars argues against identifying these objects as massive YSOs. We thus conclude that these two stars are most likely late-type foreground stars, and excluded them from further analysis.
Although nebular emission can be seen in the full image (Figure 1), it is significantly removed ($``$1′) from the cluster. Nevertheless, in order to ensure that any measured Br$`\gamma `$ (2.166 µm) is stellar in origin and not contaminated by nebular emission within the cluster, we extracted a local background spectrum for the cluster. There were no features apparent in the resulting spectrum; we thus conclude that nebular emission within the cluster is negligible. This conclusion is supported by an apparent bubble of MIR emission seen in the MSX Band A image (see Figure 5); the MIR emission avoids the cluster itself.
Figure 3 shows that source A1 has emission lines with equivalent widths stronger than -2 Å at 2.116 µm and 2.166 µm (see Table 1). The line at 2.166 µm is immediately identifiable as Br$`\gamma `$. We identify the line at 2.116µm as N III, which is consistent with the lines used in the the classification system presented in Hanson et al. (1996); the broad nature of this line is due to the multiplet nature of the transition responsible rather than broadening by stellar winds. The presence of Br$`\gamma `$ and N III 2.116µm in emission, without further information and without equivalent widths, is sufficient to identify the star as being an early to middle O supergiant; the broad Br$`\gamma `$, produced in the stellar winds, is not observed in main-sequence O stars (Hanson et al., 1996). There is a possible weak detection ($`2\sigma `$) of C IV in emission at 2.078 µm. This line only appears in O stars ranging from O5 to O6.5 (Hanson et al., 1996), and if real, significantly constrains the stellar type. Helium lines are often observed both in emission and absorption in the spectra of massive stars: He I (2.058 µm), He I (2.112 µm), and He II (2.189 µm), are all absent from the spectrum of A1. Poor removal of the telluric features near the 2.058 µmfeature prevents us from drawing any conclusions based on our non-detection. If real, the absence of the He I (2.112 µm) line restricts the spectral type to O6 or earlier. A He II line is expected in an O star; by estimating the strength of possible features dominated by the noise (as described in detail in § 3.1.2) we can place an upper limit of 0.5 Å on the equivalent width of any potential He II (2.188 µm) feature. This is consistent with the width of the feature in the stars observed by Hanson et al. (1996), so the non-detection does not rule out an O star identification for this source. Taken together, these spectral characteristics suggest a spectral type of O5Ib-O6Ib for Source A1. If the weak detection of C IV is discounted, the presence of the N III line and the limit on an He II line at 2.188 µm allows an O7-O8 identification as well. Even when present, however, the C IV line is weak, with an equivalent width weaker than -2 Å; thus, while a positive detection of this line would allow for definitive classification of this source as an O5Ib-O6Ib star, a non-detection at the given S/N does not preclude the same classification.
The intrinsic NIR colors of O and B stars range from -0.08 to -0.01 (Wegner, 1994); this small range allows an extinction to be derived even without knowing the precise spectral type of a massive star. For source A1, the extinction thus derived based on the observed $`HK`$ color is $`A_V=12`$ assuming the extinction law of Rieke & Lebofsky (1985). However, the large range in absolute $`M_K`$ for O supergiants prevents us from making a distance determination based on Source A1. We can only say the distance is greater than $`3.3`$ kpc, which would be the distance for a main-sequence O5-O6 star. Clark & Porter (2004) adopt a distance of 4 kpc to the Danks 1 and 2 clusters in the same star formation complex, calling it an upper limit to the values allowed by the radio and H$`\alpha `$ observations, and we will follow suit, acknowledging that the uncertainties in this value are $`0.5`$ kpc.
Source A2 shows a strong Br$`\gamma `$ (2.166 µm) line in absorption with an equivalent width of $`6.2\pm 1.2`$ Å and a probable weak He I (2.112 µm) line in absorption with EW = $`0.7\pm 0.2`$ Å. This combination of features occurs only in B stars; a comparison of the equivalent width of the lines with the B stars of Hanson et al. (1996) suggests an spectral type in the range of B2-B4. If the He I line is considered only as an upper limit, the classification becomes more problematic, and the star could range from B2-A2. The star has $`HK=0.68`$, which for any star in this range of spectral type excludes a foreground object. Unlike for A1, the luminosity class of these sources cannot be determined from these spectral features; as Hanson et al. (1996) points out, the $`K`$-band spectra of early B supergiants are indistinguishable from those of early B main-sequence stars about half the time, and those of late-B supergiants cannot be distinguished from early-B dwarfs.
If we assume that A2 is a cluster star, we can constrain the absolute magnitude, and thus the spectral type, by requiring the distance to be the same as for the O star. Since the intrinsic near-infrared colors vary by less than 0.1 magnitude for stars in the range of spectral types allowed by the spectrum (Wegner, 1994), we can derive a extinction for this source rather than use that derived from the O star, thus reducing the effects of differential extinction. This gives an extinction to source A2 of $`A_V=11.6`$, or $`A_K=1.3`$ using the reddening law of Rieke & Lebofsky (1985). At the distance of $`4`$ kpc, we obtain an absolute magnitude for Source A2 of $`M_K=4.0`$, roughly that expected for an O8V star. This identification is not consistent with the spectral features of A2; a smaller distance, or an identification of A2 as an early-B supergiant, could explain the spectrum of A2. If the radio distance of $`3.3\pm 0.3`$ kpc is used instead, we obtain an absolute magnitude of $`M_K=3.3`$ for Source A2, making it a B0V-B1V.
Source A3 shows only Br$`\gamma `$ in absorption with an EW of $`5.9\pm 1.3`$ Å. We place an upper limit on an He I absorption line at 2.112 µm of 0.6 Å. As discussed above, this width for Br$`\gamma `$ only constrain the classification of the star as main sequence B or early A. The observed $`K`$ magnitude is $`11.96`$, which corresponds to an absolute $`M_K2.6`$ assuming the extinction and distance of an O5Ib-O6Ib star for source A1; this is consistent with an identification of A3 as a main-sequence B1V star. The radio distance would imply a B2V identification, also consistent with the spectral features of A3. Source A3 is not among the brightest stars in the cluster region; it happened to fall in the same long slit as one of the foreground contaminants we had targeted for observation. This suggests that the other cluster members brighter than A3 are also late O or early B stars.
#### 3.1.2 G353.4-0.36 Cluster
Spectra for the three sources observed in this cluster are presented in Fig. 6.
An enlarged version of the relevant portion of Fig. 2 is presented in Fig. 7, with the positions of the spectroscopic targets indicated with arrows and labels.
The only non-nebular feature which we detect is CO absorption in Source B1; the Br $`\gamma `$ emission observed in all three spectra is contaminated by nebular emission to such a degree that we cannot disentangle any stellar component that may be present. While this line is much stronger in B1 than in the other two sources, the nebular emission is highly spatially variable in the cluster region and this does not demonstrate a stellar origin for the line. Additionally, the line width is significantly narrower than that of Source A1 and is similar to that observed in the off-source nebular spectrum (Fig. 8).
The CO absorption in Source B1 in combination with the red colors (Table 3) are similar to those associated with solar-mass young stellar objects (YSOs) (Greene & Lada, 1996), or a cool giant or supergiant. If B1 is a YSO, the CO absorption is from the circumstellar material; otherwise it is photospheric in nature. Using the radio kinematic distance (Forster & Caswell, 2000) of 3.6 kpc to the cluster, we derive an $`M_K`$ for Source B1 of -0.8 without correcting for extinction. Correcting for extinction is difficult to do accurately in this region of highly variable extinction, especially when the intrinsic colors are not known since the nature of the object is uncertain. Nevertheless, limits can be placed on the amount of extinction present, and thus the absolute magnitude of Source B1. The lower limit is given by the uncorrected value of $`M_K=0.8`$, which assumes the color observed is the intrinsic color, while the bright limit can be derived assuming an intrinsic $`HK=0.3`$, characteristic of late-type stars; this gives an extinction to source B1 of $`A_V=16.6`$ magnitudes and an extinction-corrected absolute $`M_K`$ of $`2.6`$. This is several magnitudes brighter than the expected magnitude of YSOs of approximately a solar mass at the distance and extinction of this cluster, $`M_K13`$ (Oasa, Tamura, & Sugitani, 1999), and somewhat lower than the $`M_K`$ for massive YSOs, $`M_K1`$ to $`5`$ (Ishii et al., 2001). Finally, we note that this $`M_K`$ is consistent with that for a $`7M_{}`$ YSO (Chakraborty et al., 2000). We conclude that if Source B1 is a YSO, it has a mass greater than a few solar masses based on its absolute magnitude in $`K`$, but observations of more massive YSOs are still sufficiently few that a more accurate mass determination based solely on the absolute magnitude is not possible. Given the nebular emission, seen as He I (2.058 µm), H<sub>2</sub> (2.12 µm), and Br$`\gamma `$ (2.166 µm) emission off the stellar sources (see Figure 8), G353.4-0.36 is obviously a region of current star formation; therefore, the identification as a massive YSO is more probable than a late type cool giant or supergiant located in the cluster itself.
Since Source B1 was not detected in $`J`$, it cannot be placed on a color-color diagram to determine whether a NIR excess is present, which could help to discriminate between the YSO and cool field star possibilities. For B1 to be a cool giant, it would need to be a foreground star with the appropriate color and magnitude, which falls by chance in the cluster region. Rather than use the entire 8′$`\times `$ 8′ field to determine the field star density, as we did for the G305.3+0.2 cluster (§3.2), we used only the heavily extincted region surrounding the cluster. This is because the molecular cloud in which the cluster is embedded extinguishes the background stars to such a degree that using the entire field would significantly overestimate the level of field star contamination in the immediate region of the cluster. We estimate the probability of a field source as bright as or brighter than Source B1 and red enough to satisfy the color cut falling within the cluster region to be approximately 18%. This is a conservative estimate, since at the edges of the cloud reddened sources become visible and increase the field star density, especially of red objects, over what it would be at the location of the cluster. Nevertheless, we cannot rule out either a foreground giant or a YSO explanation for Source B1.
As with Source B1, the non-detection of Sources B2 and B3 in $`J`$ prevents us from using a color-color diagram to measure NIR excess. No photospheric features are detected in the spectra of either Source B2 or B3; Source B2 shows a rising spectrum in $`K`$ suggesting a strong NIR excess, while the spectrum of B3 is essentially flat in this region. In order to determine whether the spectra were truly featureless or merely had a signal-to-noise too low to see expected features, we fit a continuum to the spectra and examined all excursions above and below the fit. 90% of these deviations had an equivalent width less than 1.7 Å. For comparison, the detected absorption lines tabulated by Greene & Lada (1996) for low-mass YSOs range in equivalent width from 0.3-5.6 Å for Na I and Ca I, with CO usually exceeding 2 Å when present. Ishii et al. (2001) conducted a similar survey of massive YSOs; the only emission lines other than Br $`\gamma `$ detected in a significant number of sources are CO (with an equivalent width exceeding 4 Å) and H<sub>2</sub> (with an EW $`>3`$ Å in all cases, and $`>5`$ Å in most cases). We thus conclude that Source B2 is genuinely featureless, but cannot classify it. The final source, Source B3, had no reliably detected features but the signal-to-noise was low enough that we cannot reliably call it featureless.
The observed $`K`$ magnitudes are consistent with a B star identification for sources B2 and B3; however, the extincted but distance-corrected $`M_K`$ magnitudes of $`0.2`$ to $`0.6`$ are also similar to those observed for the massive YSO ($`M7M_{}`$) 05361+3539 (Chakraborty et al., 2000). Thus, although these sources are massive, we cannot distinguish based on their NIR spectra or magnitudes between shrouded B stars and less-evolved YSOs. Mid-IR observations with sufficient resolution to resolve the individual sources (separated by $`5`$″) would aid in this determination; deeper J-band photometry, detecting more of the cluster stars, would also be useful. We note that although we see ionized gas suggesting the presence of O stars, we have not detected any O stars which would be the source of the ionizing radiation in this cluster.
Due to the young age of the sources observed in this cluster and the lack of photospheric features in their spectra, the spectra were unsuitable for determining a reliable distance. Thus, the kinematic distance to the associated maser and UCHII (Forster & Caswell, 2000) was used instead, adjusted to a distance to the Galactic Center of 8 kpc from the original 10 kpc. This gave a distance to the cluster of 3.6 kpc. Assuming an intrinsic $`HK=0`$, we estimate the reddening to the cluster at $`A_V=22`$ based on the narrow cluster sequence at $`HK1.3`$ and assuming the extinction law of Rieke & Lebofsky (1985). This estimate is highly uncertain due to the young age of the sources; many are likely to have a near-infrared excess leading to an overestimate of the line-of-sight extinction to the cluster.
### 3.2. Photometry
We obtained images in $`J`$, $`H`$, and $`K_s`$ of both clusters to a limiting magnitude of approximately $`J=16`$, $`H=18`$, $`K_s=18.5`$, with total integration times of 12 minutes in each band. The limiting magnitude was brighter than expected due to confusion, which is most noticeable in $`J`$ due to the slightly larger PSF and the greater sensitivity of the instrument at shorter wavelengths. Seeing was 0$`\stackrel{}{\mathrm{.}}`$7–0$`\stackrel{}{\mathrm{.}}`$8, which, since the IRIS2 platescale is 0$`\stackrel{}{\mathrm{.}}`$45/pixel, resulted in a slight undersampling of the point spread function (PSF), thus making PSF fitting more uncertain. Our individual images were taken using a random dither pattern with sub-pixel dithers employed to improve the PSF. In an effort to better understand our errors we performed both PSF fitting and aperture photometry for each source. There was no systematic offset between the two methods, but the errors were $`2`$ times larger for the aperture photometry due to the crowded fields.
Photometric calibration was performed using the 2MASS magnitudes of field stars, after correcting from the IRIS2 filter system to the 2MASS filter system as described in Carpenter (2003). The calibrated magnitudes for the stars in the cluster area are presented in Table 2. The large field of view and location in the Galactic Plane provided over 100 stars in each pointing which were bright enough to have good photometry with 2MASS, but faint enough to be unsaturated in our IRIS2 images $`(11.5<K_s<14)`$. Those stars which were relatively isolated in the IRIS2 images were used as the photometric calibration set. We chose to use a relatively large number of calibration stars rather than selecting the few most isolated stars to reduce effects of potential variability and photometric outliers among the calibration stars. The scatter in the photometric calibration derived from comparison to 2MASS is the dominant source of photometric error, contributing two to three times the measurement errors as reported by DAOPHOT. DAOPHOT errors were $`0.03`$ mag while the calibration uncertainties were ($`\mathrm{\Delta }J=\pm 0.05,\mathrm{\Delta }H=\pm 0.06`$, and $`\mathrm{\Delta }K=\pm 0.06`$ mag). Quoted errors in the 2MASS photometry were negligible, with most stars having an error of $`\pm 0.003`$ mag or less in all bands. Thus, the quoted error should be considered an overestimate when considering the *relative* photometry of stars within either cluster; the calibration errors from comparison to the 2MASS photometry will shift all our measurements by the same amount. No trend in the photometric errors, either internally or relative to the 2MASS data, was observed with location. Finally, the positions of the stars were also adjusted to agree with 2MASS by minimizing the offsets between the 2MASS and IRIS2 positions allowing for pointing offset and rotation.
#### 3.2.1 G305.3+0.2
The color composite of the full $`J`$, $`H`$, and $`K_s`$ images is presented in Fig. 1; the cluster alone is shown in Fig. 4, with the spectroscopic targets marked.
The cluster is clearly visible in the full-size image with a concentration of nebular emission to the northwest. In order to help determine whether the nebular emission is physically associated with the cluster, we overplotted the contours at 8 µm from the MSX mission<sup>2</sup><sup>2</sup>2On-line data are available from http://www.ipac.caltech.edu/ipac/msx/msx.html. (Fig. 5). The ridge of near-IR nebulosity corresponds to the brightest portion of a roughly circular structure of mid-IR emission, with the cluster located in the interior where there is no mid-IR emission present. The general appearance is that of a wind-blown bubble, and the 8 µm emission wraps entirely around the cluster at a lower level. The cluster is located off-center in this structure, near the brightest portion of the mid-IR emission, but there is no mid-IR emission and no near-IR nebulosity present in the area of the cluster itself. The cluster is dense and well-defined, with stellar density much higher than in the field.
The $`K`$ versus $`HK`$ color-magnitude diagram of the cluster region is shown in Figure 9.
At radii of approximately 30″ in the east-west direction and 20″ in the north-south direction from the cluster center the stellar density has fallen to that of the field, which we used to define the cluster region. Foreground stars are apparent in the color-magnitude diagram at $`HK0.3`$; in this cluster there is no clear separation in color between cluster and field stars, just an overdensity of redder stars in the cluster; as a result, we cannot impose a firm color cut to separate field stars from cluster stars. A color-magnitude diagram of a randomly selected control field with the same area as the cluster is shown in Fig. 10; many fewer stars are present, especially at bright magnitudes and moderately red colors.
In order to account for field star contamination within the cluster region, we determined the average number of stars per square arcminute in the image outside the cluster region in color-magnitude bins of $`\mathrm{\Delta }K=0.5`$, $`\mathrm{\Delta }(HK)=0.5`$ and randomly selected the appropriate number of stars from the cluster field for removal. This is similar to the procedure employed by, among others, Blum, Conti, & Damineli (2000) and Figuerêdo et al. (2002). In cases where less than one star was expected in the cluster field in a particular color-magnitude bin, the number expected was used as a probability for removing a star. A total of 24 “field” stars were removed, leaving 115. The main concentration of cluster stars is at about $`HK=0.8`$, with a gradually declining number present out to $`HK4`$. The resulting cluster CMD with the field stars statistically removed is shown in Fig. 11.
Given the spectroscopically confirmed presence of OB stars in the cluster, as well as the lack of an obvious color gap, we consider it more likely that these very red sources are either background sources or sources with a near-IR excess due to local dust than that they represent a separate cluster giant branch. The red sources are not concentrated toward any part of the cluster, though they may occur more frequently on the outskirts (as would be expected if they are background objects). Sources redder than $`HK=1.5`$ were excluded from analysis of the cluster KLF and IMF; they are unlikely to be main-sequence cluster members. If they are included and assumed to be on the main sequence, the resulting extinction correction would give very large values for the masses and an overly flat slope to the IMF. If these sources are cluster members, they are pre-main-sequence objects, and their masses are difficult to determine from $`H`$ and $`K`$ photometry alone. Thus, including them in the IMF determination would give an inaccurate result whether or not they are cluster members, and they have been excluded. Finally, the crowded nature of the cluster region means that these very red sources may suffer from poor photometry.
The $`JH`$ vs. $`HK`$ color-color diagram (Fig. 12) is of limited utility in identifying cluster members or determining whether some cluster members are pre-main-sequence objects.
Since many sources were undetected in $`J`$, it will not represent all cluster members, and faint red sources (where we would expect to find the relatively low-mass, pre-main-sequence objects) would be most commonly missed in the color-color diagram. A cut based solely on $`HK`$ must still be applied to exclude background sources. Fig. 12 shows few sources in the area occupied by pre-main-sequence objects. Of those sources separated from reddened main-sequence stars by more than 3$`\sigma `$, three are relatively faint sources adjacent to bright sources and one is in a particularly crowded region. The remaining three could potentially be pre-main-sequence objects. However, due to the lack of observed gaseous emission from the cluster, we consider it unlikely that these are truly pre-main-sequence stars, and exclude them from the analysis along with the objects in the unphysical blue region of the color-color diagram as likely suffering from blending or a mismatch between sources in the different bandpasses. There are few enough sources in this region that we do not expect their inclusion or exclusion to greatly affect the IMF determination.
#### 3.2.2 G353.4-0.36
The $`J`$, $`H`$, and $`K`$ color composite of G353.4-0.36 is presented in Fig. 2. The youth of this cluster is immediately apparent from its heavily embedded nature and the dense molecular cloud that surrounds it. This region has long been known to be a site of massive star formation, and it has been studied extensively in the radio and sub-mm, including continuum observations at 1.5 GHz, 5 GHz (Becker et al., 1994), and 850 µm (Carey et al., 2000) as well as molecular line observations in CS (Gardner & Whiteoak, 1978), CO (Whiteoak, Otrupcek, & Rennie, 1982), H<sub>2</sub>CO (Gardner & Whiteoak, 1978), HNCO (Zinchenko, Henkel, & Mao, 2000) (identified as a dense molecular core), and SiO (Harju et al., 1998). These signatures of ongoing star formation, combined with the strong nebular emission still present around the sources observed spectroscopically, suggest that the cluster is quite young, without main-sequence stars. Many of the continuum and molecular line observations quote slightly different positions for the source peak, and sources separated by several tens of arcseconds are all identified with the IRAS point source 17271-3439. Since the beam sizes in many instances are comparable to the size of the NIR-bright nebulosity and to the separation between sources, it is likely that the extended source measurements are observing the same complex, which may peak at different locations in different wavelengths. Many of the radio data are tabulated by Chan, Henning, & Schreyer (1996), who identify a massive YSO in the region based on the IRAS colors. It is obvious from the NIR imaging that this source is not a single point source; in addition to the NIR sources, there are at least four separate sets of masers (e.g. Caswell et al., 2000; Argon et al., 2000; Val’tts et al., 2000), one of which is associated with an UCHII (Forster & Caswell, 2000). Positions of the masers are indicated in Fig. 7. We note that the masers occur in regions which are heavily extincted in the near-IR. OH, H<sub>2</sub>O, and CH<sub>3</sub>OH masers are all known in the region; the latter in particular are indicative of ongoing massive star formation. Clearly the sources visible in the near-infrared are only the tip of the iceberg, with other massive stars still in the process of formation. Higher-resolution maps at radio and sub-mm wavelengths are necessary to obtain a full understanding of this region.
In the region of the large dark molecular cloud, only foreground stars are visible. This implies $`A_V>50`$ in order to completely obscure the stars even in $`K`$, assuming a $`K`$-band detection limit of 17 and a distribution of $`K`$ magnitudes similar to the rest of the field. The less heavily extincted region in which the cluster is visible in the near-IR must have been partially cleared out by stellar winds and ionization from massive stars. The relative position of the NIR stars and the methanol masers (which lie in regions of higher extinction) suggest that we are observing stars nearer the main sequence which are emerging from the dust, while objects at an earlier evolutionary state are offset from this region, indicating ongoing star formation.
The color-magnitude diagram of the G353.4-0.36 cluster is presented in Figure 13.
The cluster sequence is much narrower and more well-separated than in the G305.3+0.2 cluster, allowing for reliable separation of foreground objects based solely on $`HK`$. Thus, we did not carry out a statistical removal of foreground objects for this cluster, instead considering only the objects well-separated from the foreground sequence. Due to the high extinction toward this cluster, a large number of objects in the cluster area were detected only in $`K`$ (shown as limits in Figure 13). The KLF is thus likely to be more reliable than the color-magnitude and color-color diagrams in determining cluster characteristics.
## 4. The K Band Luminosity Function and the Initial Mass Function
Once field stars have been rejected as described in §3.2.1, we can compute the KLF for both clusters. For the G305.3+0.2 cluster, which has more than 100 stars remaining, we additionally compute the initial mass function (IMF) using two different techniques, the first using the KLF and the second using the color-magnitude diagram and the spectroscopy of the massive stars. The KLF is commonly used to determine the IMF even when multi-color photometry is available; we take this opportunity to test the robustness of this method and compare the results between this simple and commonly used method and the more involved method using the color-magnitude diagram. This will help to understand the uncertainties and systematic errors that may be a factor when only the KLF method can be used to derive an IMF. There were too few stars to robustly compute the IMF for the G353-0.4 cluster, so we compute only the KLF in this case.
### 4.1. The G305.3+0.2 Cluster
To provide a robust determination of the KLF and the IMF, we must determine the completeness of our data, which we established by performing artificial star tests. Five artificial stars at a time were inserted into the cluster region; the small number was chosen to avoid significantly changing the crowding characteristics. IRAF-DAOPHOT was then run on the images to determine the number of artificial stars that were successfully recovered. The procedure was repeated 50 times for each magnitude bin ($`\mathrm{\Delta }m=0.5`$), for a total of 250 artificial stars added in each bin in $`H`$ and in $`K`$. Figure 14 shows the results; completeness falls sharply to about 25% at $`H16.5`$, $`K15.5`$. We can compare these magnitudes with the turnover in the “field luminosity function”, which also probes incompleteness. The counts in the field turned over sharply at $`K16`$, in reasonable agreement with the artificial star estimate of incompleteness.
#### 4.1.1 The $`K`$ Luminosity Function
Knowing our incompleteness, we can calculate the KLF for the cluster. Figure 15 shows the uncorrected data, with the field “luminosity function” normalized to the same total number of stars overplotted for comparison. Figure 16 shows the results after correcting for incompleteness by dividing the number of stars in each magnitude bin by the recovered fraction of artificial stars.
As expected, there is an overabundance of bright stars ($`K<14.5`$) in the cluster region relative to the field. This is not an artifact of incompleteness; the completeness fraction at this magnitude is $`90\%`$, and we expect incompleteness to be higher in the cluster than the field due to the effects of crowding. Using the number counts corrected for field star contamination (as discussed in §3.2.1) and incompleteness, we fit a slope to the number counts in bins of $`\mathrm{\Delta }K=0.5`$. We excluded sources fainter than $`K=15.5`$ from the fit since errors in the incompleteness determination are likely to dominate the number counts. We derived a slope of $`0.21\pm 0.06`$ for log $`N_{}`$. This slope is somewhat flatter than the KLFs derived for more massive embedded clusters (e.g. $`0.41\pm 0.02`$ for NGC 3576 from Figuerêdo et al. (2002), $`0.40\pm 0.03`$ for W42 from Blum et al. (2000)). This suggests that this cluster is more weighted toward massive stars than the norm.
#### 4.1.2 The Initial Mass Function
In order to better compare our results with the literature, and to explore how much of a difference the use of multi-color photometry and spectra of the massive stars make in the determination of the IMF, we used two methods to derive an IMF for the G305+00.2 cluster. For both methods we use a distance to the cluster of 4.0 kpc (as discussed in §3.1.1). The first IMF-determination method, which uses only the KLF, is commonly employed even when multi-color photometry and spectra are available (e.g. Figuerêdo et al., 2002; Blum et al., 2000). This method is simply a transformation from $`K`$ magnitude bins to mass bins. To make this transformation, we first correct the observed $`K`$ for distance and extinction as discussed in §3.1.1. Using the stellar evolutionary models of Meynet & Maeder (2003) for solar metallicity, we relate the mass for each track to an absolute $`K`$ magnitude for a star on the ZAMS. We transformed $`L_{bol}`$ to $`K`$ using the bolometric corrections from Vacca, Garmany, & Shull (1996) for the early spectral types and Malagnini et al. (1986) for later spectral types. We then use the intrinsic $`VK`$ colors from Bessell & Brett (1988) for A-M stars and from Wegner (1994) for O and B stars. Finally we interpolate linearly between the masses available on the evolutionary tracks to find the masses corresponding to our magnitude bins, and fit a power law to the resulting mass function. Our resulting IMF slope is $`\mathrm{\Gamma }=1.5\pm 0.3`$, excluding the two lowest-mass bins where incompleteness is significant.
Our second method of determining the IMF made use of our multi-color photometry and spectra to estimate individual extinctions and masses for cluster members. Spectral typing of the brightest cluster stars allows their mass to be determined fairly accurately for a given stellar evolutionary model. For the models described above, the mass of an O6V star is approximately 40 M, that of a B0V star is 15 M, and that of a B2V star is 8 M. Although spectra are not available for most of the cluster stars, their masses, as well as extinctions to the individual stars, can be estimated from the accurate relative photometry. The presence of an O supergiant in the cluster suggests that, while the most massive stars have begun to evolve away from the main sequence, none have yet gone supernova, and less massive stars should still be on the zero-age main sequence. Therefore, with the exception of the few most massive stars (for which we can estimate masses from their spectral types) the cluster stars should be scattered around the zero-age main sequence (ZAMS) primarily by differential extinction and rather than the effects of stellar evolution. We can then use the same models and conversions from theoretical to observed quantities described for the KLF method, with additional transformations from $`T_{eff}`$ to $`HK`$ using intrinsic colors from from Bessell & Brett (1988) and Wegner (1994) and from $`T_{eff}`$ to spectral type from Repolust et al. (2004) or Johnson (1966).
This transformation from theoretical to observed quantities allows us to place the ZAMS on our CMD. If the cluster is sufficiently young that we can neglect the effects of stellar evolution, as discussed in the previous paragraph, we expect the ZAMS will lie in the middle of the distribution of cluster stars. The ZAMS derived from the evolutionary tracks of Meynet & Maeder (2003) is overplotted on the distance and extinction-corrected CMD in Figure 11. A significant number of stars are bluer than the ZAMS on this plot. We interpret these as stars which are less extincted than those used to determine the average cluster extinction and thus have been over-corrected by using the mean extinction. The scatter of stars around the ZAMS suggests that the extinction varies across the cluster region. To correct for this, we move the stars along the direction of the reddening vector until they lie on the ZAMS. If the resulting extinction differs from the mean cluster value by more than $`A_V=5`$ for a given star, we exclude the star from the analysis, as it probably suffers from poor photometry. Examination of the color image of the cluster region (Figure 1) suggests that the variation in internal extinction in this region is relatively small; no dust lanes or color variations across the cluster are visible to the eye. The exact value selected for the cutoff is somewhat arbitrary, but does not greatly affect the results; most of the sources thus excluded have derived extinctions that differ from the median value by $`A_V=10`$ or more.
Using the positions of the extinction-corrected photometry along the ZAMS, we are able to more accurately place stars in mass bins. The endpoints of the bins were determined by the masses for which theoretical tracks are present in the models we used. In order to have an adequate number of stars in each bin we constructed bins using alternate tracks for the endpoints, rather than every track. The analysis was repeated for three different metallicities ($`Z=0.1,0.02,0.001`$) using the evolutionary tracks of Mowlavi et al. (1998, Z = 0.1), Schaller et al. (1992, Z=0.02); Meynet & Maeder (2003, Z=0.02) and Schaller et al. (1992, Z=0.001). For the solar-metallicity case the high-mass points ($`M>9M_{}`$) are from Meynet & Maeder (2003) while the lower-mass points are from Schaller et al. (1992). The difference in $`K`$ for the two solar-metallicity tracks is always less than 0.1 magnitudes for the masses where the two sets of tracks overlap and for most masses is less than 0.03 magnitudes. The high metallicity model should be considered only as a limiting case since such a high metallicity is not expected. The use of such a wide range of metallicities allows us to estimate the importance of this parameter on the final IMF determination.
Given these sets of mass bins, for each metallicity we determine the number of stars per unit logarithmic mass interval after correcting for completeness. We then fit a power law to the data. The two lowest-mass bins ($`M<2M_{}`$), where incompleteness was significant, were excluded from the fit; uncertainty in the completeness correction applied could significantly influence the results in these mass bins. The resulting completeness-corrected IMF for the cluster is plotted in Figure 17. The solar-metallicity models yield an IMF slope $`\mathrm{\Gamma }=0.98\pm 0.2`$, where the quoted errors are only the formal fit errors and should be considered an underestimate. The low-metallicity tracks yield $`\mathrm{\Gamma }=1.01\pm 0.2`$ for the same distance, suggesting that the cluster IMF determination is insensitive to metallicity for solar and sub-solar values. The $`Z=0.1`$ tracks give $`\mathrm{\Gamma }=0.88\pm 0.15`$.
#### 4.1.3 Comparison of the IMF Methods
The IMF slopes we derive using these two methods are marginally consistent within the error bars: $`\mathrm{\Gamma }=1.5\pm 0.3`$ for the KLF method, and $`\mathrm{\Gamma }=0.98\pm 0.2`$ for the CMD + spectroscopy method assuming solar metallicity. Comparing these results individually to the Salpeter slope would lead to different conclusions, however. The KLF method produces a slope that is very close to the Salpeter value, while the slope from the CMD + spectroscopy method differs from Salpeter by about 2$`\sigma `$. While this difference in slopes could arise purely from statistical uncertainty, various systematic effects should cause the KLF-derived slope to be steeper than the CMD-derived slope, as we observe. If the more massive stars are preferentially located toward the center of the cluster, as expected due to mass segregation, and if the extinction is higher in the center of the cluster, the mean extinction used in the KLF determination would be systematically low for the more massive stars. This method would then underestimate the masses the highest mass stars, thus steepening the slope of the IMF. Evidence that this effect may be at work is provided by the six brightest cluster members, all of which lie redward of the ZAMS in Figure 11 while the fainter members are scattered more evenly. A difference in $`A_K`$ (and thus $`M_K`$) of 0.2 corresponds to 1-2 subtypes for massive stars and thus to a difference in the derived mass of at least 2 M.
However, mass segregation can only provide a partial explanation for the difference in the IMF slopes; the stars for which we obtained spectra are not in the very center of the cluster (since crowding in the 2MASS image used to select spectroscopic targets prevented us from selecting targets in the cluster core). An additional possible source of systematic error in the KLF method relative to the CMD method lies in field star rejection. In addition to the statistical field star rejection described in §3.2.1, which was done before any further analysis and thus applies to both methods, the CMD method has color-based field star rejection. The CMD method can reject foreground objects, which due to lower extinction are bluer than cluster objects, as well as background objects which are redder than the cluster. The KLF method includes these objects, which tend to be fainter on average than the cluster stars (since they are either at a greater distance or are low-mass foreground stars) and thus finds an artificially high number of low-mass stars. We find the use of $`K`$ photometry alone to derive the IMF is likely to produce an overly steep IMF in regions with significant field contamination or variable extinction.
#### 4.1.4 Comparison With Other Young Stellar Clusters
Most studies of young star clusters have found an initial mass function consistent with the Salpeter slope of $`\mathrm{\Gamma }=1.35`$, generally with uncertainties of 0.1-0.2 (e.g. Figuerêdo et al., 2002; Massey & Hunter, 1998; Hillenbrand & Carpenter, 2000; Okumura et al., 2000), including the extremely massive R136 cluster in the LMC (Massey & Hunter, 1998). A review of the results is provided in Massey (2003). In the case of NGC 6611, reanalysis of the same data by different authors has produced dramatically different results; an IMF of $`1.1\pm 0.1`$ was found by Hillenbrand et al. (1993), while a reanalysis with different treatment of extinction produced $`0.7\pm 0.2`$ (Massey et al., 1995), suggesting that the systematic effects are at work in IMF determinations that are at least as important as the statistical errors, as we see in this work. Slopes significantly flatter than Salpeter have been reported for the Arches cluster near the Galactic Center (Figer et al., 1999), though later work suggests that this result is an artifact of mass segregation; Stolte et al. (2002) found a very flat IMF in the core of the Arches Cluster with a steeper IMF at larger radii, with an overall slope consistent with a Salpeter value. The flatness we observe in both the KLF and the IMF for the G305+00.2 cluster using the CMD + spectroscopy method may similarly be due to mass segregation. In addition to the extinction effects mentioned previously, fainter stars in the outskirts of the cluster could be indistinguishable from the field star density (especially given the high field star density due to the location of the cluster in the Galactic plane) and not fall within the cluster boundaries we employ.
### 4.2. The KLF for the G353.4-0.36 cluster
Completeness tests were performed for the G353.4-0.36 cluster using artificial stars as discussed above, and the completeness-corrected KLF is plotted in Figure 18. Since the cluster is significantly less crowded and faint cluster stars less common, our detections in this cluster are nearly complete in K, even though our detection limit is brighter than in the G305.3+0.2 cluster. The turnover at $`K=15.5`$ appears to be genuine rather than an artifact of completeness. Perhaps lower-mass stars in this cluster are still more deeply embedded in the gas and dust, and thus we observe only the massive objects.
Due to the small number of stars detected in this cluster ($`N=25`$, only 7 of which were detected in $`H`$) and to the early evolutionary stage of the objects, we did not attempt to determine an IMF for this cluster or to place objects on the ZAMS. While the individual objects we observed in the G353.4-0.36 cluster were intriguing and worthy of further study, we cannot analyze the cluster as a whole because there are so few objects.
This cluster is a very promising target for study at other wavelengths more suited than the NIR to the study of YSOs and even earlier stages of star formation; the methanol masers and likely presence of massive YSOs suggest that several stages of massive star formation can be studied in this region.
## 5. Summary
We present NIR images and spectroscopy of two young stellar clusters near radio sources G353.4-0.36 and G305+00.2. Our $`K`$-band spectrum of the brightest cluster star in the G305+00.2 cluster show it to be an O5Ib-O6Ib star. Although the range of luminosities of supergiants prevents us from determining an exact distance, this identification suggests a larger distance than radio distance to the nearby methanol masers (Walsh et al., 1997) of 3.3 kpc. We also obtained spectra of early two B stars in the cluster. There was no nebular emission present in the G305+00.2 cluster, though a ridge of nebular emission, coinciding with 8 µm emission and masers, is present $``$ 1′ away and may indicate sequential star formation, with the masers and gas indicating ongoing star formation and the cluster the result of earlier star formation. We computed the KLF and IMF of this cluster, and found them to be steeper than that reported for most young clusters ($`\mathrm{\Gamma }=0.98\pm 0.2`$ for the more reliable CMD-based method) but generally consistent with the Salpeter value. We find that computing the IMF based only on a single color of photometry is prone to systematic errors when differential extinction and field-star contamination are significant.
Our $`K`$-band spectra of two of the three stars we observed in the G353.4-0.36 cluster were featureless, while the other showed CO absorption, which is consistent either with a cool foreground giant or a YSO. The absolute magnitudes derived based on the distance to the radio sources are too bright for these objects to be solar-mass YSOs. None of the objects were detected in our $`J`$-band photometry, making identification as YSOs based on NIR excess impossible. They remain candidate massive YSOs, and observations at other wavelengths are needed to make a positive identification. The images of this cluster showed a region with intense nebular emission embedded in a very dark cloud where earlier stages of star formation are progressing.
We thank the AAT and Chris Tinney for assistance with the IRIS2 instrument. We thank Phil Massey, Margaret Hanson, and the anonymous referee for comments that improved this paper. This research has made use of the SIMBAD database, operated at CDS, Strasbourg, France. This publication makes use of data products from the Two Micron All Sky Survey (2MASS), which is a joint project of the University of Massachusetts and IPAC, funded by NASA and NSF. A.C. was supported in part by NASA through the American Astronomical Society’s Small Research Grant Program.
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# 1 Introduction
## 1 Introduction
The electromagnetic response of negative refractive index materials (NIM) $`^{\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}}`$ has recently attracted unprecedented attention. Novel optical phenomena, predicted to take place in these unique systems include reversal of Snell Law, Doppler Effect, Cherenkov Effect,$`^\text{?}`$ aberration-free $`^{\text{?, ?, ?}}`$ and sub-diffraction$`^{\text{?, }\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}}`$ imaging, and excitation of the new types of surface and nonlinear waves .$`^{\text{?, }\text{?}}`$ In particular, realization of NIMs may potentially lead to fabrication of new types of lenses and prisms,$`^{\text{?, ?, ?}}`$ new lithographic techniques,$`^{\text{?, ?}}`$ novel radars, sensors, and telecommunication systems. However, despite the great advantages NIM has to offer for optical and infrared spectral range, all practical realizations of NIM are currently limited to GHz frequencies.$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}}`$
Until recently there were two major approaches for NIM design. The first one is based on the original proposal$`^\text{?}`$ that material with simultaneously negative dielectric permittivity and magnetic permeability must have a negative refraction index. This particular approach also benefits from the possibility to resonantly excite the plasmon polariton waves at the interface between NIM and surrounding media, which in turn may lead to sub-diffraction imaging.$`^{\text{?, ?, ?, ?, ?, ?, }\text{?}}`$ However, the absence of natural magnetism at high (optical or infrared) frequencies$`^\text{?}`$ requires the design and fabrication of nanostructured meta-materials, to achieve the non-trivial magnetic permeability.$`^{\text{?, }\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}}`$ As these engineered systems typically operate in close proximity to resonance, resonant losses become the dominant factor in system response, severely limiting the practicality of resonant-based systems.$`^{\text{?, ?, ?, ?, }\text{?}}`$
The second approach for NIM design involves the use of photonic crystals.$`^{\text{?, }\text{?}\text{}\text{?}\text{, ?, ?}}`$ However, the NIM response in these systems is typically associated with second or other higher-order bands and requires a complete bandgap between the band in use and the next band. The dispersion and very existence of the required bandgap are typically strongly affected by crystal disorder, unavoidable during the fabrication step. The manufacturing of the optical photonic crystals-based NIM typically requires 3D patterning with 10-$`nm`$ – accuracy, which is beyond the capabilities of modern technology.
To address the mentioned-above shortcomings of the traditional NIM schemes, we have recently introduced an alternative approach to design the NIM structure.$`^\text{?}`$ In contrast to “conventional” systems, the proposed design does not rely on either magnetism or periodicity to achieve negative refraction response. It has been shown that the combination of strong anisotropy of the dielectric constant and planar waveguide geometry yields the required negative phase velocity in the system.$`^\text{?}`$ Here we present the detailed description of NIMs proposed in Refs. \[?, ?\], study the effects related to waveguide boundaries, important for optical domain, and suggest several nanostructured materials providing the low-loss negative refraction response at optical and infrared frequencies.
The rest of the paper is organized as follows: the next Section is devoted to EM wave propagation in strongly anisotropic waveguides; Section 3 describes the proposed realizations of the structure; imagining properties of these composites are shown in Section 4; Section 5 concludes the paper.
## 2 Negative refraction in strongly anisotropic waveguides
We consider wave propagation in the $`2D`$ planar waveguide structure shown in Fig. 1. The propagation in the system is allowed in the $`y`$ and $`z`$ directions, while the waveguide walls occupy the regions $`|x|>d/2`$. The waveguide core is assumed to be a homogeneous, non-magnetic ($`\mu =1`$) material, with a uniaxial anisotropic dielectric constant with dielectric permittivities $`ϵ_{}`$ and $`ϵ_{}`$ along and perpendicular to the optical axis respectively. The optical axis of the core material ($`C`$) is assumed to be perpendicular to the direction of the wave propagation in the media ($`Cx`$). Therefore, despite the anisotropy of the system, the effective refractive index of propagating in the planar geometry waves will be completely isotropic.
Any wave propagating in the system can be represented as a linear combination of the waveguide modes.$`^{\text{?, ?}}`$ An individual mode is defined by its structure along the optical axis direction ($`C`$) and its polarization. Two different kinds of modes have to be distinguished. The modes of the first kind (known as TE waves) have their $`E`$ vector perpendicular to the optical axis. The propagation of these waves is fully described by the in-plane dielectric constant $`ϵ_{}`$. The modes of the second kind (known as TM waves) have their $`H`$ vector in the waveguide plane and are affected by both $`ϵ_{}`$ and $`ϵ_{}`$. The existence of these TM waves is crucial for the NIM described here.
In the analytical results presented below we limit ourselves to the case of single-mode propagation. We note that such a description provides complete information about the linear properties of the waveguide structure. Indeed, as mentioned above, an arbitrary wavepacket in the system can be represented as a linear combination of modes. In our numerical simulations discussed in Section 4 we utilize this property to compute the imaging performance of the system.
### 2.A Waveguide with perfectly conducting walls
As it has been shown in Ref. \[?\], the propagation of a mode in a planar waveguide can be described by the free-space-like dispersion relation:
$$k_y^2+k_z^2=ϵ\nu \mathrm{k}^2,$$
(1)
where $`ϵ`$ is $`ϵ_{}`$ for TE modes and $`ϵ_{}`$ for TM ones, $`k_y`$ and $`k_z`$ are the propagation components of the wavevector, $`\mathrm{k}=\omega /c`$ with $`\omega `$ and $`c`$ being the free-space angular frequency of the radiation, and speed of light in a vacuum; the propagation constant $`\nu `$ is given by
$$\nu =1\frac{\kappa ^2}{ϵ_{}\mathrm{k}^2},$$
(2)
and the parameter $`\kappa `$ defines the mode structure in $`x`$ direction.
As it directly follows from Eq. (1), the phase velocity of a propagating mode is equal to
$$v_p=n\mathrm{k},$$
(3)
where the effective refraction index $`n^2=ϵ\nu `$. Note that similar to the case of the plane wave propagation in free-space, the refraction index contains a product of two (mode-specific) scalar constants. A transparent structure must have both propagation constants of the same sign. The case of positive $`ϵ`$ and $`\nu `$ corresponds to “conventional” (positive refraction index) material. The case of negative $`ϵ`$ and $`\nu `$ describes NIM.$`^{\text{?, ?}}`$ The NIM behavior can be easily illustrated by comparing the Poynting vector $`S_z`$ and the wavevector $`k_z`$ as shown below.
Similar to any waveguide structure, the mode in the system described here can be related to the $`x`$ profile of the longitudinal field component (the detailed description of such a dependence is given in Ref. \[?\]). To better illustrate the physical picture behind the mode propagation, in this section we present the results for the important case of perfectly conducting waveguide walls. In this case, the EM energy is confined to the waveguide core and the longitudinal field has a $`\mathrm{cos}(\kappa x)`$ or $`\mathrm{sin}(\kappa x)`$ profile depending on the symmetry with respect to the $`x=0`$ plane, with $`\kappa =(2j+1)\pi /d`$ for symmetric and $`\kappa =2\pi j/d`$ for anti-symmetric modes respectively with $`j`$ being the integer mode number. The deviation from this idealized picture due to finite conductance of the waveguide material does not change the physical picture described in this section, and for the practical case of “good” metals (Ag,Al,Au) at near-IR to THz frequencies can be treated perturbatively. Results of such a perturbation approach are presented in the Section 2b.
The electric ($`U_E`$) and magnetic ($`U_H`$) field contribution to the energy density of a mode in weakly-dispersive material ($`|ϵ/\omega ||\mathrm{d}ϵ/\mathrm{d}\omega |`$) is given by $`U_E=\frac{1}{8\pi d}(𝐃𝐄^{})𝑑x`$ and $`U_H=\frac{1}{8\pi d}(𝐇𝐇^{})𝑑x`$ respectively$`^\text{?}`$ (the asterisk () denotes the complex conjugation). Using the explicit mode structure for TE and TM waves (see Ref. \[?\]) we arrive to:
$`U_E^{(TM)}`$ $`=`$ $`U_H^{(TM)}={\displaystyle \frac{1}{16\pi }}{\displaystyle \frac{ϵ_{}^2\mathrm{k}^2}{\kappa ^2}}|A_0|^2;U^{(TM)}=U_E^{(TM)}+U_H^{(TM)}={\displaystyle \frac{ϵ_{}^2\mathrm{k}^2}{8\pi \kappa ^2}}|A_0|^2;`$ (4)
$`U_E^{(TE)}`$ $`=`$ $`U_H^{(TE)}={\displaystyle \frac{ϵ_{}}{16\pi }}|A_0|^2;U^{(TE)}=U_E^{(TE)}+U_H^{(TE)}={\displaystyle \frac{ϵ_{}}{8\pi }}|A_0|^2,`$ (5)
where $`A_0`$ is the mode amplitude. Thus, extending the similarity between the waveguide system described here and the free-space propagation, the EM energy of any propagating wave is always positive and contains equal contributions from the electric and magnetic components of the field.
It is also seen that the TE mode is in some sense very similar to the conventional plane wave propagating in the isotropic homogeneous dielectric. Namely, (i) energy density of the TE waves is exactly equal to that of the plane waves; (ii) there is no wave propagation in material with $`ϵ_{}<0`$. In contrast to this behavior, the sign of the dielectric permittivity alone does not impose limitations on the propagation of TM modes.
Another important characteristic of the energy transport in the EM system is the average energy flux given by the propagating component of the Poynting vector $`𝐒=\frac{c}{4\pi }[𝐄\times 𝐇]`$. Selecting the direction of the wave propagation as $`z`$ axis, we obtain:
$`S_z^{(\{TE;TM\})}=c{\displaystyle \frac{k_z}{ϵ_{\{;\}}\mathrm{k}}}U^{(\{TE;TM\})}`$ (6)
It is clearly seen from Eq. 6 that the relation between the direction of the phase velocity and energy flux is defined by the sign of the dielectric constant (for a given mode polarization) – positive $`ϵ`$ means $`n>0`$ propagation, while $`ϵ<0`$ signifies the NIM case. Of course, for this relation to take place, we must require the medium to be transparent – both propagation constants $`ϵ`$ and $`\nu `$ must be of the same sign. As it can be seen from Eq. (1), the NIM condition can be satisfied only for TM wave and only in the case of extreme anisotropy of the dielectric constant of the core material ($`ϵ_{}ϵ_{}<0`$). The feasibility of the fabrication of such unusual materials will be discussed in the Section 3.
### 2.B The effect of finite wall conductance
In this section we consider the practical realization of the system described above, in which the anisotropic core material is surrounded by metallic walls. The electromagnetic properties of metals at high (GHz to optical) frequencies are dominated by the dynamics of the free-electron plasma-like gas. Following the approach described in e.g. Ref. \[?\] it is possible to write down the high-frequency effective permittivity of metal in Drude form:
$$ϵ_m(\omega )=ϵ_{\mathrm{}}\frac{\mathrm{\Omega }_{\mathrm{pl}}^2}{\omega (\omega +i\tau )},$$
(7)
where the constant term $`ϵ_{\mathrm{}}`$ describes the contribution of the bound electrons, $`\tau `$ is responsible for EM losses due to (inelastic) processes, and $`\mathrm{\Omega }_{\mathrm{pl}}=\frac{N_ee^2}{m_{\mathrm{eff}}}`$ is the plasma frequency with $`N_e,e`$, and $`m_{\mathrm{eff}}`$ being the free-electron concentration, charge, and effective mass respectively. Note that for $`\omega <\mathrm{\Omega }_{\mathrm{pl}}/\sqrt{ϵ_{\mathrm{}}}`$ the permittivity of the metal becomes negative $`ϵ_m^{}<0`$ (here and below single and double prime ($`{}_{}{}^{};_{}^{\prime \prime }`$) denote the real and imaginary parts respectively). For most of “good” metals (Ag,Al,Au) the plasma frequency is of the order of $`10eV`$ and $`ϵ_{\mathrm{}}1`$, which means that $`ϵ_m^{}`$ is negative from optical to GHz frequencies. The losses, given by the parameter $`ϵ_m^{\prime \prime }/|ϵ_m^{}|1`$ are typically small in these spectral ranges.
Similar to the case of perfectly conducting waveguide walls, the structure of the modes in the system can be still derived from the dependence of the longitudinal ($`z`$) field component on the $`x`$ coordinate, which has $`\mathrm{cos}(\kappa x)`$ or $`\mathrm{sin}(\kappa x)`$ behavior depending on its symmetry. The exact value of the mode parameter $`\kappa `$ is given by the requirement of the in-plane ($`y,z`$) field components continuity throughout $`x=\pm d/2`$ planes. For symmetric ($`\mathrm{cos}`$) mode profile, we obtain:
$`\mathrm{tan}\left({\displaystyle \frac{\kappa ^{(TM)}d}{2}}\right)`$ $`=`$ $`{\displaystyle \frac{ϵ_m\kappa ^{(TM)}}{\sqrt{\mathrm{k}^2ϵ_{}^2(ϵ_{}ϵ_m)\kappa ^{(TM)^2}ϵ_{}ϵ_{}}}}`$ (8)
$`\mathrm{tan}\left({\displaystyle \frac{\kappa ^{(TE)}d}{2}}\right)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\mathrm{k}^2(ϵ_{}ϵ_m)\kappa ^{(TE)^2}}}{\kappa ^{(TE)}}}`$
In the limit of $`ϵ_m\mathrm{}`$, these equations yield the values $`\kappa _0=\pi (2j+1)/d`$, used in the previous Section. As we previously noted, these values correspond to the well-known condition of zero mode magnitude at the waveguide boundary. In the limit of sufficiently large $`|ϵ_m|`$ it is possible to find the correction to the above values of the mode parameter $`\kappa `$. Specifically,
$`\kappa ^{(TM)}`$ $``$ $`\kappa _0\left(1{\displaystyle \frac{2\mathrm{k}ϵ_{}}{\kappa _0^2d\sqrt{ϵ_m}}}\right)`$ (9)
$`\kappa ^{(TE)}`$ $``$ $`\kappa _0\left(1{\displaystyle \frac{2}{\mathrm{k}d\sqrt{ϵ_m}}}\right)`$
As the mode parameter $`\kappa `$ plays a role of an inverse confinement length of the mode in $`x`$ direction, the negative $`\kappa `$ correction signifies the “mode expansion” into the waveguide wall region. Such a mode expansion is illustrated in Fig. 2.
The immediate effect of such a change in the mode structure is the change of the effective phase velocity, given by the refraction index:
$`n^{(TM)}\pm \sqrt{ϵ_{}\nu _0}\left(1+{\displaystyle \frac{2}{\mathrm{k}d\nu _0\sqrt{ϵ_m}}}\right)`$ (10)
$`n^{(TE)}\sqrt{ϵ_{}\nu _0}\left(1+{\displaystyle \frac{2\kappa _0^2}{\mathrm{k}^3dϵ_{}\nu _0\sqrt{ϵ_m}}}\right),`$
where $`\nu _0=1\kappa _0^2/(ϵ_{}\mathrm{k}^2)`$. As it has been described above, the sign of the refraction index for the $`TM`$ polarization has to be selected positive for $`ϵ_{}>0;\nu >0`$, and negative for $`ϵ_{}<0;\nu <0`$.
Penetration of the mode into the waveguide wall region has another effect on the wave propagation. Namely, the finite value of the $`ϵ_m^{\prime \prime }`$ introduces an additional (with respect to the core material) absorption into the system. As a result, the magnitude of a mode will exponentially decay as it propagates through the system. Such an attenuation can be related to the imaginary part of the effective refractive index through $`E\mathrm{exp}(n^{\prime \prime }\mathrm{k}z)`$. In the limit of small absorption in the metal ($`ϵ_m^{\prime \prime }/|ϵ_m^{}|1`$) the “waveguide-induced” mode decay is described by:
$`n^{(TM)^{\prime \prime }}{\displaystyle \frac{1}{\mathrm{k}d}}\sqrt{{\displaystyle \frac{ϵ_{}}{\nu _0|ϵ_m|}}}{\displaystyle \frac{ϵ_m^{\prime \prime }}{|ϵ_m^{}|}}`$ (11)
$`n^{(TE)^{\prime \prime }}{\displaystyle \frac{\kappa _0^2}{\mathrm{k}^3d\sqrt{ϵ_{}\nu _0|ϵ_m|}}}{\displaystyle \frac{ϵ_m^{\prime \prime }}{|ϵ_m^{}|}}`$
Note that in agreement with causality principle$`^{\text{?, ?}}`$ the losses in the system are positive, regardless of the sign of the refractive index.
Using Eq. (11) we estimate that for wavelengths $`\lambda 850`$ nm, the losses introduced by a silver waveguide walls are substantially small ($`n^{\prime \prime }/n0.01`$).
## 3 Anisotropic nanoplasmonic composites
We now consider the fabrication perspectives of the material with strong optical anisotropy required for NIM waveguide core region. A number of naturally occurring materials with the required properties exist at THz or far IR frequencies. Some examples include Bi and Sapphire.$`^\text{?}`$ Unfortunately, no known material exhibits anisotropy exceeding 30% at optical or infrared spectral range. Here we propose to take advantage of a new class of nano-engineered media, known as meta-materials.$`^\text{?}`$ In these composites, nanostructured particles are used as meta-atoms to achieve the desired EM properties.
To realize the strong optical anisotropy we propose to use a combination of plasmonic or polar particles (providing the negative permittivity) and dielectric media (having $`ϵ>0`$). If the characteristic size of inhomogeneities and their typical separation are much smaller than the wavelength of incident radiation, the EM response of the composite structure can be described in terms of the effective dielectric constant $`ϵ_{\mathrm{eff}}`$:$`^\text{?}`$
$$<D(r)>_\alpha =<ϵ(r)_{\alpha ,\beta }E(r)_\beta >=ϵ_{\mathrm{eff}_{\alpha ,\beta }}<E(r)>_\beta ,$$
(12)
where the brackets ($`<>`$) denote the averaging over the microscopically-large (multi-particle), macroscopically small (subwavelength) spatial area, Greek indices denote Cartesian components, and summation over repeated indices is assumed. Since the size of a particle $`a`$ typically enters Maxwell equations in the combination $`\mathrm{k}a`$,$`^\text{?}`$ all size-related effects play a minor role in the considered here “quasi-static” averaging process. Therefore, we note that the composites proposed below are highly tolerant with respect to size variation. Also, since the effects described here are obtained in effective medium approximation, the desired response does not require any periodicity of the particle arrangement and only the average concentration has to be controlled during the fabrication step.
Below we present two meta-material designs of the strongly anisotropic composite for optical and infrared spectrum ranges.
### 3.A Layered system
We first consider the permittivity of a stack of interlacing plasmonic (Ag,Au,Al,$`\mathrm{}`$) or polar (SiC) ($`ϵ_{pl}<0`$) and dielectric (Si,GaAs,$`\mathrm{}`$) ($`ϵ_d>0`$) layers. We assume that the layers are aligned in the waveguide ($`y,z`$) plane (see Fig. 3). As noted above, the absolute thickness of the layers is not important (as long as it is subwavelength), and only the average concentration of plasmonic layers $`N_{pl}`$ plays the role.
To compute $`ϵ_{\mathrm{eff}}`$ we note that the $`E_y,E_z`$, and $`ϵE_x`$ have to be continuous throughout the system,$`^{\text{?, }\text{?}\text{}\text{?}}`$ leading to:
$`ϵ_{}=ϵ_{\mathrm{eff}_{y,z}}=N_{pl}ϵ_{pl}+(1N_{pl})ϵ_d`$ (13)
$`ϵ_{}=ϵ_{\mathrm{eff}_x}={\displaystyle \frac{ϵ_{pl}ϵ_d}{(1N_{pl})ϵ_{pl}+N_{pl}ϵ_d}}`$
The effective permittivities for several layered composites are shown in Fig. 4. We note that while the strong anisotropy $`ϵ_{}ϵ_{}<0`$ can be easily achieved in the layered system, the actual realizations of the materials with $`ϵ_{}>0,ϵ_{}<0`$ required for the high-frequency NIM described here typically have substantial absorption<sup>1</sup><sup>1</sup>1This particular realization of layered NIM structure for IR frequencies has been earlier proposed in Ref. \[?\], and therefore have a limited range of applications.$`^{\text{?, ?}}`$
On the contrary, the materials with $`ϵ_{}<0,ϵ_{}>0`$ (achieved, for example by a repeated deposition of Ag-Si layers) form low-loss media. While this configuration has a positive refraction index, it may be potentially used to concentrate propagating modes in subwavelength areas.$`^{\text{?, }\text{?}}`$
### 3.B Aligned wire structure
The array of aligned $`ϵ_{pl}<0`$ nanowires embedded in the dielectric ($`ϵ_d>0`$) host, schematically shown in Fig. 5 is in some sense a counterpart of the layered system described above. In fact, the boundary conditions now require the continuity of the $`E_x`$ field, along with the solution of quasi-static equations in the $`yz`$ plane. While in the general case the analytical solution of this problem is complicated, the case of small plasmonic material concentration is adequately described by the Maxwell-Garnett approximation:$`^{\text{?, }\text{?}\text{}\text{?}\text{}\text{?}}`$
$`ϵ_{}=ϵ_{\mathrm{eff}_{y,z}}={\displaystyle \frac{N_{pl}ϵ_{pl}E_{in}+(1N_{pl})ϵ_dE_0}{N_{pl}E_{in}+(1N_{pl})E_0}}`$ (14)
$`ϵ_{}=ϵ_{\mathrm{eff}_x}=N_{pl}ϵ_{pl}+(1N_{pl})ϵ_d,`$
where $`E_{in}=\frac{2ϵ_d}{ϵ_d+ϵ_{pl}}E_0`$ is the field inside the plasmonic inclusion and $`E_0`$ is the excitation field.
To illustrate the validity of the MG approximation, we numerically solve the Maxwell equations in the planar geometry using the coupled-dipole approach (CDA), described in detail in Refs. \[?, ?, ?\]. In these calculations the composite is represented by a large number of interacting point dipoles, and the resulting dipole moment distribution is related to the effective dielectric constant. Fig. 5 shows the excellent agreement between the numerical simulations and the analytical result \[Eq. (14)\].
The effective dielectric constants for some composite materials are presented in Fig. 6. Note that in contrast to the layered system described above, these wired composites have extremely low absorption in the near-IR-NIM regime - in a way solving the major problem with the “conventional” design of optical LHMs.
## 4 Imaging Properties of non-magnetic optical NIMs
To illustrate the imaging performance of the proposed system we calculate the propagation of a wavepacket formed by a double-slit source through the $`5\mu m`$-long planar layer of $`5\%`$ $`Ag`$, $`95\%`$ $`SiO_2`$ wire-based NIM-core described above \[see Fig. 6(a,b)\], embedded in the $`Si`$ waveguide. We select the thickness of the dielectric core to be $`d=0.3\mu m`$, and assume the excitation by the telecom-wavelength $`\lambda =1.5\mu m`$. The Eqs. (10,14) yield the following values of the refraction index: $`n^{(+)}2.6`$, $`n^{(LHM)}2.6+0.05i`$.
To calculate the resulting field distribution we first represent the wavepacket at the $`z=0`$ plane as a linear combination of the waveguide modes.$`^{\text{?, ?}}`$ We then use the boundary conditions at the front and back interfaces of the NIM region to calculate the reflection and transmission of individual mode. The solutions of Maxwell equations are then represented as a sum of solutions for the individual modes.
To better illustrate the imaging properties of the system and distinguish between the effects of negative refractive index and material absorption, we first neglect losses in the NIM core. The resulting intensity distribution in the system is shown in Fig. 7(a). The image formation in the focal plane ($`z=10\mu m`$) of the far-field planar NIM lens is clearly seen.
In Fig. 7(b) we compare the imaging through the planar NIM lens with and without the material absorption and demonstrate that the presence of weak loss, although it reduces the magnitude of the signal, it does not destroy the far-field imaging. Similar to any far-field imaging system,$`^{\text{?, ?, ?}}`$ the resolution $`\mathrm{\Delta }`$ of the non-magnetic NIM structure presented here is limited by the internal wavelength: $`\mathrm{\Delta }\lambda _{in}/2=\lambda /(2n)0.3\mu m`$ \[see Fig. 7(c)\].
## 5 Conclusions
In conclusion, we presented a non-magnetic non-periodic design of a system with negative index of refraction. We have further proposed several low-loss nanoplasmonic-based realizations of the proposed structure for optical and infrared frequencies. We have presented analytical description of the effective dielectric permittivity of strongly anisotropic nanostructured composites, and showed the excellent agreement of the developed theory with results of numerical solution of Maxwell equations. Finally, we have demonstrated the low-loss far-field planar NIM lens for $`\lambda =1.5\mu m`$ with resolution $`\mathrm{\Delta }0.3\mu m`$.
The authors would like to thank E.E. Mishchenko and A.L. Efros for fruitful discussions. The research was partially supported by NSF grants DMR-0134736, ECS-0400615 and Oregon State University.
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# The impact of lens galaxy environments on the image separation distribution
## 1 Introduction
Strong gravitational lensing offers a unique probe of cosmology and the physical properties of early-type galaxies. For example, the ensemble properties of a well-defined sample of strong lens systems can constrain the cosmological constant (Turner, 1990; Fukugita et al., 1990; Kochanek, 1996b; Chae, 2003; Mitchell et al., 2005), the density profile and evolution of early-type galaxies (Kochanek et al., 2000; Rusin et al., 2003a, b; Rusin & Kochanek, 2005), and the amount of substructure in lens galaxies (Metcalf & Madau, 2001; Dalal & Kochanek, 2002; Kochanek & Dalal, 2004). In addition, the image separation distribution of lensed quasars (Turner et al., 1984) probes the efficiency of baryon cooling inside dark matter halos of different masses, and thereby constrains models of galaxy formation (Kochanek & White, 2001; Keeton, 2001a; Oguri, 2002). So far, more than 80 lens systems have been discovered, and the largest homogeneous subsample suitable for statistical studies is the Cosmic Lens All-Sky Survey (CLASS), containing 13 lensed radio sources at image separations $`0\stackrel{}{.}3<\theta <3^{\prime \prime }`$ (Myers et al., 2003; Browne et al., 2003). A larger statistical sample is being constructed from the Sloan Digital Sky Survey (SDSS), which has already discovered more than 10 new lenses including the largest-separation lensed quasar known to date (Inada et al., 2003; Oguri et al., 2004).
Lensing probability distributions are often computed using a clean spherical lens object without external fields. However, in reality, most observed lenses have non-spherical galaxies, and many require a significant external tidal shear (Keeton et al., 1997). While a portion of the shear may come from fluctuations of matter along the line of sight (Seljak, 1994; Bar-Kana, 1996; Momcheva et al., 2005), much of it is thought to be associated with mass in the immediate environment of the lens galaxy. Most lens galaxies are early-type, which lie preferentially in dense environments. For instance, Keeton et al. (2000, see also ) used galaxy demographics (specifically, the galaxy luminosity function as a function of type and environment) to predict that at least $``$25% of lens galaxies lie in groups or clusters of galaxies. Indeed, more than 20 lens galaxies are already known to or suspected to be surrounded by groups or clusters (Young et al., 1981; Kundić et al., 1997a, b; Tonry, 1998; Fischer et al., 1998; Tonry & Kochanek, 1999, 2000; Fassnacht et al., 1999; Hagen & Reimers, 2000; Fassnacht et al., 2004; Kneib et al., 2000; Soucail et al., 2001; Faure et al., 2002, 2004; Fassnacht & Lubin, 2002; Johnston et al., 2003; Inada et al., 2003; Morgan et al., 2005; Oguri et al., 2005; McKean et al., 2005; Momcheva et al., 2005).
Conventional wisdom holds that ellipticity and environment mainly affect the relative numbers of double and quadruple lenses and do not significantly modify the total lensing probability (Kochanek, 1996b; Keeton et al., 1997; Rusin & Tegmark, 2001; Chae, 2003, but see Keeton & Zabludoff 2004). Huterer et al. (2005) showed explicitly that ellipticity and shear change the lensing probability by only a few percent for most source luminosity functions (see also Premadi & Martel, 2004). They also found that shear (and to a lesser extent ellipticity) shifts and broadens the distribution of image separations for a given lens galaxy. However, they did not put the two pieces together and discuss changes in the overall image separation distribution. Also, Huterer et al. (and everyone else, for that matter) neglected the effects of external convergence from the environment. External convergence is often omitted from models of individual lenses because the mass sheet degeneracy renders it unmeasurable (Gorenstein et al., 1988; Saha, 2000). Nevertheless, theoretical studies show that external convergence affects lensing analyses in many important ways (Keeton & Zabludoff, 2004).
There is already some observational evidence that lens galaxy environments affect the image separation distribution. Oguri et al. (2005) found an overabundance of lensed quasars with image separations $`\theta 3^{\prime \prime }`$. They argued that the excess arises from lens galaxy environments, since many of the $`\theta 3^{\prime \prime }`$ lenses appear to lie in groups or clusters. We extend this consideration to all lensed quasars: Figure 1 shows the image separation distribution for all lensed quasars and for the subset whose lens galaxies lie in groups or clusters. Lens galaxy environment is clearly correlated with image separation such that lenses with larger separations tend to lie in groups and clusters. The correlation is almost tautological at large separations (say $`\theta >4^{\prime \prime }`$), because it is hard to achieve a large deflection angle without some environmental convergence. But the figure shows a gradual increase of the fraction even at $`\theta <4^{\prime \prime }`$. This trend provides further evidence that environmental convergence and shear are important: it contradicts the prediction by Keeton et al. (2000) that, if external convergence and shear are irrelevant, lenses in dense environments should have a smaller mean image separation than lenses in the field (because dense environments tend to have a larger ratio of dwarf to giant galaxies than the field).
While these data certainly suggest a connection between environment and lensing, the observational evidence is incomplete because systematic surveys of lens environments have only just begun (see Momcheva et al., 2005). It is therefore valuable to undertake a theoretical analysis of how environments affect lensing probabilities, and what range of environments is expected for lenses with various image separations. In this paper, we compute the lensing effects of shear and convergence associated with matter near the lens galaxy (see also Möller et al., 2002, for the effect of groups on lens statistics). We use a realistic model for the distribution of galaxy environments, which is based on $`N`$-body simulations and the halo occupation distribution and calibrated by observations of galaxy–galaxy lensing and number counts of massive elliptical galaxies (Dalal & Watson, 2005). For the lens objects, we assume a singular isothermal sphere (SIS) because it is both simple and a reasonable model for the density profile of lens galaxies (e.g., Cohn et al., 2001; Treu & Koopmans, 2002; Koopmans et al., 2003; Rusin & Kochanek, 2005). We neglect the ellipticity of lens galaxies for simplicity, but note that the effects of ellipticity are generally smaller than those of external shear (Huterer et al., 2005).
This paper is organized as follows. Section 2 describes the joint probability distribution function of convergence and shear that we use throughout the paper. Section 3 reviews the calculation of lensing probabilities and the image separation distribution. In §4 we examine the image separation distribution from various viewpoints in order to understand how lens environments affect the distribution. Section 5 summarizes our main conclusions. Throughout the paper we assume a concordance cosmology with mass density $`\mathrm{\Omega }_M=0.3`$, cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, and dimensionless Hubble parameter $`h=0.7`$.
## 2 Joint PDF of convergence and shear
We consider the probability distribution functions (PDFs) of convergence and shear originating from lens galaxy environments, using the model derived by Dalal & Watson (2005). They placed galaxies in an $`N`$-body simulation using a halo occupation distribution calibrated to match number counts and tangential shear profiles measured for massive elliptical galaxies in the SDSS (see Sheldon et al., 2004), and determined the convergence and shear at each galaxy position. Figure 3 of their paper shows the resulting distributions. The convergence $`\kappa _{\mathrm{ext}}`$ and shear $`|\gamma _{\mathrm{ext}}|`$ have similar means $`0.03`$, but very different tails: large convergences ($`\kappa _{\mathrm{ext}}0.1`$) are more common than large shears ($`|\gamma _{\mathrm{ext}}|0.1`$). There are two reasons for the difference. First, for a typical group or cluster dark matter halo, we expect that a large shear is always accompanied by an even larger convergence. Figure 2 illustrates this idea by showing the convergence and shear profiles for a dark matter halo modeled with an NFW density profile (Navarro et al., 1997); in the central region ($`\kappa ,|\gamma |0.1`$), we see that $`\kappa >|\gamma |`$. <sup>1</sup><sup>1</sup>1More generally, for a power-law density profile $`\rho (r)r^\eta `$, the convergence and shear are related as $`|\gamma |/\kappa =|\eta 1|/|3\eta |`$. Therefore the argument holds for different mass profiles, as long as they have inner profiles that are shallower than the isothermal sphere ($`\eta =2`$). Second, convergence and shear sum differently when there are multiple perturbers. The probability distribution of convergence is highly skewed, with a tail extending to large positive values but no corresponding tail at large negative values. In contrast, the distributions of both shear components are symmetric about zero. Thus, a large convergence provided by one perturber cannot be canceled by contributions from other perturbers, while large shear values can be canceled. As we shall see, the longer tail of convergence is one reason that convergence is more important than shear.
Figure 2 also implies that convergence and shear are not independent but rather correlated. Therefore in this paper we consider the joint PDF of convergence and shear, which is shown in Figure 3. The correlation between $`\kappa _{\mathrm{ext}}`$ and $`|\gamma _{\mathrm{ext}}|`$ is quite evident, particularly for large values.
Although the environment distribution was derived for a lens redshift $`z_l=0.45`$, we assume little evolution in the elliptical galaxy population over the range $`0z1`$ (see Schade et al., 1999; Im et al., 2002; Ofek et al., 2003; Chae & Mao, 2003). Thus, we simply extrapolate the surface mass density associated with environment to all redshifts. Since convergence and shear depend on the surface mass density in units of the critical density for lensing, $`\mathrm{\Sigma }_{\mathrm{crit}}D_{os}/(D_{ol}D_{ls})`$, the redshift dependence of the PDF for convergence and shear is straightforward. The PDF was derived for galaxies with a velocity dispersion $`\sigma 216`$ km s<sup>-1</sup>, which is typical for lenses with image separation $`\theta 1^{\prime \prime }`$, but we apply it to all galaxies. In other words, we neglect any correlation between a galaxy’s velocity dispersion and environment. However, we know that galaxies with higher velocity dispersion will tend to reside in denser environments; accordingly, our estimate of environmental effects on wide-splitting lenses is strictly a lower limit. We exclude extreme cases with $`\kappa _{\mathrm{ext}}+|\gamma _{\mathrm{ext}}|1`$, because those would be categorized as lensing by groups/clusters rather than by a galaxy perturbed by its environment.
## 3 Lens Statistics with External Convergence and Shear
We model each lens object as an SIS galaxy with external convergence and shear. The Einstein radius of an SIS with velocity dispersion $`\sigma `$ is
$$\theta _\mathrm{E}(\sigma )=4\pi \left(\frac{\sigma }{c}\right)^2\frac{D_{ls}}{D_{os}},$$
(1)
where $`D_{os}`$ and $`D_{ls}`$ are, respectively, angular diameter distances from the observer and lens to the source.
The important properties of the lens are the image separation as a function of source position, the lensing cross section, and the magnification bias. We define the image separation $`\theta `$ to be the maximum distance between multiple images, which is a well-defined and observable quantity. It is useful to define the normalized separation $`\widehat{\theta }`$ such that
$$\theta (𝐮)=\theta _\mathrm{E}\widehat{\theta }(𝐮).$$
(2)
Note that $`\widehat{\theta }(𝐮)=2`$ for all $`𝐮`$ in the absence of convergence and shear.
We combine the cross section and magnification bias to compute the “biased cross section”,
$$BA=\frac{\mathrm{\Phi }(L/\mu )}{\mathrm{\Phi }(L)}𝑑𝐮=\mu ^{\beta 1}𝑑𝐮,$$
(3)
where the integral is over the multiply-imaged region of the source plane, and we assume a power law source luminosity function $`\varphi _L(L)L^\beta `$. Unless otherwise specified, we adopt the source luminosity function of the CLASS survey, $`\beta =2.1`$ (Myers et al., 2003). We take $`\mu `$ to be the total magnification of all images, which is appropriate for surveys in which lenses are identified from high-resolution follow-up observations of unresolved targets (including both CLASS and SDSS). We choose coordinates such that the source position $`𝐮`$ is in units of $`\theta _\mathrm{E}`$, which means that $`BA`$ is naturally (and conveniently) in units of $`\theta _\mathrm{E}^2`$.
We compute image separations and biased cross sections using the public software gravlens by Keeton (2001b). For each set of ($`\kappa _{\mathrm{ext}}`$, $`|\gamma _{\mathrm{ext}}|`$) we place $`10^5`$ random sources in the smallest circle enclosing the lensing caustics. We solve the lens equation to find the images, and then compute the image separation and total magnification. Finally, we sum over the multiply-imaged sources to compute the biased cross section.
The next step is to integrate over lens galaxy populations to obtain the total lensing probability. This involves integrating over appropriate distributions of galaxy masses (or velocity dispersions), redshifts, and environments:
$`P`$ $`=`$ $`{\displaystyle 𝑑z_l\frac{cdt}{dz_l}(1+z_l)^3𝑑\kappa _{\mathrm{ext}}d|\gamma _{\mathrm{ext}}|p(\kappa _{\mathrm{ext}},|\gamma _{\mathrm{ext}}|)}`$ (4)
$`\times {\displaystyle }d\sigma {\displaystyle \frac{dn}{d\sigma }}\left(D_{ol}\theta _\mathrm{E}\right)^2BA.`$
Here $`D_{ol}`$ is the angular diameter distance from the observer to the lens, and $`p(\kappa _{\mathrm{ext}},|\gamma _{\mathrm{ext}}|)`$ is the joint PDF of external convergence and shear from §2. We specify the distribution of galaxy velocity dispersions using the velocity function $`dn/d\sigma `$ of early-type galaxies determined from $``$30,000 galaxies at $`0.01<z<0.3`$ in the SDSS (Sheth et al., 2003; Mitchell et al., 2005).
The image separation distribution of lensed quasars can then be obtained by differentiating equation (4):
$`{\displaystyle \frac{dP}{d\theta }}`$ $`=`$ $`{\displaystyle 𝑑z_l\frac{cdt}{dz_l}(1+z_l)^3𝑑\kappa _{\mathrm{ext}}d|\gamma _{\mathrm{ext}}|p(\kappa _{\mathrm{ext}},|\gamma _{\mathrm{ext}}|)}`$ (6)
$`\times {\displaystyle }d\widehat{\theta }{\displaystyle \frac{1}{\widehat{\theta }}}{\displaystyle \frac{d\sigma }{d\theta _\mathrm{E}}}{\displaystyle \frac{dn}{d\sigma }}\left(D_{ol}\theta _\mathrm{E}\right)^2{\displaystyle \frac{d(BA)}{d\widehat{\theta }}}`$
$`=`$ $`{\displaystyle 𝑑z_l\frac{cdt}{dz_l}(1+z_l)^3𝑑\kappa _{\mathrm{ext}}d|\gamma _{\mathrm{ext}}|p(\kappa _{\mathrm{ext}},|\gamma _{\mathrm{ext}}|)}`$
$`\times {\displaystyle }d\widehat{\theta }{\displaystyle \frac{\sigma }{2\theta }}{\displaystyle \frac{dn}{d\sigma }}{\displaystyle \frac{\left(D_{ol}\theta \right)^2}{\widehat{\theta }^2}}{\displaystyle \frac{d(BA)}{d\widehat{\theta }}}.`$
We fix the source redshift to $`z_s=2.0`$ which is a typical redshift for lensed quasars.
## 4 Results
### 4.1 Dependence on convergence and shear
Before presenting results that account for full distribution of convergence and shear, it is useful to study how fixed values of $`\kappa _{\mathrm{ext}}`$ and/or $`|\gamma _{\mathrm{ext}}|`$ affect the image separation distribution. These results are shown in Figure 4.
Convergence and shear both modify the shape of image separation distribution, particularly at relatively large image separations. For instance, at $`\theta =5^{\prime \prime }`$ a convergence $`\kappa _{\mathrm{ext}}=0.2`$ increases the lensing probability by one order of magnitude, and a shear of $`|\gamma _{\mathrm{ext}}|=0.2`$ enhances it by a factor 2–3. For the same values, convergence clearly has much more effect on lensing probabilities than shear. Combined with the fact that large values of convergence are more likely to occur than large values of shear (see §2), we expect that convergence has the greater impact on image separation distributions.
The results can be understood as follows. First, external convergence magnifies the image plane with respect to the source plane, which increases the image separation by a factor $`(1\kappa _{\mathrm{ext}})^1`$. Second, while external convergence does not affect the caustics, it does change the biased cross section by giving additional magnification to the images. Specifically, the biased cross section increases by a factor $`(1\kappa _{\mathrm{ext}})^{2(\beta 1)}`$ (see Keeton & Zabludoff, 2004). As a result of these two effects, convergence shifts the image separation distribution up and to the right in Figure 4. The increase is significant where the image separation distribution is a deceasing function ($`\theta 1^{\prime \prime }`$), but negligible where it is an increasing function ($`\theta 1^{\prime \prime }`$). By contrast, shear mainly broadens the image separation distribution. It modestly enhances the lensing probability at $`\theta 1^{\prime \prime }`$, and suppresses it at $`\theta 1^{\prime \prime }`$.
### 4.2 Full result
We are now ready to consider the image separation distribution with the full joint PDF of convergence and shear, which is shown in Figure 5. The enhancement of the lensing probability, which is plotted in Figure 6, is small at $`\theta 1^{\prime \prime }`$, but significant at larger image separations: the increase is $``$30% at $`\theta =3^{\prime \prime }`$ and $``$200% at $`\theta =5^{\prime \prime }`$. Thus, lens galaxy environments are very important when we discuss the shape of the image separation distribution. Indeed, they can account for the excess of lensing probabilities reported by Oguri et al. (2005).
In Figure 6, we also plot the enhancement of the cumulative lens probability $`P(>\theta )`$. We find that the enhancement of the total lensing probability integrated over the all image separation is $`15\%`$. This enhancement is equivalent to that obtained by shifting the cosmological constant by $`\mathrm{\Delta }\mathrm{\Omega }_\mathrm{\Lambda }0.05`$, therefore it cannot be ignored in accurate determination of cosmological parameters with lens statistics.
As discussed above, some of the effects of convergence and shear are mediated by magnification bias: convergence directly magnifies images, while shear generates and lengthens the tangential caustic and thus increases the cross section for high magnifications. The enhancement of the lensing probability is therefore expected to depend on the source luminosity function. Indeed, Figure 7 shows that the enhancement at $`\theta =3^{\prime \prime }`$ increases strongly with the luminosity function slope $`\beta `$ (recall $`\mathrm{\Phi }L^\beta `$). This result implies that bright optical lensed quasars with large image separations are more likely to lie in dense environments.
### 4.3 What enhances the lensing probability?
Dissecting the results further helps us understand the changes to the image separation in more detail. As discussed in §4.1, we believe that convergence is more important than shear because of its stronger effect on the lensing probability (see Fig. 4) and the longer tail to high $`\kappa _{\mathrm{ext}}`$ in the joint PDF (see Fig. 3). To check this hypothesis, we project the joint PDF to a distribution of convergence or shear alone (the same distributions shown in Fig. 3 Dalal & Watson 2005) and recompute the image separation distribution. The results, shown in the upper panel of Figure 8, confirm that the enhancement of the lensing probability is in fact driven by convergence.
In that case, it is useful to understand what values of $`\kappa _{\mathrm{ext}}`$ contribute most to the lensing probability. We see this by decomposing the image separation distribution into contributions from different $`\kappa _{\mathrm{ext}}`$ intervals, as shown in the bottom panel of Figure 8. The interpretation depends on the image separation. At small image separations ($`\theta <3^{\prime \prime }`$), most lenses are expected to have small convergences ($`\kappa _{\mathrm{ext}}0.05`$), and larger and larger convergences are more and more rare. At intermediate separations ($`3^{\prime \prime }\theta 4^{\prime \prime }`$), the largest convergences ($`\kappa _{\mathrm{ext}}0.35`$) become more important than intermediate values. Finally at large separations ($`\theta 5^{\prime \prime }`$), the largest convergences become dominant. In other words, while dense environments with large convergences may be rare, nearly all of the largest separation lenses will be found there.
### 4.4 Environments of lens galaxies
Figure 8 suggests that the typical environments of lenses are very sensitive to image separation. Put another way, the distribution of convergence and shear for actual lens galaxies will differ from the distribution shown in Figure 3 (which applies to normal elliptical galaxies), by a factor of the lensing probability. To quantify this effect, we derive the posterior PDF of convergence and shear for lenses at fixed image separations of $`\theta =1^{\prime \prime }`$ or $`3^{\prime \prime }`$. Figure 9 shows that the distributions differ from each other, and from the original (unweighted) distributions. The distributions for lenses with $`\theta =1^{\prime \prime }`$ are similar to the original distributions, just shifted to slightly larger values. However, the distributions for lenses with $`\theta =3^{\prime \prime }`$ show particularly large increases in the tail of high convergence or shear values.
To see this more clearly, in Figure 10 we show the fraction of lenses with a “strong” environment (defined to be $`\kappa _{\mathrm{ext}}`$ or $`|\gamma _{\mathrm{ext}}|>0.1`$), as a function of image separation $`\theta `$. The fraction increases monotonically as the image separation grows. For instance, at $`\theta =1^{\prime \prime }`$ the fraction of lenses with a strong convergence is 11%, while the fraction with a strong shear is 6%. By $`\theta =3^{\prime \prime }`$ the fractions have increased to 24% and 13%, respectively. (For comparison, the fractions of non-lens elliptical galaxies with strong convergence or shear are just 6% and 3%; see Dalal & Watson 2005.) In other words, the distribution of lens galaxy environments is a strong function of image separation, and this effect should be taken into account when comparing observed environments with theoretical predictions. The increasing fraction of “strong” environments with image separation is qualitatively consistent with the data shown in Figure 1.
## 5 Summary and Discussion
We have studied how the convergence and shear from lens galaxy environments affect lens statistics, in particular the distribution of lens image separations. We find that the external field enhances the lensing probability, especially at large image separations, and the effect increases with the slope of the source luminosity function. We argue that the enhancement is driven mainly by the external convergence, which has been neglected in previous studies of this sort. Our results mesh with those of Keeton & Zabludoff (2004) to indicate that it is essential to include convergence from the environment in order to obtain correct results from both lens modeling and lens statistics. For example, since the external field changes the shape of the image separation distribution rather than its overall amplitude, it will bias attempts to constrain the velocity dispersion function of early-type galaxies using the lens image separation distribution (e.g., Chae, 2005).
The environmental boost in the lensing probability depends on the image separation, which means that the posterior distribution of lens environments does as well. The fraction of lenses with a “strong” environment (a convergence or shear larger than 0.1, say) increases monotonically with image separation. Large separation lenses are more likely to be found in dense environments. This bias will be important when comparing the observed distribution of lens environments to theoretical predictions (e.g., Momcheva et al., 2005).
One puzzle remains: even when we take the separation/environment correlation into account, we find that the predicted fraction of lens galaxies with a large shear $`|\gamma _{\mathrm{ext}}|>0.1`$ is $``$10%. However, models of most observed 4-image lenses require shears of $``$0.1 or larger (Keeton et al., 1997). While 4-image lenses are certainly biased toward dense environments (see Holder & Schechter 2003 for a theoretical perspective, and Momcheva et al. 2005 for intriguing observational results), it does not appear that the bias is strong enough to reconcile our predicted shear distribution with the “observed” values. At present, it is not clear whether the problem lies with the lens models or the predicted environment distributions. If the latter, then our conclusions regarding the importance of lens galaxy environments should clearly be revisited.
It is also possible that relaxing some simplifying assumptions in our analysis will affect the quantitative results. The most important simplification is the assumption that all galaxies have the same distribution of convergence and shear. As shown in Sheldon et al. (2004), however, the tangential shear signal on scales $`1h^1`$ Mpc increases with increasing velocity dispersion, suggesting that more massive ellipticals reside in denser environments than less massive ellipticals. We expect this effect to enhance the importance of environment for wide-splitting lenses relative to the estimates presented here.
Another simplification was our neglect of lens galaxy ellipticity. The effects of ellipticity are somewhat similar to those of shear (see Huterer et al., 2005), which has a modest effect on the lensing probability. Therefore we expect that our results would not change significantly with the addition of ellipticity. Nevertheless, since ellipticity is important in predictions of the relative numbers of 4-image and 2-image lenses, it will need to be considered more carefully when comparing observed and predicted environment distributions for 4-image lenses alone.
It seems to be too early to say whether there is conflict between the observed and predicted distribution of lens galaxy environments. But it is an intriguing question that deserves to be studied both observationally and theoretically.
## Acknowledgments
We thank the referee, David Rusin, for helpful comments and suggestions. M. O. is supported by JSPS through JSPS Research Fellowship for Young Scientists. N. D. acknowledges the support of NASA through Hubble Fellowship grant #HST-HF-01148.01-A awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA, under contract NAS 5-26555.
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# Spin rotation and birefringence effect for a particle in a high energy storage ring and measurement of the real part of the coherent elastic zero-angle scattering amplitude, electric and magnetic polarizabilities
## 1 INTRODUCTION
Investigation of spin-dependent interactions of elementary particles at high energies is a very important part of program of scientific research has been preparing for carry out at storage rings (RHIC, CERN, COSY, GSI). It is well known in experimental particle physics how to measure a total spin-dependent cross-section of proton-proton (pp) and proton-deuteron (pd) or proton-nucleus (pN) and deuteron-nucleus (dN) interactions.
Through analicity we can get dispersion relations between the real and imaginary parts of the forward scattering amplitude. These relations are very valuable for analyzing interactions, especially if we know both real and imaginary parts of the forward scattering amplitude in a broad energy range through independent experimental measurements.
There are several experimental possibilities for the indirect measurement of the real part of the forward scattering amplitude .
Since no scattering experiment is possible in the forward direction, the determination of the real part of the forward amplitudes has always consisted in the measurement of well chosen elastic scattering observables at small angles and then in the extrapolation of these observables towards zero angle . All of these methods, however, contain discrete ambiguities in the reconstruction of the forward scattering matrix, which can be removed only by new independent measurements. Consequently,what is needed is a direct reconstruction of the real part of the forward scattering matrix such we have in the case of the imaginary part through the measurement of a total cross section.
It has been shown in - that there is an unambiguous method which makes the direct measurement of the real part of the spin-dependent forward scattering amplitude in the high energy range possible. This technique is based on the effect proton (deuteron, antiproton) beam spin rotation in a polarized nuclear target and on the phenomenon of deuteron spin rotation and oscillation in a nonpolarized target. This technique uses the measurement of angle of spin rotation of high energy proton (deuteron, antiproton) in conditions of transmission experiment - the so-called spin rotation experiment.
The analogous phenomenon for thermal neutrons was theoretically predicted in and experimentally observed in - (the phenomena of nuclear precession of neutron spin in a nuclear pseudomagnetic field of a target).
Spin rotation and oscillation experiments as well as investigation of spin dichroism (i.e. investigation of dependence of beam absorption on spin orientation) also allow to carry out new experiments to study P- and T-odd interactions . Deuteron spin rotation and oscillation experiments allow to measure the tensor electric polarizability and, as it is shown below, the tensor magnetic polarizability, too. Change of spin state of a particle at passing deep into target can influence the experiments, studying nonelastic processes at collisions of polarized nucleons and nuclei. This impels to investigate possible influence of spin rotation on the cross-section of such processes.
Observation of particle spin rotation and birefringence effect with a storage ring requires to cancel influence of $`(g2)`$ precession ($`g`$ is the gyromagnetic ratio). This precession appears due to interaction of the particle magnetic moment with an external electromagnetic field. The requirement for $`(g2)`$ precession influence cancellation also arises when searching for a deuteron electric dipole moment (EDM) by the deuteron spin precession in an electric field in a storage ring .
In the present paper it is shown that the influence of particle spin precession in a magnetic field on the process of measurement of spin rotation of a particle passed through a polarized target can be eliminated with the aid of making the vector (tensor) polarization of the target rotating (oscillating) with the frequency coinciding with the frequency of particle spin rotation due to the particle magnetic moment interacting with a magnetic field i.e. providing for the paramagnetic resonance under the action of periodic in time pseudomagnetic (pseudoelectric) field of the target.
## 2 Interactions contributing to the spin motion of a particle in a storage ring
Considering evolution of the spin of a particle in a storage ring one should take into account several interactions:
1. interactions of the magnetic and electric dipole moments with an electromagnetic field;
2. interaction of the particle with the electric field due to the tensor electric polarizability;
3. interaction of the particle with the magnetic field due to the tensor magnetic polarizability;
4. interaction of the particle with the pseudoelectric and pseudomagnetic nuclear fields of matter.
The equation for the particle spin wavefunction considering all these interactions is as follows:
$$i\mathrm{}\frac{\mathrm{\Psi }(t)}{t}=\left(\widehat{H}_0+\widehat{V}_{EDM}+\widehat{V}_\stackrel{}{E}+\widehat{V}_\stackrel{}{B}+\widehat{V}_E^{nucl}+\widehat{V}_B^{nucl}\right)\mathrm{\Psi }(t)$$
(1)
where $`\mathrm{\Psi }(t)`$ is the particle spin wavefunction,
$`\widehat{H}_0`$ is the Hamiltonian describing the spin behavior caused by interaction of the magnetic moment with the electromagnetic field (equation (1) with the only $`\widehat{H}_0`$ summand converts to the Bargman-Myshel-Telegdy equation),
$`\widehat{V}_{EDM}`$ describes interaction of the particle EDM with the electric field,
$`\widehat{V}_{EDM}`$ $`=`$ $`d\left(\stackrel{}{\beta }\times \stackrel{}{B}+\stackrel{}{E}\right)\stackrel{}{S},`$ (2)
$`\widehat{V}_\stackrel{}{E}`$ describes interaction of the particle with the electric field due to the tensor electric polarizability:
$$\widehat{V}_\stackrel{}{E}=\frac{1}{2}\widehat{\alpha }_{ik}(E_{eff})_i(E_{eff})_k,$$
(3)
where $`\widehat{\alpha }_{ik}`$ is the electric polarizability tensor of the particle , $`\stackrel{}{E}_{eff}=(\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B})`$ is the effective electric field; the expression (3) can be rewritten as follows:
$`\widehat{V}_\stackrel{}{E}=\alpha _SE_{eff}^2\alpha _TE_{eff}^2\left(\stackrel{}{S}\stackrel{}{n}_E\right)^2,\stackrel{}{n}_E={\displaystyle \frac{\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B}}{|\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B}|}}`$ (4)
where $`\alpha _S`$ is the scalar electric polarizability and $`\alpha _T`$ is the tensor electric polarizability of the particle.
A particle with the spin $`S1`$ also has the magnetic polarizability which is described by the magnetic polarizability tensor $`\widehat{\beta }_{ik}`$ and interaction of the particle with the magnetic field due to the tensor magnetic polarizability is as follows:
$$\widehat{V}_\stackrel{}{B}=\frac{1}{2}\widehat{\beta }_{ik}(B_{eff})_i(B_{eff})_k,$$
(5)
where $`(B_{eff})_i`$ are the components of the effective magnetic field $`\stackrel{}{B}_{eff}=(\stackrel{}{B}\stackrel{}{\beta }\times \stackrel{}{E})`$; $`\widehat{V}_\stackrel{}{B}`$ (5) could be expressed as:
$$\widehat{V}_\stackrel{}{B}=\beta _SB_{eff}^2\beta _TB_{eff}^2\left(\stackrel{}{S}\stackrel{}{n}_B\right)^2,\stackrel{}{n}_B=\frac{\stackrel{}{B}\stackrel{}{\beta }\times \stackrel{}{E}}{|\stackrel{}{B}\stackrel{}{\beta }\times \stackrel{}{E}|}.$$
(6)
where $`\beta _S`$ is the scalar magnetic polarizability and $`\beta _T`$ is the tensor magnetic polarizability of the particle.
$`\widehat{V}_B^{nucl}`$ describes the effective potential energy of particle magnetic moment interaction the pseudomagnetic field of the target -,.
$`\widehat{V}_E^{nucl}`$ describes the effective potential energy of particle electric moment interaction the pseudoelectric field of the target -,.
It should be emphasized that $`\widehat{V}_B^{nucl}`$ and $`\widehat{V}_E^{nucl}`$ include contributions from strong interactions as well as those caused by weak interaction violating P (space) and T (time) invariance.
## 3 The equations describing the spin evolution of a particle in a storage ring
Let us consider particles moving in a storage ring with low pressure of residual gas ($`10^{10}`$ Torr) and without targets inside the storage ring. In this case we can omit the effects caused by the interactions $`\widehat{V}_B^{nucl}`$ and $`\widehat{V}_E^{nucl}`$.
Let us consider a particle with $`S=1`$ (for example, deuteron) moving in a storage ring. According to the above analysis spin behavior of such a particle can not be described by the Bargman-Myshel-Telegdy equation. The equations for particle spin motion including contribution from the tensor electric polarizability were obtained in . Considering that deuteron possesses also the tensor magnetic polarizability and adding the terms caused by it to the equations obtained in finally we get:
$`\{\begin{array}{c}\frac{d\stackrel{}{P}}{dt}=\frac{e}{mc}\left[\stackrel{}{P}\times \left\{\left(a+\frac{1}{\gamma }\right)\stackrel{}{B}a\frac{\gamma }{\gamma +1}\left(\stackrel{}{\beta }\stackrel{}{B}\right)\stackrel{}{\beta }\left(\frac{g}{2}\frac{\gamma }{\gamma +1}\right)\stackrel{}{\beta }\times \stackrel{}{E}\right\}\right]+\hfill \\ +\frac{d}{\mathrm{}}\left[\stackrel{}{P}\times \left(\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B}\right)\right]\frac{2}{3}\frac{\alpha _TE_{eff}^2}{\mathrm{}}[\stackrel{}{n}_E\times \stackrel{}{n}_E^{}]\frac{2}{3}\frac{\beta _TB_{eff}^2}{\mathrm{}}[\stackrel{}{n}_B\times \stackrel{}{n}_B^{}],\hfill \\ \\ \frac{dP_{ik}}{dt}=\left(\epsilon _{jkr}P_{ij}\mathrm{\Omega }_r+\epsilon _{jir}P_{kj}\mathrm{\Omega }_r\right)\hfill \\ \frac{3}{2}\frac{\alpha _TE_{eff}^2}{\mathrm{}}\left([\stackrel{}{n}_E\times \stackrel{}{P}]_in_{E,k}+n_{E,i}[\stackrel{}{n}_E\times \stackrel{}{P}]_k\right)\hfill \\ \frac{3}{2}\frac{\beta _TB_{eff}^2}{\mathrm{}}\left([\stackrel{}{n}_B\times \stackrel{}{P}]_in_{B,k}+n_{B,i}[\stackrel{}{n}_B\times \stackrel{}{P}]_k\right),\hfill \end{array}`$ (13)
where $`m`$ is the mass of the particle, $`e`$ is its charge, $`\stackrel{}{P}`$ is the spin polarization vector, $`P_{xx}+P_{yy}+P_{zz}=0`$, $`\gamma `$ is the Lorentz-factor, $`\stackrel{}{\beta }=\stackrel{}{v}/c`$, $`\stackrel{}{v}`$ is the particle velocity, $`a=(g2)/2`$, $`g`$ is the gyromagnetic ratio, $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ are the electric and magnetic fields in the point of particle location, $`\stackrel{}{E}_{eff}=(\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B})`$, $`\stackrel{}{B}_{eff}=(\stackrel{}{B}\stackrel{}{\beta }\times \stackrel{}{E})`$, $`\stackrel{}{n}=\stackrel{}{k}/k`$, $`\stackrel{}{n}_E=\frac{\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B}}{|\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B}|}`$, $`\stackrel{}{n}_B=\frac{\stackrel{}{B}\stackrel{}{\beta }\times \stackrel{}{E}}{|\stackrel{}{B}\stackrel{}{\beta }\times \stackrel{}{E}|}`$, $`n_i^{}=P_{ik}n_k`$, $`n_{E,i}^{}=P_{ik}n_{E,k}`$, $`n_{Bi}^{}=P_{il}n_{Bl}=P_{i3}`$, $`\mathrm{\Omega }_r(d)`$ are the components of the vector $`\stackrel{}{\mathrm{\Omega }}(d)`$ ($`r=1,2,3`$ correspond to $`x,y,z`$, respectively).
$`\stackrel{}{\mathrm{\Omega }}(d)`$ $`=`$ $`{\displaystyle \frac{e}{mc}}\left\{\left(a+{\displaystyle \frac{1}{\gamma }}\right)\stackrel{}{B}a{\displaystyle \frac{\gamma }{\gamma +1}}\left(\stackrel{}{\beta }\stackrel{}{B}\right)\stackrel{}{\beta }\left({\displaystyle \frac{g}{2}}{\displaystyle \frac{\gamma }{\gamma +1}}\right)\stackrel{}{\beta }\times \stackrel{}{E}\right\}+`$ (14)
$`+`$ $`{\displaystyle \frac{d}{\mathrm{}}}\left(\stackrel{}{E}+\stackrel{}{\beta }\times \stackrel{}{B}\right).`$
## 4 Deuteron birefringence effect in electromagnetic field
When omitting contribution from interaction of the particle EDM with the electric field $`\widehat{V}_{EDM}`$ we can rewrite the equations for particle spin motion (13) as follows:
$`{\displaystyle \frac{d\stackrel{}{P}}{dt}}=[\stackrel{}{P}\times \stackrel{}{\mathrm{\Omega }}]+\mathrm{\Omega }_T[\stackrel{}{n}_E\times \stackrel{}{n}_E^{}]+\mathrm{\Omega }_T^\mu [\stackrel{}{n}_B\times \stackrel{}{n}_B^{}],`$
$`{\displaystyle \frac{d\stackrel{}{P_{ik}}}{dt}}=(ϵ_{jkr}P_{ij}\mathrm{\Omega }_r+ϵ_{jir}P_{kj}\mathrm{\Omega }_r)+\mathrm{\Omega }_T^{}([\stackrel{}{n}_E\times \stackrel{}{P}]_in_{Ek}+n_{Ei}[\stackrel{}{n}_E\times \stackrel{}{P}]_k)+`$
$`+\mathrm{\Omega }_T^\mu ([\stackrel{}{n}_B\times \stackrel{}{P}]_in_{Bk}+n_{Bi}[\stackrel{}{n}_B\times \stackrel{}{P}]_k)`$ (15)
where
$`\begin{array}{c}\stackrel{}{\mathrm{\Omega }}=\frac{e}{mc}\left[\left(a+\frac{1}{\gamma }\right)\stackrel{}{B}a\frac{\gamma }{\gamma +1}\left(\stackrel{}{\beta }\stackrel{}{B}\right)\stackrel{}{\beta }\left(\frac{g}{2}\frac{\gamma }{\gamma +1}\right)\stackrel{}{\beta }\times \stackrel{}{E}\right],\hfill \\ \\ \mathrm{\Omega }_T=\frac{2}{3}\frac{\alpha _TE_{eff}^2}{\mathrm{}},\mathrm{\Omega }_T^{}=\frac{3}{2}\frac{\alpha _TE_{eff}^2}{\mathrm{}},\mathrm{\Omega }_T^{}=\frac{2}{3}\mathrm{\Omega }_T,\hfill \\ \\ \mathrm{\Omega }_T^\mu =\frac{2}{3}\frac{\beta _TB_{eff}^2}{\mathrm{}},\mathrm{\Omega }_T^\mu =\frac{3}{2}\frac{\beta _TB_{eff}^2}{\mathrm{}},\mathrm{\Omega }_T^\mu =\frac{2}{3}\mathrm{\Omega }_T^\mu .\hfill \end{array}`$ (17)
Thus presence of the electric and magnetic tensor polarizabilities makes impossible to describe the spin evolution of a particle in a by the Bargman-Myshel-Telegdy equation
$$\frac{d\stackrel{}{P}}{dt}=[\stackrel{}{P}\times \stackrel{}{\mathrm{\Omega }}]$$
(18)
but requires considering of the system (15).
Let us consider the coordinate system and vectors $`\stackrel{}{v},\stackrel{}{E},\stackrel{}{B}`$ as shown in figure and denote the axes by $`x,y,z`$ (or $`1,2,3`$, respectively).
Suppose that an electric field is absent and the particle initial polarization coincides with $`\stackrel{}{v}`$ direction, therefore, the components of the vectors are:
$`\begin{array}{c}\stackrel{}{P}=(P_1,P_2,P_3),\stackrel{}{P}_0=(0,P,0),\hfill \\ \\ \stackrel{}{n}_E=(1,0,0),n_{Ei}^{}=P_{il}n_{El}=P_{i1}\hfill \\ \\ [\stackrel{}{n}_E\times \stackrel{}{n}_E^{}]_1=0,[\stackrel{}{n}_E\times \stackrel{}{n}_E^{}]_2=P_{31},[\stackrel{}{n}_E\times \stackrel{}{n}_E^{}]_3=P_2,\hfill \\ \\ [\stackrel{}{P}\times \stackrel{}{\mathrm{\Omega }}]_1=\mathrm{\Omega }P_2,[\stackrel{}{P}\times \stackrel{}{\mathrm{\Omega }}]_2=\mathrm{\Omega }P_1,[\stackrel{}{P}\times \stackrel{}{\mathrm{\Omega }}]_3=P_2,\hfill \\ \\ [\stackrel{}{n}_E\times \stackrel{}{P}]_1=0,[\stackrel{}{n}_E\times \stackrel{}{P}]_2=P_3,[\stackrel{}{n}_E\times \stackrel{}{P}]_3=P_2,\hfill \end{array}`$ (28)
$`\begin{array}{c}\stackrel{}{\mathrm{\Omega }}=\frac{e}{mc}\left(a+\frac{1}{\gamma }\right)\stackrel{}{B}=(0,0,\mathrm{\Omega }),\hfill \\ \\ \stackrel{}{n}_B=(0,0,1),n_{Bi}^{}=P_{il}n_{Bl}=P_{i3}\hfill \\ \\ [\stackrel{}{n}_B\times \stackrel{}{n}_B^{}]_1=P_{23},[\stackrel{}{n}_B\times \stackrel{}{n}_B^{}]_2=P_{13},[\stackrel{}{n}_B\times \stackrel{}{n}_B^{}]_3=0,\hfill \\ \\ [\stackrel{}{n}_B\times \stackrel{}{P}]_1=P_2,[\stackrel{}{n}_B\times \stackrel{}{P}]_2=P_1,[\stackrel{}{n}_B\times \stackrel{}{P}]_3=0.\hfill \end{array}`$ (36)
Substituting (28,36) to the system (13) we obtain:
$`\begin{array}{c}\frac{dP_1}{dt}=\mathrm{\Omega }P_2\mathrm{\Omega }_T^\mu P_{23},\hfill \\ \\ \frac{dP_2}{dt}=\mathrm{\Omega }P_1+(\mathrm{\Omega }_T^\mu \mathrm{\Omega }_T)P_{13},\hfill \\ \\ \frac{dP_3}{dt}=\mathrm{\Omega }_TP_{12}\hfill \end{array}`$ (42)
$`\begin{array}{c}\frac{dP_{11}}{dt}=2\mathrm{\Omega }_3P_{12},\hfill \\ \\ \frac{dP_{22}}{dt}=2\mathrm{\Omega }_3P_{12},\hfill \\ \\ \frac{dP_{33}}{dt}=0,\hfill \end{array}`$ (48)
$`\begin{array}{c}\frac{dP_{12}}{dt}=\mathrm{\Omega }\left(P_{11}P_{22}\right)\mathrm{\Omega }_T^{}P_3,\hfill \\ \\ \frac{dP_{13}}{dt}=\mathrm{\Omega }P_{23}+\mathrm{\Omega }_T^{}P_2\mathrm{\Omega }_T^\mu P_2,\hfill \\ \\ \frac{dP_{23}}{dt}=\mathrm{\Omega }P_{13}+\mathrm{\Omega }_T^\mu P_1\hfill \end{array}`$ (54)
remembering that $`P_{11}+P_{22}+P_{33}=0`$ and $`P_{ik}=P_{ki}`$, then getting $`P_{33}=const`$ from the last equation in (48)we can conclude that $`P_{11}+P_{22}=const`$
### 4.1 Contribution from the tensor electric polarizability to deuteron spin oscillation
From the system (48) it follows
$`\begin{array}{c}\frac{d(P_{11}P_{22})}{dt}=4\mathrm{\Omega }P_{12},\hfill \\ \\ \frac{d^2P_{12}}{dt^2}=\mathrm{\Omega }\frac{d(P_{11}P_{22})}{dt}\mathrm{\Omega }_T^{}\frac{dP_3}{dt}=(4\mathrm{\Omega }^2+\mathrm{\Omega }_T\mathrm{\Omega }_T^{})P_{12}.\hfill \end{array}`$ (58)
Thus we have the equation
$$\frac{d^2P_{12}}{dt^2}+\omega _{12}^2P_{12}=0$$
(59)
where $`\omega _{12}=\sqrt{4\mathrm{\Omega }^2+\mathrm{\Omega }_T\mathrm{\Omega }_T^{}}2\mathrm{\Omega }`$, because $`\mathrm{\Omega }_T\mathrm{\Omega }_T^{}\mathrm{\Omega }^2`$.
The solution for this equation can be found in the form:
$$P_{12}=c_1\mathrm{cos}\omega _{12}t+c_2\mathrm{sin}\omega _{12}t$$
(60)
Let us find coefficients $`c_1`$ and $`c_2`$: when $`t=0`$ the equation (60) gives $`c_1=P_{12}(0)`$. The coefficient $`c_2`$ can be found from
$$\frac{d(P_{12})}{dt}(t0)=\omega _{12}c_2,$$
(61)
therefore
$$c_2=\frac{1}{\omega _{12}}\frac{d(P_{12})}{dt}(t0),$$
(62)
From the equation (54)
$$\frac{dP_{12}}{dt}(t0)=\mathrm{\Omega }\left(P_{11}(t0)P_{22}(t0)\right),$$
(63)
that
$$c_2=\frac{P_{11}P_{22}}{2},$$
(64)
and
$$P_{12}=P_{12}(0)\mathrm{cos}\omega _{12}t\frac{P_{11}P_{22}}{2}\mathrm{sin}\omega _{12}t$$
(65)
As a result we can write the following equation for the vertical component of the spin $`P_3`$:
$$\frac{dP_3}{dt}=\mathrm{\Omega }_TP_{12}(t)=\mathrm{\Omega }_T[P_{12}(0)\mathrm{cos}2\mathrm{\Omega }t\frac{P_{11}(0)P_{22}(0)}{2}\mathrm{sin}2\mathrm{\Omega }t]$$
(66)
As it can be seen the vertical component of the spin oscillates with the frequency $`2\mathrm{\Omega }`$.
But it should be mentioned that according to the equations (13) interaction of the EDM with an electric field causes oscillations of the vertical component of the spin with the frequency $`\mathrm{\Omega }`$. According to the idea these oscillations can be eliminated if the deuteron velocity is modulated with the frequency $`\mathrm{\Omega }`$:
$$v=v_0+\delta v\mathrm{sin}(\mathrm{\Omega }t+\phi )$$
(67)
As a result $`E_{eff}`$ depends on $`\stackrel{}{\beta }=\stackrel{}{v}/c`$ it also appears modulated:
$$E_{eff}=E_{eff}^0+\delta E_{eff}\mathrm{sin}(\mathrm{\Omega }t+\phi )$$
(68)
here $`\phi `$ is a phase. Therefore,
$$\frac{dP_3}{dt}=\mathrm{\Omega }_TP_{12}dE_{eff}P_2$$
(69)
as $`P_2`$ also oscillate with $`\mathrm{\Omega }`$ frequency, then in the product $`E_{eff}P_2`$ there non-oscillating terms and $`P_3`$ linearly grows with time.
It is important that modulation of the velocity $`v=v_0+\delta v\mathrm{sin}(\mathrm{\Omega }t+\phi )`$ results in oscillation of $`E_{eff}^2`$ also oscillates with time and appears proportional to $`\mathrm{sin}^2(\mathrm{\Omega }_T+\phi )`$. As a result
$$\frac{dP_3}{dt}\mathrm{\Delta }\mathrm{\Omega }_T\mathrm{sin}^2(\mathrm{\Omega }_Tt+\phi )[P_{12}(0)\mathrm{cos}2\mathrm{\Omega }t\frac{P_{11}(0)P_{22}(0)}{2}\mathrm{sin}2\mathrm{\Omega }t]$$
(70)
i.e.
$$\frac{dP_3}{dt}\frac{1}{2}\mathrm{\Delta }\mathrm{\Omega }_T\mathrm{cos}(2\mathrm{\Omega }t+2\phi )[P_{12}(0)\mathrm{cos}2\mathrm{\Omega }t\frac{P_{11}(0)P_{22}(0)}{2}\mathrm{sin}2\mathrm{\Omega }t]$$
(71)
According to (71) if the phase $`\phi =0`$ then the contribution to the linear growth of $`P_3`$ is provided by the term $`P_{12}(0)`$. If $`\phi =\pi /4`$ then linear growth is due to the second term proportional to $`(P_{11}(0)P_{22}(0))`$.
Measurement of these contribution provides to measure the tensor electric polarizability.
According to the evaluations $`\alpha _T10^{40}`$ cm<sup>3</sup> for the field $`E_{eff}=B10^4`$ gauss, therefore the frequency $`\mathrm{\Omega }_T10^5`$ sec<sup>-1</sup>. When considering modulation we should estimate $`\mathrm{\Delta }\mathrm{\Omega }_T\mathrm{\Omega }_T(\frac{\delta }{v_0})^2`$, then suppose $`(\frac{\delta }{v_0})^210^210^3`$ we obtain $`\mathrm{\Delta }\mathrm{\Omega }_T10^710^8`$ sec<sup>-1</sup>.
### 4.2 Contribution from the tensor magnetic polarizability to deuteron spin oscillation
Let us consider now contributions caused by the tensor magnetic polarizability $`\beta _T`$. Let we omit the terms proportional to the tensor electric polarizability in the system (48):
$`\begin{array}{c}\frac{dP_1}{dt}=\mathrm{\Omega }P_2\mathrm{\Omega }_T^\mu P_{23},\hfill \\ \\ \frac{dP_2}{dt}=\mathrm{\Omega }P_1+\mathrm{\Omega }_T^\mu P_{13},\hfill \\ \\ \frac{dP_{13}}{dt}=\mathrm{\Omega }P_{23}\mathrm{\Omega }_T^\mu P_2,\hfill \\ \\ \frac{dP_{23}}{dt}=\mathrm{\Omega }P_{13}+\mathrm{\Omega }_T^\mu P_1\hfill \end{array}`$ (79)
Introducing new variables $`P_+=P_1+iP_2`$ and $`G_+=P_{13}+iP_{23}`$ and recomposing equations (79) to determine $`P_+`$ and $`G_+`$ we obtain:
$`\begin{array}{c}\frac{dP_+}{dt}=i\mathrm{\Omega }P_++i\mathrm{\Omega }_T^\mu G_+,\hfill \\ \\ \frac{dG_+}{dt}=i\mathrm{\Omega }G_++i\mathrm{\Omega }_T^\mu P_+,\hfill \end{array}`$ (82)
or
$`\begin{array}{c}i\frac{dP_+}{dt}=\mathrm{\Omega }P_+\mathrm{\Omega }_T^\mu G_+,\hfill \\ \\ i\frac{dG_+}{dt}=\mathrm{\Omega }G_+\mathrm{\Omega }_T^\mu P_+,\hfill \end{array}`$ (85)
Let us search $`P_+,G_+e^{i\omega t}`$ then (85) transforms as follows:
$`\begin{array}{c}\omega \stackrel{~}{P}_+=\mathrm{\Omega }\stackrel{~}{P}_+\mathrm{\Omega }_T^\mu \stackrel{~}{G}_+,\hfill \\ \\ \omega \stackrel{~}{G}_+=\mathrm{\Omega }\stackrel{~}{G}_+\mathrm{\Omega }_T^\mu \stackrel{~}{P}_+.\hfill \end{array}`$ (88)
The solution of this system can be easily find:
$`\begin{array}{c}(\omega \mathrm{\Omega })^2\mathrm{\Omega }_T^\mu \mathrm{\Omega }_T^\mu =0\hfill \end{array}`$ (90)
that finally gives
$$\omega _{1,2}=\mathrm{\Omega }\pm \sqrt{\mathrm{\Omega }_T^\mu \mathrm{\Omega }_T^\mu }$$
(91)
Rewriting the solution
$$P_+(t)=c_1e^{i\omega _1t}+c_2e^{i\omega _2t}=|c_1|e^{i(\omega _1t+\delta _1)}+|c_2|e^{i(\omega _2t+\delta _2)}$$
(92)
Therefore,
$$P_1(t)=|c_1|\mathrm{cos}(\omega _1t+\delta _1)+|c_2|\mathrm{cos}(\omega _2t+\delta _2)$$
(93)
This means that spin rotates with two frequencies $`\omega _1`$ and $`\omega _2`$ instead of $`\mathrm{\Omega }`$ and, therefore, experiences beating with the frequency $`\mathrm{\Delta }\omega =\omega _1\omega _2=2\sqrt{\mathrm{\Omega }_T^\mu \mathrm{\Omega }_T^\mu }=\frac{\beta _TB_{eff}^2}{\mathrm{}}`$.
According to the evaluation the tensor magnetic polarizability $`\beta _T210^{40}`$, therefore for the beating frequency $`\mathrm{\Delta }\omega 10^5`$ in the field $`B10^4`$ gauss.
Measurement of the frequency of this beating makes possible to measure the tensor magnetic polarizability of the deuteron (nuclei).
Thus, due to the presence of tensor magnetic polarizability the the horizontal component of spin rotates around $`\stackrel{}{B}`$ with two frequencies $`\omega _1,\omega _2`$ instead of expected rotation with the frequency $`\mathrm{\Omega }`$. The resulting motion of the spin is beating: $`P_1(t)\mathrm{cos}\mathrm{\Omega }t\mathrm{sin}\mathrm{\Delta }\omega t`$.
This is the reason for the component $`P_3`$ caused by the EDM to experience the similar beating. Therefore, particle velocity modulation with the frequency $`\mathrm{\Omega }`$ ($`v=v_0+\delta v\mathrm{sin}(\mathrm{\Omega }t+\phi )`$) provides for eliminating oscillation with $`\mathrm{\Omega }`$ frequency, but oscillations with the frequency $`\mathrm{\Delta }\omega `$ rest.
## 5 Spin rotation of proton (deuteron, antiproton) in a storage ring with a polarized target and paramagnetic resonance in the nuclear pseudoelectric and pseudomagnetic fields
Another class of experiments deals with the use of polarized targets. Preparing such experiment one should remember that density of polarized gas target is lower than nonpolarized that and for example for COSY density of polarized target is $`j=10^{14}cm^2`$.
In 1964 it was shown that while slow neutrons are propagating through the target with polarized nuclei a new effect of nucleon spin precession occurred. It is stipulated by the fact that in a polarized target the neutrons are characterized by two refraction indices ($`N_{}`$ for neutrons with the spin parallel to the target polarization vector and $`N_{}`$ for neutrons with the opposite spin orientation , $`N_{}N_{})`$. According to the , in the target with polarized nuclei there is a nuclear pseudomagnetic field and the interaction of an incident neutron with this field results in neutron spin rotation. The results obtained in , initiated experiments which proved the existence of this effect -.
The effective potential energy of a particle in the pseudomagnetic nuclear field $`\stackrel{}{G}`$ of matter can be written as:
$$\widehat{V}_B^{nucl}=\stackrel{}{\mu }\stackrel{}{G},$$
(94)
where $`\stackrel{}{\mu }`$ is the magnetic moment of the particle and $`\stackrel{}{G}`$ can be expressed as -
$`\begin{array}{c}\stackrel{}{G}=\stackrel{}{G}_s+\stackrel{}{G}_w,\\ \\ \stackrel{}{G}_s=\frac{2\pi \mathrm{}^2}{\mu m}\rho [A_1\stackrel{}{J}+A_2\stackrel{}{n}(\stackrel{}{n}\stackrel{}{J})+\mathrm{}],\\ \\ \stackrel{}{G}_w=\frac{2\pi \mathrm{}^2}{\mu m}\rho [b\stackrel{}{n}+b_1[\stackrel{}{J}\times \stackrel{}{n}]+b_2\stackrel{}{n}_1+b_3\stackrel{}{n}(\stackrel{}{n}\stackrel{}{n}_1)+b_5[\stackrel{}{n}\times \stackrel{}{n}_1]+\mathrm{}]\end{array}`$ (100)
where $`\stackrel{}{n}=\stackrel{}{v}/v`$, $`\stackrel{}{J}`$ is the spin of nuclei of matter, $`\stackrel{}{J}=\text{Sp}\rho _{nucl}\stackrel{}{J}`$ is the average value of nuclear spin, $`\stackrel{}{n}_1`$ has the components $`n_{1j}=Q_{ij}n_j`$, where $`Q_{ij}=\text{Sp}\rho _{nucl}Q_{ij}`$ is the polarization tensor
$$Q_{ij}=\frac{1}{2J(2J1)}\left\{J_iJ_j+J_jJ_i\frac{2}{3}J(J+1)\delta _{ij}\right\},$$
(101)
It is easy to see that interaction (94) looks like the interaction of a magnetic moment with a magnetic field, thus the field $`\stackrel{}{G}`$ contributes to the change of the particle polarization similar a magnetic field does. It should be especially mentioned that $`\widehat{V}_B^{nucl}`$ contains both the real part, which is responsible for spin rotation, and imaginary part, which contributes to spin dichroism (i.e. beam absorption dependence on spin orientation). The detailed analysis of the effects caused by the nuclear pseudoelectric field was done in .
Interaction with the field $`\stackrel{}{G}=\stackrel{}{G}_s+\stackrel{}{G}_w`$ contains two summands: the first $`\stackrel{}{G}_s`$ corresponds to the strong interaction, which is T,P-even, while the second $`\stackrel{}{G}_w`$ describes spin rotation by the weak interaction, which has both T,P-odd (the term containing the constant $`b_1`$) and T-odd, P-even (the term containing the constant $`b_5`$) terms.
If either vector or tensor polarization of a target rotates then the effects provided by $`\stackrel{}{G}_s`$, $`\stackrel{}{G}_w`$ periodically depend on time i.e. equation (1) converts to:
$$i\mathrm{}\frac{\mathrm{\Psi }(t)}{t}=\left(\widehat{H}_0+\widehat{V}_{EDM}+\widehat{V}_\stackrel{}{E}+\widehat{V}_\stackrel{}{B}+\widehat{V}_E^{nucl}(t)+\widehat{V}_B^{nucl}(t)\right)\mathrm{\Psi }(t).$$
(102)
This equation coincides with the well-known equation for the paramagnetic resonance. Really, if we have the strong field orthogonal to the weak one (in this case $`\stackrel{}{B}\stackrel{}{G}`$) and $`\stackrel{}{G}`$ either rotates or oscillates with the frequency corresponding to the splitting, caused by the field $`\stackrel{}{B}`$, the resonance occurs. In our case this leads to the conversion of horizontal spin component to the vertical one with the frequency determined by the frequency of spin precession in the field $`\stackrel{}{G}`$. Thus we can measure all the constants containing in $`\stackrel{}{G}_s`$ and $`\stackrel{}{G}_w`$: constants $`A_i`$ give the spin-dependent part of elastic coherent forward scattering amplitude of proton (deuteron, antiproton) that is important for the projects at GSI and COSY; amplitudes $`b_i`$ provides to measure the constants of T-,P-odd interactions.
First of all we should pay attention to the effects caused by the T-odd nucleon-nucleon interaction of protons (antoprotons) and deuterons with polarized nuclei and, in particular, interaction described by $`V_{P,T}\stackrel{}{S}\left[\stackrel{}{p}_N\times \stackrel{}{n}\right]`$, where $`\stackrel{}{P}_N(t)`$ is the polarization vector of target. The interaction $`V_{P,T}`$ leads to the spin rotation around the axis determined by the unit vector $`\stackrel{}{n}_T`$ parallel to the vector $`\left[\stackrel{}{P}_N(t)\times \stackrel{}{n}\right]`$. Spin dichroism also appears with respect to this vector $`\stackrel{}{n}_T`$ i.e. a proton (deuteron) beam with the spin parallel to $`\stackrel{}{n}_T`$ has the absorption cross-section different from the absorption cross-section for a proton (deuteron) beam with the opposite spin direction.
P-even T-odd spin rotation, oscillation and dichroism of deuterons (nuclei with $`S1`$) caused by the interaction either $`V_T(\stackrel{}{S}[\stackrel{}{P}_N(t)\times \stackrel{}{n}])(\stackrel{}{S}\stackrel{}{n})`$ could be observed ; P-even T-odd spin rotation and dichroism for a proton, deuteron (nucleus with the spin $`S1/2`$) $`V_T^{}b_5[\stackrel{}{n}\times \stackrel{}{n}_1(t)]`$ could be observed in paramagnetic resonance conditions, too.
## 6 Conclusion
In the present paper the equations for the spin evolution of a particle in a storage ring are analyzed considering contributions from the tensor electric and magnetic polarizabilities of the particle. Study of spin rotation and birefringence effect for a particle in a high energy storage ring provides for measurement as the real part of the coherent elastic zero-angle scattering amplitude as well as tensor electric and magnetic polarizabilities.
We proposed the method for measurement the real part of the elastic coherent zero-angle scattering amplitude of particles and nuclei in a storage ring by the paramagnetic resonance in the periodical in time nuclear pseudoelectric and pseudomagnetic fields.
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# Effective critical behaviour of diluted Heisenberg-like magnets
## 1 Introduction
Relevance of structural disorder for the critical behaviour remains to be an important problem of modern condensed matter physics. Even a weak disorder may change drastically the behaviour near the critical point and in this respect may be related to the global characteristics of a physical system, such as the space dimension, order parameter symmetry and the origin of interparticle interaction. In this paper, we are going to discuss some peculiarities of a paramagnetic-ferromagnetic phase transition in magnets, where the randomness of structure has the form of substitutional random-site or random-bond quenched disorder. Solid solutions of magnets with small concentration of non-magnetic component as well as amorphous magnets with large relaxation times may serve as an example of such systems.
Intuitively, it is clear that for a weak enough disorder the ferromagnetic phase persists in such systems. Obviously, intuition fails to predict whether the critical exponents characterizing phase transition into ferromagnetic state will differ in a disordered system and in a “pure” one. The answer here is given by the Harris criterion which states that the critical exponents of the disordered system are changed only if the heat capacity critical exponent of a pure system is positive, otherwise the critical exponents of a disordered system coincide with those of a pure one. Returning to $`d=3`$ dimensional magnets with $`O(m)`$ symmetric spontaneous magnetization one is lead to the conclusion, that here only the critical exponents of uniaxial magnets described by the $`d=3`$ Ising model ($`m=1`$) are the subject of influence by weak quenched disorder. Indeed, the heat capacity diverges $`\alpha =0.109\pm 0.004>0`$ for $`m=1`$, whereas it does not diverge for the easy-plane and Heisenberg-like magnets: $`\alpha =0.011\pm 0.004`$ and $`\alpha =0.122\pm 0.010`$ for $`m=2`$ and $`m=3`$, respectively .
Note however that the Harris criterion tells about the scaling behaviour at the critical point $`T_c`$. In other words it predicts (possible) changes in the asymptotic values of the critical exponents defined at $`T_c`$. In real situations one often deals with the effective critical exponents governing scaling when $`T_c`$ still is not reached . These are non-universal. As far as in our study of particular interest will be the isothermal magnetic susceptibility $`\chi _T`$ let us define the corresponding effective exponent by :
$$\gamma _{\mathrm{eff}}(\tau )=\frac{\mathrm{d}\mathrm{ln}\chi (\tau )}{\mathrm{d}\mathrm{ln}\tau },\text{with}\tau =|TT_c|/T_c.$$
(1)
In the limit $`TT_c`$ the effective exponent coincides with the asymptotic one $`\gamma _{\mathrm{eff}}=\gamma `$.
Already in the first experimental studies of weakly diluted uniaxial (Ising-like) $`d=3`$ random magnets the asymptotic values of critical exponents were found. For the solid solutions, the exponents do not depend on the concentration of non-magnetic component and belong to the new universality class as predicted by the Harris criterion. We do not know analogous experiments where an influence of disorder on criticality of easy-plane magnets was examined. However its irrelevance was experimentally proven for the superfluid phase transition in $`\mathrm{He}^4`$ which belongs to the same $`O(2)`$ universality class as the ferromagnetic phase transition in easy-plane magnets.
As far as the disorder should be irrelevant for the asymptotic critical behaviour of the Heisenberg magnets, the diluted $`d=3`$ Heisenberg magnets should belong to the same $`O(3)`$ universality class as the pure ones. Theoretically predicted values of the isothermal magnetic susceptibility, correlation length, heat capacity, pair correlation function, and the order parameter asymptotic critical exponents in this universality class read :
$`\gamma =1.3895\pm 0.0050,\nu =0.7073\pm 0.0035,\alpha =0.122\pm 0.009,`$ (2)
$`\eta =0.0355\pm 0.0025,\beta =0.3662\pm 0.0025.`$
The experimental picture is more controversial. The bulk of experiments on critical behaviour of disordered Heisenberg-like magnets performed up to middle 80-ies is discussed in the comprehensive reviews . More recent experiments may be found in and references therein. We show typical results of measurements of the isothermal magnetic susceptibility effective critical exponent $`\gamma _{\mathrm{eff}}`$ (1) in Figs. 1. As it is seen from the pictures, the behaviour of $`\gamma _{\mathrm{eff}}`$ is non-monotonic. The exponent differs from its value predicted in the asymptotic limit (2) and is a subject of a wide crossover behaviour. Before reaching asymptotics $`\gamma _{\mathrm{eff}}`$ possess maximum (except of the fig. 1.d), the value of the maximum is system dependent: it differs for different magnets.
It is standard now to rely on the renormalization group (RG) method to get a reliable quantitative description of the behaviour in the vicinity of critical point. Namely in this way the cited above values (2) of the critical exponents of $`d=3`$ Heisenberg model were obtained. The RG approach appeared to be a powerful tool to describe asymptotic and effective critical behaviour of disordered Ising-like magnets as well. The purpose of the present paper is to describe the crossover behaviour of disordered Heisenberg-like magnets in frames of the field-theoretical RG technique. In particular we want to calculate theoretically the isothermal magnetic susceptibility effective critical exponent and to explain in this way the appearance of the peak in its typical experimental dependencies. The rest of the paper is organized as follows. In the Section 2 we formulate the model and review main theoretical results obtained for it so far by means of the RG technique, effective critical behaviour is analyzed in the Section 3, we end by conclusions and outlook in the Section 4.
## 2 The model and its RG analysis
The model of a random quenched magnet we are going to consider is described by the following Hamiltonian:
$$H=\frac{1}{2}\underset{𝐑,𝐑^{}}{}J(|𝐑𝐑^{}|)\stackrel{}{S}_𝐑\stackrel{}{S}_𝐑^{}c_𝐑c_𝐑^{}.$$
(3)
Here, the sum spans over all sites $`𝐑`$ of $`d`$-dimensional hypercubic lattice, $`J(|𝐑𝐑^{}|)`$ is a short-range (ferro)magnetic interaction between classical “spins” $`\stackrel{}{S}_𝐑`$ and $`\stackrel{}{S}_𝐑^{}`$. We consider the spins $`\stackrel{}{S}_𝐑`$ to be $`m`$-component vectors and the Hamiltonian (3) contains their scalar product. Obviously, for the particular case of Heisenberg spins we will put later $`m=3`$. The randomness is introduced into the Hamiltonian (3) by the occupation numbers $`c_𝐑`$ which are equal 1 if the site $`𝐑`$ is occupied by a spin and $`0`$ if the site is empty. Considering the case when occupied sites are distributed without any correlation and fixed in certain configuration one obtains so-called uncorrelated quenched $`m`$-vector model.
In principle, the above information is enough to apply the RG approach for a study of the critical behaviour of the model (3). One should obtain an effective Hamiltonian corresponding to the model under consideration and then one analyzes its long-distance properties by analyzing appropriate RG equations . But already on this step there are at least two different possibilities to proceed and both were exploited for the model (3). On one hand, to get the free energy of the model one can average the logarithm of configuration-dependent partition function over different possible configurations of disorder . Then, making use of the replica trick one arrives to the familiar effective Hamiltonian :
$$H_{\mathrm{eff}}=d^dR\{\frac{1}{2}\underset{\alpha =1}{\overset{n}{}}\left[\mu _{0}^{}{}_{}{}^{2}|\stackrel{}{\varphi }^\alpha |^2+|\stackrel{}{}\stackrel{}{\varphi }^\alpha |^2\right]+\frac{u_0}{4!}\underset{\alpha =1}{\overset{n}{}}|\stackrel{}{\varphi }^\alpha |^4+\frac{v_0}{4!}(\underset{\alpha =1}{\overset{n}{}}|\stackrel{}{\varphi }^\alpha |^2)^2\}$$
(4)
describing in the replica limit $`n0`$ critical properties of the model (3). Here, $`\mu _0`$ is a bare mass, $`u_0>0`$ and $`v_00`$ are bare couplings and $`\stackrel{}{\varphi }^\alpha \stackrel{}{\varphi }^\alpha (𝐑)`$ is an $`\alpha `$-replica of $`m`$-component vector field. The prevailing amount of RG studies of the critical behaviour of quenched $`m`$-vector model was performed on the base of the effective Hamiltonian (4) .
However, one more effective Hamiltonian corresponding to the model (3) is discussed in the literature . It is obtained exploiting the idea that a quenched disordered system can be described as an equilibrium system with additional forces of constraints . In such approach both variables $`\stackrel{}{S}_𝐑`$ and $`c_𝐑`$ are treated equivalently and one ends up with the effective Hamiltonian which differs from (4) and, consequently, leads to different results for the critical behaviour of the model (3) . Whereas the effective Hamiltonian (4) was used in the wide context of general $`m`$-vector models , the approach of Refs. was mainly used in explanations of crossover behaviour in Heisenberg-like systems . Below, we will discuss our results, based on the effective Hamiltonian (4) for $`m=3`$ and compare them with those derived in .
As it is well known, the renormalization group (RG) approach makes use of the scaling symmetry of the system in the asymptotic limit to extract the universal content and at the same time removes divergencies which occur for the evaluation of the bare functions in this limit . A change in the renormalized couplings $`u`$, $`v`$ of the effective Hamiltonian (3) under the RG transformation is described by the flow equations:
$$\mathrm{}\frac{\mathrm{d}}{\mathrm{d}\mathrm{}}u(\mathrm{})=\beta _u(u(\mathrm{}),v(\mathrm{})),\mathrm{}\frac{\mathrm{d}}{\mathrm{d}\mathrm{}}v(\mathrm{})=\beta _v(u(\mathrm{}),v(\mathrm{})).$$
(5)
Here, $`\mathrm{}`$ is the flow parameter related to the distance $`\tau `$ to the critical point. The fixed points ($`u^{},v^{}`$) of the system of differential equations (5) are given by:
$$\beta _u(u^{},v^{})=0,\beta _v(u^{},v^{})=0.$$
(6)
A fixed point is said to be stable if the stability matrix
$$B_{ij}\beta _{u_i}/u_j,i,j=1,2;u_i=\{u,v\},$$
(7)
possess in this point eigenvalues $`\omega _1,\omega _2`$ with positive real parts. In the limit $`\mathrm{}0`$, $`u(\mathrm{})`$ and $`v(\mathrm{})`$ attain the stable fixed point values $`u^{},v^{}`$. If the stable fixed point is reachable from the initial conditions (let us recall that for the effective Hamiltonian (3) they read $`u>0,v0`$) it corresponds to the critical point of the system. The asymptotic critical exponents values are defined by the fixed point values of the RG $`\gamma `$-functions. In particular the isothermal magnetic susceptibility exponent $`\gamma `$ is expressed in terms of the RG functions $`\gamma _\varphi `$ and $`\overline{\gamma }_{\varphi ^2}`$ describing renormalization of the field $`\varphi `$ and of the two-point vertex function with a $`\varphi ^2`$ insertion correspondingly :
$$\gamma ^1=1\frac{\overline{\gamma }_{\varphi ^2}}{2\gamma _\varphi }.$$
(8)
In Eq. (8), the functions $`\gamma _\varphi \gamma _\varphi (u,v)`$, $`\overline{\gamma }_{\varphi ^2}\overline{\gamma }_{\varphi ^2}(u,v)`$ are calculated in the stable fixed point $`u^{},v^{}`$. In the RG scheme, the effective critical exponents are calculated in the region, where couplings $`u(\mathrm{}),v(\mathrm{})`$ have not reached their fixed point values and depend on $`\mathrm{}`$. In particular for the exponent $`\gamma _{\mathrm{eff}}`$ one gets:
$$\gamma _{\mathrm{eff}}^1(\tau )=1\frac{\overline{\gamma }_{\varphi ^2}[u\{\mathrm{}(\tau )\},v\{\mathrm{}(\tau )\}]}{2\gamma _\varphi [u\{\mathrm{}(\tau )\},v\{\mathrm{}(\tau )\}]}+\mathrm{}.$$
(9)
In (9) the part denoted by dots is proportional to the $`\beta `$–functions (5) and comes from the change of the amplitude part of the susceptibility. In the subsequent calculations we will neglect this part, taking the contribution of the amplitude function to the crossover to be small .
For the effective Hamiltonian (4), the fixed point structure is well established . It is schematically shown in Figs. 2.a, 2.b. Two qualitatively different scenarios are observed: for $`m>m_c`$ the critical behaviour of the disordered magnet is governed by the fixed point of the pure magnet ($`u^{}>0`$, $`v^{}=0`$), whereas for $`m<m_c`$ the new stable fixed point ($`u^{}>0`$, $`v^{}<0`$) governs the asymptotic critical behaviour of the disordered magnet. At the marginal dimensionality $`m_c`$ which separates these two regimes, the $`\alpha `$ exponent of the pure magnet equals zero in agreement with the Harris criterion.
Best theoretical estimates of $`m_c`$ definitely support $`m_c<2`$: $`m_c=1.942\pm 0.026`$ , $`m_c=1.912\pm 0.004`$ . Consequently, the fixed point structure of the model of diluted Heisenberg-like magnet ($`m=3`$) is given by Fig. 2.a: the stable reachable fixed points of the diluted and pure Heisenberg-like magnets do coincide ($`u^{}0,v^{}=0`$), hence their asymptotic critical exponents do coincide as well. However the last statement does not concern the effective exponents. These are defined by the running values of the couplings $`u(\mathrm{})0,v(\mathrm{})0`$ and will be calculated in the next section.
## 3 The RG flows and the effective critical behaviour
The RG functions of the model (4) are known by now in pretty high orders of the perturbation theory . For the purpose of present study we will restrict ourselves by the first approximation where the described crossover phenomena manifests itself for the Heisenberg-like disordered magnets in non-trivial way. Within the two loop approximation in the minimal subtraction RG scheme the RG-functions read :
$`\beta _u(u,v)`$ $`=`$ $`u(\epsilon {\displaystyle \frac{m+8}{6}}u2v+{\displaystyle \frac{3m+14}{12}}u^2+{\displaystyle \frac{5mn+82}{36}}v^2+`$ (10)
$`{\displaystyle \frac{11m+58}{18}}uv),`$
$`\beta _v(u,v)`$ $`=`$ $`v(\epsilon {\displaystyle \frac{m+2}{3}}u{\displaystyle \frac{mn+8}{6}}v+{\displaystyle \frac{5(m+2)}{36}}u^2+`$ (11)
$`{\displaystyle \frac{3mn+14}{12}}v^2+{\displaystyle \frac{11(m+2)}{18}}uv),`$
$`\gamma _\varphi (u,v)`$ $`=`$ $`{\displaystyle \frac{m+2}{72}}u^2+{\displaystyle \frac{mn+2}{72}}v^2+{\displaystyle \frac{m+2}{36}}uv,`$ (12)
$`\overline{\gamma }_{\varphi ^2}(u,v)`$ $`=`$ $`{\displaystyle \frac{m+2}{6}}u+{\displaystyle \frac{mn+2}{6}}v{\displaystyle \frac{m+2}{12}}u^2{\displaystyle \frac{mn+2}{12}}v^2{\displaystyle \frac{m+2}{6}}uv.`$ (13)
Here, $`\epsilon =4d`$ and replica limit $`n=0`$ is to be taken.
Starting form the expressions (10)–(13) one can either develop the $`\epsilon `$-expansion, or work directly at $`d=3`$ putting in (10), (11) $`\epsilon =1`$ and considering renormalized couplings $`u,v`$ as the expansion parameters . However, such RG perturbation theory series with several couplings are known to be asymptotic at best . One should apply appropriate resummation technique to improve their convergence to get reliable numerical data on their basis. We used several different resummation schemes for this purpose. Here we will give the results obtained by the method which allowed to analyze the largest region in the parametric $`uv`$ space. The method was proposed in Ref. and was successfully applied to study random $`d=3`$ Ising model . Moreover, it was shown that the RG functions of the $`d=0`$ random Ising model are Borel-summable by this method . The main idea proposed in Ref. is to consider resummation in variables $`u`$ and $`v`$ separately. Taken that the RG function $`f(u,v)`$ is given to the order of $`p`$ loops, one first rewrites it as a power series in $`v`$:
$$f(u,v)=\underset{k=0}{\overset{p}{}}A_k(u)v^k.$$
(14)
Then each coefficient $`A_k(u)`$ is considered as power series in $`u`$ and resummed as a function of a single variable $`u`$ thus obtaining the resummed functions $`A_k^{res}(u)`$. Next one substitutes these functions into (14) and resums the RG function $`f`$ in single variable $`v`$. For the resummation in a single variable one may use any of familiar methods. Our results are obtained by making use of the Padé-Borel-Leroy method .
First, applying the above described resummation procedure to the $`\beta `$-functions (10), (11) we get the pure Heisenberg fixed point coordinates $`u^{}=0.8956`$, $`v^{}=0`$. The stability matrix (7) eigenvalues are positive at this fixed point ($`\omega _1=0.577`$, $`\omega _2=0.147`$) providing its stability. Then for the resummed values of the asymptotic critical exponents we get :
$$\gamma =1.382,\nu =0.701,\alpha =0.104,\eta =0.030,\beta =0.361.$$
(15)
We do not give the confidence intervals in (15), as far as they can be estimated only by comparison of changes introduced by different orders of perturbation theory. Note however that the results (15) are in a good agreement with the most accurate estimates of the exponents in the $`O(3)`$ universality class (2). This brings about that both the considered here two-loop approximation as well as the chosen resummation technique give an adequate description of asymptotic critical phenomena.
Before passing to the effective critical exponents let us first analyze the corrections to scaling. For the pure Heisenberg magnet, taking into account the leading correction to scaling results in the following formula for the isothermal susceptibility:
$$\chi (\tau )=\mathrm{\Gamma }_0\tau ^\gamma (1+\mathrm{\Gamma }_1\tau ^\mathrm{\Delta }),$$
(16)
where the correction-to-scaling exponent is given by $`\mathrm{\Delta }=\omega \nu `$ with $`\omega =\beta _u(u)/u|_{u=u^{}}`$ and non-universal critical amplitudes $`\mathrm{\Gamma }_0,\mathrm{\Gamma }_1`$. For the diluted Heisenberg magnet the corresponding formula includes two leading corrections $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$ (see e.g. ):
$$\chi _(\tau )=\mathrm{\Gamma }^{}_0\tau ^\gamma (1+\mathrm{\Gamma }_1^{}\tau ^{\mathrm{\Delta }_1}+\mathrm{\Gamma }_2^{}\tau ^{\mathrm{\Delta }_2}),$$
(17)
with critical amplitudes $`\mathrm{\Gamma }_0^{},\mathrm{\Gamma }_1^{},\mathrm{\Gamma }_2^{}`$. The exponents $`\mathrm{\Delta }_i`$ are expressed in terms of the stability matrix (7) eigenvalues $`\omega _i`$ in the pure Heisenberg fixed point: $`\mathrm{\Delta }_i=\nu \omega _i`$. At this fixed point, the eigenvalues of the stability matrix (7) read:
$$\omega _1=\frac{\beta _u(u,v)}{u}|_{u^{}0,v^{}=0},\omega _2=\frac{\beta _v(u,v)}{v}|_{u^{}0,v^{}=0}.$$
(18)
It is straightforward to see that the value $`\omega _1`$ coincides with the exponent $`\omega `$ of the pure model whereas it may be shown (see e.g. ) that the exponent $`\omega _2=|\alpha |/\nu `$ where $`\alpha `$ and $`\nu `$ are the heat capacity and correlation length critical exponent of the pure Heisenberg model. On the base of the numerical values of the exponents (15) we get:
$$\mathrm{\Delta }_1=0.405,\mathrm{\Delta }_2=0.104.$$
(19)
Again, obtained by us in the two-loop approximation numbers (19) can be compared with those in the six-loop approximation making use of the data (2) together with the value of $`\omega `$ of pure 3d Heisenberg model $`\omega =0.782\pm 0.0013`$ . As we have noted above, in order to get the numerical values of the correction-to-scaling exponents of diluted Heisenberg model it is no need to consider the RG functions (10)–(13) in the whole region of couplings $`u,v`$: it is enough to know them for the case of the pure model (i.e. for $`u0,v=0`$). However, to get the effective exponents it is necessary to study complete set of the RG functions (10)–(13) working also in the region where both couplings $`u`$ and $`v`$ differ from zero.
To this end we use the above described resummation technique in order to restore the convergence of the RG expansions in couplings $`u`$, $`v`$. First we solve the system of differential equations (5) and get the running values of couplings $`u(\mathrm{})`$, $`v(\mathrm{})`$ (10)–(13). They define the flow in the parametric space $`u,v`$ and in the limit $`\mathrm{}0`$ attain the stable fixed point value (shown by the filled box in Fig. 3).
Character of the flow depends on the initial conditions $`u_0,v_0`$ for solving the system of differential equations (5). For the model (3), the coupling $`v`$ is proportional to variance of disorder thus one can use the ratio $`|v_0/u_0|`$ to define the degree of dilution. Typical flows which are obtained for different ratios $`|v_0/u_0|`$ are shown in Fig. 3 by curves 1-3. We choose the starting values in the region with the appropriate signs of couplings $`u>0,v<0`$ near the origin (in the vicinity of the Gaussian fixed point $`u^{}=v^{}=0`$ shown by the filled circle in the figure). The flow No 1 is obtained for $`v_0=0`$, it corresponds to the pure Heisenberg model. The flow No 2 results from the small ratio $`|v_0/u_0|`$ and corresponds to the weak disorder whereas the flow No 3 is obtained for large $`|v_0/u_0|`$ and corresponds to the stronger dilution.
Obtained running values of coupling constant presented by flows in Fig. 3 allow one to get the effective critical exponents. Calculating resummed expression for the effective exponent $`\gamma _{\mathrm{eff}}`$ (9) along the flows 1-3 we get the results shown in the Fig. 4.
Again, the curve 1 corresponds to the effective critical exponent of the pure Heisenberg model, whereas curves 2 and 3 provide two possible scenarios for the effective exponents of the disordered Heisenberg model. Curve 2 corresponds to the weak dilution region: here, the exponent increases with approach to the critical point, although the crossover region is larger in comparison with the pure magnet (compare curves 1 and 2 in Fig. 4). This may lead to the peculiar situation that the asymptotic value of the exponent is reached earlier than the asymptotic values of the coupling. The effective exponents for the flows originating from non-zero ratio $`|v_0/u_0|`$ always attain the value which are larger than the asymptotic one. But the absolute value of this “overshooting” for small enough $`|v_0/u_0|`$ is too small to be observed experimentally. An experimental observation of such type of $`\gamma _{\mathrm{eff}}`$ behaviour of the disordered Heisenberg-like magnet is provided e.g. by Fig. 1d. Different behaviour of $`\gamma _{\mathrm{eff}}`$ is demonstrated by the curve 3 in Fig. 4. Here, before reaching the asymptotic region the exponent possess a distinct peak. Such behaviour is in agreement with observed experimental data presented by Figs. 1a1c, 1e1h. The value of maximum depends on the initial values for the RG flows. Larger ratio $`|v_0/u_0|`$ (i.e. stronger disorder) leads to the larger maximum. Thus, within unique approach one may explain both scenarios observed in the diluted Heisenberg-like magnets effective critical exponent $`\gamma _{\mathrm{eff}}`$ behaviour.
As we have noticed in the section 2, the crossover behaviour of random Heisenberg-like magnets was analyzed by means of an alternative approach in . There, the quenched disordered magnet was described as an equilibrium one with additional forces of constraints . This resulted in an effective Hamiltonian which differs from (4). The fixed point structure of this Hamiltonian differs from those given in Fig. 2 and, for different concentrations, leads to different crossover regimes. In particular, it predicts that there exists a limiting value of concentration where the critical behaviour is governed by Fisher-renormalized tricritical exponents which coincide with those of a $`d=3`$ spherical model: $`\gamma =2`$, $`\nu =1`$, $`\alpha =1`$, $`\eta =0`$, $`\beta =1/2`$. There exist two more fixed points which may be stable in the weak dilution regime. Their stability differs in different orders of the perturbation theory (compare and ) but the numerical values of the critical exponents do not differ essentially at these fixed points. The maximal possible value of the effective critical exponent $`\gamma _{\mathrm{eff}}`$ has been estimated as $`\gamma _{\mathrm{eff}}2.6`$ . However, the distinct feature of the behaviour of $`\gamma _{\mathrm{eff}}(\tau )`$ obtained in is its monotonic dependence. Hence, the experimentally observed peaks (see Fig. 1) can not be explained within such approach.
## 4 Conclusions
In the present paper we used the field-theoretical RG technique to study the effective critical behaviour of diluted Heisenberg-like magnets. The question of particular interest was to explain the peak in the exponent $`\gamma _{\mathrm{eff}}`$ as function of distance from $`T_c`$ observed in some experiments. Our two-loop calculations refined by the resummation of the perturbation theory series resulted in typical behaviour of diluted Heisenberg-like magnets $`\gamma _{\mathrm{eff}}`$ exponent represented by curves 2 and 3 in Fig. 4. The exponent can either reach it asymptotic value without demonstrating distinct maximum or it can first reach the peak and then cross-over to the asymptotic value from above. The strength of disorder is a physical reason which discriminates between these two regimes.
Our calculations are quite general and do not specify any particular object. In order to fit our curves to certain experiment one should include into consideration non-universal parameters to specify the magnetic system. The same concerns the flow parameter $`\mathrm{}`$ which as we have already noted is related to the distance to the critical point $`\tau `$. In principle such calculations may be done. However we want to emphasize that our analysis shows the reason of the peak in $`\gamma _{\mathrm{eff}}(\tau )`$ dependence for different disordered magnets which belong to the $`O(3)`$ universality class and this reason may be explained within the traditional RG approach. This concerns not only the magnetic susceptibility effective critical exponent. One more example is given by the order parameter effective exponents $`\beta _{\mathrm{eff}}`$ which has minimum when $`\tau `$ goes to zero (see e.g. ). Interpretation of this effect will be the goal of a separate study
In conclusion we want to note that similar peculiarities of the effective critical behaviour may be observed in studies of disordered easy-plane magnets which belong to the $`O(2)`$ universality class. Since the heat capacity does not diverge in such systems, the RG fixed point scenario is given by the Fig. 2a as for the Heisenberg-like disordered magnets. Up to our knowledge such experiments have not been performed yet and we hope that our calculations may stimulate them.
M. D. acknowledges the Ernst Mach research fellowship of the Österreichisher Austauschdienst. This work was supported in part by Österreichische Nationalbank Jubiläumsfonds through grant No 7694.
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# Testing the dynamics of 𝑩→𝝅𝝅 and constraints on 𝜶
## I Introduction
Nonleptonic $`B`$ decays to light hadrons provide information about $`CP`$ violation. In particular, the decays to $`\pi \pi `$, $`\rho \pi `$ and $`\rho \rho `$ can determine the weak phase $`\alpha `$. The theoretical challenge is to disentangle the strong interaction physics from the weak phase one would like to determine. For the decay $`B^0\pi ^+\pi ^{}`$ the $`B`$ factories study the $`CP`$ asymmetry,
$`{\displaystyle \frac{\mathrm{\Gamma }[\overline{B}{}_{}{}^{0}(t)\pi ^+\pi ^{}]\mathrm{\Gamma }[B^0(t)\pi ^+\pi ^{}]}{\mathrm{\Gamma }[\overline{B}{}_{}{}^{0}(t)\pi ^+\pi ^{}]+\mathrm{\Gamma }[B^0(t)\pi ^+\pi ^{}]}}`$ (1)
$`=S_+\mathrm{sin}(\mathrm{\Delta }mt)C_+\mathrm{cos}(\mathrm{\Delta }mt),`$
with the present world averages pipiCP ; Group:2004cx
$$S_+=0.50\pm 0.12,C_+=0.37\pm 0.10.$$
(2)
If the $`B\pi ^+\pi ^{}`$ amplitude were dominated by contributions with a single weak phase, the observable
$$\mathrm{sin}(2\alpha _{\mathrm{eff}})=S_+/\sqrt{1C_+^2},$$
(3)
would be equal to $`\mathrm{sin}2\alpha `$ and $`C_+`$ would be zero. The data indicate that this is not a good approximation. An isospin analysis Gronau:1990ka still allows a theoretically clean determination of $`\alpha `$ if the $`B^0\pi ^0\pi ^0`$ and $`\overline{B}{}_{}{}^{0}\pi ^0\pi ^0`$ rates are precisely measured. Since this requires very large data samples, several strategies have been proposed to extract $`\alpha `$ from $`\alpha _{\mathrm{eff}}`$ relying on theoretical inputs.
In the last few years the theory of $`B\pi \pi `$ decays has advanced considerably. Using the heavy quark limit, factorization theorems have been proven for the decay amplitudes at leading order in $`\mathrm{\Lambda }/m_b`$. The amplitudes in Eq. (I) arise from the matrix element of the effective Hamiltonian,
$`H_{\mathrm{eff}}={\displaystyle \frac{4G_F}{\sqrt{2}}}[\lambda _u(C_1O_1^u+C_2O_2^u+{\displaystyle \underset{i3}{}}C_i^cO_i)`$
$`+\lambda _c\left(C_1O_1^c+C_2O_2^c+{\displaystyle \underset{i3}{}}C_i^cO_i\right)`$
$`+\lambda _t{\displaystyle \underset{i3}{}}C_i^tO_i],`$ (4)
where CKM-unitarity was not used, and $`i=3,\mathrm{},6,8`$. (In the usual notation one has $`C_i=C_i^cC_i^t`$.) Its $`\overline{B}\pi \pi `$ matrix element can be parameterized as
$`\overline{A}(\overline{B}{}_{}{}^{0}\pi ^+\pi ^{})`$ $`=`$ $`\lambda _u(T+P_u)\lambda _cP_c\lambda _tP_t`$
$`=`$ $`e^{i\gamma }T_{\pi \pi }+e^{i\varphi }P_{\pi \pi },`$
$`\sqrt{2}\overline{A}(\overline{B}{}_{}{}^{0}\pi ^0\pi ^0)`$ $`=`$ $`\lambda _u(C+P_u)+\lambda _cP_c+\lambda _tP_t`$
$`=`$ $`e^{i\gamma }C_{\pi \pi }e^{i\varphi }P_{\pi \pi },`$
$`\sqrt{2}\overline{A}(B^{}\pi ^{}\pi ^0)`$ $`=`$ $`\lambda _u(T+C)=e^{i\gamma }T_0,`$ (5)
where $`\lambda _q=V_{qb}V_{qd}^{}`$. (We neglect isospin breaking isospinbreak and the contributions of electroweak penguins, the dominant part of which can be included model independently ewp .) In Eq. (I) $`T+P_u`$ and $`CP_u`$ are the $`B\pi ^+\pi ^{}`$ and $`B\pi ^0\pi ^0`$ matrix elements of the terms in the first line in Eq. (I), while $`P_c`$ and $`P_t`$ are the matrix elements of the second and third lines, respectively. This implies that each of the $`T+P_u`$, $`CP_u`$, $`P_c`$ and $`P_t`$ terms are separately renormalization group invariant.
There is an ambiguity in Eq. (I) related to the freedom in choosing the weak phase $`\varphi `$, in terms of which the amplitudes are written. There are two widely used conventions corresponding to eliminating either $`\lambda _t`$ or $`\lambda _c`$ using unitarity (some aspects of this were discussed in Refs. conv ). In the t-convention one eliminates $`\lambda _t`$ from Eq. (I), while in the c-convention one eliminates $`\lambda _c`$. Table 1 shows the expressions for the amplitudes and $`\varphi `$ in these conventions. Once a choice is made, $`T_{\pi \pi }`$, $`C_{\pi \pi }`$, $`P_{\pi \pi }`$, and $`T_0`$ can be extracted from the data, while further theoretical input is needed to determine $`T`$, $`C`$ and $`P_{u,c,t}`$.
The amplitudes in Eq. (I) (and their $`CP`$ conjugates) satisfy the isospin relation
$$\frac{1}{\sqrt{2}}\overline{A}(\overline{B}{}_{}{}^{0}\pi ^+\pi ^{})+\overline{A}(\overline{B}{}_{}{}^{0}\pi ^0\pi ^0)=\overline{A}(B^{}\pi ^{}\pi ^0).$$
(6)
The “tree” amplitudes also satisfy the relation
$$T_{\pi \pi }+C_{\pi \pi }=T_0,$$
(7)
which will play an important role in this paper, and we refer to it as the “tree triangle” (TT).
Expanding the amplitudes in soft-collinear effective theory (SCET) scet , one can define the leading (in $`\mathrm{\Lambda }/m_b`$) parts of $`T`$, $`C`$, and $`P_u`$ separately in terms of matrix elements of distinct SCET operators bp , which we denote with $`(0)`$ superscripts. The relative strong phase of $`T^{(0)}`$ and $`C^{(0)}`$ is suppressed by $`\alpha _s`$ bprs ; qcdf , and therefore
$$\varphi _T\mathrm{arg}\left(\frac{T^{(0)}+P_u^{(0)}}{T+C}\right)=𝒪[\alpha _s(m_b),\mathrm{\Lambda }_{\mathrm{QCD}}/m_b].$$
(8)
The numerator includes $`P_u^{(0)}`$ so that $`\varphi _T`$ is scale independent. The denominator could be defined to contain $`T^{(0)}+C^{(0)}`$, and our choice is for later convenience. Neither of these affect the right-hand side of Eq. (8) \[recall: $`P_u^{(0)}/T^{(0)}=𝒪(\alpha _s)`$\]. We define $`T^{(0)}T^{(0)}+P_u^{(0)}`$ and $`T+P_uT^{(0)}+P_u^{}`$, and in the rest of this paper the primes will be dropped. Thus, hereafter, $`P_u`$ contains the power suppressed corrections to $`T+P_u`$ (including weak annihilation).
The implications of Eq. (8) for the determination of $`\alpha `$ are obscured by the fact that $`T`$ and $`C`$ are not directly observable. The amplitudes $`T_{\pi \pi }`$ and $`C_{\pi \pi }`$ in Eq. (I) that can be extracted from the data include contributions from $`P_{u,c,t}`$. The heavy quark limit also determines the power counting for the penguin amplitudes, however, the convergence of the expansion for the penguins is less clear than it is for the trees. At leading order in $`\mathrm{\Lambda }/m_b`$ the calculable parts of $`P_{u,c,t}`$ are suppressed by $`\alpha _s`$ or the small Wilson coefficients $`C_{3,4}`$. At subleading order, the QCD factorization (QCDF) formula for $`P_t`$ contains sizeable “chirally enhanced” corrections, comparable to the leading order term qcdf . The possible size of nonperturbative contributions to $`P_c`$ has also been the subject of debate bprs ; pcdedbate . A large $`P_c`$ amplitude was found in fits using the leading order factorization results in SCET bprs , or adding a free parameter to the leading order QCDF result cpengs . In QCDF $`P_c`$ is claimed to be computable at leading order without nonperturbative inputs, while $`P_t`$ receives sizable “chirally enhanced” $`𝒪(\mathrm{\Lambda }/m_b)`$ corrections. Equation (8) and allowing for large long distance contribution to $`P_c`$ was used in Ref. brs to determine $`\alpha `$ without using the measurement of $`C_{00}`$ (the direct $`CP`$ asymmetry in $`B\pi ^0\pi ^0`$).
The penguin amplitudes $`P_c`$ and $`P_t`$ introduce a difference between the TTs in the two conventions. The $`P_u`$ amplitude is common to $`T_{\pi \pi }`$ in the t- and c-conventions, but $`P_c`$ enters $`T_{\pi \pi }`$ in the c-convention and $`P_t`$ enters $`T_{\pi \pi }`$ in the t-convention. Understanding the relative hierarchy of the three penguin amplitudes, $`P_{u,c,t}`$, is important if one is to use Eq. (8) for the determination of $`\alpha `$. In addition, it may also shed light on the $`\mathrm{\Lambda }/m_b`$ power counting for the penguin amplitudes. In this paper we show that by comparing the shapes of the TT in the c and t-conventions we can gain empirical knowledge about the relative sizes of $`P_u`$, $`P_c`$ and $`P_t`$.
## II Isospin analysis and tree triangle
The isospin relation in Eq. (6) holds for both the $`\overline{B}`$ and $`B`$ decay amplitudes, denoted by $`\overline{A}`$ and $`A`$, respectively. It is convenient to define $`\stackrel{~}{A}^{ij}=e^{2i\gamma }\overline{A}^{ij}`$, so that $`A^{0+}=\stackrel{~}{A}^0`$. Figure 1 shows the resulting two isospin triangles, $`WZX`$ and $`WZY`$, where the tree triangle, $`WZV`$, is also drawn. We follow the notation of Ref. GLSS , but normalize $`A(B^+\pi ^0\pi ^+)=\overline{WZ}=1`$.
To determine the TT from the data, recall that the $`WZX`$ and $`WZY`$ isospin triangles can be obtained from the direct $`CP`$ asymmetries $`C_+`$ and $`C_{00}`$, and the ratios of branching fractions
$`R_+`$ $`=`$ $`{\displaystyle \frac{(B^0\pi ^+\pi ^{})}{2(B^+\pi ^+\pi ^0)}}{\displaystyle \frac{\tau _{B^+}}{\tau _{B^0}}}=0.44_{0.06}^{+0.07},`$
$`R_{00}`$ $`=`$ $`{\displaystyle \frac{(B^0\pi ^0\pi ^0)}{(B^+\pi ^+\pi ^0)}}{\displaystyle \frac{\tau _{B^+}}{\tau _{B^0}}}=0.29_{0.06}^{+0.07},`$ (9)
where we used the experimental inputs from pipiBr ; Group:2004cx . Taking the ratios eliminates an arbitrary overall normalization parameter. To determine the coordinates of $`V`$, however, the measurement of $`S_+`$ is also needed.
It is convenient to define the coordinates of $`X`$ and $`Y`$ to be $`(\pm \mathrm{},0)`$, with
$$\mathrm{}^2=\frac{1}{2}R_+\left[1\sqrt{1C_+^2}\mathrm{cos}2\mathrm{\Delta }\alpha \right],$$
(10)
where $`\mathrm{\Delta }\alpha \alpha \alpha _{\mathrm{eff}}`$ and $`\alpha _{\mathrm{eff}}`$ is defined in Eq. (3). The four coordinates of $`W`$ and $`Z`$ and the phase $`\mathrm{\Delta }\alpha `$ are given by the solutions of the five equations GLSS
$`1`$ $`=`$ $`(x_Zx_W)^2+(y_Zy_W)^2,`$
$`R_{00}`$ $`=`$ $`x_Z^2+y_Z^2+\mathrm{}^2,`$
$`R_+`$ $`=`$ $`x_W^2+y_W^2+\mathrm{}^2,`$
$`R_+C_+`$ $`=`$ $`2\mathrm{}x_W,`$
$`R_{00}C_{00}`$ $`=`$ $`2\mathrm{}x_Z.`$ (11)
The $`XVY`$ angle is $`2(\varphi +\gamma )`$, so that the $`y`$ coordinate of $`V(0,y_V)`$ is
$$y_V=\{\begin{array}{cc}\mathrm{}\mathrm{cot}\gamma ,\hfill & \text{in the t-convention}\text{ ,}\hfill \\ \mathrm{}\mathrm{cot}\alpha ,\hfill & \text{in the c-convention}\text{ .}\hfill \end{array}$$
(12)
Equations (II) can be solved for $`\mathrm{\Delta }\alpha `$ and the coordinates of $`W`$ and $`Z`$. Because of the relative orientation of the amplitudes $`A^+`$ and $`\stackrel{~}{A}^+`$ adopted in Fig. 1, the solution must also satisfy $`\text{sgn}(\mathrm{\Delta }\alpha )=\text{sgn}(y_W)`$.
Some important properties of the solutions are apparent. First, $`x_W=0`$ if and only if $`C_+=0`$ (similarly, $`x_Z=0`$ if and only if $`C_{00}=0`$). Second, the sign of $`x_W`$ ($`x_Z`$) is opposite of that of $`C_+`$ ($`C_{00}`$). Thus, $`WZ`$ crosses the $`y`$ axis if and only if the direct $`CP`$ asymmetries in the charged and neutral modes have opposite signs.
In the rest of this section, we treat the simplified case where $`C_{00}`$ is not known. The first four equations in (II) can be used to solve for the coordinates of $`W`$ and $`Z`$ as functions of $`\mathrm{\Delta }\alpha `$. For any given value of $`\mathrm{\Delta }\alpha `$, $`W`$ and $`Z`$ are determined up to a two-fold ambiguity, corresponding to the reflection of $`Z`$ about the $`WO`$ line. These equations also place bounds on $`\mathrm{}`$ and $`\mathrm{\Delta }\alpha `$ Grossman:1997jr ; GLSS
$`\mathrm{}^2`$ $``$ $`R_+R_{00}{\displaystyle \frac{(1R_+R_{00})^2}{4}},`$
$`\mathrm{cos}(2\mathrm{\Delta }\alpha )`$ $``$ $`{\displaystyle \frac{(1+R_+R_{00})^22R_+}{2R_+\sqrt{1C_+^2}}}.`$ (13)
We refer to these inequalities as the isospin bound, and define $`\alpha _{\mathrm{bound}}\alpha _{\mathrm{eff}}\pm \mathrm{\Delta }\alpha _{\mathrm{max}}`$, which can be obtained from Eqs. (3) and (II), and $`\gamma _{\mathrm{bound}}\pi \beta \alpha _{\mathrm{bound}}`$. (Here, and in what follows $`\beta `$ is treated as known.) The coordinates of $`W`$ and $`Z`$ at the isospin bound satisfy
$$\frac{x_Z}{x_W}|_{\mathrm{bound}}=\frac{y_Z}{y_W}|_{\mathrm{bound}}=\frac{1+R_{00}R_+}{1R_{00}+R_+}.$$
(14)
This means that at the isospin bound $`W`$, $`Z`$, and $`O`$ are on one line and that at the bound
$$C_{00}|_{\mathrm{bound}}=\frac{R_+}{R_{00}}\frac{1+R_{00}R_+}{1R_{00}+R_+}C_+|_{\mathrm{bound}}.$$
(15)
The present data gives at the isospin bound $`C_{00}=(1.1\pm 0.1)C_+`$, which is almost $`2\sigma `$ from the measurements of $`C_+`$ in Eq. (2) and $`C_{00}=0.28_{0.40}^{+0.39}`$ c00 ; Group:2004cx .
In general, and even at the isospin bound, the $`V`$ vertex of the TT depends on $`S_+`$ via Eq. (12). Thus, the shape of the TT at the bound is not fixed, but depends on the experimental results. This dependence enters through $`\alpha _{\mathrm{eff}}+\mathrm{\Delta }\alpha `$ and implies that if one uses a constraint on the shape of the TT to extract $`\alpha `$, then i) the solution is not invariant under $`\mathrm{\Delta }\alpha \mathrm{\Delta }\alpha `$, and ii) the allowed values of $`\mathrm{\Delta }\alpha `$ are not the same for each discrete ambiguity of $`\alpha _{\mathrm{eff}}`$. Both of these points are different from the well-known symmetry properties of the usual isospin analysis.
The theory prediction of a small strong phase in Eq. (8) implies that the TT should be nearly flat, up to penguin contributions, small $`\alpha _s`$ and unknown $`\mathrm{\Lambda }/m_b`$ corrections. While the penguin contamination makes the definition of the TT itself convention dependent, it is interesting to consider under what conditions the TT can be flat, and its relation to the isospin bound. Since at the isospin bound $`W`$, $`Z`$, and $`O`$ are on a line, unless $`y_V=0`$, the TT is flat at the isospin bound if and only if $`x_W=x_Z=0`$. This implies that if any two of the following statements hold, then the other three follow:
| 1. The t-convention TT is flat for generic $`\alpha `$; |
| --- |
| 2. The c-convention TT is flat for generic $`\alpha `$; |
| 3. $`\alpha `$ is at the isospin bound; |
| 4. $`C_+=0`$; |
| 5. $`C_{00}=0`$. |
(16)
Equivalently, when one of the statements in (16) holds, the other four are either all true or all false. This shows that whether the TT is flat near the isospin bound or not depends on the value of $`\alpha `$; i.e., the TT being flat and $`\alpha `$ (or $`\gamma `$) being close to the isospin bound are in principle unrelated.
## III Constraints on $`𝜶`$
In Ref. brs , the predicted smallness of $`\varphi _T`$ and $`P_{ut}`$ was used to imply that the TT in the t-convention is (near) flat, which, in turn, was used to extract $`\gamma `$ without the insufficiently known $`C_{00}`$. In this section we discuss the implications of knowing an angle in the TT for the determination of $`\alpha `$, using a method which makes transparent the dependence of the constraints on $`\alpha `$ on the data.
For given $`R_+`$, $`R_{00}`$, and $`C_+`$, the first four equations in (II) together with (10) determine the coordinates of $`W`$ and $`Z`$ as functions of $`\mathrm{\Delta }\alpha `$. If, in addition, an angle in the TT is also known, then the position of the point $`V`$ is determined. We find it simplest to discuss the constraints in terms of the (convention dependent) observable phase,
$$\tau ^{(q)}\mathrm{arg}\left(\frac{T_{\pi \pi }^{(q)}}{T_0}\right)=\mathrm{arg}\left(1+\frac{P_{uq}}{T^{(0)}}\right)+\varphi _T,$$
(17)
where $`q=c`$ or $`t`$. The TT is near flat in either convention if $`|\tau |1`$. Note that if the penguin amplitudes vanished, then $`\tau ^{(t)}=\tau ^{(c)}=\varphi _T`$. We can determine the coordinates of $`V`$ as a function of $`\mathrm{\Delta }\alpha `$ in two ways: from the value of $`\tau `$ and the coordinates of $`W`$ and $`Z`$
$$y_V(\mathrm{\Delta }\alpha )=y_Wx_W\frac{y_Zy_W(x_Zx_W)\mathrm{tan}\tau }{x_Zx_W+(y_Zy_W)\mathrm{tan}\tau },$$
(18)
and from Eq. (12) if $`\beta `$, $`S_+`$ and $`C_+`$ are measured
$$y_V(\mathrm{\Delta }\alpha )=\{\begin{array}{cc}\mathrm{}\mathrm{cot}(\beta +\alpha _{\mathrm{eff}}+\mathrm{\Delta }\alpha ),\hfill & \text{t-convention}\text{,}\hfill \\ \mathrm{}\mathrm{cot}(\alpha _{\mathrm{eff}}+\mathrm{\Delta }\alpha ),\hfill & \text{c-convention}\text{.}\hfill \end{array}$$
(19)
The expression in (19) is convention dependent, because so is the definition of $`\tau `$ that enters in (18). These two equations form an implicit equation for $`\mathrm{\Delta }\alpha `$.
Figure 2 illustrates this method for the central values of the data. The solid curves show the solution for $`y_V(\mathrm{\Delta }\alpha )`$ vs. $`\mathrm{\Delta }\alpha `$ from Eq. (19): the darker (blue) curve corresponds to the t-convention and $`\alpha _{\mathrm{eff}}106^{}`$, while the lighter (red) curves correspond to the c-convention (the upper one for $`\alpha _{\mathrm{eff}}106^{}`$, the lower one for its mirror solution $`\alpha _{\mathrm{eff}}164^{}`$). The dashed curve shows $`y_V`$ vs. $`\mathrm{\Delta }\alpha `$ from Eq. (18) for $`\tau =0`$, and its intersections with the solid curves determine the value of $`\mathrm{\Delta }\alpha `$, which together with $`\alpha _{\mathrm{eff}}`$ gives $`\alpha `$. For the purpose of illustration the dotted curves show $`\tau =+10^{}`$ (lower curve) and $`10^{}`$ (up-most curve).
The $`\tau =0`$ curve goes to $`y_V=0`$ at the isospin bound (see Fig. 2), in accordance with our result in Sec. II that if $`\mathrm{\Delta }\alpha `$ is at the isospin bound and the TT is flat, then $`y_V=0`$. The right-hand side of Eq. (19) is small in this region of $`\mathrm{\Delta }\alpha `$, since the argument of the cotangent is close to $`90^{}`$ (the central values of the $`\pi \pi `$ data give $`\alpha _{\mathrm{eff}}106^{}`$, so that at the smallest value of $`\mathrm{\Delta }\alpha 28^{}`$, $`\beta +\alpha _{\mathrm{eff}}+\mathrm{\Delta }\alpha 102^{}`$ and $`\alpha _{\mathrm{eff}}+\mathrm{\Delta }\alpha 79^{}`$). These two facts imply that there is a solution for $`\mathrm{\Delta }\alpha `$ near the isospin bound with a flat TT; however, this is a coincidence and not a necessity.
In Ref. brs it was found that for small $`\tau ^{(t)}`$ the solution for $`\mathrm{\Delta }\alpha `$ was close to the isospin bound. This can be easily seen from Fig. 2. The dashed and dotted curves are steep near the bound for negative $`\mathrm{\Delta }\alpha `$, so changing $`\tau `$ hardly changes the solution for $`\mathrm{\Delta }\alpha `$. However, for the other solution (corresponding to positive $`\mathrm{\Delta }\alpha `$, and a value of $`\alpha `$ disfavored by the global CKM fit Charles:2004jd ), the error is significantly larger, since the dependence of $`\mathrm{\Delta }\alpha `$ on $`\tau `$ is stronger. The allowed region of $`\mathrm{\Delta }\alpha `$ is particularly sensitive to $`R_{00}`$; for example, for $`R_{00}=0.2`$ (which is a bit more than $`1\sigma `$ lower than its present central value) the $`|\tau |<10^{}`$ constraint would include almost all values of $`\mathrm{\Delta }\alpha `$ that are allowed by the isospin analysis. Note that with the current data the error of $`\alpha `$ extracted using the constraint of a small $`\tau `$ increases with decreasing $`R_{00}`$, contrary to the isospin analysis.
The confidence level (CL) of $`\alpha `$ obtained by imposing a constraint on $`\tau `$ is shown in Fig. 3 using the CKMfitter package Charles:2004jd . In the left plot the curves show (see the labels) the CL of $`\alpha `$ imposing $`\tau =0`$ in both the t- and c-conventions without using the $`C_{00}`$ measurement in the fit. For comparison, we also show the result of the usual isospin analysis with and without using $`C_{00}`$. The plot on the right-hand side shows the CL of $`\alpha `$ imposing $`\tau =0`$ in the t-convention with and without using $`C_{00}`$, and the constraint in the t-convention imposing $`|\tau |<5^{}`$, $`10^{}`$, and $`20^{}`$. The restriction on $`\alpha `$ from a constraint $`|\tau |<\tau _0`$ becomes quite weak as $`\tau _0`$ increases in the range $`10^{}<\tau _0<20^{}`$. We can compare our results with those of brs , which use as theory input an upper bound on $`ϵ=|\mathrm{Im}(C_{\pi \pi }^{(t)}/T_{\pi \pi }^{(t)})|`$. Assuming $`\{\gamma ,|\mathrm{arg}(P_{\pi \pi }/T_{\pi \pi })|\}<90^{}`$, we find $`\mathrm{sin}\tau ^{(t)}<ϵ\sqrt{R_+}`$, i.e., $`\tau ^{(t)}<15.5^{}(7.8^{})`$ for the bounds considered in brs , $`ϵ<0.4(0.2)`$.
Imposing $`\tau =0`$ gives only two solutions with $`\chi ^2=0`$ with the current data, around $`\alpha 78^{}`$ and $`132^{}`$. The first one, which is consistent with the Standard Model (SM) CKM fit, is disfavored by the measurement of $`C_{00}`$. While the two solutions have comparable errors for $`\tau =0`$, allowing a finite range of $`\tau `$ to account for subleading effects increases the error of the $`\alpha 132^{}`$ solution more rapidly. Imposing a bound on $`|\mathrm{Im}(C/T)|`$ brs allows, in addition to $`\tau `$ being near 0, that $`\tau `$ is near $`\pi `$ (mod $`2\pi `$); however, the theory disfavors the latter possibility. It is constraining $`|\tau |`$ modulo $`2\pi `$ and not $`\pi `$ that makes some of the CL curves not periodic with a period of $`\pi `$.
These results for $`\alpha `$ should not be taken at face value, because in the next Section we find that extracting $`\tau `$ using the SM CKM fit as an input gives significantly larger values of $`|\tau |`$ than considered here. The implications of this are discussed below.
## IV The penguin hierarchy problem
If the penguin amplitudes were small then the statements in (16) would all hold to a good precision, and $`\alpha `$ could be extracted simply from $`S_+`$. This is known not to be the case, so the question is to determine which penguins are large or small. This is complicated by the fact that, as explained in Sec. II, the amplitudes $`T`$, $`C`$, $`P_{uc}`$, and $`P_{ut}`$ are not separately observable from the $`B\pi \pi `$ data alone. They can be disentangled using $`SU(3)`$ flavor symmetry and data on $`BK\pi `$, $`K\overline{K}`$, etc.
In this section we propose to use the theory expectation for $`\varphi _T`$ in Eq. (8) to test the magnitude of the penguins. (Another test of corrections to factorization in $`B\pi \pi `$ was proposed in Feldmann:2004mg .) We assume $`\varphi _T=0`$, although we may learn from other data that power corrections to tree amplitudes are sizable. For example, a power suppressed strong phase around $`30^{}`$ is observed in $`BD\pi `$ decays Mantry:2003uz .
In the t-convention $`P_{ut}`$ (recall, $`P_{ij}P_iP_j`$) contributes to the TT in Eq. (7), while in the c-convention it is $`P_{uc}`$. (We choose, for convenience, the pure tree amplitude $`T_0`$ to be real.) Thus, comparing the TT in the two conventions teaches us about the relative size of $`P_{ut}`$ and $`P_{uc}`$. (The same information can in principle be obtained from the fit in any one convention; this comparison makes the results more transparent.) We use the SM global fit to the CKM matrix that determines the weak phase $`\gamma =(59.0_{4.9}^{+6.4})^{}`$ Charles:2004jd . This allows the construction of the tree triangles in both conventions, as explained in Sec. II. Comparing how flat they are, i.e., how small the angle $`\tau `$ of the TT is, the following outcomes are possible:
* $`|\tau ^{(t)}||\tau ^{(c)}|`$. This would imply $`\mathrm{Im}(P_{ut})\mathrm{Im}(P_{uc})`$, and the likely explanation would be $`|P_c||P_u||P_t|`$.
* $`|\tau ^{(t)}||\tau ^{(c)}|`$. This would imply $`\mathrm{Im}(P_{ut})\mathrm{Im}(P_{uc})`$, and the likely explanation would be $`|P_t||P_u||P_c|`$.
* $`|\tau ^{(t)}||\tau ^{(c)}|1`$. This would imply that both $`\mathrm{Im}(P_{ut}/T^{(0)})`$ and $`\mathrm{Im}(P_{uc}/T^{(0)})`$ are small. In this case the likely explanation would be that $`P_q/T^{(0)}`$ is small for each of the penguin amplitudes.
* $`|\tau ^{(t)}||\tau ^{(c)}|=𝒪(1)`$ and $`|\tau ^{(t)}\tau ^{(c)}|1`$. This would imply that $`\mathrm{Im}(P_{ut}/T^{(0)})`$ and $`\mathrm{Im}(P_{uc}/T^{(0)})`$ are both much larger than $`\mathrm{Im}(P_{ct}/T^{(0)})`$. There appears to be no single plausible explanation for such a case. It may indicate that $`P_u`$ (that includes weak annihilation) is large, while $`P_c`$ and $`P_t`$ are small or have small phases. Another, fine tuned, possibility is that both $`P_c`$ and $`P_t`$ have large but nearly equal phases. Last, it might be that $`\varphi _T=𝒪(1)`$, indicating large corrections to the heavy quark limit.
* $`|\tau ^{(t)}||\tau ^{(c)}|=𝒪(1)`$ and $`|\tau ^{(t)}\tau ^{(c)}|=𝒪(1)`$. This would imply that $`\mathrm{Im}(P_{ut}/T^{(0)})`$, $`\mathrm{Im}(P_{uc}/T^{(0)})`$, and $`\mathrm{Im}(P_{ct}/T^{(0)})`$ are all large. In this case the likely explanation would be that all penguins are large and comparable to $`T^{(0)}`$.
Note that the $`\tau ^{(t)}\tau ^{(c)}`$ difference is related to the penguin-to-tree ratio,
$$\tau ^{(t)}\tau ^{(c)}=\mathrm{arg}\left(1\frac{|\lambda _u|}{|\lambda _c|}\frac{P_{\pi \pi }^{(t)}}{T_{\pi \pi }^{(t)}}\right),$$
(20)
and can be determined with better precision than $`\tau ^{(t,c)}`$ separately.
### IV.1 $`𝑩\mathbf{}𝝅𝝅`$
Using the experimental data we can determine $`\tau `$ in the t- and c-conventions. The results for the confidence levels of $`\tau ^{(t,c)}`$ are shown in the left plot in Fig. 4. At the one sigma level only one solution is allowed (because $`C_{00}`$ disfavors one of the solutions at a near $`2\sigma `$ level). Including $`C_{00}`$ in the fit drives $`|\tau |`$ to larger values
$$\tau =\{\begin{array}{cc}\left(36_8^{+6}\right)^{},\hfill & \text{t-convention,}\hfill \\ \left(30_8^{+6}\right)^{},\hfill & \text{c-convention.}\hfill \end{array}$$
(21)
Note that the central values indicate rather large values for $`\tau `$ in both conventions. Their difference is more accurately determined by Eq. (20), where the fit gives
$$\tau ^{(t)}\tau ^{(c)}=\left(5.7_{1.7}^{+2.0}\right)^{}.$$
(22)
Eqs. (21) and (22) favor scenario (iv). While this may have several reasons as explained above, the least fine-tuned one, i.e., a large $`P_u`$ (including weak annihilation) and smaller $`P_{c,t}`$ penguins (or that the $`\varphi _T1`$ prediction receives large corrections), would be puzzling for any approach to factorization. At present, this is not a very firm conclusion yet. (Note that a similar enhancement of the $`u`$-penguin amplitude is observed in $`\overline{B}K\pi `$ and $`b(s\overline{s})s`$ decays, if the apparent anomalies therein are interpreted within the SM.)
### IV.2 $`𝑩\mathbf{}𝝆𝝆`$
Since $`B\rho \rho `$ decays are dominantly longitudinally polarized, the determination of $`\alpha `$ from this mode is very similar to that from $`B\pi \pi `$, except that at the few percent level an $`I=1`$ amplitude may be present Falk:2003uq . Using dynamical input to reduce the uncertainty of $`\alpha `$ from $`B\rho \rho `$ has received little attention so far, because the isospin bound puts tight constraints on $`\alpha \alpha _{\mathrm{eff}}`$. However, this bound may become worse in the future, since the strong present bound is a consequence of the fact that the isospin triangles do not close with the central values of the current world averages. This is a consequence of both the branching ratios, whose central values in units of $`10^3`$ are $`\sqrt{(B\rho ^+\rho ^0)}=5.14`$, $`\sqrt{(B\rho ^+\rho ^{})/2}=3.87`$, and $`\sqrt{(B\rho ^0\rho ^0)}<1.05`$ (90% CL), and the smallness of $`C_{\rho ^+\rho ^{}}=0.03\pm 0.20`$ Group:2004cx ; rhorho . Therefore, although at present imposing $`|\tau |<10^{}`$ does not improve the constraint on $`\alpha \alpha _{\mathrm{eff}}`$ in this mode, such a dynamical input may become useful in the future.
In this case, the $`\tau `$ values in the two conventions differ by less than a degree as shown in the right plot in Fig. 4, giving $`\tau =(0\pm 12)^{}`$. This may tend towards the above scenario (iii). If in the future the measured value of the $`B\rho ^+\rho ^0`$ branching ratio decreases (or that of $`\rho ^0\rho ^0`$ increases) then the pure isospin bound will become worse, and the fit results for $`\tau `$ will also change. If that fit still favors $`|\tau ^{(t)}||\tau ^{(c)}|`$ or $`|\tau ^{(t)}||\tau ^{(c)}|1`$ \[cases (i) or (iii)\] then we would feel comfortable imposing a constraint on the magnitude of $`\tau ^{(t)}`$ to improve the determination of the CKM angle $`\alpha `$.
## V Conclusions
The tree amplitudes in $`B\pi \pi `$ decays can be computed in an expansion of $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ using factorization. In the heavy quark limit the strong phase between the tree amplitudes is suppressed, which may help to improve the determination of the weak phase $`\alpha `$. Using this theory input as an additional constraint in the fit for $`\alpha `$, requires some understanding of the power corrections and penguin amplitudes.
While the present measurement of $`C_{00}`$ does not provide a significant determination of $`\alpha `$ from the $`B\pi \pi `$ isospin analysis, it provides useful information about the hadronic amplitudes. The determination of $`\alpha `$ using the central values of the present data with $`C_{00}`$ replaced by the assumption of a flat TT gives a solution near the isospin bound. While a $`|\tau ^{(t)}|<5^{}`$ or $`10^{}`$ theoretical bound is quite powerful to constrain $`\alpha `$, allowing for larger deviations from the heavy quark limit $`(|\tau ^{(t)}|<20^{})`$ reduces significantly the predictive power of the constraint on $`\alpha `$. The present $`C_{00}`$ result, however, disfavors being at the isospin bound at about the $`2\sigma `$ level. This observation is exhibited by the like-sign $`C_+`$ and $`C_{00}`$ measurements, whereas the opposite signs of the $`P_{\pi \pi }`$ terms in the $`\pi ^+\pi ^{}`$ and $`\pi ^0\pi ^0`$ amplitudes would imply opposite signs for $`C_+`$ and $`C_{00}`$ if the tree triangle was flat.
We proposed a comparison of fits that can give information about the relative size of the penguins, using only $`\pi \pi `$ data and the global fit for $`\gamma `$. While the present data is not yet precise enough to give firm conclusions, its most likely implication is that not only the charm (nor the top) penguins in $`B\pi \pi `$ are large, but so are the up penguins (including terms proportional to $`V_{ub}`$ that are power suppressed in the heavy quark limit), thus one may not be able to use theory instead of $`C_{00}`$. On the other hand, for $`B\rho \rho `$ decay, it may well be the case that the data will continue to favor $`|\tau ^{(t)}||\tau ^{(c)}|1`$ or $`|\tau ^{(t)}||\tau ^{(c)}|`$, in which case the theory can be useful to reduce the error on $`\alpha `$ without a measurement of $`C_{00}`$.
###### Acknowledgements.
We thank Christian Bauer, Andy Cohen, Marco Ciuchini, and Iain Stewart for useful discussions. Special thanks to Jérôme Charles for helpful comments and for pointing out a mistake in our earlier numerical results. We thank the Institute for Nuclear Theory at the University of Washington for its hospitality and partial support while some of this work was completed. Z.L. thanks the particle theory group at Boston University for its hospitality while part of this work was completed. This work was supported in part by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of High Energy Physics, of the U.S. Department of Energy under Contract DE-AC02-05CH11231 and by a DOE Outstanding Junior Investigator award (Z.L.); and by the U.S. Department of Energy under cooperative research agreement DOE-FC02-94ER40818 (D.P.).
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# On the spectrum of Jacobi operators with quasi-periodic algebro-geometric coefficients
## 1. Introduction
It is well-known from the work of Date and Tanaka , , Dubrovin, Matveev, and Novikov , Flaschka , McKean , McKean and van Moerbeke , , Mumford , Novikov, Manakov, Pitaevskii, and Zakharov , Teschl \[53, Chs. 9,13\], Toda \[54, Ch. 4\], \[55, Chs. 26-30\], van Moerbeke , van Moerbeke and Mumford , that the self-adjoint Jacobi operator
$$H=aS^++a^{}S^{}+b,\text{dom}(H)=\mathrm{}^2(),$$
(1.1)
in $`\mathrm{}^2()`$ with real-valued periodic, or more generally, algebro-geometric quasi-periodic and real-valued coefficients $`a`$ and $`b`$ (i.e., coefficients that satisfy one (and hence infinitely many) equation(s) of the stationary Toda hierarchy), leads to a finite-gap, or perhaps more appropriately, to a finite-band spectrum $`\sigma (H)`$ of the form
$$\sigma (H)=\underset{m=1}{\overset{p+1}{}}[E_{2m2},E_{2m1}],E_0<E_1<\mathrm{}<E_{2p+1}.$$
(1.2)
Compared to the real-valued case, the corresponding spectral properties of Jacobi operators with periodic and complex-valued coefficients $`a`$ and $`b`$, to the best of our knowledge, have been studied rather sparingly in the literature. We are only aware of two papers by Naĭman , , in which it is shown that the spectrum consists of a set of piecewise analytic arcs which may have common endpoints. (This is in anology to the case of (non-self-adjoint) one-dimensional periodic Schrödinger operators, cf. , .) It seems plausible that the latter case is connected with (complex-valued) stationary solutions of equations of the Toda hierarchy. In particular, the next scenario in line, the determination of the spectrum of $`H`$ in the case of quasi-periodic and complex-valued solutions of the stationary Toda hierarchy apparently has never been clarified. The latter spectral problems have been open since the late seventies and it is the purpose of this paper to provide a comprehensive solution of them. For theta function representations of $`a`$ and $`b`$ in the complex-valued algebro-geometric case (without addressing spectral theoretic questions) we refer, for instance, to Aptekarev , Dubrovin, Krichever, and Novikov , Krichever (cf. also the appendix written by Krichever in ).
To describe our results, a bit of preparation is needed. Let
$$G(z,n,n^{})=(Hz)^1(n,n^{}),z\backslash \sigma (H),n,n^{},$$
(1.3)
be the Green’s function of $`H`$ (here $`\sigma (H)`$ denotes the spectrum of $`H`$) and denote by $`g(z,n)`$ the corresponding diagonal Green’s function of $`H`$ defined by
$`g(z,n)=G(z,n,n)={\displaystyle \frac{_{j=1}^p[z\mu _j(n)]}{R_{2p+2}(z)^{1/2}}},`$ (1.4)
$`R_{2p+2}(z)={\displaystyle \underset{m=0}{\overset{2p+1}{}}}(zE_m),\{E_m\}_{m=0}^{2p+1},`$ (1.5)
$`E_mE_m^{}\text{ for }mm^{}\text{}m,m^{}=0,1,\mathrm{},2p+1\text{.}`$ (1.6)
For any quasi-periodic (in fact, Bohr (uniformly) almost periodic) sequence $`f=\{f(k)\}_k`$ the mean value $`f`$ of $`f`$ is defined by
$$f=\underset{N\mathrm{}}{lim}\frac{1}{2N+1}\underset{k=N}{\overset{N}{}}f(k).$$
(1.7)
Moreover, we introduce the set $`\mathrm{\Sigma }`$ by
$$\mathrm{\Sigma }=\left\{\lambda \right|\text{Re}\left(\text{ln}\left(\frac{G_{p+1}(\lambda ,)y}{G_{p+1}(\lambda ,)+y}\right)\right)=0\},$$
(1.8)
where $`y^2=R_{2p+2}(z)`$, and the polynomial $`G_{p+1}(z,n)`$ of degree $`p+1`$ in $`z`$ will be defined in (2.21). Here we observe that $`G_{p+1}`$ is given in terms of the off-diagonal Green’s function $`G(z,n+1,n)`$ by
$$\frac{G_{p+1}(z,n)}{R_{2p+2}(z)^{1/2}}=12a(n)G(z,n+1,n)$$
(1.9)
(cf. also (3.13)). In addition, we note that
$$g(z,)=\frac{_{j=1}^p\left(z\stackrel{~}{\lambda }_j\right)}{R_{2p+2}(z)^{1/2}}$$
(1.10)
for some constants $`\{\stackrel{~}{\lambda }_j\}_{j=1}^p`$.
Finally, we denote by $`\sigma _\mathrm{p}(T)`$, $`\sigma _\mathrm{r}(T)`$, $`\sigma _\mathrm{c}(T)`$, $`\sigma _\mathrm{e}(T)`$, and $`\sigma _{\mathrm{ap}}(T)`$, the point spectrum (i.e., the set of eigenvalues), the residual spectrum, the continuous spectrum, the essential spectrum (cf. (4.16)), and the approximate point spectrum of a densely defined closed operator $`T`$ in a complex Hilbert space, respectively.
Our principal new results, to be proved in Section 4, then read as follows:
###### Theorem 1.1.
Assume that $`a`$ and $`b`$ are quasi-periodic $`(`$complex-valued $`)`$ solutions of the $`p`$th stationary Toda equation associated with the hyperelliptic curve $`y^2=R_{2p+2}(z)`$ subject to (1.5) and (1.6). Then the following assertions hold:
$`(i)`$ The point spectrum and residual spectrum of $`H`$ are empty and hence the spectrum of $`H`$ is purely continuous,
$`\sigma _\mathrm{p}(H)=\sigma _\mathrm{r}(H)=\mathrm{},`$ (1.11)
$`\sigma (H)=\sigma _\mathrm{c}(H)=\sigma _\mathrm{e}(H)=\sigma _{\mathrm{ap}}(H).`$ (1.12)
$`(ii)`$ The spectrum of $`H`$ coincides with $`\mathrm{\Sigma }`$ and equals the conditional stability set of $`H`$,
$`\sigma (H)`$ $`=\left\{\lambda \right|\text{Re}\left(\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(\lambda ,)y}{G_{p+1}(\lambda ,)+y}}\right)\right)=0\}`$ (1.13)
$`=\{\lambda |\text{there exists at least one bounded solution}`$
$`0\psi \mathrm{}^{\mathrm{}}()\text{ of }H\psi =\lambda \psi \}.`$ (1.14)
$`(iii)`$ $`\sigma (H)`$ is bounded,
$$\sigma (H)\{z|\text{Re}(z)[M_1,M_2],\text{Im}(z)[M_3,M_4]\},$$
(1.15)
where
$`\begin{array}{cc}& M_1=2\underset{n}{sup}[|\text{Re}(a(n))|]+\underset{n}{inf}[\text{Re}(b(n))],\hfill \\ & M_2=2\underset{n}{sup}[|\text{Re}(a(n))|]+\underset{n}{sup}[\text{Re}(b(n))],\hfill \\ & M_3=2\underset{n}{sup}[|\text{Im}(a(n))|]+\underset{n}{inf}[\text{Im}(b(n))],\hfill \\ & M_4=2\underset{n}{sup}[|\text{Im}(a(n))|]+\underset{n}{sup}[\text{Im}(b(n))].\hfill \end{array}`$ (1.16)
$`(iv)`$ $`\sigma (H)`$ consists of finitely many simple analytic arcs. These analytic arcs may only end at the points $`\stackrel{~}{\lambda }_1,\mathrm{},\stackrel{~}{\lambda }_p`$, $`E_0,\mathrm{},E_{2p+1}`$.
$`(v)`$ Each $`E_m`$, $`m=0,\mathrm{},2p+1`$, is met by at least one of these arcs. More precisely, a particular $`E_{m_0}`$ is hit by precisely $`2N_0+1`$ analytic arcs, where $`N_0\{0,\mathrm{},p\}`$ denotes the number of $`\stackrel{~}{\lambda }_j`$ that coincide with $`E_{m_0}`$. Adjacent arcs meet at an angle $`2\pi /(2N_0+1)`$ at $`E_{m_0}`$. $`(`$Thus, generically, $`N_0=0`$ and precisely one arc hits $`E_{m_0}`$.$`)`$
$`(vi)`$ Crossings of spectral arcs are permitted and take place precisely when
$`\begin{array}{cc}& \text{Re}\left(\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(\stackrel{~}{\lambda }_{j_0},)y}{G_{p+1}(\stackrel{~}{\lambda }_{j_0},)+y}}\right)\right)=0\hfill \\ & \text{for some }j_0\{1,\mathrm{},p\}\text{ with }\stackrel{~}{\lambda }_{j_0}\{E_m\}_{m=0}^{2p+1}.\hfill \end{array}`$ (1.17)
In this case $`2M_0+2`$ analytic arcs are converging toward $`\stackrel{~}{\lambda }_{j_0}`$, where $`M_0\{1,\mathrm{},p\}`$ denotes the number of $`\stackrel{~}{\lambda }_j`$ that coincide with $`\stackrel{~}{\lambda }_{j_0}`$. Adjacent arcs meet at an angle $`\pi /(M_0+1)`$ at $`\stackrel{~}{\lambda }_{j_0}`$. $`(`$Thus, if crossings occur, generically, $`M_0=1`$ and two arcs cross at a right angle.$`)`$
$`(vii)`$ The resolvent set $`\backslash \sigma (H)`$ of $`H`$ is path-connected.
Naturally, Theorem 1.1 applies to the special case where $`a`$ and $`b`$ are periodic complex-valued solutions of the $`p`$th stationary Toda equation associated with a nonsingular hyperelliptic curve. Even in this special case, Theorem 1.1 yields new facts which go beyond the previous results by Naĭman , .
For analogous results in the context of one-dimensional Schrödinger operators with quasi-periodic algebro-geometric KdV potentials we refer to .
Theorem 1.1 focuses on stationary quasi-periodic solutions of the Toda hierarchy for the following reasons. First of all, the class of algebro-geometric solutions of the (time-dependent) Toda hierarchy is defined as the class of all solutions of some (and hence infinitely many) equations of the stationary Toda hierarchy. Secondly, time-dependent algebro-geometric solutions of a particular equation of the (time-dependent) Toda hierarchy just represent isospectral deformations (the deformation parameter being the time variable) of fixed stationary algebro-geometric Toda solutions (the latter can be viewed as the initial condition at a fixed time $`t_0`$). In the present case of quasi-periodic algebro-geometric solutions of the $`p`$th Toda equation, the isospectral manifold of such given solutions is a complex $`p`$-dimensional torus, and time-dependent solutions trace out a path in that isospectral torus (cf. the discussions in and ).
Finally, we give a brief discussion of the contents of each section. In Section 2 we provide the necessary background material including a quick construction of the Toda hierarchy of nonlinear evolution equations and its Lax pairs using a polynomial recursion formalism. We also discuss the hyperelliptic Riemann surface underlying the stationary Toda hierarchy, the corresponding Baker–Akhiezer function, and the necessary ingredients to describe the analog of the Its–Matveev formula for stationary Toda solutions. Section 3 focuses on the Green’s function of the Jacobi operator $`H`$, a key ingredient in our characterization of the spectrum $`\sigma (H)`$ of $`H`$ in Section 4 (cf. (1.13)). Our principal Section 4 is then devoted to a proof of Theorem 1.1. Appendix A provides the necessary summary of tools needed from elementary algebraic geometry (most notably the theory of compact (hyperelliptic) Riemann surfaces) and sets the stage for some of the notation used in Sections 24. Appendix B provides additional insight into one ingredient of the theta function representation of the coefficients $`a`$ and $`b`$.
## 2. The Toda Hierarchy, Hyperelliptic Curves, and Theta Function Representations of the Coefficients $`a`$ and $`b`$
In this section we briefly review the recursive construction of the Toda hierarchy and associated Lax pairs following and . Moreover, we discuss the class of algebro-geometric solutions of the Toda hierarchy corresponding to the underlying hyperelliptic curve and recall the analog of the Its–Matveev formula for such solutions. The material in this preparatory section is known and detailed accounts with proofs can be found, for instance, in . For the notation employed in connection with elementary concepts in algebraic geometry (more precisely, the theory of compact Riemann surfaces), we refer to Appendix A.
Throughout this section we assume that
$$a,b\mathrm{}^{\mathrm{}}(),a(n)0\text{ for all }n,$$
(2.1)
and consider the second-order Jacobi difference expression
$$L=aS^++a^{}S^{}+b,$$
(2.2)
where $`S^\pm `$ denote the shift operators
$$(S^\pm f)(n)=f^\pm (n)=f(n\pm 1),n,f\mathrm{}^{\mathrm{}}().$$
(2.3)
To construct the stationary Toda hierarchy we need a second difference expression of order $`2p+2,p_0,`$ defined recursively in the following. We take the quickest route to the construction of $`P_{2p+2}`$, and hence to the Toda hierarchy, by starting from the recursion relations (2.4)–(2.6) below.
Define $`\{f_j\}_{j_0}`$ and $`\{g_j\}_{j_0}`$ recursively by
$`f_0=1,g_0=c_1,`$ (2.4)
$`2f_{j+1}+g_j+g_j^{}2bf_j=0,j_0,`$ (2.5)
$`g_{j+1}g_{j+1}^{}+2\left(a^2f_j^+(a^{})^2f_j^{}\right)b(g_jg_j^{})=0,j_0.`$ (2.6)
Explicitly, one finds
$`f_0=1,`$
$`f_1=b+c_1,`$
$`f_2=a^2+(a^{})^2+b^2+c_1b+c_2,\text{ etc.,}`$ (2.7)
$`g_0=c_1,`$
$`g_1=2a^2c_2,`$
$`g_2=2a^2(b^++b)+c_1(2a^2)c_3,\text{ etc.}`$
Here $`\{c_j\}_j`$ denote undetermined summation constants which naturally arise when solving (2.4)–(2.6).
Subsequently, it will be convenient to also introduce the corresponding homogeneous coefficients $`\widehat{f}_j`$ and $`\widehat{g}_j`$, defined by vanishing of the constants $`c_k,k`$,
$`\widehat{f}_0=1,\widehat{f}_j=f_j|_{c_k=0,k=1,\mathrm{},j},j,`$
$`\widehat{g}_j=g_j|_{c_k=0,k=1,\mathrm{},j+1},j_0.`$ (2.8)
Hence,
$$f_j=\underset{k=0}{\overset{j}{}}c_{jk}\widehat{f}_k,g_j=\underset{k=1}{\overset{j}{}}c_{jk}\widehat{g}_kc_{j+1},j_0,$$
(2.9)
introducing
$$c_0=1.$$
(2.10)
Next we define difference expressions $`P_{2p+2}`$ of order $`2p+2`$ by
$$P_{2p+2}=L^{p+1}+\underset{j=0}{\overset{p}{}}\left(g_j+2af_jS^+\right)L^{pj}+f_{p+1},p_0.$$
(2.11)
Using the recursion relations (2.4)–(2.6), the commutator of $`P_{2p+2}`$ and $`L`$ can be explicitly computed and one obtains
$`[P_{2p+2},L]=`$ $`a\left(g_p^++g_p+f_{p+1}^++f_{p+1}2b^+f_p^+\right)S^+`$
$`+2\left(b(g_p+f_{p+1})+a^2f_p^+(a^{})^2f_p^{}+b^2f_p\right)`$
$`a^{}\left(g_p+g_p^{}+f_{p+1}+f_{p+1}^{}2bf_p\right)S^{},p_0.`$ (2.12)
In particular, $`(L,P_{2p+2})`$ represents the celebrated Lax pair of the Toda hierarchy. Varying $`p_0`$, the stationary Toda hierarchy is then defined in terms of the vanishing of the commutator of $`P_{2p+2}`$ and $`L`$ in (2.12) by,
$$[P_{2p+2},L]=\mathrm{s}\mathrm{Tl}_p(a,b)=0,p_0.$$
(2.13)
Thus one finds
$`g_p+g_p^{}+f_{p+1}+f_{p+1}^{}2bf_p`$ $`=0,`$ (2.14)
$`b(g_p+f_{p+1})+a^2f_p^+(a^{})^2f_g^{}+b^2f_p`$ $`=0.`$ (2.15)
Using (2.5) with $`j=p`$ one concludes that (2.14) reduces to
$$f_{p+1}=f_{p+1}^{},$$
(2.16)
that is, $`f_{p+1}`$ is a lattice constant. Similarly, one infers by subtracting $`b`$ times (2.14) from twice (2.15) and using (2.6) with $`j=p`$, that $`g_{p+1}`$ is a lattice constant as well, that is,
$$g_{p+1}=g_{p+1}^{}.$$
(2.17)
Equations (2.16) and (2.17) constitute the $`p`$th stationary equation in the Toda hierarchy, which will be denoted by
$$\mathrm{s}\mathrm{Tl}_p(a,b)=\left(\begin{array}{c}f_{p+1}^+f_{p+1}\\ g_{p+1}g_{p+1}^{}\end{array}\right)=0,p_0.$$
(2.18)
Explicitly,
$`\mathrm{s}\mathrm{Tl}_0(a,b)=\left(\begin{array}{c}b^+b\\ 2\left((a^{})^2a^2\right)\end{array}\right)=0,`$
$`\mathrm{s}\mathrm{Tl}_1(a,b)=\left(\begin{array}{c}(a^+)^2(a^{})^2+(b^+)^2b^2\\ 2(a^{})^2(b+b^{})2a^2(b^++b)\end{array}\right)`$ (2.19)
$`+c_1\left(\begin{array}{c}b^+b\\ 2\left((a^{})^2a^2\right)\end{array}\right)=0,\text{etc.,}`$
represent the first few equations of the stationary Toda hierarchy. By definition, the set of solutions of (2.13), with $`p`$ ranging in $`_0`$ and $`c_k`$, $`k`$, represents the class of algebro-geometric Toda solutions.
In the following we will frequently assume that $`a`$ and $`b`$ satisfy the $`p`$th stationary Toda equation. By this we mean they satisfy one of the $`p`$th stationary Toda equations after a particular choice of integration constants $`c_k`$, $`k=1,\mathrm{},p`$, $`p`$, has been made.
Next, we introduce polynomials $`F_p(z,n)`$ and $`G_{p+1}(z,n)`$ of degree $`p`$ and $`p+1`$ with respect to the spectral parameter $`z`$ by
$`F_p(z,n)`$ $`={\displaystyle \underset{j=0}{\overset{p}{}}}z^jf_{pj}(n),`$ (2.20)
$`G_{p+1}(z,n)`$ $`=z^{p+1}+{\displaystyle \underset{j=0}{\overset{p}{}}}z^jg_{pj}(n)+f_{p+1}(n).`$ (2.21)
Explicitly, one obtains
$`F_0`$ $`=1,`$
$`F_1`$ $`=z+b+c_1,`$
$`F_2`$ $`=z^2+bz+a^2+(a^{})^2+b^2+c_1(z+b)+c_2,\text{ etc.,}`$ (2.22)
$`G_1`$ $`=z+b,`$
$`G_2`$ $`=z^2+(a^{})^2a^2+b^2+c_1(z+b),\text{ etc.}`$
Next, we study the restriction of the difference expression $`P_{2p+2}`$ to the two-dimensional kernel (i.e., the formal null space in an algebraic sense as opposed to the functional analytic one) of $`(Lz)`$. More precisely, let
$`\text{ker}(Lz)=\{\psi :\{\mathrm{}\}(Lz)\psi =0\}.`$ (2.23)
Then (2.11) implies
$$P_{2p+2}_{\text{ker}(Lz)}=\left(2aF_p(z)S^++G_{p+1}(z)\right)|_{\text{ker}(Lz)}.$$
(2.24)
Therefore, the Lax relation (2.12) becomes
$`0`$ $`=[P_{2p+2},L]_{\text{ker}(Lz)}=(a(2(zb^+)F_p^+2(zb)F_p+G_{p+1}^{}G_{p+1}^+)S^+`$
$`+(2(a^{})^2F_p^{}2a^2F_p^++(zb)(G_{p+1}^{}G_{p+1}))\left)\right|_{\text{ker}(Lz)},`$ (2.25)
or, equivalently,
$`2(zb^+)F_p^+2(zb)F_p+G_{p+1}^+G_{p+1}^{}`$ $`=0,`$ (2.26)
$`2a^2F_p^+2(a^{})^2F_p^{}+(zb)(G_{p+1}G_{p+1}^{})`$ $`=0.`$ (2.27)
Upon summing (2.26) one infers
$$2(zb^+)F_p^++G_{p+1}^++G_{p+1}=0,p_0,$$
(2.28)
and inserting (2.26) into (2.27) then implies
$$(zb)^2F_p+(zb)G_{p+1}+a^2F_p^+(a^{})^2F_p^{}=0,p_0.$$
(2.29)
Combining equations (2.27) and (2.28) one concludes that the quantity
$$R_{2p+2}(z)=G_{p+1}(z,n)^24a(n)^2F_p(z,n)F_p^+(z,n)$$
(2.30)
is a lattice constant, and hence depends on $`z`$ only. Thus, one can write
$$R_{2p+2}(z)=\underset{m=0}{\overset{2p+1}{}}(zE_m),\{E_m\}_{m=0}^{2p+1}.$$
(2.31)
One can show that equation (2.30) leads to an explicit determination of the integration constants $`c_1,\mathrm{},c_p`$ in
$$\mathrm{s}\mathrm{Tl}_p(a,b)=0$$
(2.32)
in terms of the zeros $`E_0,\mathrm{},E_{2p+1}`$ of the associated polynomial $`R_{2p+2}`$ in (2.31). In fact, one can prove that
$$c_k=c_k(\underset{¯}{E}),k=1,\mathrm{},p,$$
(2.33)
where
$`c_k(\underset{¯}{E})`$ $`={\displaystyle \underset{\begin{array}{c}j_0,\mathrm{},j_{2p+1}=0\\ j_0+\mathrm{}+j_{2p+1}=k\end{array}}{\overset{k}{}}}{\displaystyle \frac{(2j_0)!\mathrm{}(2j_{2p+1})!}{2^{2k}(j_0!)^2\mathrm{}(j_{2p+1}!)^2(2j_01)\mathrm{}(2j_{2p+1}1)}}E_0^{j_0}\mathrm{}E_{2p+1}^{j_{2p+1}},`$
$`k=1,\mathrm{},p`$ (2.34)
are symmetric functions of $`\underset{¯}{E}=(E_0,\mathrm{},E_{2p+1})`$.
###### Remark 2.1.
Since by (2.5), (2.6), (2.20) and (2.21), $`a`$ enters quadratically in $`F_p`$ and $`G_{p+1}`$, the Toda hierarchy (2.18) is invariant under the substitution
$$a(n)a_ϵ(n)=ϵ(n)a(n),ϵ(n)\{+1,1\},n.$$
(2.35)
We emphasize that the result (2.24) is valid independently of whether or not $`P_{2p+2}`$ and $`L`$ commute. However, the fact that the two difference expressions $`P_{2p+2}`$ and $`L`$ commute implies the existence of an algebraic relationship between them. This gives rise to the Burchnall–Chaundy polynomial for the Toda hierarchy first discussed in the discrete context by Naĭman , .
###### Theorem 2.2.
Fix $`p_0`$ and assume that $`P_{2p+2}`$ and $`L`$ commute, $`[P_{2p+2},L]=0`$, or equivalently, suppose $`\mathrm{s}\mathrm{Tl}_p(a,b)=0`$. Then $`L`$ and $`P_{2p+2}`$ satisfy an algebraic relationship of the type $`(`$cf. (2.31)$`)`$
$`\begin{array}{cc}& _p(L,P_{2p+2})=P_{2p+2}^2R_{2p+2}(L)=0,\hfill \\ & R_{2p+2}(z)={\displaystyle \underset{m=0}{\overset{2p+1}{}}}(zE_m),z.\hfill \end{array}`$ (2.36)
The expression $`_p(L,P_{2p+2})`$ is called the Burchnall–Chaundy polynomial of the Lax pair $`(L,P_{2p+2})`$. Equation (2.36) naturally leads to the hyperelliptic curve $`𝒦_p`$ of (arithmetic) genus $`p_0`$, where
$`\begin{array}{cc}& 𝒦_p:_p(z,y)=y^2R_{2p+2}(z)=0,\hfill \\ & R_{2p+2}(z)={\displaystyle \underset{m=0}{\overset{2p+1}{}}}(zE_m),\{E_m\}_{m=0}^{2p+1}.\hfill \end{array}`$ (2.37)
The curve $`𝒦_p`$ is compactified by joining two points $`P_\mathrm{}_\pm `$, $`P_\mathrm{}_+P_{\mathrm{}_{}},`$ at infinity. For notational simplicity, the resulting curve is still denoted by $`𝒦_p`$. Points $`P`$ on $`𝒦_p\backslash P_\mathrm{}_\pm `$ are represented as pairs $`P=(z,y)`$, where $`y()`$ is the meromorphic function on $`𝒦_p`$ satisfying $`_p(z,y)=0`$. The complex structure on $`𝒦_p`$ is then defined in the usual way, see Appendix A for the case of nonsingular curves. Hence, $`𝒦_p`$ becomes a two-sheeted hyperelliptic Riemann surface of (arithmetic) genus $`p_0`$ (possibly with a singular affine part) in a standard manner.
We also emphasize that by fixing the curve $`𝒦_p`$ (i.e., by fixing $`E_0,\mathrm{},E_{2p+1}`$), the integration constants $`c_1,\mathrm{},c_p`$ in the corresponding stationary $`\mathrm{s}\mathrm{Tl}_p`$ equation are uniquely determined as is clear from (2.33) and (2.34), which establish the integration constants $`c_k`$ as symmetric functions of $`E_0,\mathrm{},E_{2p+1}`$.
For notational simplicity we will usually tacitly assume that $`p`$. The trivial case $`p=0`$, which leads to $`a(n)^2=(E_1E_0)^2/16`$, $`b(n)=(E_0+E_1)/2`$, $`n`$, is of no interest to us in this paper.
In the following, the zeros<sup>1</sup><sup>1</sup>1If $`a,b\mathrm{}^{\mathrm{}}()`$, these zeros are the Dirichlet eigenvalues of a bounded operator on $`\mathrm{}^2()`$ associated with the difference expression $`L`$ and a Dirichlet boundary condition at $`n`$. of the polynomial $`F_p(,n)`$ (cf. (2.20)) will play a special role. We denote them by $`\{\mu _j(n)\}_{j=1}^p`$ and write
$$F_p(z,n)=\underset{j=1}{\overset{p}{}}[z\mu _j(n)].$$
(2.38)
The next step is crucial; it permits us to “lift” the zeros $`\mu _j`$ of $`F_p`$ from $``$ to the curve $`𝒦_p`$. From (2.30) and (2.38) one infers
$$R_{2p+2}(z)G_{p+1}(z)^2=0,z\{\mu _j,\mu _k^+\}_{j,k=1,\mathrm{},p}.$$
(2.39)
We now introduce $`\{\widehat{\mu }_j(n)\}_{j=1,\mathrm{},p}𝒦_p`$ by
$`\widehat{\mu }_j(n)=(\mu _j(n),G_{p+1}(\mu _j(n),n)),j=1,\mathrm{},p,n.`$ (2.40)
Next, we recall equation (2.30) and define the fundamental meromorphic function $`\varphi (,n)`$ on $`𝒦_p`$ by
$`\varphi (P,n)`$ $`={\displaystyle \frac{yG_{p+1}(z,n)}{2a(n)F_p(z,n)}}`$ (2.41)
$`={\displaystyle \frac{2a(n)F_p(z,n+1)}{y+G_{p+1}(z,n)}},`$ (2.42)
$`P=(z,y)𝒦_p,n`$
with divisor $`(\varphi (,n))`$ of $`\varphi (,n)`$ given by
$$\left(\varphi (,n)\right)=𝒟_{P_\mathrm{}_+\underset{¯}{\overset{^}{\mu }}(n+1)}𝒟_{P_{\mathrm{}_{}}\underset{¯}{\overset{^}{\mu }}(n)},$$
(2.43)
using (2.38) and (2.40). Here we abbreviated
$$\underset{¯}{\overset{^}{\mu }}=\{\widehat{\mu }_1,\mathrm{},\widehat{\mu }_p\}\mathrm{Sym}^p(𝒦_p)$$
(2.44)
(cf. the notation introduced in Appendix A). The stationary Baker–Akhiezer function $`\psi (,n,n_0)`$ on $`𝒦_p`$ is then defined in terms of $`\varphi (,n)`$ by
$`\psi (P,n,n_0)=\{\begin{array}{cc}_{m=n_0}^{n1}\varphi (P,m)\hfill & \text{for}nn_0+1,\hfill \\ 1\hfill & \text{for}n=n_0,\hfill \\ _{m=n}^{n_01}\varphi (P,m)^1\hfill & \text{for}nn_01\hfill \end{array}`$ (2.45)
with divisor $`\left(\psi (,n,n_0)\right)`$ of $`\psi (P,n,n_0)`$ given by
$$\left(\psi (,n,n_0)\right)=𝒟_{\underset{¯}{\overset{^}{\mu }}(n)}𝒟_{\underset{¯}{\overset{^}{\mu }}(n_0)}+(nn_0)(𝒟_{P_\mathrm{}_+}𝒟_P_{\mathrm{}_{}}).$$
(2.46)
Basic properties of $`\varphi `$ and $`\psi `$ are summarized in the following result. We denote by $`W(f,g)(n)=a(fg^+f^+g)`$ the Wronskian of two complex-valued sequences $`f`$ and $`g`$, and denote $`P^{}=(z,y)`$ for $`P=(z,y)𝒦_p`$.
###### Lemma 2.3.
Assume (2.1) and suppose $`a,b`$ satisfy the $`p`$th stationary Toda equation (2.18). Moreover, let $`P=(z,y)𝒦_p\backslash \{P_\mathrm{}_\pm \}`$ and $`(n,n_0)^2`$. Then $`\varphi `$ satisfies the Riccati-type equation
$`a\varphi (P)+a^{}\varphi ^{}(P)^1=zb,`$ (2.47)
as well as
$`\varphi (P)\varphi (P^{})={\displaystyle \frac{F_p^+(z)}{F_p(z)}},`$ (2.48)
$`\varphi (P)+\varphi (P^{})={\displaystyle \frac{G_{p+1}(z)}{aF_p(z)}},`$ (2.49)
$`\varphi (P)\varphi (P^{})={\displaystyle \frac{y(P)}{aF_p(z)}}.`$ (2.50)
Moreover, $`\psi `$ satisfies
$`\left(Lz(P)\right)\psi (P)=0,\left(P_{2p+2}y(P)\right)\psi (P)=0,`$ (2.51)
$`\psi (P,n,n_0)\psi (P^{},n,n_0)={\displaystyle \frac{F_p(z,n)}{F_p(z,n_0)}},`$ (2.52)
$`a(n)\left[\psi (P,n,n_0)\psi (P^{},n+1,n_0)+\psi (P^{},n,n_0)\psi (P,n+1,n_0)\right]`$
$`={\displaystyle \frac{G_{p+1}(z,n)}{F_p(z,n_0)}},`$ (2.53)
$`W(\psi (P,,n_0),\psi (P^{},,n_0))={\displaystyle \frac{y(P)}{F_p(z,n_0)}}.`$ (2.54)
Combining the polynomial recursion approach with (2.38) readily yields trace formulas for the Toda invariants, which are polynomial expressions of $`a`$ and $`b`$, in terms of the zeros $`\mu _j`$ of $`F_p`$.
###### Lemma 2.4.
Assume (2.1) and suppose $`a,b`$ satisfy the $`p`$th stationary Toda equation (2.18). Then,
$`a(n)^2={\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{p}{}}}R_{2p+2}^{1/2}(\widehat{\mu }_j(n)){\displaystyle \underset{\begin{array}{c}k=1\\ kj\end{array}}{\overset{p}{}}}[\mu _j(n)\mu _k(n)]^1+{\displaystyle \frac{1}{4}}[b^{(2)}(n)b(n)^2],`$ (2.55)
$`b(n)={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=0}{\overset{2p+1}{}}}E_m{\displaystyle \underset{j=1}{\overset{p}{}}}\mu _j(n),`$ (2.56)
$`b^{(k)}(n)={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=0}{\overset{2p+1}{}}}E_m^k{\displaystyle \underset{j=1}{\overset{p}{}}}\mu _j(n)^k,k.`$ (2.57)
Strictly speaking, (2.55) is only valid in the case where for all $`n`$, $`\mu _j(n)\mu _k(n)`$ for $`jk`$. The case where some of the $`\mu _j`$ coincide requires a limiting argument that will be omitted as (2.55) will play no further role in this paper. The details of this limiting procedure can be found in .
From this point on we assume that the affine part of $`𝒦_p`$ is nonsingular, that is,
$$E_mE_m^{}\text{ for }mm^{}\text{}m,m^{}=0,1,\mathrm{},2p+1.$$
(2.58)
Since nonspecial divisors play a fundamental role in this context we also recall the following fact.
###### Lemma 2.5.
Assume (2.1) and suppose $`a,b`$ satisfy the $`p`$th stationary Toda equation (2.18). In addition, assume that the affine part of $`𝒦_p`$ is nonsingular. Let $`𝒟_{\underset{¯}{\overset{^}{\mu }}}`$, $`\underset{¯}{\overset{^}{\mu }}=\{\widehat{\mu }_1,\mathrm{},\widehat{\mu }_p\}\mathrm{Sym}^p(𝒦_p)`$, be the Dirichlet divisor of degree $`p`$ associated with $`a`$, $`b`$ defined according to (2.40), that is,
$$\widehat{\mu }_j(n)=(\mu _j(n),G_{p+1}(\mu _j(n),n)),j=1,\mathrm{},p,n.$$
(2.59)
Then $`𝒟_{\underset{¯}{\overset{^}{\mu }}(n)}`$ is nonspecial for all $`n`$. Moreover, there exists a constant $`C_\mu >0`$ such that
$$|\mu _j(n)|C_\mu ,j=1,\mathrm{},p,n.$$
(2.60)
We remark that if $`a,b\mathrm{}^{\mathrm{}}()`$ satisfy the $`p`$th stationary Toda equation (2.18), then automatically $`a(n)0`$ for all $`n`$ (cf. ).
We continue with the theta function representation for $`\psi `$, $`a`$, and $`b`$. For general background information and the notation employed we refer to Appendix A.
Let $`\theta `$ denote the Riemann theta function associated with $`𝒦_p`$ (whose affine part is assumed to be nonsingular) and a fixed homology basis $`\{a_j,b_j\}_{j=1}^p`$. Next, choosing a base point $`Q_0(𝒦_p)`$ in the set of branch points of $`𝒦_p`$, we recall that the Abel maps $`\underset{¯}{A}_{Q_0}`$ and $`\underset{¯}{\alpha }_{Q_0}`$ are defined by (A.40) and (A.43), and the Riemann vector $`\underset{¯}{\mathrm{\Xi }}_{Q_0}`$ is given by (A.55). Then Abel’s theorem (cf. (A.53)) and (2.46) yields
$`\begin{array}{cc}\hfill \underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(n)})& =\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(n_0)})\underset{¯}{A}_P_{\mathrm{}_{}}(P_\mathrm{}_+)(nn_0)\hfill \\ & =\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(n_0)})2\underset{¯}{A}_{Q_0}(P_\mathrm{}_+)(nn_0).\hfill \end{array}`$ (2.61)
Next, let $`\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}`$ denote the normalized differential of the third kind defined by
$`\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}={\displaystyle \frac{1}{y}}{\displaystyle \underset{j=1}{\overset{p}{}}}(z\lambda _j)dz\underset{\zeta 0}{=}\pm \left(\zeta ^1+O(1)\right)d\zeta \text{ as }PP_\mathrm{}_\pm ,`$ (2.62)
$`\zeta =1/z,`$
where the constants $`\lambda _j`$, $`j=1,\mathrm{},p`$, are determined by employing the normalization
$$_{a_j}\omega _{P_\mathrm{}+,P_{\mathrm{}}}^{(3)}=0,j=1,\mathrm{},p.$$
(2.63)
One then infers
$$_{Q_0}^P\omega _{P_\mathrm{}+,P_{\mathrm{}}}^{(3)}\underset{\zeta 0}{=}\pm \mathrm{ln}(\zeta )+e_0^{(3)}(Q_0)+O(\zeta )\text{ as }PP_{\mathrm{}}$$
(2.64)
for some constant $`e_0^{(3)}(Q_0)`$. The vector of $`b`$-periods of $`\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}/(2\pi i)`$ is denoted by
$$\underset{¯}{U}_0^{(3)}=(U_{0,1}^{(3)},\mathrm{},U_{0,p}^{(3)}),U_{0,j}^{(3)}=\frac{1}{2\pi i}_{b_j}\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)},j=1,\mathrm{},p.$$
(2.65)
Since $`Q_0`$ is a branch point, $`Q_0(𝒦_p)`$, one concludes by (A.42) that
$$U_0^{(3)}=\underset{¯}{A}_P_{\mathrm{}}(P_\mathrm{}+)=2\underset{¯}{A}_{Q_0}(P_\mathrm{}+).$$
(2.66)
In the following it will be convenient to introduce the abbreviation
$`\underset{¯}{z}(P,\underset{¯}{Q})=\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{Q}}),`$ (2.67)
$`P𝒦_p,\underset{¯}{Q}=\{Q_1,\mathrm{},Q_p\}\mathrm{Sym}^p(𝒦_p).`$
We note that $`\underset{¯}{z}(,\underset{¯}{Q})`$ is independent of the choice of base point $`Q_0`$.
The zeros and the poles of $`\psi `$ as recorded in (2.46) suggest consideration of the following expression involving $`\theta `$-functions (cf. (A.29))
$$\frac{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n))\right)}{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n_0))\right)}\text{exp}\left(_{Q_0}^P\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}\right).$$
(2.68)
Here we agree to use the same path of integration from $`Q_0`$ to $`P`$ on $`𝒦_p`$ in the Abel map $`\underset{¯}{\overset{^}{A}}_{Q_0}(P)`$ in $`\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n))`$ and in the integral of $`\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}`$ in the exponent of (2.68). With this convention the expression (2.68) is well-defined on $`𝒦_p`$ (cf. Remark A.4, however) and one concludes
$$\psi (P,n,n_0)=C(n,n_0)\frac{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n))\right)}{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n_0))\right)}\text{exp}\left((nn_0)_{Q_0}^P\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}\right).$$
(2.69)
To determine $`C(n,n_0)`$ one can use (2.52) for $`P=P_\mathrm{}_+`$ and $`P^{}=P_{\mathrm{}_{}}`$. Hence,
$$C(n,n_0)^2=\frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n_0))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n_01))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n1))\right)}.$$
(2.70)
Thus, one obtains the following well-known result.
###### Theorem 2.6.
Assume (2.1) and suppose $`a,b`$ satisfy the $`p`$th stationary Toda equation (2.18). In addition, assume the affine part of $`𝒦_p`$ to be nonsingular and let $`P𝒦_p\backslash \{P_\mathrm{}_\pm \}`$ and $`n,n_0`$. Then $`𝒟_{\underset{¯}{\overset{^}{\mu }}(n)}`$ is nonspecial for $`n`$. Moreover,<sup>2</sup><sup>2</sup>2To avoid multi-valued expressions in formulas such as (2.71), etc., we agree to always choose the same path of integration connecting $`Q_0`$ and $`P`$ and refer to Remark A.4 for additional tacitly assumed conventions.
$$\psi (P,n,n_0)=C(n,n_0)\frac{\theta (\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n)))}{\theta (\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n_0)))}\mathrm{exp}\left((nn_0)_{Q_0}^P\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}\right),$$
(2.71)
where
$$C(n,n_0)=\left[\frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n_0))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n_01))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n1))\right)}\right]^{1/2},$$
(2.72)
with the linearizing property of the Abel map,
$$\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(n)})=\left(\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(n_0)})2\underset{¯}{A}_{Q_0}(P_\mathrm{}_+)(nn_0)\right)(modL_p).$$
(2.73)
The coefficients $`a`$ and $`b`$ are given by
$`a(n)=\stackrel{~}{a}\left[{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n1))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n+1))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n))\right)^2}}\right]^{1/2},n,`$ (2.74)
$`b(n)={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=0}{\overset{2p+1}{}}}E_m{\displaystyle \underset{j=1}{\overset{p}{}}}\lambda _j+{\displaystyle \underset{j=1}{\overset{p}{}}}c_j(p){\displaystyle \frac{}{\omega _j}}\mathrm{ln}\left[{\displaystyle \frac{\theta \left(\underset{¯}{\omega }+\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n))\right)}{\theta \left(\underset{¯}{\omega }+\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(n1))\right)}}\right]|_{\underset{¯}{\omega }=0},`$
$`n,`$ (2.75)
where the constant $`\stackrel{~}{a}`$ depends only on $`𝒦_p`$ and $`c_j(p)`$ is given by (A.23).
Combining (2.73) and (2.75), one observes the remarkable linearity of the theta function with respect to $`n`$ in formulas (2.74), (2.75). In fact, one can rewrite (2.75) as
$$b(n)=\mathrm{\Lambda }_0+\underset{j=1}{\overset{p}{}}c_j(p)\frac{}{\omega _j}\mathrm{ln}\left(\frac{\theta (\underset{¯}{\omega }+\underset{¯}{A}\underset{¯}{B}n)}{\theta (\underset{¯}{\omega }+\underset{¯}{C}\underset{¯}{B}n)}\right)|_{\underset{¯}{\omega }=0},$$
(2.76)
where
$`\underset{¯}{A}`$ $`=\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P_\mathrm{}+)+\underset{¯}{U}_0^{(3)}n_0+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(n_0)}),`$ (2.77)
$`\underset{¯}{B}`$ $`=\underset{¯}{U}_0^{(3)},`$ (2.78)
$`\underset{¯}{C}`$ $`=\underset{¯}{A}+\underset{¯}{B},`$ (2.79)
$`\mathrm{\Lambda }_0`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=0}{\overset{2p+1}{}}}E_m{\displaystyle \underset{j=1}{\overset{p}{}}}\lambda _j.`$ (2.80)
Hence, the constants $`\mathrm{\Lambda }_0`$ and $`\underset{¯}{B}^p`$ are uniquely determined by $`𝒦_p`$ (and its homology basis), and the constant $`\underset{¯}{A}^p`$ is in one-to-one correspondence with the Dirichlet data $`\underset{¯}{\overset{^}{\mu }}(n_0)=(\widehat{\mu }_1(n_0),\mathrm{},\widehat{\mu }_p(n_0))\mathrm{Sym}^p𝒦_p`$ at the point $`n_0`$.
###### Remark 2.7.
If one assumes $`a`$, $`b`$ in (2.74) and (2.75) to be quasi-periodic (cf. (3.6) and (3.7)), then there exists a homology basis $`\{\stackrel{~}{a}_j,\stackrel{~}{b}_j\}_{j=1}^p`$ on $`𝒦_p`$ such that $`\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}`$ satisfies the constraint
$$\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}^p.$$
(2.81)
This is studied in detail in Appendix B.
## 3. The Green’s Function of $`H`$
In this section we focus on the properties of the Green’s function of $`H`$ and derive a variety of results to be used in our principal Section 4.
Introducing
$`\text{G}(P,m,n)={\displaystyle \frac{1}{W(\psi (P^{},,n_0),\psi (P,,n_0))}}\{\begin{array}{cc}\psi (P^{},m,n_0)\psi (P,n,n_0),\hfill & mn,\hfill \\ \psi (P,m,n_0)\psi (P^{},n,n_0),\hfill & mn,\hfill \end{array}`$
$`P𝒦_p\backslash \{P_\mathrm{}_\pm \},n,n_0,`$ (3.1)
and
$$g(P,n)=G(P,n,n)=\frac{\psi (P,n,n_0)\psi (P^{},n,n_0)}{W(\psi (P^{},,n_0),\psi (P,,n_0))},$$
(3.2)
equations (2.52) and (2.54) then imply
$$g(P,n)=\frac{F_p(z,n)}{y(P)},P=(z,y)𝒦_p\backslash \{P_\mathrm{}_\pm \},n.$$
(3.3)
Together with $`g(P,n)`$ we also introduce its two branches $`g_\pm (z,n)`$ defined on the upper and lower sheets $`\mathrm{\Pi }_\pm `$ of $`𝒦_p`$ (cf. (A.3), (A.4), and (A.14))
$$g_\pm (z,n)=\frac{F_p(z,n)}{R_{2p+2}(z)^{1/2}},z\mathrm{\Pi },n$$
(3.4)
with $`\mathrm{\Pi }=\backslash 𝒞`$ the cut plane introduced in (A.4).
For convenience we shall focus on $`g_{}`$ whenever possible and use the simplified notation
$$g(z,n)=g_{}(z,n),z\mathrm{\Pi },n$$
(3.5)
from now on.
Next we briefly review a few properties of quasi-periodic and almost-periodic discrete functions.
We denote by $`\mathrm{QP}()`$ and $`\mathrm{AP}()`$ the sets of quasi-periodic and almost periodic sequences on $``$, respectively.
In particular, a sequence $`f`$ is called quasi-periodic with fundamental periods $`(\mathrm{\Omega }_1,\mathrm{},\mathrm{\Omega }_N)(0,\mathrm{})^N`$ if the frequencies $`2\pi /\mathrm{\Omega }_1,\mathrm{},2\pi /\mathrm{\Omega }_N`$ are linearly independent over $``$ and if there exists a continuous function $`FC(^N)`$, periodic of period $`1`$ in each of its arguments,
$$F(x_1,\mathrm{},x_j+1,\mathrm{},x_N)=F(x_1,\mathrm{},x_N),x_j,j=1,\mathrm{},N,$$
(3.6)
such that
$$f(n)=F(\mathrm{\Omega }_1^1n,\mathrm{},\mathrm{\Omega }_N^1n),n.$$
(3.7)
Any quasi-periodic sequence on $``$ is almost periodic on $``$. Moreover, a sequence $`f=\{f(k)\}_k`$ is almost periodic on $``$ if and only if there exists a Bohr almost periodic function $`g`$ on $``$ such that $`f(k)=g(k)`$ for all $`k`$ (see, e.g., \[11, p. 47\]).
For any almost periodic sequence $`f=\{f(k)\}_k`$, the mean value $`f`$ of $`f`$, defined by
$$f=\underset{N\mathrm{}}{lim}\frac{1}{2N+1}\underset{k=n_0N}{\overset{n_0+N}{}}f(k),$$
(3.8)
exists and is independent of $`n_0`$. Moreover, we recall the following facts for almost periodic sequences that can be deduced from corresponding properties of Bohr almost periodic functions, see, for instance, \[5, Ch. I\], \[6, Sects. 39–92\], \[11, Ch. I\], \[22, Chs. 1,3,6\], , \[38, Chs. 1,2,6\], and .
###### Theorem 3.1.
Assume $`f,g\mathrm{AP}()`$ and $`n_0,n`$. Then the following assertions hold:
$`(i)`$ $`f\mathrm{}^{\mathrm{}}()`$.
$`(ii)`$ $`\overline{f}`$, $`cf`$, $`c`$, $`f(+n)`$, $`f(n)`$, $`n`$, $`|f|^\alpha `$, $`\alpha 0`$ are all in $`\mathrm{AP}()`$.
$`(iii)`$ $`f+g,fg\mathrm{AP}()`$.
$`(iv)`$ $`1/g\mathrm{AP}()`$ if and only if $`1/g\mathrm{}^{\mathrm{}}()`$.
$`(v)`$ Let $`G`$ be uniformly continuous on $``$ and $`f(n)`$ for all $`n`$. Then
$`G(f)\mathrm{AP}()`$.
$`(vi)`$ Let $`f=0`$, then $`_{k=n_0}^nf(k)\underset{|n|\mathrm{}}{=}o(|n|)`$.
$`(vii)`$ Let $`F(n)=_{k=n_0}^nf(k)`$. Then $`F\mathrm{AP}()`$ if and only if $`F\mathrm{}^{\mathrm{}}()`$.
$`(viii)`$ If $`0f\mathrm{AP}()`$, $`f0`$, then $`f>0`$.
$`(ix)`$ If $`1/f\mathrm{}^{\mathrm{}}()`$ and $`f=|f|\mathrm{exp}(i\phi )`$, then $`|f|\mathrm{AP}()`$ and $`\phi `$ is of the type
$`\phi (n)=cn+\psi (n)`$, where $`c`$ and $`\psi \mathrm{AP}()`$ $`(`$and real-valued $`)`$.
$`(x)`$ If $`F(n)=\mathrm{exp}\left(_{k=n_0}^nf(k)\right)`$, then $`F\mathrm{AP}()`$ if and only if $`f(n)=i\beta +\psi (n)`$,
where $`\beta `$, $`\psi \mathrm{AP}()`$, and $`\mathrm{\Psi }\mathrm{}^{\mathrm{}}()`$, where $`\mathrm{\Psi }(n)=_{k=n_0}^n\psi (k)`$.
For the rest of this paper it will be convenient to introduce the following hypothesis:
###### Hypothesis 3.2.
Assume the affine part of $`𝒦_p`$ to be nonsingular. Moreover, suppose that $`a,b\mathrm{QP}()`$ satisfy the $`p`$th stationary Toda equation (2.18) on $``$.
Next, we note the following result.
###### Lemma 3.3.
Assume Hypothesis 3.2. Then all $`z`$-derivatives of $`F_p(z,)`$ and $`G_{p+1}(z,)`$, $`z`$, and $`g(z,)`$, $`z\mathrm{\Pi }`$, are quasi-periodic. Moreover, taking limits to points on $`𝒞`$, the last result extends to either side of cuts in the set $`𝒞\backslash \{E_m\}_{m=0}^{2p+1}`$ $`(`$cf. (A.3)$`)`$ by continuity with respect to $`z`$.
###### Proof.
Since $`f_{\mathrm{}}`$ and $`g_{\mathrm{}}`$ are polynomials with respect to $`a`$ and $`b`$, $`f_{\mathrm{}}`$ and $`g_{\mathrm{}}`$, $`\mathrm{}`$, are quasi-periodic by Theorem 3.1. The corresponding assertion for $`F_p(z,)`$ then follows from (2.20) and that for $`g(z,)`$ follows from (3.4). ∎
In the following we represent $`G_{p+1}(z,n)+G_{p+1}^+(z,n)`$ as
$$G_{p+1}(z,n)+G_{p+1}^+(z,n)=2\underset{k=1}{\overset{p+1}{}}[z\nu _k(n)],z,n,$$
(3.9)
and note that the roots $`\nu _k`$ are bounded,
$$\nu _k_{\mathrm{}}\stackrel{~}{C},k=1,\mathrm{},p+1$$
(3.10)
for some constant $`\stackrel{~}{C}>0`$, since the coefficients of $`G_{p+1}(z,n)`$ are defined in terms of bounded coefficients $`a`$ and $`b`$ by(2.6). For future purposes we introduce the set
$`\mathrm{\Pi }_C`$ $`=\mathrm{\Pi }\backslash (\{z||z|C+1\}`$
$`\{z|\underset{m=0,\mathrm{},2p+1}{\mathrm{min}}[\text{Re}(E_m)]1\text{Re}(z)\underset{m=0,\mathrm{},2p+1}{\mathrm{max}}[\text{Re}(E_m)]+1,`$
$`\underset{m=0,\mathrm{},2p+1}{\mathrm{min}}[\text{Im}(E_m)]1\text{Im}(z)\underset{m=0,\mathrm{},2p+1}{\mathrm{max}}[\text{Im}(E_m)]+1\}),`$ (3.11)
where $`C=\mathrm{max}\{C_\mu ,b_{\mathrm{}},\stackrel{~}{C}\}`$ and $`C_\mu `$ is the constant in (2.60). Without loss of generality, we may also assume that $`\mathrm{\Pi }_C`$ contains no cuts, that is,
$$\mathrm{\Pi }_C𝒞=\mathrm{}.$$
(3.12)
Next, we derive a fundamental equation for the mean value of the diagonal Green’s function $`g(z,)`$ that will allow us to analyze the spectrum of the Jacobi operator $`H`$. First, we note that by (2.49), (2.50), (2.54), and (3.1) one obtains
$`{\displaystyle \frac{\text{G}(P,n,n+1)}{\text{G}(P^{},n,n+1)}}={\displaystyle \frac{G_{p+1}(z,n)y}{G_{p+1}(z,n)+y}},P=(z,y)𝒦_p,n.`$ (3.13)
Differentiating the logarithm of the expression on the right-hand side of (3.13) with respect to $`z`$ and using (2.30), one infers
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dz}}\text{ln}\left({\displaystyle \frac{G_{p+1}(z,n)y}{G_{p+1}(z,n)+y}}\right)={\displaystyle \frac{\frac{R_{2p+2}^{}(z)}{2y}G_{p+1}(z,n)yG_{p+1}^{}(z,n)}{4a(n)^2F_p(z,n)F_p^+(z,n)}},z\mathrm{\Pi }_C.`$ (3.14)
Here $``$ abbreviates $`d/dz`$. We note that the left-hand side of (3.14) is well-defined since by (2.30), (2.38), and (3.11),
$`[G_{p+1}(z,n)y][G_{p+1}(z,n)+y]=G_{p+1}(z,n)^2R_{2p+2}(z)`$
$`=4a(n)^2F_p(z,n)F_p^+(z,n)`$
$`=4a(n)^2{\displaystyle \underset{j=1}{\overset{p}{}}}[z\mu _j(n)][z\mu _j(n+1)]0,z\mathrm{\Pi }_C.`$ (3.15)
Adding and subtracting $`g(z,n)`$ on the right-hand side of (3.14) yields
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dz}}\text{ln}\left({\displaystyle \frac{G_{p+1}(z,n)y}{G_{p+1}(z,n)+y}}\right)=g(z,n)+{\displaystyle \frac{K(z,n)}{y}},z\mathrm{\Pi }_C,`$ (3.16)
where
$$K(z,n)=\frac{1}{2}G_{p+1}(z,n)\left(\frac{F_p^{}(z,n)}{F_p(z,n)}+\frac{(F_p^+)^{}(z,n)}{F_p^+(z,n)}\right)G_{p+1}^{}(z,n)F_p(z,n).$$
(3.17)
Next we prove that the mean value of $`K(z,)`$ equals zero.
###### Lemma 3.4.
Assume Hypothesis 3.2. Then
$$K(z,)=0,z\mathrm{\Pi }_C.$$
(3.18)
###### Proof.
Let $`z\mathrm{\Pi }_C`$. Using (2.28) we rewrite (3.17) as
$`K(z,n)=`$ $`{\displaystyle \frac{1}{2}}G_{p+1}(z,n)[{\displaystyle \frac{d}{dz}}\mathrm{ln}(G_{p+1}(z,n)+G_{p+1}^{}(z,n))`$
$`+{\displaystyle \frac{d}{dz}}\mathrm{ln}(G_{p+1}^+(z,n)+G_{p+1}(z,n))]`$
$`{\displaystyle \frac{d}{dz}}G_{p+1}(z,n)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{G_{p+1}^{}(z,n)}{zb(n)}}{\displaystyle \frac{G_{p+1}(z,n)}{zb^+(n)}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{2}}G_{p+1}(z,n)[{\displaystyle \frac{G_{p+1}^{}(z,n)+(G_{p+1}^{})^{}(z,n)}{G_{p+1}(z,n)+G_{p+1}^{}(z,n)}}`$
$`+{\displaystyle \frac{(G_{p+1}^+)^{}(z,n)+G_{p+1}^{}(z,n)}{G_{p+1}^+(z,n)+G_{p+1}(z,n)}}]`$
$`G_{p+1}^{}(z,n)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{G_{p+1}^{}(z,n)}{zb(n)}}{\displaystyle \frac{G_{p+1}(z,n)}{zb^+(n)}}\right)`$
$`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{(G_{p+1}^+)^{}(z,n)G_{p+1}(z,n)G_{p+1}^{}(z,n)G_{p+1}^+(z,n)}{G_{p+1}^+(z,n)+G_{p+1}(z,n)}}`$
$`{\displaystyle \frac{G_{p+1}^{}(z,n)G_{p+1}^{}(z,n)(G_{p+1}^{})^{}(z,n)G_{p+1}(z,n)}{G_{p+1}(z,n)+G_{p+1}^{}(z,n)}}]`$
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{G_{p+1}^{}(z,n)}{zb(n)}}{\displaystyle \frac{G_{p+1}(z,n)}{zb^+(n)}}\right),z\mathrm{\Pi }_C.`$ (3.19)
Since $`K(z,)`$ is a sum of two difference expressions and $`G_{p+1}(z,)`$ and $`G_{p+1}^{}(z,)`$ are bounded for fixed $`z\mathrm{\Pi }_C`$, one obtains
$$K(z,)=0,z\mathrm{\Pi }_C.$$
(3.20)
Using (3.16) and Lemma 3.4, one derives the following result that will subsequently play a crucial role in this paper.
###### Lemma 3.5.
Assume Hypothesis 3.2 and let $`z,z_0\mathrm{\Pi }`$. Then
$$<\mathrm{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>=2_{z_0}^z𝑑z^{}g(z^{},)+<\mathrm{ln}\left(\frac{G_{p+1}(z_0,)y}{G_{p+1}(z_0,)+y}\right)>,$$
(3.21)
where the path connecting $`z_0`$ and $`z`$ is assumed to lie in the cut plane $`\mathrm{\Pi }`$. Moreover, by taking limits to points on $`𝒞`$ in (3.21), the result (3.21) extends to either side of the cuts in the set $`𝒞`$ by continuity with respect to $`z`$.
###### Proof.
Let $`z,z_0\mathrm{\Pi }_C`$. Integrating equation (3.16) from $`z_0`$ to $`z`$ along a smooth path in $`\mathrm{\Pi }_C`$ yields
$`\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(z,n)y}{G_{p+1}(z,n)+y}}\right)\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(z_0,)y}{G_{p+1}(z_0,)+y}}\right)`$ $`=2{\displaystyle _{z_0}^z}𝑑z^{}g(z^{},n)+`$
$`+2{\displaystyle _{z_0}^z}𝑑z^{}{\displaystyle \frac{K(z^{},n)}{y}}.`$ (3.22)
By Lemma 3.3, $`K(z,)`$ is quasi-periodic. Consequently, also
$$_{z_0}^z𝑑z^{}\frac{K(z^{},)}{y},z\mathrm{\Pi }_C,$$
(3.23)
is a family of uniformly almost periodic functions for $`z`$ varying in compact subsets of $`\mathrm{\Pi }_C`$ as discussed in \[22, Sect. 2.7\]. By Lemma 3.4 one thus obtains
$$\left[_{z_0}^z𝑑z^{}\frac{K(z^{},)}{y}\right]=0.$$
(3.24)
Hence, taking mean values in (3.22) (taking into account (3.24)), proves (3.21) for $`z\mathrm{\Pi }_C`$. Since $`f_{\mathrm{}}`$, $`\mathrm{}_0`$, are quasi-periodic by Lemma 3.3 (we recall that $`f_0=1`$), (2.20) and (3.4) yield
$$_{z_0}^z𝑑z^{}g(z^{},)=\underset{\mathrm{}=0}{\overset{p}{}}f_p\mathrm{}_{z_0}^z𝑑z^{}\frac{(z^{})^{\mathrm{}}}{R_{2p+2}(z^{})^{1/2}}.$$
(3.25)
Thus, $`_{z_0}^z𝑑z^{}g(z^{},)`$ has an analytic continuation with respect to $`z`$ to all of $`\mathrm{\Pi }`$ and consequently, (3.21) for $`z\mathrm{\Pi }_C`$ extends by analytic continuation to $`z\mathrm{\Pi }`$. By continuity this extends to either side of the cuts in $`𝒞`$. Interchanging the role of $`z`$ and $`z_0`$, analytic continuation with respect to $`z_0`$ then yields (3.21) for $`z,z_0\mathrm{\Pi }`$. ∎
###### Remark 3.6.
For $`z\mathrm{\Pi }_C`$, the sequence $`\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)`$ is quasi-periodic and hence $`<\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>`$ is well-defined. But if one analytically continues $`\text{ln}\left(\frac{G_{p+1}(z,n)y}{G_{p+1}(z,n)+y}\right)`$ with respect to $`z`$, then $`(G_{p+1}(z,n)y)`$ and $`(G_{p+1}(z,n)+y)`$ may acquire zeros for some $`n`$ and hence $`\text{ln}\left(\frac{G_{p+1}(z,n)y}{G_{p+1}(z,n)+y}\right)\mathrm{QP}()`$. Nevertheless, as shown by the right-hand side of (3.21), $`<\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>`$ admits an analytic continuation in $`z`$ from $`\mathrm{\Pi }_C`$ to all of $`\mathrm{\Pi }`$, and from now on, $`\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)`$, $`z\mathrm{\Pi }`$, always denotes that analytic continuation (cf. also (3.27)).
Next, we will invoke the Baker-Akhiezer function $`\psi (P,n,n_0)`$ and analyze the expression $`<\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>`$ in more detail.
###### Theorem 3.7.
Assume Hypothesis 3.2, let $`P=(z,y)\mathrm{\Pi }_\pm `$, and $`n,n_0`$. Moreover, select a homology basis $`\{\stackrel{~}{a}_j,\stackrel{~}{b}_j\}_{j=1}^p`$ on $`𝒦_p`$ such that $`\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}`$, with $`\underset{¯}{\overset{~}{U}}_0^{(3)}`$ the vector of $`\stackrel{~}{b}`$-periods of the normalized differential of the third kind, $`\stackrel{~}{\omega }_{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}`$, satisfies the constraint
$$\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}^p$$
(3.26)
$`(`$cf. Appendix B$`)`$. Then,
$$\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)\right)=2\text{Re}\left(_{Q_0}^P\stackrel{~}{\omega }_{P_\mathrm{}_+,\mathrm{}_{}}^{(3)}\right).$$
(3.27)
###### Proof.
Using (2.41), (2.45) and (2.46) one obtains the following representation of the Baker-Akhiezer function $`\psi (P,n,n_0)`$ for $`n>n_0`$, $`n,n_0`$, $`P𝒦_p`$,
$`\psi (P,n,n_0)`$ $`={\displaystyle \underset{m=n_0}{\overset{n1}{}}}\varphi (P,m)=\left[{\displaystyle \underset{m=n_0}{\overset{n1}{}}}{\displaystyle \frac{yG_{p+1}(z,m)}{yG_{p+1}(z,m)}}{\displaystyle \frac{F_p(z,m+1)}{F_p(z,m)}}\right]^{1/2}`$
$`=\left({\displaystyle \frac{F_p(z,n)}{F_p(z,n_0)}}\right)^{1/2}\left[{\displaystyle \underset{m=n_0}{\overset{n1}{}}}{\displaystyle \frac{G_{p+1}(z,m)y}{G_{p+1}(z,m)+y}}\right]^{1/2}`$
$`=\left({\displaystyle \frac{F_p(z,n)}{F_p(z,n_0)}}\right)^{1/2}`$
$`\times \text{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{m=n_0}{\overset{n1}{}}}\left[\text{ln}\left({\displaystyle \frac{G_{p+1}(z,m)y}{G_{p+1}(z,m)+y}}\right)<\text{ln}\left({\displaystyle \frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}}\right)>\right]\right)`$
$`\times \text{exp}\left({\displaystyle \frac{1}{2}}(nn_0)<\text{ln}\left({\displaystyle \frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}}\right)>\right),`$ (3.28)
$`P=(z,y)\mathrm{\Pi }_\pm ,z\mathrm{\Pi }_C,n,n_0.`$
A similar representation can be written for $`\psi (P,n,n_0)`$ if $`n<n_0`$, $`n,n_0`$, $`P𝒦_p`$.
Since $`\left[\text{ln}\left(\frac{G_{p+1}(z,m)y}{G_{p+1}(z,m)+y}\right)<\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>\right]`$ has mean zero,
$`\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{m=n_0}{\overset{n1}{}}}\left[\text{ln}\left({\displaystyle \frac{G_{p+1}(z,m)y}{G_{p+1}(z,m)+y}}\right)<\text{ln}\left({\displaystyle \frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}}\right)>\right]\right)\underset{|n|\mathrm{}}{=}o(|n|),`$
$`z\mathrm{\Pi }_C,`$ (3.29)
by Theorem 3.1$`(vi)`$. In addition, the factor $`F_p(z,n)/F_p(z,n_0)`$ in (3.28) is quasi-periodic and hence bounded on $``$.
On the other hand, (2.71) yields
$`\psi (P,n,n_0)`$ $`=C(n,n_0){\displaystyle \frac{\theta (\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n)))}{\theta (\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n_0)))}}\mathrm{exp}\left((nn_0){\displaystyle _{Q_0}^P}\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}\right)`$
$`=\mathrm{\Theta }(P,n,n_0)\mathrm{exp}\left((nn_0){\displaystyle _{Q_0}^P}\stackrel{~}{\omega }_{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}\right),`$ (3.30)
$`P𝒦_p\backslash \left(\{P_\mathrm{}_\pm \}\{\widehat{\mu }_j(n_0)\}_{j=1}^p\right).`$
Taking into account (2.67), (2.73), (2.81), (A.29), and the fact that by (2.60) no $`\widehat{\mu }_j(n)`$ can reach $`P_\mathrm{}_\pm `$ as $`n`$ varies in $``$, one concludes that
$$\mathrm{\Theta }(P,,n_0)\mathrm{}^{\mathrm{}}(),P𝒦_p\backslash \{\widehat{\mu }_j(n_0)\}_{j=1}^p.$$
(3.31)
A comparison of (3.28) and (3.30) then shows that the $`o(|n|)`$-term in (3.29) must actually be bounded on $``$ and hence the left-hand side of (3.29) is almost periodic (in fact, quasi-periodic). In addition, the term
$$\text{exp}\left(\frac{1}{2}\underset{m=n_0}{\overset{n1}{}}\left[\text{ln}\left(\frac{G_{p+1}(z,m)y}{G_{p+1}(z,m)+y}\right)<\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>\right]\right),z\mathrm{\Pi }_C,$$
(3.32)
is then almost periodic (in fact, quasi-periodic) by Theorem 3.1$`(x)`$. A further comparison of (3.28) and (3.30) then yields (3.27) for $`z\mathrm{\Pi }_C`$. Analytic continuation with respect to $`z`$ then implies (3.27) for $`z\mathrm{\Pi }`$. By continuity with respect to $`z`$, taking boundary values to either side of the cuts in the set $`𝒞`$, this then extends to $`z𝒞`$ (cf. (A.3), (A.4)) and hence proves (3.27) for $`P=(z,y)𝒦_p\backslash \{P_\mathrm{}_\pm \}`$. ∎
## 4. Spectra of Jacobi Operators with Quasi-Periodic Algebro-Geometric Coefficients
In this section we establish the connection between the algebro-geometric formalism of Section 2 and the spectral theoretic description of Jacobi operators $`H`$ in $`\mathrm{}^2()`$ with quasi-periodic algebro-geometric coefficients. In particular, we introduce the conditional stability set of $`H`$ and prove our principal result, the characterization of the spectrum of $`H`$. Finally, we provide a qualitative description of the spectrum of $`H`$ in terms of analytic spectral arcs.
Suppose that $`a`$, $`b\mathrm{}^{\mathrm{}}()\mathrm{QP}()`$ satisfy the $`p`$th stationary Toda equation (2.18) on $``$. The corresponding Jacobi operator $`H`$ in $`\mathrm{}^2()`$ is then defined by
$$H=aS^++a^{}S^{}+b,\text{dom}(H)=\mathrm{}^2().$$
(4.1)
Thus, $`H`$ is a bounded operator on $`\mathrm{}^2()`$ (it is self-adjoint if and only if $`a`$ and $`b`$ are real-valued).
Before we turn to the spectrum of $`H`$ in the general non-self-adjoint case, we briefly mention the following result on the spectrum of $`H`$ in the self-adjoint case with quasi-periodic (or almost periodic) real-valued coefficients $`a`$ and $`b`$. We denote by $`\sigma (A)`$, $`\sigma _\mathrm{e}(A)`$, and $`\sigma _\mathrm{d}(A)`$ the spectrum, essential spectrum, and discrete spectrum of a self-adjoint operator $`A`$ in a complex Hilbert space, respectively.
###### Theorem 4.1 (See, e.g., in the continuous context).
Let $`a,b\mathrm{QP}()`$ be real-valued. Define the self-adjoint Jacobi operator $`H`$ in $`\mathrm{}^2()`$ as in (4.1). Then,
$`\sigma (H)=\sigma _\mathrm{e}(H)`$ (4.2)
$`[2\underset{n}{sup}\left(|\text{Re}(a(n))|\right)+\underset{n}{inf}\left(\text{Re}(b(n))\right),2\underset{n}{sup}\left(|\text{Re}(a(n))|+\underset{n}{sup}\text{Re}(b(n))\right)],`$
$`\sigma _\mathrm{d}(H)=\mathrm{}.`$ (4.3)
Moreover, $`\sigma (H)`$ contains no isolated points, that is, $`\sigma (H)`$ is a perfect set.
In the special periodic case where $`a,b`$ are real-valued, the spectrum of $`H`$ is purely absolutely continuous and a finite union of some compact intervals (see, e.g., , , , ,).
Next, we turn to the analysis of the generally non-self-adjoint operator $`H`$ in (4.1). Assuming Hypothesis 3.2 we introduce the set $`\mathrm{\Sigma }`$ by
$$\mathrm{\Sigma }=\left\{\lambda \right|\text{Re}\left(\text{ln}\left(\frac{G_{p+1}(\lambda ,)y}{G_{p+1}(\lambda ,)+y}\right)\right)=0\}.$$
(4.4)
Below we will show that $`\mathrm{\Sigma }`$ plays the role of the conditional stability set of $`H`$, familiar from the spectral theory of one-dimensional periodic differential and difference operators.
###### Lemma 4.2.
Assume Hypothesis 3.2. Then $`\mathrm{\Sigma }`$ coincides with the conditional stability set of $`H`$, that is,
$`\mathrm{\Sigma }`$ $`=\{\lambda |\text{there exists at least one bounded solution}`$
$`0\psi \mathrm{}^{\mathrm{}}()\text{ of }H\psi =\lambda \psi \}.`$ (4.5)
###### Proof.
By (2.71) and (2.72),
$`\psi (P,n,n_0)`$ $`=C(n,n_0){\displaystyle \frac{\theta (\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n)))}{\theta (\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(n_0)))}}\mathrm{exp}\left((nn_0){\displaystyle _{E_0}^z}\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}\right),`$ (4.6)
$`P=(z,y)\mathrm{\Pi }_\pm ,`$
is a solution of $`H\psi =z\psi `$ which is bounded on $``$ if and only if the exponential function in (4.6) is bounded on $``$. By (3.27), the latter holds if and only if
$$\text{Re}\left(\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)\right)=0.$$
(4.7)
###### Remark 4.3.
At first sight our a priori choice of cuts $`𝒞`$ for $`R_{2p+2}()^{1/2}`$, as described in Appendix A, might seem unnatural as they completely ignore the actual spectrum of $`H`$. However, the spectrum of $`H`$ is not known from the outset, and in the case of complex-valued potentials, spectral arcs of $`H`$ may actually cross each other (cf. Theorem 4.7$`(iv)`$) which renders them unsuitable for cuts of $`R_{2p+2}()^{1/2}`$.
Before we state our first principal result on the spectrum of $`H`$, we find it convenient to recall a number of basic definitions and well-known facts in connection with the spectral theory of non-self-adjoint operators (we refer to \[18, Chs. I, III, IX\], \[27, Sects. 1, 21–23\], \[31, Sects. IV.5.6, V.3.2\], and \[49, p. 178–179\] for more details). Let $`S`$ be a densely defined closed operator in complex separable Hilbert space $``$. Denote by $`()`$ the Banach space of all bounded linear operators on $``$ and by $`\mathrm{ker}(T)`$ and $`\text{ran}(T)`$ the kernel (null space) and range of a linear operator $`T`$ in $``$. The resolvent set, $`\rho (S)`$, spectrum, $`\sigma (S)`$, point spectrum (the set of eigenvalues), $`\sigma _\mathrm{p}(S)`$, continuous spectrum, $`\sigma _\mathrm{c}(S)`$, residual spectrum, $`\sigma _\mathrm{r}(S)`$, field of regularity, $`\pi (S)`$, approximate point spectrum, $`\sigma _{\mathrm{ap}}(S)`$, two kinds of essential spectra, $`\sigma _\mathrm{e}(S)`$, and $`\stackrel{~}{\sigma }_\mathrm{e}(S)`$, the numerical range of $`S`$, $`\mathrm{\Theta }(S)`$, and the sets $`\mathrm{\Delta }(S)`$ and $`\stackrel{~}{\mathrm{\Delta }}(S)`$ are defined as follows:
$`\rho (S)`$ $`=\{z|(SzI)^1()\},`$ (4.8)
$`\sigma (S)`$ $`=\backslash \rho (S),`$ (4.9)
$`\sigma _\mathrm{p}(S)`$ $`=\{\lambda |\mathrm{ker}(S\lambda I)\{0\}\},`$ (4.10)
$`\sigma _\mathrm{c}(S)`$ $`=\{\lambda |\mathrm{ker}(S\lambda I)=\{0\}\text{ and }\text{ran}(S\lambda I)\text{ is dense in }`$
$`\text{but not equal to }\},`$ (4.11)
$`\sigma _\mathrm{r}(S)`$ $`=\{\lambda |\mathrm{ker}(S\lambda I)=\{0\}\text{ and }\text{ran}(S\lambda I)\text{ is not dense in }\},`$ (4.12)
$`\pi (S)`$ $`=\{z|\text{there exists }k_z>0\text{ s.t. }(SzI)u_{}k_zu_{}`$
$`\text{for all }u\text{dom}(S)\},`$ (4.13)
$`\sigma _{\mathrm{ap}}(S)`$ $`=\backslash \pi (S),`$ (4.14)
$`\mathrm{\Delta }(S)`$ $`=\{z|dim(\mathrm{ker}(SzI))<\mathrm{}\text{ and }\text{ran}(SzI)\text{ is closed}\},`$ (4.15)
$`\sigma _\mathrm{e}(S)`$ $`=\backslash \mathrm{\Delta }(S),`$ (4.16)
$`\stackrel{~}{\mathrm{\Delta }}(S)`$ $`=\{z|dim(\mathrm{ker}(SzI))<\mathrm{}\text{ or }dim(\mathrm{ker}(S^{}\overline{z}I))<\mathrm{}\},`$ (4.17)
$`\stackrel{~}{\sigma }_\mathrm{e}(S)`$ $`=\backslash \stackrel{~}{\mathrm{\Delta }}(S),`$ (4.18)
$`\mathrm{\Theta }(S)`$ $`=\{(f,Sf)|f\text{dom}(S),f_{}=1\},`$ (4.19)
respectively. One then has
$`\sigma (S)`$ $`=\sigma _\mathrm{p}(S)\sigma _\mathrm{c}(S)\sigma _\mathrm{r}(S)\text{(disjoint union)}`$ (4.20)
$`=\sigma _\mathrm{p}(S)\sigma _\mathrm{e}(S)\sigma _\mathrm{r}(S),`$ (4.21)
$`\sigma _\mathrm{c}(S)`$ $`\sigma _\mathrm{e}(S)\backslash (\sigma _\mathrm{p}(S)\sigma _\mathrm{r}(S)),`$ (4.22)
$`\sigma _\mathrm{r}(S)`$ $`=\sigma _\mathrm{p}(S^{})^{}\backslash \sigma _\mathrm{p}(S),`$ (4.23)
$`\sigma _{\mathrm{ap}}(S)`$ $`=\{\lambda |\text{there exists a sequence }\{f_n\}_n\text{dom}(S)`$
$`\text{with }f_n_{}=1\text{}n\text{, and }lim_n\mathrm{}(S\lambda I)f_n_{}=0\},`$ (4.24)
$`\stackrel{~}{\sigma }_\mathrm{e}(S)`$ $`\sigma _\mathrm{e}(S)\sigma _{\mathrm{ap}}(S)\sigma (S)\text{ (all four sets are closed)},`$ (4.25)
$`\rho (S)`$ $`\pi (S)\mathrm{\Delta }(S)\stackrel{~}{\mathrm{\Delta }}(S)\text{ (all four sets are open),}`$ (4.26)
$`\stackrel{~}{\sigma }_\mathrm{e}(S)`$ $`\overline{\mathrm{\Theta }(S)},\mathrm{\Theta }(S)\text{ is convex,}`$ (4.27)
$`\stackrel{~}{\sigma }_\mathrm{e}(S)`$ $`=\sigma _\mathrm{e}(S)\text{ if }S=S^{}\text{.}`$ (4.28)
Here $`\sigma ^{}`$ in the context of (4.23) denotes the complex conjugate of the set $`\sigma `$, that is,
$$\sigma ^{}=\{\overline{\lambda }|\lambda \sigma \}.$$
(4.29)
We note that there are several other versions of the concept of the essential spectrum in the non-self-adjoint context (cf. \[18, Ch. IX\]) but we will only use the two in (4.16) and in (4.18) in this paper.
We start with the following elementary result.
###### Lemma 4.4.
Let $`H`$ be defined as in (4.1). Then,
$$\sigma _\mathrm{e}(H)=\stackrel{~}{\sigma }_\mathrm{e}(H)\overline{\mathrm{\Theta }(H)}.$$
(4.30)
###### Proof.
Since $`H`$ and $`H^{}`$ are second-order difference operators on $``$,
$$dim(\mathrm{ker}(HzI))2,dim(\mathrm{ker}(H^{}\overline{z}I))2.$$
(4.31)
Moreover, we note that $`S`$ closed and densely defined and $`dim(\mathrm{ker}(S^{}\overline{z}I))<\mathrm{}`$ implies that $`\text{ran}(SzI)`$ is closed (cf. \[18, Theorem I.3.2\]). Equations (4.15)–(4.18) and (4.27) then prove (4.30). ∎
###### Theorem 4.5.
Assume Hypothesis 3.2. Then the point spectrum and residual spectrum of $`H`$ are empty and hence the spectrum of $`H`$ is purely continuous,
$`\sigma _\mathrm{p}(H)=\sigma _\mathrm{r}(H)=\mathrm{},`$ (4.32)
$`\sigma (H)=\sigma _\mathrm{c}(H)=\sigma _\mathrm{e}(H)=\sigma _{\mathrm{ap}}(H).`$ (4.33)
###### Proof.
First we prove the absence of the point spectrum of $`H`$. Suppose $`z\mathrm{\Pi }\backslash (\mathrm{\Sigma }\{\mu _j(n_0)\}_{j=1}^p)`$. Then $`\psi (P,,n_0)`$ and $`\psi (P^{},,n_0)`$ are linearly independent solutions of $`H\psi =z\psi `$ which are unbounded at $`+\mathrm{}`$ or $`\mathrm{}`$. This argument extends to all $`z\mathrm{\Pi }\backslash \mathrm{\Sigma }`$ by multiplying $`\psi (P,,n_0)`$ and $`\psi (P^{},,n_0)`$ with an appropriate function of $`z`$ and $`n_0`$ (independent of $`n`$). It also extends to either side of the cut $`𝒞\backslash \mathrm{\Sigma }`$ by continuity with respect to $`z`$. On the other hand, any solution $`\psi (z,)\mathrm{}^2()`$ of $`H\psi =z\psi `$, $`z`$, is necessarily bounded (since any sequence in $`\mathrm{}^2()`$ is bounded). Thus,
$$\{\backslash \mathrm{\Sigma }\}\sigma _\mathrm{p}(H)=\mathrm{}.$$
(4.34)
Hence, it remains to rule out eigenvalues located in $`\mathrm{\Sigma }`$. We consider a fixed $`\lambda \mathrm{\Sigma }`$ and note that by (2.52), there exists at least one solution $`\psi _1(\lambda ,)\mathrm{}^{\mathrm{}}()`$ of $`H\psi =\lambda \psi `$. Actually, a comparison of (3.28) and (4.4) shows that one can choose $`\psi _1(\lambda ,)`$ such that $`|\psi _1(\lambda ,)|\mathrm{QP}()`$ and hence $`\psi _1(\lambda ,)\mathrm{}^2()`$.
Next, suppose there exists a second solution $`\psi _2(\lambda ,)\mathrm{}^2()`$ of $`H\psi =\lambda \psi `$ which is linearly independent of $`\psi _1(\lambda ,)`$. Then one concludes that the Wronskian of $`\psi _1(\lambda ,)`$ and $`\psi _2(\lambda ,)`$ lies in $`\mathrm{}^2()`$,
$$W(\psi _1(\lambda ,),\psi _2(\lambda ,))\mathrm{}^2().$$
(4.35)
However, by hypothesis, $`W(\psi _1(\lambda ,),\psi _2(\lambda ,))=c(\lambda )0`$ is a nonzero constant. This contradiction proves that
$$\mathrm{\Sigma }\sigma _\mathrm{p}(H)=\mathrm{}$$
(4.36)
and hence $`\sigma _\mathrm{p}(H)=\mathrm{}`$.
Next, we note that the same argument yields that $`H^{}`$ also has no point spectrum,
$$\sigma _\mathrm{p}(H^{})=\mathrm{}.$$
(4.37)
Indeed, if $`a,b\mathrm{}^{\mathrm{}}()\mathrm{QP}()`$ satisfy the $`p`$th stationary Toda equation (2.18) on $``$, then $`\overline{a},\overline{b}`$ also satisfy one of the $`p`$th stationary Toda equation (2.18) associated with a hyperelliptic curve of genus $`p`$ with $`\{E_m\}_{m=0}^{2p+1}`$ replaced by $`\{\overline{E}_m\}_{m=0}^{2p+1}`$, etc. Since by general principles (cf. (4.29)),
$$\sigma _\mathrm{r}(B)\sigma _\mathrm{p}(B^{})^{}$$
(4.38)
for any densely defined closed linear operator $`B`$ in some complex separable Hilbert space (see, e.g., \[28, p. 71\]), one obtains $`\sigma _\mathrm{r}(H)=\mathrm{}`$ and hence (4.32). This proves that the spectrum of $`H`$ is purely continuous, $`\sigma (H)=\sigma _\mathrm{c}(H)`$. The remaining equalities in (4.33) then follow from (4.22) and (4.25). ∎
The following result is a fundamental one:
###### Theorem 4.6.
Assume Hypothesis 3.2. Then the spectrum of $`H`$ coincides with $`\mathrm{\Sigma }`$ and hence equals the conditional stability set of $`H`$,
$`\sigma (H)`$ $`=\left\{\lambda \right|\text{Re}\left(\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(\lambda ,)y}{G_{p+1}(\lambda ,)+y}}\right)\right)=0\}`$ (4.39)
$`=\{\lambda |\text{there exists at least one bounded solution}`$
$`0\psi \mathrm{}^{\mathrm{}}()\text{ of }H\psi =\lambda \psi \}.`$ (4.40)
In particular,
$$\{E_m\}_{m=0}^{2p+1}\sigma (H),$$
(4.41)
and $`\sigma (H)`$ contains no isolated points.
###### Proof.
First we will prove that
$$\sigma (H)\mathrm{\Sigma }$$
(4.42)
by adapting a method due to Chisholm and Everitt (in the context of differential operators). For this purpose we temporarily choose $`z\mathrm{\Pi }\backslash (\mathrm{\Sigma }\{\mu _j(n_0)\}_{j=1}^p)`$ and construct the resolvent of $`H`$ as follows. Introducing the two branches $`\psi _\pm (P,n,n_0)`$ of the Baker–Akhiezer function $`\psi (P,n,n_0)`$ by
$$\psi _\pm (P,n,n_0)=\psi (P,n,n_0),P=(z,y)\mathrm{\Pi }_\pm ,n,n_0,$$
(4.43)
we define
$`\widehat{\psi }_+(z,n,n_0)`$ $`=\{\begin{array}{cc}\psi _+(z,n,n_0)\hfill & \text{if }\psi _+(z,,n_0)\mathrm{}^2(n_0,\mathrm{})\text{,}\hfill \\ \psi _{}(z,n,n_0)\hfill & \text{if }\psi _{}(z,,n_0)\mathrm{}^2(n_0,\mathrm{})\text{,}\hfill \end{array}`$ (4.44)
$`\widehat{\psi }_{}(z,n,n_0)`$ $`=\{\begin{array}{cc}\psi _{}(z,n,n_0)\hfill & \text{if }\psi _{}(z,,n_0)\mathrm{}^2(\mathrm{},n_0)\text{,}\hfill \\ \psi _+(z,n,n_0)\hfill & \text{if }\psi _+(z,,n_0)\mathrm{}^2(\mathrm{},n_0)\text{,}\hfill \end{array}`$ (4.45)
$`z\mathrm{\Pi }\backslash \mathrm{\Sigma },n,n_0,`$
and
$`G(z,n,n^{})`$ $`={\displaystyle \frac{1}{W(\widehat{\psi }_{}(z,n,n_0),\widehat{\psi }_+(z,n,n_0))}}\{\begin{array}{cc}\widehat{\psi }_{}(z,n^{},n_0)\widehat{\psi }_+(z,n,n_0),\hfill & nn^{},\hfill \\ \widehat{\psi }_{}(z,n,n_0)\widehat{\psi }_+(z,n^{},n_0),\hfill & nn^{},\hfill \end{array}`$
$`z\mathrm{\Pi }\backslash \mathrm{\Sigma },n,n_0.`$ (4.46)
Due to the homogeneous nature of $`G`$, (4.46) extends to all $`z\mathrm{\Pi }`$. Moreover, we extend (4.44)–(4.46) to either side of the cut $`𝒞`$ except at possible points in $`\mathrm{\Sigma }`$ (i.e., to $`𝒞\backslash \mathrm{\Sigma }`$) by continuity with respect to $`z`$, taking limits to $`𝒞\backslash \mathrm{\Sigma }`$. Next, we introduce the operator $`R(z)`$ in $`\mathrm{}^2()`$ defined by
$`(R(z)f)(n)={\displaystyle \underset{n^{}}{}}G(z,n,n^{})f(n^{}),f\mathrm{}_0^{\mathrm{}}(),z\mathrm{\Pi },`$ (4.47)
where $`\mathrm{}_0^{\mathrm{}}()`$ denotes the linear space of compactly supported (i.e., finite) complex-valued sequences, and extend it to $`z𝒞\backslash \mathrm{\Sigma }`$, as discussed in connection with $`G(,n,n^{})`$. The explicit form of $`\widehat{\psi }_\pm (z,n,n_0)`$, inferred from (3.30) by restricting $`P`$ to $`\mathrm{\Pi }_\pm `$, then yields the estimates
$$|\widehat{\psi }_\pm (z,n,n_0)|C_\pm (z,n_0)e^{\kappa (z)n},z\mathrm{\Pi }\backslash \mathrm{\Sigma },n$$
(4.48)
for some constants $`C_\pm (z,n_0)>0`$, $`\kappa (z)>0`$, $`z\mathrm{\Pi }\backslash \mathrm{\Sigma }`$. One can follow the second part of the proof of Theorem 5.3.2 in line by line and prove that $`R(z)`$, $`z\backslash \mathrm{\Sigma }`$, extends from $`\mathrm{}_0^{\mathrm{}}()`$ to a bounded linear operator defined on all of $`\mathrm{}^2()`$. A straightforward computation then proves
$$(HzI)R(z)f=f,f\mathrm{}^2(),z\backslash \mathrm{\Sigma }$$
(4.49)
and hence also
$$R(z)(HzI)g=g,g\mathrm{}^2(),z\backslash \mathrm{\Sigma }.$$
(4.50)
Thus, $`R(z)=(HzI)^1`$, $`z\backslash \mathrm{\Sigma }`$, and hence (4.42) holds.
Next we will prove that
$$\sigma (H)\mathrm{\Sigma }.$$
(4.51)
We will adapt a strategy of proof applied by Eastham in the continuous case of (real-valued) periodic potentials (reproduced in the proof of Theorem 5.3.2 of ) to the (complex-valued) quasi-periodic discrete case at hand. Suppose $`\lambda \mathrm{\Sigma }`$. By the characterization (4.5) of $`\mathrm{\Sigma }`$, there exists a bounded solution $`\psi (\lambda ,)`$ of $`H\psi =\lambda \psi `$. A comparison with the Baker-Akhiezer function (3.30) then shows that one can assume, without loss of generality, that
$$|\psi (\lambda ,)|\mathrm{QP}().$$
(4.52)
By Theorem 3.1$`(i)`$, one obtains
$$\psi (\lambda ,)\mathrm{}^{\mathrm{}}().$$
(4.53)
Next, we pick $`\mathrm{\Omega }`$ and consider $`g(n),n=0,1,\mathrm{},\mathrm{\Omega }`$, satisfying
$`\begin{array}{cc}& g(0)=0,g(\mathrm{\Omega })=1,\hfill \\ & 0g(n)1,n=1,\mathrm{},\mathrm{\Omega }1.\hfill \end{array}`$ (4.54)
Moreover, we introduce the sequence $`\{h_k\}_k\mathrm{}^2()`$ by
$$h_k(n)=\{\begin{array}{cc}1,\hfill & |n|(k1)\mathrm{\Omega },\hfill \\ g(k\mathrm{\Omega }|n|),\hfill & (k1)\mathrm{\Omega }|n|k\mathrm{\Omega },\hfill \\ 0,\hfill & |n|k\mathrm{\Omega }\hfill \end{array}$$
(4.55)
and the sequence $`\{f_k(\lambda )\}_k\mathrm{}^2()`$ by
$$f_k(\lambda ,n)=d_k(\lambda )\psi (\lambda ,n)h_k(n),n,d_k(\lambda )>0,k.$$
(4.56)
Here $`d_k(\lambda )`$ is determined by the normalization requirement
$$f_k(\lambda )_2=1,k.$$
(4.57)
Of course,
$$f_k(\lambda ,)\mathrm{}^2(),k,$$
(4.58)
since $`f_k(\lambda ,)`$ is finitely supported. Next, we note that as a consequence of Theorem 3.1$`(viii)`$,
$$\underset{N}{\overset{N}{}}|\psi (\lambda ,n)|^2\underset{N\mathrm{}}{=}(2N+1)|\psi (\lambda ,)|^2+o(N)$$
(4.59)
with
$$|\psi (\lambda ,)|^2>0.$$
(4.60)
Thus, one computes
$`1`$ $`=f_k(\lambda )_2^2=d_k(\lambda )^2{\displaystyle \underset{n}{}}|\psi (\lambda ,n)|^2h_k(n)^2`$
$`=d_k(\lambda )^2{\displaystyle \underset{|n|k\mathrm{\Omega }}{}}|\psi (\lambda ,n)|^2h_k(n)^2d_k(\lambda )^2{\displaystyle \underset{|n|(k1)\mathrm{\Omega }}{}}|\psi (\lambda ,n)|^2`$
$`d_k(\lambda )^2\left[|\psi (\lambda ,)|^2(k1)\mathrm{\Omega }+o(k)\right].`$ (4.61)
Consequently,
$$d_k(\lambda )\underset{k\mathrm{}}{=}O\left(k^{1/2}\right).$$
(4.62)
Next, one computes
$`(H\lambda I)f_k(\lambda ,n)=`$ $`d_k(\lambda )[a(n)\psi (\lambda ,n)[h_k(n+1)h_k(n)]`$
$`+a(n1)\psi (\lambda ,n1)[h_k(n1)h_k(n)]]`$ (4.63)
and hence
$`(H\lambda I)f_k_2`$ $`2d_k(\lambda )a_{\mathrm{}}\psi (\lambda )(h_k^+h_k)_2,k.`$ (4.64)
Using (4.53) and (4.55) one estimates
$`\psi (\lambda )\left[h_k^+h_k\right]_2^2`$ $`={\displaystyle \underset{(k1)\mathrm{\Omega }|n|k\mathrm{\Omega }}{}}|\psi (\lambda ,n)|^2|h_k(n+1)h_k(n)|^2`$
$`2\psi (\lambda )_{\mathrm{}}^2\left(\mathrm{\Omega }+1\right).`$ (4.65)
Thus, combining (4.62) and (4.64)–(4.65) one infers
$$\underset{n\mathrm{}}{lim}(H\lambda I)f_k_2=0$$
(4.66)
and hence $`\lambda \sigma _{\mathrm{ap}}(H)=\sigma (H)`$ by (4.24) and (4.33).
Relation (4.41) follows from (4.5) and the fact that by (2.52) there exists a solution $`\psi ((E_m,0),,n_0)\mathrm{}^{\mathrm{}}()`$ of $`H\psi =E_m\psi `$ for all $`m=0,\mathrm{},2p+1`$.
Finally, $`\sigma (H)`$ contains no isolated points since those would necessarily be essential singularities of the resolvent of $`H`$, as $`H`$ has no eigenvalues by (4.32) (cf. \[31, Sect. III.6.5\]). An analysis of the Green’s function of $`H`$ reveals at most an algebraic singularity at the points $`\{E_m\}_{m=0}^{2p+1}`$ and hence excludes the possibility of an essential singularity of $`(HzI)^1`$. ∎
In the special self-adjoint case where $`a,b`$ are real-valued, the result (4.39) is equivalent to the vanishing of the Lyapunov exponent of $`H`$ which characterizes the (purely absolutely continous) spectrum of $`H`$ as discussed by Carmona and Lacroix \[9, Chs. IV, VII\] (cf. also , , , ).
The explicit formula for $`\mathrm{\Sigma }`$ in (4.4) permits a qualitative description of the spectrum of $`H`$ as follows. We recall (3.16) and (3.25) and write
$$\frac{1}{2}\frac{d}{dz}\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)=g(z,)=\frac{_{j=1}^p(z\stackrel{~}{\lambda }_j)}{\left(_{j=0}^{2p+1}(zE_m)\right)^{1/2}},z\mathrm{\Pi },$$
(4.67)
for some constants
$$\{\stackrel{~}{\lambda }_j\}_{j=1}^p.$$
(4.68)
As in similar situations before, (4.67) extends to either side of the cuts in $`𝒞`$ by continuity with respect to $`z`$.
###### Theorem 4.7.
Assume Hypothesis 3.2. Then the spectrum $`\sigma (H)`$ of $`H`$ has the following properties:
$`(i)`$ $`\sigma (H)`$ is bounded,
$$\sigma (H)\{z|\text{Re}(z)[M_1,M_2],\text{Im}(z)[M_3,M_4]\},$$
(4.69)
where
$`\begin{array}{cc}& M_1=2\underset{n}{sup}[|\text{Re}(a(n))|]+\underset{n}{inf}[\text{Re}(b(n))],\hfill \\ & M_2=2\underset{n}{sup}[|\text{Re}(a(n))|]+\underset{n}{sup}[\text{Re}(b(n))],\hfill \\ & M_3=2\underset{n}{sup}[|\text{Im}(a(n))|]+\underset{n}{inf}[\text{Im}(b(n))],\hfill \\ & M_4=2\underset{n}{sup}[|\text{Im}(a(n))|]+\underset{n}{sup}[\text{Im}(b(n))].\hfill \end{array}`$ (4.70)
$`(ii)`$ $`\sigma (H)`$ consists of finitely many simple analytic arcs $`(`$cf. Remark 4.8$`)`$. These analytic arcs may only end at the points $`\stackrel{~}{\lambda }_1,\mathrm{},\stackrel{~}{\lambda }_p`$, $`E_0,\mathrm{},E_{2p+1}`$.
$`(iii)`$ Each $`E_m`$, $`m=0,\mathrm{},2p+1`$, is met by at least one of these arcs. More precisely, a particular $`E_{m_0}`$ is hit by precisely $`2N_0+1`$ analytic arcs, where $`N_0\{0,\mathrm{},p\}`$ denotes the number of $`\stackrel{~}{\lambda }_j`$ that coincide with $`E_{m_0}`$. Adjacent arcs meet at an angle $`2\pi /(2N_0+1)`$ at $`E_{m_0}`$. $`(`$Thus, generically, $`N_0=0`$ and precisely one arc hits $`E_{m_0}`$.$`)`$
$`(iv)`$ Crossings of spectral arcs are permitted. This phenomenon takes place precisely when for a particular $`j_0\{1,\mathrm{},p\}`$, $`\stackrel{~}{\lambda }_{j_0}\sigma (H)`$ such that
$`\begin{array}{cc}& \text{Re}\left(\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(\stackrel{~}{\lambda }_{j_0},)y}{G_{p+1}(\stackrel{~}{\lambda }_{j_0},)+y}}\right)\right)=0\hfill \\ & \text{ for some }j_0\{1,\mathrm{},p\}\text{ with }\stackrel{~}{\lambda }_{j_0}\{E_m\}_{m=0}^{2p+1}.\hfill \end{array}`$ (4.71)
In this case $`2M_0+2`$ analytic arcs are converging toward $`\stackrel{~}{\lambda }_{j_0}`$, where $`M_0\{1,\mathrm{},p\}`$ denotes the number of $`\stackrel{~}{\lambda }_j`$ that coincide with $`\stackrel{~}{\lambda }_{j_0}`$. Adjacent arcs meet at an angle $`\pi /(M_0+1)`$ at $`\stackrel{~}{\lambda }_{j_0}`$. $`(`$Thus, if crossings occur, generically, $`M_0=1`$ and two arcs cross at a right angle.$`)`$
$`(v)`$ The resolvent set $`\backslash \sigma (H)`$ of $`H`$ is path-connected.
###### Proof.
Item $`(i)`$ follows from (4.30) and (4.33) upon noticing that
$$(f,Hf)=2\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}a(k)\text{Re}[f(k+1)\overline{f(k)}]+(f,\text{Re}(b)f)+i(f,\text{Im}(b)f),f\mathrm{}^2().$$
(4.72)
To prove $`(ii)`$ we first introduce the meromorphic differential of the third kind
$`\mathrm{\Omega }^{(3)}=g(P,)dz={\displaystyle \frac{F_p(z,)dz}{y}}={\displaystyle \frac{_{j=1}^p\left(z\stackrel{~}{\lambda }_j\right)dz}{R_{2p+2}(z)^{1/2}}},`$
$`P=(z,y)𝒦_p\backslash \{P_\mathrm{}_\pm \}`$ (4.73)
(cf. (4.68)). Then, by Lemma 3.5,
$$<\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)>=2_{Q_0}^P\mathrm{\Omega }^{(3)}+<\text{ln}\left(\frac{G_{p+1}(z_0,)y}{G_{p+1}(z_0,)+y}\right)>,$$
(4.74)
for some fixed $`Q_0=(z_0,y_0)𝒦_p\backslash \{P_\mathrm{}_\pm \}`$, is holomorphic on $`𝒦_p\backslash \{P_\mathrm{}_\pm \}`$. By (4.67), (4.68), the characterization (4.39) of the spectrum,
$$\sigma (H)=\left\{\lambda \right|\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(\lambda ,)y}{G_{p+1}(\lambda ,)+y}\right)\right)=0\},$$
(4.75)
and the fact that $`\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)\right)`$ is a harmonic function on the cut plane $`\mathrm{\Pi }`$, the spectrum $`\sigma (H)`$ of $`H`$ consists of analytic arcs which may only end at the points $`\stackrel{~}{\lambda }_1,\mathrm{},\stackrel{~}{\lambda }_p`$, $`E_0,\mathrm{},E_{2p+1}`$. (Since $`\sigma (H)`$ is independent of the chosen set of cuts, if a spectral arc crosses or runs along a part of one of the cuts in $`𝒞`$, one can slightly deform the original set of cuts to extend an analytic arc along or across such an original cut.)
To prove $`(iii)`$ one first recalls that by Theorem 4.6 the spectrum of $`H`$ contains no isolated points. On the other hand, since $`\{E_m\}_{m=0}^{2p+1}\sigma (H)`$ by (4.41), one concludes that at least one spectral arc meets each $`E_m`$, $`m=0,\mathrm{},2p+1`$. Choosing $`Q_0=(E_{m_0},0)`$ in (4.74) one obtains
$`\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}}\right)=2{\displaystyle _{E_{m_0}}^z}𝑑z^{}g(z^{},)+\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(E_{m_0},)y}{G_{p+1}(E_{m_0},)+y}}\right)`$
$`=2{\displaystyle _{E_{m_0}}^z}𝑑z^{}{\displaystyle \frac{_{j=1}^p\left(z^{}\stackrel{~}{\lambda }_j\right)}{\left(_{m=0}^{2p+1}(z^{}E_m)\right)^{1/2}}}+\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(E_{m_0},)y}{G_{p+1}(E_{m_0},)+y}}\right)`$
$`\underset{zE_{m_0}}{=}{\displaystyle _{E_{m_0}}^z}𝑑z^{}(z^{}E_{m_0})^{N_0(1/2)}[C+O(z^{}E_{m_0})]`$
$`+\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(E_{m_0},)y}{G_{p+1}(E_{m_0},)+y}}\right)`$
$`\underset{zE_{m_0}}{=}{\displaystyle \frac{(zE_{m_0})^{N_0+(1/2)}}{N_0+(1/2)}}[C+O(zE_{m_0})]+\mathrm{ln}\left({\displaystyle \frac{G_{p+1}(E_{m_0},)y}{G_{p+1}(E_{m_0},)+y}}\right)`$ (4.76)
for some $`C=|C|e^{i\phi _0}\backslash \{0\}`$. Using
$$\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(E_m,)y}{G_{p+1}(E_m,)+y}\right)\right)=0,m=0,\mathrm{},2p+1,$$
(4.77)
as a consequence of (4.41), $`\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)\right)=0`$ and $`z=E_{m_0}+\rho e^{i\phi }`$ imply
$$0\underset{\rho 0}{=}\mathrm{cos}[(N_0+(1/2))\phi +\phi _0]\rho ^{N_0+(1/2)}[|C|+O(\rho )].$$
(4.78)
This proves the assertions made in item $`(iii)`$.
In order to prove $`(iv)`$ it suffices to refer to (4.67) and observe that locally $`\frac{1}{2}\frac{d}{dz}\text{ln}\left(\frac{G_{p+1}(z,)y}{G_{p+1}(z,)+y}\right)`$ behaves like $`C_0(z\stackrel{~}{\lambda }_{j_0})^{M_0}`$ for some $`C_0\backslash \{0\}`$ in a sufficiently small neighborhood of $`\stackrel{~}{\lambda }_{j_0}`$.
Finally we will show that all arcs are simple (i.e., do not self-intersect each other). Assume that the spectrum of $`H`$ contains a simple closed loop $`\gamma `$, $`\gamma \sigma (H)`$. Then
$$\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(z(P),)y(P)}{G_{p+1}(z(P),)+y(P)}\right)\right)=0,P\mathrm{\Gamma },$$
(4.79)
where the closed simple curve $`\mathrm{\Gamma }𝒦_p`$ denotes an appropriate lift of $`\gamma `$ to $`𝒦_p`$, yields the contradiction
$$\text{Re}\left(\mathrm{ln}\left(\frac{G_{p+1}(z(P),)y(P)}{G_{p+1}(z(P),)+y(P)}\right)\right)=0\text{ for all }P\text{ in the interior of }\mathrm{\Gamma }$$
(4.80)
by Corollary 8.2.5 in . Therefore, since there are no closed loops in $`\sigma (H)`$ and no analytic arc tends to infinity, the resolvent set of $`H`$ is connected and hence path-connected, proving $`(v)`$. ∎
###### Remark 4.8.
Here $`\sigma `$ is called an arc if there exists a parameterization $`\gamma C([0,1])`$ such that $`\sigma =\{\gamma (t)|t[0,1]\}`$. The arc $`\sigma `$ is called simple if there exists a parameterization $`\gamma `$ such that $`\gamma :[0,1]`$ is injective.
## Appendix A Hyperelliptic Curves and their Theta Functions
We provide a brief summary of some of the fundamental notations needed from the theory of hyperelliptic Riemann surfaces. More details can be found in some of the standard textbooks and as well as in monographs dedicated to integrable systems such as \[4, Ch. 2\], \[24, App. A, B\].
Fix $`p`$. We intend to describe the hyperelliptic Riemann surface $`𝒦_p`$ of genus $`p`$ of the Toda-type curve (2.36), associated with the polynomial
$`\begin{array}{cc}& _p(z,y)=y^2R_{2p+2}(z)=0,\hfill \\ & R_{2p+2}(z)={\displaystyle \underset{m=0}{\overset{2p+1}{}}}(zE_m),\{E_m\}_{m=0}^{2p+1}.\hfill \end{array}`$ (A.1)
To simplify the discussion we will assume that the affine part of $`𝒦_p`$ is nonsingular, that is, we assume that
$$E_mE_m^{}\text{ for }mm^{},m,m^{}=0,\mathrm{},2p+1$$
(A.2)
throughout this appendix. Next we introduce an appropriate set of (nonintersecting) cuts $`𝒞_j`$ joining $`E_{m(j)}`$ and $`E_{m^{}(j)}`$, $`j=1,\mathrm{},p+1`$, and denote
$$𝒞=\underset{j=1}{\overset{p+1}{}}𝒞_j,𝒞_j𝒞_k=\mathrm{},jk.$$
(A.3)
Define the cut plane
$$\mathrm{\Pi }=\backslash 𝒞,$$
(A.4)
and introduce the holomorphic function
$$R_{2p+2}()^{1/2}:\mathrm{\Pi },z\left(\underset{m=0}{\overset{2p+1}{}}(zE_m)\right)^{1/2}$$
(A.5)
on $`\mathrm{\Pi }`$ with an appropriate choice of the square root branch in (A.5). Next we define the set
$$_p=\{(z,\sigma R_{2p+2}(z)^{1/2})z,\sigma \{1,1\}\}\{P_\mathrm{}_+,P_{\mathrm{}_{}}\}$$
(A.6)
by extending $`R_{2p+2}()^{1/2}`$ to $`𝒞`$. The hyperelliptic curve $`𝒦_p`$ is then the set $`_p`$ with its natural complex structure obtained upon gluing the two sheets of $`_p`$ crosswise along the cuts. Moreover, we introduce the set of branch points
$$(𝒦_p)=\{(E_m,0)\}_{m=0}^{2p+1}.$$
(A.7)
Points $`P𝒦_p\backslash \{P_\mathrm{}_\pm \}`$ are denoted by
$$P=(z,\sigma R_{2p+2}(z)^{1/2})=(z,y),$$
(A.8)
where $`y(P)`$ denotes the meromorphic function on $`𝒦_p`$ satisfying $`_p(z,y)=y^2R_{2p+2}(z)=0`$ and
$$y(P)\underset{\zeta 0}{=}\left(1\frac{1}{2}\left(\underset{m=0}{\overset{2p+1}{}}E_m\right)\zeta +O(\zeta ^2)\right)\zeta ^{p1}\text{ as }PP_\mathrm{}_\pm \text{,}\zeta =1/z.$$
(A.9)
In addition, we introduce the holomorphic sheet exchange map (involution)
$$:𝒦_p𝒦_p,P=(z,y)P^{}=(z,y),P_\mathrm{}_\pm P_\mathrm{}_\pm ^{}=P_{\mathrm{}_{}}$$
(A.10)
and the two meromorphic projection maps
$$\stackrel{~}{\pi }:𝒦_p\{\mathrm{}\},P=(z,y)z,P_\mathrm{}_\pm \mathrm{}$$
(A.11)
and
$$y:𝒦_p\{\mathrm{}\},P=(z,y)y,P_\mathrm{}_\pm \mathrm{}.$$
(A.12)
Thus the map $`\stackrel{~}{\pi }`$ has a pole of order 1 at $`P_\mathrm{}_\pm `$ and $`y`$ has a pole of order $`p+1`$ at $`P_\mathrm{}_\pm `$. Moreover,
$$\stackrel{~}{\pi }(P^{})=\stackrel{~}{\pi }(P),y(P^{})=y(P),P𝒦_p.$$
(A.13)
As a result, $`𝒦_p`$ is a two-sheeted branched covering of the Riemann sphere $`^1`$ ($`\{\mathrm{}\}`$) branched at the $`2p+4`$ points $`\{(E_m,0)\}_{m=0}^{2p+1},P_\mathrm{}_\pm `$. $`𝒦_p`$ is compact since $`\stackrel{~}{\pi }`$ is open and $`^1`$ is compact. Therefore, the compact hyperelliptic Riemann surface resulting in this manner has topological genus $`p`$.
Next we introduce the upper and lower sheets $`\mathrm{\Pi }_\pm `$ by
$$\mathrm{\Pi }_\pm =\{(z,\pm R_{2p+2}(z)^{1/2})_pz\mathrm{\Pi }\}$$
(A.14)
and the associated charts
$$\zeta _\pm :\mathrm{\Pi }_\pm \mathrm{\Pi },Pz.$$
(A.15)
Let $`\{a_j,b_j\}_{j=1}^p`$ be a homology basis for $`𝒦_p`$ with intersection matrix of the cycles satisfying
$$a_jb_k=\delta _{j,k},a_ja_k=0,b_jb_k=0,j,k=1,\mathrm{},p.$$
(A.16)
Associated with the homology basis $`\{a_j,b_j\}_{j=1}^p`$ we also recall the canonical dissection of $`𝒦_p`$ along its cycles yielding the simply connected interior $`\widehat{𝒦}_p`$ of the fundamental polygon $`\widehat{𝒦}_p`$ given by
$$\widehat{𝒦}_p=a_1b_1a_1^1b_1^1a_2b_2a_2^1b_2^1\mathrm{}a_p^1b_p^1.$$
(A.17)
Let $`(𝒦_p)`$ and $`^1(𝒦_p)`$ denote the set of meromorphic functions (0-forms) and meromorphic differentials (1-forms) on $`𝒦_p`$, respectively. The residue of a meromorphic differential $`\nu ^1(𝒦_p)`$ at a point $`Q𝒦_p`$ is defined by
$$\text{res}_Q(\nu )=\frac{1}{2\pi i}_{\gamma _Q}\nu ,$$
(A.18)
where $`\gamma _Q`$ is a counterclockwise oriented smooth simple closed contour encircling $`Q`$ but no other pole of $`\nu `$. Holomorphic differentials are also called Abelian differentials of the first kind. Abelian differentials of the second kind $`\omega ^{(2)}^1(𝒦_p)`$ are characterized by the property that all their residues vanish. Any meromorphic differential $`\omega ^{(3)}`$ on $`𝒦_p`$ not of the first or second kind is said to be of the third kind. A differential of the third kind $`\omega ^{(3)}^1(𝒦_p)`$ is usually normalized by vanishing of its $`a`$-periods, that is,
$$_{a_j}\omega ^{(3)}=0,j=1,\mathrm{},p.$$
(A.19)
A normal differential $`\omega _{P_1,P_2}^{(3)}`$, associated with two distinct points $`P_1,P_2\widehat{𝒦}_p`$, by definition, has simple poles at $`P_1`$ and $`P_2`$ with residues $`+1`$ at $`P_1`$ and $`1`$ at $`P_2`$ and vanishing $`a`$-periods. If $`\omega _{P,Q}^{(3)}`$ is a normal differential of the third kind associated with $`P,Q\widehat{𝒦}_p`$, holomorphic on $`𝒦_p\backslash \{P,Q\}`$, then
$$_{b_j}\omega _{P,Q}^{(3)}=2\pi i_P^Q\omega _j,j=1,\mathrm{},p.$$
(A.20)
We shall always assume (without loss of generality) that all poles of $`\omega ^{(3)}`$ on $`𝒦_p`$ lie on $`\widehat{𝒦}_p`$ (i.e., not on $`\widehat{𝒦}_p`$).
Using local charts one infers that $`dz/y`$ is a holomorphic differential on $`𝒦_p`$ with zeros of order $`p1`$ at $`P_\mathrm{}_\pm `$ and hence
$$\eta _j=\frac{z^{j1}dz}{y},j=1,\mathrm{},p,$$
(A.21)
form a basis for the space of holomorphic differentials on $`𝒦_p`$. Introducing the invertible matrix $`C`$ in $`^p`$
$`C`$ $`=\left(C_{j,k}\right)_{j,k=1,\mathrm{},p},C_{j,k}={\displaystyle _{a_k}}\eta _j,`$ (A.22)
$`\underset{¯}{c}(k)`$ $`=(c_1(k),\mathrm{},c_p(k)),c_j(k)=\left(C^1\right)_{j,k},j,k=1,\mathrm{},p,`$ (A.23)
the normalized differentials $`\omega _j`$ for $`j=1,\mathrm{},p`$,
$$\omega _j=\underset{\mathrm{}=1}{\overset{p}{}}c_j(\mathrm{})\eta _{\mathrm{}},_{a_k}\omega _j=\delta _{j,k},j,k=1,\mathrm{},p,$$
(A.24)
form a canonical basis for the space of holomorphic differentials on $`𝒦_p`$.
In the chart $`(U_{P_\mathrm{}_\pm },\zeta _{P_\mathrm{}_\pm })`$ induced by $`1/\stackrel{~}{\pi }`$ near $`P_\mathrm{}_\pm `$ one infers,
$`\underset{¯}{\omega }`$ $`=(\omega _1,\mathrm{},\omega _p)={\displaystyle \underset{j=1}{\overset{p}{}}}{\displaystyle \frac{\underset{¯}{c}(j)\zeta ^{pj}d\zeta }{\left(_{m=0}^{2p+1}(1\zeta E_m)\right)^{1/2}}}`$ (A.25)
$`=\pm \left(\underset{¯}{c}(p)+\zeta \left[{\displaystyle \frac{1}{2}}\underset{¯}{c}(p){\displaystyle \underset{m=0}{\overset{2p+1}{}}}E_m+\underset{¯}{c}(p1)\right]+O(\zeta ^2)\right)d\zeta \text{ as }PP_\mathrm{}_\pm \text{,}`$
$`\zeta =1/z.`$
The matrix $`\tau =\left(\tau _{j,\mathrm{}}\right)_{j,\mathrm{}=1}^p`$ in $`^{p\times p}`$ of $`b`$-periods defined by
$$\tau _{j,\mathrm{}}=_{b_j}\omega _{\mathrm{}},j,\mathrm{}=1,\mathrm{},p$$
(A.26)
satisfies
$$\text{Im}(\tau )>0\text{ and }\tau _{j,\mathrm{}}=\tau _{\mathrm{},j},j,\mathrm{}=1,\mathrm{},p.$$
(A.27)
Associated with the matrix $`\tau `$ one introduces the period lattice
$$L_p=\{\underset{¯}{z}^p\underset{¯}{z}=\underset{¯}{m}+\underset{¯}{n}\tau ,\underset{¯}{m},\underset{¯}{n}^p\}$$
(A.28)
and the Riemann theta function associated with $`𝒦_p`$ and the given homology basis $`\{a_j,b_j\}_{j=1,\mathrm{},p}`$,
$$\theta (\underset{¯}{z})=\underset{\underset{¯}{n}^p}{}\mathrm{exp}\left(2\pi i(\underset{¯}{n},\underset{¯}{z})+\pi i(\underset{¯}{n},\underset{¯}{n}\tau )\right),\underset{¯}{z}^p,$$
(A.29)
where $`(\underset{¯}{u},\underset{¯}{v})=\overline{\underset{¯}{u}}\underset{¯}{v}^{}=_{j=1}^p\overline{u}_jv_j`$ denotes the scalar product in $`^p`$. It has the following fundamental properties
$`\theta (z_1,\mathrm{},z_{j1},z_j,z_{j+1},\mathrm{},z_n)=\theta (\underset{¯}{z}),`$ (A.30)
$`\theta (\underset{¯}{z}+\underset{¯}{m}+\underset{¯}{n}\tau )=\mathrm{exp}\left(2\pi i(\underset{¯}{n},\underset{¯}{z})\pi i(\underset{¯}{n},\underset{¯}{n}\tau )\right)\theta (\underset{¯}{z}),\underset{¯}{m},\underset{¯}{n}^p.`$ (A.31)
Next we briefly describe some consequences of a change of homology basis. Let
$$\{a_1,\mathrm{},a_p,b_1,\mathrm{},b_p\}$$
(A.32)
be a canonical homology basis on $`𝒦_p`$ with intersection matrix satisfying (A.16) and let
$$\{a_1^{},\mathrm{},a_p^{},b_1^{},\mathrm{},b_p^{}\}$$
(A.33)
be a homology basis on $`𝒦_p`$ related to (A.32) by
$$\left(\begin{array}{c}\underset{¯}{a}_{}^{}{}_{}{}^{}\\ \underset{¯}{b}_{}^{}{}_{}{}^{}\end{array}\right)=X\left(\begin{array}{c}\underset{¯}{a}^{}\\ \underset{¯}{b}^{}\end{array}\right),$$
(A.34)
where
$`\underset{¯}{a}^{}`$ $`=(a_1,\mathrm{},a_p)^{},\underset{¯}{b}^{}=(b_1,\mathrm{},b_p)^{},`$
$`\underset{¯}{a}_{}^{}{}_{}{}^{}`$ $`=(a_1^{},\mathrm{},a_p^{})^{},\underset{¯}{b}_{}^{}{}_{}{}^{}=(b_1^{},\mathrm{},b_p^{})^{},`$ (A.35)
$`X`$ $`=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ (A.36)
with $`A,B,C`$, and $`D`$ being $`p\times p`$ matrices with integer entries. Then (A.33) is also a canonical homology basis on $`𝒦_p`$ with intersection matrix satisfying (A.16) if and only if
$$X\text{Sp}(p,),$$
(A.37)
where
$$\text{Sp}(p,)=\left\{X=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)\right|X\left(\begin{array}{cc}0& I_p\\ I_p& 0\end{array}\right)X^{}=\left(\begin{array}{cc}0& I_p\\ I_p& 0\end{array}\right),det(X)=1\}$$
(A.38)
denotes the symplectic modular group (here $`A,B,C,D`$ in $`X`$ are again $`p\times p`$ matrices with integer entries). If $`\{\omega _j\}_{j=1}^p`$ and $`\{\omega _j^{}\}_{j=1}^p`$ are the normalized bases of holomorphic differentials corresponding to the canonical homology bases (A.32) and (A.33), with $`\tau `$ and $`\tau ^{}`$ the associated $`b`$ and $`b^{}`$-periods of $`\underset{¯}{\omega }=\omega _1,\mathrm{},\omega _p`$ and $`\underset{¯}{\omega }^{}=\omega _1^{},\mathrm{},\omega _p^{}`$, respectively, then one computes
$$\underset{¯}{\omega }^{}=\underset{¯}{\omega }(A+B\tau )^1,\tau ^{}=(C+D\tau )(A+B\tau )^1.$$
(A.39)
Fixing a base point $`Q_0𝒦_p\backslash \{P_\mathrm{}_\pm \}`$, one denotes by $`J(𝒦_p)=^p/L_p`$ the Jacobi variety of $`𝒦_p`$, and defines the Abel map $`\underset{¯}{A}_{Q_0}`$ by
$$\underset{¯}{A}_{Q_0}:𝒦_nJ(𝒦_p),\underset{¯}{A}_{Q_0}(P)=(_{Q_0}^P\omega _1,\mathrm{},_{Q_0}^P\omega _p)(modL_p),P𝒦_p.$$
(A.40)
Next, consider the vector $`\underset{¯}{U}_0^{(3)}`$ of $`b`$-periods of $`\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}/(2\pi i)`$, the normalized differential of the third kind, holomorphic on $`𝒦_p\backslash \{P_\mathrm{}_\pm \}`$,
$$\underset{¯}{U}_0^{(3)}=(U_{0,1}^{(3)},\mathrm{},U_{0,p}^{(3)}),U_{0,j}^{(3)}=\frac{1}{2\pi i}_{b_j}\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)},j=1,\mathrm{},p.$$
(A.41)
Using (A.20) one then computes
$$\underset{¯}{U}_0^{(3)}=\underset{¯}{A}_P_{\mathrm{}}(P_\mathrm{}+)=2\underset{¯}{A}_{Q_0}(P_\mathrm{}+),$$
(A.42)
where $`Q_0`$ is chosen to be a branch point of $`𝒦_p`$, $`Q_0(𝒦_p)`$, in the last part of (A.42).
Similarly, one introduces
$$\underset{¯}{\alpha }_{Q_0}:\mathrm{Div}(𝒦_p)J(𝒦_p),𝒟\underset{¯}{\alpha }_{Q_0}(𝒟)=\underset{P𝒦_p}{}𝒟(P)\underset{¯}{A}_{Q_0}(P),$$
(A.43)
where $`\mathrm{Div}(𝒦_p)`$ denotes the set of divisors on $`𝒦_p`$. Here a map $`𝒟:𝒦_p`$ is called a divisor on $`𝒦_p`$ if $`𝒟(P)0`$ for only finitely many $`P𝒦_p`$. (In the main body of this paper we will choose $`Q_0`$ to be one of the branch points, i.e., $`Q_0(𝒦_p)`$, and we will always choose the same path of integration from $`Q_0`$ to $`P`$ in all Abelian integrals.) For subsequent use in Remark A.4 we also introduce
$`\underset{¯}{\overset{^}{A}}_{Q_0}`$ $`:\widehat{𝒦}_p^p,`$ (A.44)
$`P\underset{¯}{\overset{^}{A}}_{Q_0}(P)=(\widehat{A}_{Q_0,1}(P),\mathrm{},\widehat{A}_{Q_0,p}(P))=({\displaystyle _{Q_0}^P}\omega _1,\mathrm{},{\displaystyle _{Q_0}^P}\omega _p)`$
and
$$\underset{¯}{\overset{^}{\alpha }}_{Q_0}:\mathrm{Div}(\widehat{𝒦}_p)^p,𝒟\underset{¯}{\overset{^}{\alpha }}_{Q_0}(𝒟)=\underset{P\widehat{𝒦}_p}{}𝒟(P)\underset{¯}{\overset{^}{A}}_{Q_0}(P).$$
(A.45)
In connection with divisors on $`𝒦_p`$ we will employ the following (additive) notation,
$`𝒟_{Q_0\underset{¯}{Q}}=𝒟_{Q_0}+𝒟_{\underset{¯}{Q}},𝒟_{\underset{¯}{Q}}=𝒟_{Q_1}+\mathrm{}+𝒟_{Q_m},`$ (A.46)
$`\underset{¯}{Q}=\{Q_1,\mathrm{},Q_m\}\mathrm{Sym}^m𝒦_p,Q_0𝒦_p,m,`$
where for any $`Q𝒦_p`$,
$$𝒟_Q:𝒦_p_0,P𝒟_Q(P)=\{\begin{array}{cc}1\hfill & \text{for }P=Q,\hfill \\ 0\hfill & \text{for }P𝒦_p\backslash \{Q\},\hfill \end{array}$$
(A.47)
and $`\mathrm{Sym}^m𝒦_p`$ denotes the $`m`$th symmetric product of $`𝒦_p`$. In particular, $`\mathrm{Sym}^m𝒦_p`$ can be identified with the set of nonnegative divisors $`0𝒟\mathrm{Div}(𝒦_p)`$ of degree $`m`$.
For $`f(𝒦_p)\backslash \{0\}`$, $`\omega ^1(𝒦_p)\backslash \{0\}`$ the divisors of $`f`$ and $`\omega `$ are denoted by $`(f)`$ and $`(\omega )`$, respectively. Two divisors $`𝒟`$, $`\mathrm{Div}(𝒦_p)`$ are called equivalent, denoted by $`𝒟`$, if and only if $`𝒟=(f)`$ for some $`f(𝒦_p)\backslash \{0\}`$. The divisor class $`[𝒟]`$ of $`𝒟`$ is then given by $`[𝒟]=\{\mathrm{Div}(𝒦_p)𝒟\}`$. We recall that
$$\mathrm{deg}((f))=0,\mathrm{deg}((\omega ))=2(p1),f(𝒦_p)\backslash \{0\},\omega ^1(𝒦_p)\backslash \{0\},$$
(A.48)
where the degree $`\mathrm{deg}(𝒟)`$ of $`𝒟`$ is given by $`\mathrm{deg}(𝒟)=_{P𝒦_p}𝒟(P)`$. It is customary to call $`(f)`$ (respectively, $`(\omega )`$) a principal (respectively, canonical) divisor.
Introducing the complex linear spaces
$`(𝒟)`$ $`=\{f(𝒦_p)f=0\text{ or }(f)𝒟\},r(𝒟)=dim_{}(𝒟),`$ (A.49)
$`^1(𝒟)`$ $`=\{\omega ^1(𝒦_p)\omega =0\text{ or }(\omega )𝒟\},i(𝒟)=dim_{}^1(𝒟)`$ (A.50)
(with $`i(𝒟)`$ the index of specialty of $`𝒟`$), one infers that $`\mathrm{deg}(𝒟)`$, $`r(𝒟)`$, and $`i(𝒟)`$ only depend on the divisor class $`[𝒟]`$ of $`𝒟`$. Moreover, we recall the following fundamental facts.
###### Theorem A.1.
Let $`𝒟\mathrm{Div}(𝒦_p)`$, $`\omega ^1(𝒦_p)\backslash \{0\}`$. Then,
$$i(𝒟)=r(𝒟(\omega )),p_0.$$
(A.51)
The Riemann-Roch theorem reads
$$r(𝒟)=\mathrm{deg}(𝒟)+i(𝒟)p+1,n_0.$$
(A.52)
By Abel’s theorem, $`𝒟\mathrm{Div}(𝒦_p)`$, $`p`$, is principal if and only if
$$\mathrm{deg}(𝒟)=0\text{ and }\underset{¯}{\alpha }_{Q_0}(𝒟)=\underset{¯}{0}.$$
(A.53)
Finally, assume $`p`$. Then $`\underset{¯}{\alpha }_{Q_0}:\mathrm{Div}(𝒦_p)J(𝒦_p)`$ is surjective $`(`$Jacobi’s inversion theorem$`)`$.
###### Theorem A.2.
Let $`𝒟_{\underset{¯}{Q}}\mathrm{Sym}^p𝒦_p`$, $`\underset{¯}{Q}=\{Q_1,\mathrm{},Q_p\}`$. Then,
$$1i(𝒟_{\underset{¯}{Q}})=s$$
(A.54)
if and only if there are $`s`$ pairs of the type $`\{P,P^{}\}\{Q_1,\mathrm{},Q_p\}`$ $`(`$this includes, of course, branch points for which $`P=P^{}`$$`)`$. One has $`sp/2`$.
Next, we denote by $`\underset{¯}{\mathrm{\Xi }}_{Q_0}=(\mathrm{\Xi }_{Q_{0,1}},\mathrm{},\mathrm{\Xi }_{Q_{0,p}})`$ the vector of Riemann constants,
$$\mathrm{\Xi }_{Q_{0,j}}=\frac{1}{2}(1+\tau _{j,j})\underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{p}{}}_a_{\mathrm{}}\omega _{\mathrm{}}(P)_{Q_0}^P\omega _j,j=1,\mathrm{},p.$$
(A.55)
###### Theorem A.3.
Let $`\underset{¯}{Q}=\{Q_1,\mathrm{},Q_p\}\mathrm{Sym}^p𝒦_p`$ and assume $`𝒟_{\underset{¯}{Q}}`$ to be nonspecial, that is, $`i(𝒟_{\underset{¯}{Q}})=0`$. Then,
$$\theta \left(\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\alpha _{Q_0}(𝒟_{\underset{¯}{Q}})\right)=0\text{ if and only if }P\{Q_1,\mathrm{},Q_p\}.$$
(A.56)
###### Remark A.4.
In Section 2 we dealt with theta function expressions of the type
$$\psi (P)=\frac{\theta (\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_1))}{\theta (\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_2))}\mathrm{exp}\left(c_{Q_0}^P\mathrm{\Omega }^{(3)}\right),P𝒦_p,$$
(A.57)
where $`𝒟_j\mathrm{Sym}^p𝒦_p`$, $`j=1,2`$, are nonspecial positive divisors of degree $`p`$, $`c`$ is a constant, and $`\mathrm{\Omega }^{(3)}`$ is a normalized differential of the third kind with a prescribed singularity at $`P_\mathrm{}_\pm `$. Even though we agree to always choose identical paths of integration from $`P_0`$ to $`P`$ in all Abelian integrals (A.57), this is not sufficient to render $`\psi `$ single-valued on $`𝒦_p`$. To achieve single-valuedness one needs to replace $`𝒦_p`$ by its simply connected canonical dissection $`\widehat{𝒦}_p`$ and then replace $`\underset{¯}{A}_{Q_0}`$ and $`\underset{¯}{\alpha }_{Q_0}`$ in (A.57) with $`\underset{¯}{\overset{^}{A}}_{Q_0}`$ and $`\underset{¯}{\overset{^}{\alpha }}_{Q_0}`$ as introduced in (A.44) and (A.45). In particular, one regards $`a_j,b_j`$, $`j=1,\mathrm{},p`$, as curves (being a part of $`\widehat{𝒦}_p`$, cf. (A.17)) and not as homology classes. Similarly, one then replaces $`\underset{¯}{\mathrm{\Xi }}_{Q_0}`$ by $`\underset{¯}{\overset{^}{\mathrm{\Xi }}}_{Q_0}`$ (replacing $`\underset{¯}{A}_{Q_0}`$ by $`\underset{¯}{\overset{^}{A}}_{Q_0}`$ in (A.55), etc.). Moreover, in connection with $`\psi `$, one introduces the vector of $`b`$-periods $`\underset{¯}{U}^{(3)}`$ of $`\mathrm{\Omega }^{(3)}`$ by
$$\underset{¯}{U}^{(3)}=(U_1^{(3)},\mathrm{},U_p^{(3)}),U_j^{(3)}=\frac{1}{2\pi i}_{b_j}\mathrm{\Omega }^{(3)},j=1,\mathrm{},p,$$
(A.58)
and then renders $`\psi `$ single-valued on $`\widehat{𝒦}_p`$ by requiring
$$\underset{¯}{\overset{^}{\alpha }}_{Q_0}(𝒟_1)\underset{¯}{\overset{^}{\alpha }}_{Q_0}(𝒟_2)=c\underset{¯}{U}^{(3)}$$
(A.59)
$`(`$as opposed to merely $`\underset{¯}{\alpha }_{Q_0}(𝒟_1)\underset{¯}{\alpha }_{Q_0}(𝒟_2)=c\underset{¯}{U}^{(3)}(modL_p)`$$`)`$. Actually, by (A.31),
$$\underset{¯}{\overset{^}{\alpha }}_{Q_0}(𝒟_1)\underset{¯}{\overset{^}{\alpha }}_{Q_0}(𝒟_2)c\underset{¯}{U}^{(3)}^p,$$
(A.60)
suffices to guarantee single-valuedness of $`\psi `$ on $`\widehat{𝒦}_p`$. Without the replacement of $`\underset{¯}{A}_{Q_0}`$ and $`\underset{¯}{\alpha }_{Q_0}`$ by $`\underset{¯}{\overset{^}{A}}_{Q_0}`$ and $`\underset{¯}{\overset{^}{\alpha }}_{Q_0}`$ in (A.57) and without the assumption (A.59) $`(`$or (A.60)$`)`$, $`\psi `$ is a multiplicative $`(`$multi-valued$`)`$ function on $`𝒦_p`$, and then most effectively discussed by introducing the notion of characters on $`𝒦_p`$ $`(`$cf. \[21, Sect. III.9\]$`)`$. For simplicity, we decided to avoid the latter possibility and throughout this paper will always tacitly assume (A.59) or (A.60).
## Appendix B Restrictions on $`\underset{¯}{B}=\underset{¯}{U}_0^{(3)}`$
The purpose of this appendix is to prove the result (2.81), $`\underset{¯}{B}=\underset{¯}{U}_0^{(3)}^p`$, for some choice of homology basis $`\{a_j,b_j\}_{j=1}^p`$ on $`𝒦_p`$ as recorded in Remark 2.7.
To this end we first recall a few notions in connection with periodic meromorphic functions of $`p`$ complex variables.
###### Definition B.1.
Let $`p`$ and $`F:^p\{\mathrm{}\}`$ be meromorphic (i.e., a ratio of two entire functions of $`p`$ complex variables). Then,
(i) $`\underset{¯}{\omega }=(\omega _1,\mathrm{},\omega _p)^p\backslash \{0\}`$ is called a period of $`F`$ if
$$F(\underset{¯}{z}+\underset{¯}{\omega })=F(\underset{¯}{z})$$
(B.1)
for all $`\underset{¯}{z}^p`$ for which $`F`$ is analytic. The set of all periods of $`F`$ is denoted by $`𝒫_F`$.
(ii) $`F`$ is called degenerate if it depends on less than $`p`$ complex variables; otherwise, $`F`$ is called nondegenerate.
###### Theorem B.2.
Let $`p`$, $`F:^p\{\mathrm{}\}`$ be meromorphic, and $`𝒫_F`$ be the set of all periods of $`F`$. Then either
$`(i)`$ $`𝒫_F`$ has a finite limit point,
or
$`(ii)`$ $`𝒫_F`$ has no finite limit point.
In case $`(i)`$, $`𝒫_F`$ contains infinitesimal periods $`(`$i.e., sequences of nonzero periods converging to zero$`)`$. In addition, in case $`(i)`$ each period is a limit point of periods and hence $`𝒫_F`$ is a perfect set.
Moreover, $`F`$ is degenerate if and only if $`F`$ admits infinitesimal periods. In particular, for nondegenerate functions $`F`$ only alternative $`(ii)`$ applies.
Next, let $`\underset{¯}{\omega }_q^p\backslash \{0\}`$, $`q=1,\mathrm{},r`$ for some $`r`$. Then $`\underset{¯}{\omega }_1,\mathrm{},\underset{¯}{\omega }_r`$ are called linearly independent over $``$ $`(`$resp. $``$$`)`$ if
$`\nu _1\underset{¯}{\omega }_1+\mathrm{}+\nu _r\underset{¯}{\omega }_r=0,\nu _q\text{ (resp., }\nu _q\text{)},q=1,\mathrm{},r,`$
$`\text{implies }\nu _1=\mathrm{}=\nu _r=0.`$ (B.2)
Clearly, the maximal number of vectors in $`^p`$ linearly independent over $``$ equals $`2p`$.
###### Theorem B.3.
Let $`p`$.
$`(i)`$ If $`F:^p\{\mathrm{}\}`$ is a nondegenerate meromorphic function with periods $`\underset{¯}{\omega }_q^p\backslash \{0\}`$, $`q=1,\mathrm{},r`$, $`r`$, linearly independent over $``$, then $`\underset{¯}{\omega }_1,\mathrm{},\underset{¯}{\omega }_r`$ are also linearly independent over $``$. In particular, $`r2p`$.
$`(ii)`$ A nondegenerate entire function $`F:^p`$ cannot have more than $`p`$ periods linearly independent over $``$ $`(`$or $``$$`)`$.
For $`p=1`$, $`\mathrm{exp}(z)`$, $`\mathrm{sin}(z)`$ are examples of entire functions with precisely one period. Any non-constant doubly periodic meromorphic function of one complex variable is elliptic (and hence indeed has poles).
###### Definition B.4.
Let $`p,r`$. A system of periods $`\underset{¯}{\omega }_q^p\backslash \{0\}`$, $`q=1,\mathrm{},r`$, of a nondegenerate meromorphic function $`F:^p\{\mathrm{}\}`$, linearly independent over $``$, is called fundamental or a basis of periods for $`F`$ if every period $`\underset{¯}{\omega }`$ of $`F`$ is of the form
$$\underset{¯}{\omega }=m_1\underset{¯}{\omega }_1+\mathrm{}+m_r\underset{¯}{\omega }_r\text{ for some }m_q\text{}q=1,\mathrm{},r\text{.}$$
(B.3)
The representation of $`\underset{¯}{\omega }`$ in (B.3) is unique since by hypothesis $`\underset{¯}{\omega }_1,\mathrm{},\underset{¯}{\omega }_r`$ are linearly independent over $``$. In addition, $`𝒫_F`$ is countable in this case. (This rules out case $`(i)`$ in Theorem B.2 since a perfect set is uncountable. Hence, one does not have to assume that $`F`$ is nondegenerate in Definition B.4.)
This material is standard and can be found, for instance, in \[39, Ch. 2\].
Next, returning to the Riemann theta function $`\theta (\underset{¯}{})`$ in (A.29), we introduce the vectors $`\{\underset{¯}{e}_j\}_{j=1}^p,\{\underset{¯}{\tau }_j\}_{j=1}^p^p\backslash \{0\}`$ by
$$\underset{¯}{e}_j=(0,\mathrm{},0,\underset{j}{\underset{}{1}},0,\mathrm{},0),\underset{¯}{\tau }_j=\underset{¯}{e}_j\tau ,j=1,\mathrm{},p.$$
(B.4)
Then
$$\{\underset{¯}{e}_j\}_{j=1}^p$$
(B.5)
is a basis of periods for the entire (nondegenerate) function $`\theta (\underset{¯}{}):^p`$. Moreover, fixing $`k\{1,\mathrm{},p\}`$, then
$$\{\underset{¯}{e}_j,\underset{¯}{\tau }_j\}_{j=1}^p$$
(B.6)
is a basis of periods for the meromorphic function $`_{z_k}\mathrm{ln}\left(\frac{\theta (\underset{¯}{})}{\theta (\underset{¯}{}+\underset{¯}{V})}\right):^p\{\mathrm{}\}`$, $`\underset{¯}{V}^p`$ (cf. (A.31) and \[21, p. 91\]).
Next, let $`\underset{¯}{A}^p`$, $`\underset{¯}{D}=(D_1,\mathrm{},D_p)^p`$, $`D_j\backslash \{0\}`$, $`j=1,\mathrm{},p`$, and consider
$`\begin{array}{cc}\hfill f_k:,f_k(n)& =_{z_k}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{A}\underset{¯}{z})}{\theta (\underset{¯}{C}\underset{¯}{z})}}\right)|_{\underset{¯}{z}=\underset{¯}{D}n}\hfill \\ & =_{z_k}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{A}\underset{¯}{z}\mathrm{diag}(\underset{¯}{D}))}{\theta (\underset{¯}{C}\underset{¯}{z}\mathrm{diag}(\underset{¯}{D}))}}\right)|_{\underset{¯}{z}=(n,\mathrm{},n)}.\hfill \end{array}`$ (B.7)
Here $`\mathrm{diag}(\underset{¯}{D})`$ denotes the diagonal matrix
$$\mathrm{diag}(\underset{¯}{D})=\left(D_j\delta _{j,j^{}}\right)_{j,j^{}=1}^p.$$
(B.8)
Then the quasi-periods $`D_j^1`$, $`j=1,\mathrm{},p`$, of $`f_k`$ are in a one-to-one correspondence with the periods of
$$F_k:^p\{\mathrm{}\},F_k(\underset{¯}{z})=_{z_k}\mathrm{ln}\left(\frac{\theta (\underset{¯}{A}\underset{¯}{z}\mathrm{diag}(\underset{¯}{D}))}{\theta (\underset{¯}{C}\underset{¯}{z}\mathrm{diag}(\underset{¯}{D}))}\right)$$
(B.9)
of the special type
$$\underset{¯}{e}_j\left(\mathrm{diag}(\underset{¯}{D})\right)^1=(0,\mathrm{},0,\underset{j}{\underset{}{D_j^1}},0,\mathrm{},0).$$
(B.10)
Moreover,
$$f_k(n)=F_k(\underset{¯}{z})|_{\underset{¯}{z}=(n,\mathrm{},n)},n.$$
(B.11)
###### Theorem B.5.
Suppose $`a`$ and $`b`$ in (2.75) to be quasi-periodic. Then there exists a homology basis $`\{\stackrel{~}{a}_j,\stackrel{~}{b}_j\}_{j=1}^p`$ on $`𝒦_p`$ such that the vector $`\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}`$ with $`\underset{¯}{\overset{~}{U}}_0^{(3)}`$ the vector of $`\stackrel{~}{b}`$-periods of the corresponding normalized differential of the third kind, $`\stackrel{~}{\omega }_{P_\mathrm{}+,P_{\mathrm{}}}^{(3)}`$, satisfies the constraint
$$\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}^p.$$
(B.12)
###### Proof.
By (A.41), the vector of $`b`$-periods $`\underset{¯}{U}_0^{(3)}`$ associated with a given homology basis $`\{a_j,b_j\}_{j=1}^p`$ on $`𝒦_p`$ and the normalized differential of the third kind, $`\omega _{P_\mathrm{}+,P_{\mathrm{}}}^{(3)}`$, is continuous with respect to $`E_0,\mathrm{},E_{2p+1}`$. Hence, we may assume in the following that
$$B_j0,j=1,\mathrm{},p,\underset{¯}{B}=(B_1,\mathrm{},B_p)$$
(B.13)
by slightly altering $`E_0,\mathrm{},E_{2p+1}`$, if necessary. Using (2.76), we may write
$`\begin{array}{cc}\hfill b(n)& =\mathrm{\Lambda }_0{\displaystyle \underset{j=1}{\overset{p}{}}}c_j(p){\displaystyle \frac{}{\omega _j}}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\omega }+\underset{¯}{A}\underset{¯}{B}n)}{\theta (\underset{¯}{\omega }+\underset{¯}{C}\underset{¯}{B}n)}}\right)|_{\underset{¯}{\omega }=0}\hfill \\ & =\mathrm{\Lambda }_0{\displaystyle \underset{j=1}{\overset{p}{}}}c_j(p)_{z_j}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{A}\underset{¯}{z})}{\theta (\underset{¯}{C}\underset{¯}{z})}}\right)|_{\underset{¯}{z}=\underset{¯}{B}n},\hfill \end{array}`$ (B.14)
where by (2.78),
$$\underset{¯}{B}=\underset{¯}{U}_0^{(3)}.$$
(B.15)
Introducing the meromorphic (nondegenerate) function $`𝒱:^p\{\mathrm{}\}`$ by
$$𝒱(\underset{¯}{z})=\mathrm{\Lambda }_0\underset{j=1}{\overset{n}{}}c_j(p)_{z_j}\mathrm{ln}\left(\frac{\theta (\underset{¯}{A}\underset{¯}{z}\mathrm{diag}(\underset{¯}{B}))}{\theta (\underset{¯}{C}\underset{¯}{z}\mathrm{diag}(\underset{¯}{B}))}\right),$$
(B.16)
one observes that
$$b(n)=𝒱(\underset{¯}{z})|_{\underset{¯}{z}=(n,\mathrm{},n)}.$$
(B.17)
In addition, $`𝒱`$ has a basis of periods
$$\{\underset{¯}{e}_j\left(\mathrm{diag}(\underset{¯}{B})\right)^1,\underset{¯}{\tau }_j\left(\mathrm{diag}(\underset{¯}{B})\right)^1\}_{j=1}^p$$
(B.18)
by (B.6), where
$`\underset{¯}{e}_j\left(\mathrm{diag}(\underset{¯}{B})\right)^1`$ $`=(0,\mathrm{},0,\underset{j}{\underset{}{B_j^1}},0,\mathrm{},0),j=1,\mathrm{},p,`$ (B.19)
$`\underset{¯}{\tau }_j\left(\mathrm{diag}(\underset{¯}{B})\right)^1`$ $`=(\tau _{j,1}B_1^1,\mathrm{},\tau _{j,p}B_p^1),j=1,\mathrm{},p.`$ (B.20)
By hypothesis, $`b`$ in (B.14) is quasi-periodic and hence has $`p`$ real (scalar) quasi-periods. The latter are not necessarily linearly independent over $``$ from the outset, but by slightly changing the locations of branchpoints $`\{E_m\}_{m=0}^{2p+1}`$ into, say, $`\{\stackrel{~}{E}_m\}_{m=0}^{2p+1}`$, one can assume they are. In particular, since the period vectors in (B.18) are linearly independent and the (scalar) quasi-periods of $`b`$ are in a one-one correspondence with vector periods of $`𝒱`$ of the special form (B.19) (cf. (B.9), (B.10)), there exists a homology basis $`\{\stackrel{~}{a}_j,\stackrel{~}{b}_j\}_{j=1}^p`$ on $`𝒦_p`$ such that the vector $`\underset{¯}{\overset{~}{B}}=\underset{¯}{\overset{~}{U}}_0^{(3)}`$, corresponding to the normalized differential of the third kind, $`\stackrel{~}{\omega }_{P_\mathrm{}+,P_{\mathrm{}}}^{(3)}`$ and this particular homology basis, is real-valued. By continuity of $`\underset{¯}{\overset{~}{U}}_0^3`$ with respect to $`\stackrel{~}{E}_0,\mathrm{},\stackrel{~}{E}_{2p+1}`$, this proves (B.12). ∎
Acknowledgments. We are grateful to Leonid Golinskii for pointing out references and to us.
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# Two Body Relaxation in Simulated Cosmological Haloes
## 1 Introduction
$`N`$-body modelling of the dynamical evolution of the dark matter component involves the approximation whereby the system to be simulated is represented by a surrogate one consisting of many order of magnitudes fewer particles. Discreteness noise is thus necessarilly greatly enhanced. Furthermore, within the context of the cold dark matter scenario, material rapidly collapses into haloes; in a hierarchical process that ensures that the first structures are badly resolved, with particle noise propagating through the hierarchy, further enhancing discreteness effects and weakening the $`N`$-dependence of the effective relaxation time (Binney & Knebe 2002; Diemand et al. 2004).
Relaxation is expected to play its most prominent role in the very central regions of collapsed structures — precisely the place where there remains considerable controversy as to whether the slope of the density profile persists as a power law up to the resolution limit (Diemand et al. 2005), or instead gradually flattens into a smooth core (Navarro et al. 2004). From the difference between these two situations follow important implications concerning the compatibility of the cold dark matter scenario with the observed inner mass distribution in galaxies. Relaxation also affects the symmetry of a system in physical and velocity space, rendering it more isotropic; and triaxial haloes may play an important role in determining the dynamics of the baryonic component (e.g., El-Zant & Shlosman 2002; El-Zant et al. 2003).
Since substructure haloes are resolved with far fewer particles, relaxation is expected to have a more dominant effect on their internal structure and (by modifying the efficiency of stripping) their spatial distribution. If accretion and merging are are dominant processes in determining them (???), this will have important implications for the parent halo properties too.
Estimating the importance of particle noise in $`N`$-body simulations is an essential step in evaluating the role it may play in determining the dynamical state of collapsed structures. Our focus here is on the generalisation of the standard formulation of the problem, leading to the familiar formulas for the ‘relaxation time’ in stellar systems, to the cosmologically relevant case, when the physical extent and the mass distribution of the object under consideration continuously vary with time. The details of the procedure, along with the assumptions involved, are outlined in Section 2. In our treatment we will distinguish between ‘parent’ haloes, representing the most massive progenitors and ‘subhaloes’, representing the substructure. The former are assumed to continue to grow in mass and size, while the latter’s growth stops once they are incorporated into a larger structure. In both cases we derive simple formulas for calculating characteristic relaxation times (Section 3) and, in addition, undertake a more detailed calculation of the expected root mean square (RMS) perturbation to trajectories’ velocities for haloes evolving in the context of the ‘concordance’ CDM model; giving simple empirical fits to the results (Section 4). Conclusions are summarised in Section 5.
## 2 The model
We are interested in the description of relaxation due to discreteness (particle) noise in evolving cosmological haloes. The model through which this is done is described below. First the time evolution of the halo mass distribution is parametrised; the standard description of relaxation in gravitational systems is then generalised to encompass the case of these evolving systems.
### 2.1 Evolution of the halo mass distribution
Describing the temporal evolution of the halo mass distribution is simplified by the fact that their densities can always be approximately parametrized by double power law functions, independently of the mass of the halo, or the redshift at which it is identified. Two profiles have been extensively used: that introduced by ? (NFW) has density $`1/r`$ in the inner region and $`1/r^3`$ in the outer one; the associated mass fraction within radius $`r`$ is
$$\frac{M(r)}{M(r_{\mathrm{vir}})}=\frac{\mathrm{ln}(1+r/r_s)\frac{r/r_s}{1+r/r_s}}{\mathrm{ln}(1+c)\frac{c}{1+c}}.$$
(1)
The profile put forward by ? (M99) has similar asymptotic density variation at large $`r`$, but a steeper central cusp ($`1/r^{3/2}`$). It has mass distribution
$$\frac{M(r)}{M(r_{\mathrm{vir}})}=\frac{\mathrm{ln}[1+(r/r_s)^{3/2}]}{\mathrm{ln}\left(1+c^{3/2}\right)}.$$
(2)
In these formulas, $`r_s`$ is the characteristic halo scale length (separating the $`1/r`$ from the $`1/r^3`$ region) and $`r_{\mathrm{vir}}`$ is the virial radius. Their ratio $`c=r_{\mathrm{vir}}/r_s`$ determines the mean density contrast inside radius $`r`$: $`\mathrm{\Delta }=\frac{M(r)}{M(r_{\mathrm{vir}})}\frac{r_{\mathrm{vir}}^3}{r^3}`$ (cf. Fig 1).
The virial radius is usually defined as the radius inside which the average density is $`200`$ times the critical density for closure $`\rho _c`$. Relative to $`\rho _c`$, the mean density contrast inside $`r`$ is
$$\nu =\frac{\rho (r)}{\rho _c}=200\times \mathrm{\Delta }=200\frac{M(r)}{M(r_{\mathrm{vir}})}\frac{r_{\mathrm{vir}}^3}{r^3}.$$
(3)
In a universe with matter and vacuum energy densities such that $`\mathrm{\Omega }_m(z)+\mathrm{\Omega }_V(z)=1`$, the virial radius varies with redshift as
$$\frac{r_{vir(z)}}{r_{vir(0)}}=\frac{(M(r_{\mathrm{vir}},z)/M_0)^{1/3}}{(\mathrm{\Omega }_m(1+z)^3+\mathrm{\Omega }_V)^{1/3}},$$
(4)
where $`M_0=M(r_{\mathrm{vir}},z=0)`$. ? suggested the empirical relation whereby
$$M_{\mathrm{vir}}(z)=M_0e^{\alpha z}.$$
(5)
This defines a characteristic halo growth time<sup>1</sup><sup>1</sup>1The exponential form reflects a two phased evolution; a period of rapid accretion followed by slower growth. Zhao et al. (2003) have confirmed its basic premises for halos growing through almost four orders of magnitudes in mass and three in radius (although they do not use this particular form to fit the data).
$$\tau _s\left(\frac{d\mathrm{ln}M}{dt}\right)^1=\frac{1}{\alpha dz/dt}.$$
(6)
The parameter $`\alpha 1`$ varies systematically with halo mass (but appears constrained to the range $`0.4\alpha 1.6`$). In this paper we fix $`\alpha =1`$, noting that our results were found to be quite insensitive as to the exact choice of $`\alpha `$. Wechsler et al. also find that, for NFW type fits to the density profile, $`c(1+z)^1`$, with an average value of order $`10`$ at $`z=0`$. Their results thus generally confirmed the prescription (due to Bullock et al. 2001)
$$c[\mathrm{NFW},z]=\frac{9}{1+z},$$
(7)
which is used here. We assign a minimal value of $`c=9/5`$ at high $`z`$ (because structures with $`\rho 1/r`$ in their outer regions are necessarily far from virial equilibrium; the simulations of Bullock et al. 2001; Zhao et al. 2003 and Tasitsiomi et al. 2004 confirm the presence of a minimal $`c`$). The concentration parameter of an M99 fit relates to to that of an NFW fit through $`c[\mathrm{M99}](c[\mathrm{NFW}]/1.7)^{0.9}`$.
For any given $`z`$ there is a mass dependent spread in $`c`$; it is given, in terms of the typical mass scale $`M_{}(z)`$ collapsing at redshift $`z`$, via $`c(M_{\mathrm{vir}}/M_{})^{0.13}`$. Nevertheless, this relation shows that, even for a halo that is ten orders of magnitude more massive than $`M_{}`$, $`c`$ only changes by a factor of ten. From Fig. 1, this leads to a corresponding change in the density contrast of about the same factor or less. But since the relaxation time (cf, Eq. 15 below) $`N(r)\tau _D(r)\sqrt{\rho (r)}`$ inside any fixed radius; it changes by a factor of only a few even for such a highly unlikely mass deviation. Trials with different normalisations in Eq. (7), and with the prescription proposed by (Eq. 10 of) Zhao et al. (2003), have confirmed this. In fact it turns out (Section 4.2) that the relaxation rate is quite a strong function of radius, and will be important at small radii for any reasonable set of $`c`$ values during the evolution. The results presented in this paper assume the quoted form for this relation, ignoring the weak mass dependence. Remarkably, our general conclusions should hold for haloes in any mass range
### 2.2 Generalisation of the standard formulation
The standard approach to relaxation in gravitational systems has developed along the foundation set by ?. In analogy with the case of dilute gases, it is assumed that perturbations experienced by a test particle arise from independent and local two body encounters. Since the mean field along a particle’s trajectory changes on a timescale comparable to its dynamical time $`\tau _D`$, while the timescale of an average encounter is $`\tau _D/N`$, this assumption is generally satisfied; encounters with duration $`\tau _D`$ being relatively rare. As the number of particles in the representation of a self-gravitating system is increased, the effect of strong encounters becomes increasingly unimportant; in softened systems, the subject of cosmological simulations, strong encounters are completely suppressed if the softening length is of the order of the interparticle distance. The standard formulation of relaxation theory therefore assumes that the effect of discreteness noise can be modelled in the form of weak random encounters. Numerical tests (e.g., ?), seem to vindicate these basic premises.<sup>2</sup><sup>2</sup>2Note nevertheless that such work generally focused on energy changes along particle trajectories. The response of the trajectories themselves (and thus the effect on quantities like angular momentum) can exhibit significantly stronger sensitivity to noise — even in the simplest case of deflections caused by fixed background particles (??).
The product of weak, local and independent encounters is a diffusion process, whereby a particle’s dynamical variables undergo a random walk around their unperturbed values. For a spherical system with isotropic velocities that can approximated in terms of a Maxwellian, the mean diffusion coefficient within radius $`r`$
$$(\mathrm{\Delta }v)^2=\frac{G^2\rho m\mathrm{ln}\mathrm{\Lambda }}{K\sigma }$$
(8)
describes the average rate at which the square velocities of particles deviate from those of unperturbed trajectories in a corresponding smooth system (with $`N\mathrm{}`$). Here $`\mathrm{\Lambda }`$ is the ratio of maximum to minimum impact parameters and $`K=1/15.4`$ (according to Spitzer & Hart 1971). In quasiequilibrium, at any given instant, the average velocity dispersion inside radius $`r`$ can be calculated from
$$\sigma ^2(r)=\frac{3}{r}_0^r\overline{v_r^2}𝑑r=\gamma ^2v_c^2=\gamma ^2GM/r=\gamma ^2G\frac{4}{3}\pi \rho r^2.$$
(9)
The local one dimensional velocity dispersion $`\overline{v_r^2}`$ can be evaluated by solving the Jeans equation (cf. Appendix B). The above relations thus define the weak function of radius $`\gamma `$ that is used in the calculations below. This definition should hold because the virialised region of a cosmological halo remains near equilibrium through most of its evolution — major mergers being relatively rare, and even when they do occur the system is out of equilibrium for a time $`\tau _D`$ generally much smaller than the (relaxation) timescales of interest.
When the system parameters are changing, as during the growth of a cosmological halo, the mass inside any given radius will be time dependent. As a consequence, there will be corresponding variations in $`\rho =\frac{M(r,t)}{4/3\pi r^3}`$ and $`\gamma `$. Nevertheless, if the timescale between encounters giving rise to the relaxation process ($`\tau _e\tau _D/N`$) is much smaller than the timescale for variation of the system parameters $`\tau _s`$, then local (in time) averages are allowed, and the diffusion coefficient becomes a well defined function of time. From Eq. (6), and because
$$|\dot{z}|=(1+z)H$$
(10)
and
$$\rho _c=\nu \frac{3H^2}{8\pi G},$$
(11)
one finds that $`t_s\frac{1}{1+z}\sqrt{\frac{3}{8\pi G\rho _c}}`$. Therefore, unless $`z`$ is very large, $`\tau _s>\tau _D1/\sqrt{\nu G\rho _c}`$ (recall that $`\nu 200`$); with the implication that the temporal locality condition ($`\tau _e\tau _s`$) is in fact weaker than that for spatial locality ($`\tau _e\tau _D`$), required for the validation of the diffusion approach.
The expression for the mean diffusion coefficient, inside radius $`r`$ and at time $`t`$, can then be written as
$$\frac{(\mathrm{\Delta }v)^2(r,t)}{\sigma ^2(r,t)}=\frac{\sqrt{G\rho (r,t)}}{K\gamma ^3(r,t)\sqrt{4\pi /3}}\frac{m}{M(r,t)}\mathrm{ln}\mathrm{\Lambda }(r,t).$$
(12)
The relative mean square perturbation to particle velocities due to encounters with other particles during the halo evolution is
$$v_p^2/\sigma ^2=_{t_f}^{t_0}\frac{(\mathrm{\Delta }v)^2(r,t)}{\sigma ^2(r,t)}𝑑t,$$
(13)
where $`t_f`$ refers to some chosen initial, reference, formation time (for any given $`r`$ measured at $`z=0`$, we will take the redshift to be the time when $`r_{\mathrm{vir}}(t_f)=r`$) and $`t_0`$ refers to the end of the simulation (assumed to correspond to $`z=0`$).
Using (12) and (13) one can define the relaxation time in a simulated cosmological halo implicitly:
$$v_p^2/\sigma ^2(t_{\mathrm{relax}})=_{t_f}^{t_{\mathrm{relax}}}\frac{\sqrt{G\rho }}{K\gamma ^3\sqrt{4\pi /3}}\frac{m}{M}\mathrm{ln}\mathrm{\Lambda }dt=1.$$
(14)
If $`t_{\mathrm{relax}}t_0`$ the dynamics can be expected to be completely dominated by discreteness noise. This is in line with the standard application to time independent systems, where $`t_f=0`$ and $`v_p^2/\sigma ^2`$ is assumed to be time independent; the the result are the familiar expressions for the relation time
$$t_{\mathrm{relax}}=K\frac{\sigma ^3}{G^2\rho m\mathrm{ln}\mathrm{\Lambda }}=\frac{K\gamma ^3\sqrt{\frac{4}{3}\pi }}{\mathrm{ln}\mathrm{\Lambda }\sqrt{G\rho }}\frac{M}{m}0.1\frac{N}{\mathrm{ln}\mathrm{\Lambda }}\tau _D$$
(15)
(where $`N(r)`$ and $`\tau _D(r)`$ are the particle numbers and mean dynamical time inside $`r`$).
To fix the value of $`\mathrm{\Lambda }`$, we assume that the resolution of the simulations in question corresponds to the local interparticle distance. The average of this quantity within radius $`r`$ will define the minimum impact parameter $`b_{\mathrm{min}}=(\frac{4}{3}\pi r^3/N(r))^{1/3}`$. The maximum impact parameter will be taken to correspond to the virial radius. We therefore have
$$\mathrm{\Lambda }(r,t)=\frac{b_{max}}{b_{min}}=\left(\frac{3}{4\pi }\right)^{1/3}N^{1/3}(r,t)\frac{r_{vir}(t)}{r}.$$
(16)
## 3 Simple estimates
In the next section we will numerically evaluate the integral in Eq. (14), invoking the time dependence appropriate in the currently favoured $`\mathrm{\Lambda }`$CDM cosmology. Some insight can however be gained by defining characteristic relaxation times for an evolving cosmological halo and for its substructure. This is the subject of this section.
### 3.1 A characteristic relaxation time for cosmological haloes
We consider first the evolution of the parent halo, or most massive progenitor, and assume that all mass is added to it in the form of a smooth component. Of course, in reality, accretion of material takes the form of clumpy subhaloes, but most of their mass is rapidly stripped, so that, at any given time, there is usually only $`1020\%`$ of the mass of the halo in the substructure. Their internal relaxation is dealt with separately below.
If the relaxation time within some radius $`r`$ is smaller than the timescale $`\tau _s`$ for the mass distribution to significantly evolve, it is likely that the system is heavilly affected by discreteness noise within radius $`r`$ — the rationale being that if the local relaxation time is small compared to the time it takes for more particles to be added to the system (thus decreasing the relaxation rate) then discreteness noise will have a significant effect.
A local relaxation time can be obtained by freezing the system at some time $`t=t(z)`$ and using (15). Comparing the two timescales one can write
$$R=t_{\mathrm{relax}}/\tau _\mathrm{s}=\frac{K\alpha \gamma ^3\sqrt{\frac{4}{3}\pi }}{\mathrm{ln}\mathrm{\Lambda }}\frac{|dz/dt|}{\sqrt{G\rho }}\frac{M}{m},$$
(17)
and eliminate $`dz/dt`$ by invoking (3), (10) and (11) to get
$$R(r,z)=\frac{0.087\sqrt{2/\nu }\pi \alpha \gamma ^3}{\mathrm{ln}\mathrm{\Lambda }}(1+z)N(r,z),$$
(18)
where (by Eq. 5) the number of particles inside $`r`$ at redshift $`z`$ can be expressed in terms of the final number of particles with which the halo is resolved $`N_0=M_0/m`$:
$$N(r,z)=\frac{M}{M_0}N_0e^{\alpha z}.$$
(19)
The mass ratio entering into this last relation can be calculated via either Eq. (1) or Eq. (2).
The dimensionless relaxation time $`R`$ has to be significantly greater than unity, at all $`z`$, if the simulated halo can be considered free of artificial relaxation inside radius $`r`$.<sup>3</sup><sup>3</sup>3Note that, because the virial radius and concentration are time dependent, accreted mass is not deposited uniformly; there is a preference, especially at later times, for mass increase in the outer regions (e.g., Zhao et al. 2003). However this implies that the characteristic time for mass change in the inner regions, where most of the relaxation occurs, is larger than that predicted by Eq. (6). The condition $`R>1`$ thus constitutes a minimal requirement. What is required therefore is that
$$N_0>11.5\frac{\sqrt{\nu /2}}{\pi \alpha \gamma ^3}\mathrm{ln}\mathrm{\Lambda }\frac{M_0}{M}\frac{e^{\alpha z}}{1+z}.$$
(20)
In Fig. 2 we show the variation of the of the quantity $`(\frac{\gamma ^3}{\mathrm{ln}\mathrm{\Lambda }})N_0`$ as a function of radius at different redshifts. The radii are expressed in terms of the final virial radius (at $`z=0`$), by transforming them using Eq. (4). The cutoffs in the curves correspond to the virial radius at the denoted redshift (that is the maximum size of the halo considered at that redshift).
From this figure it is apparent that, provided $`\gamma ^3\mathrm{ln}\mathrm{\Lambda }`$, that artificial relaxation may have a dominant effect inside the inner $`1\%`$ of the final virial radius for $`N_010^6`$. Note that the demand for larger particle number is most stringent at larger times (smaller $`z`$). This is because the characteristic growth time $`\tau _s`$ is a steeply increasing function of $`z`$. The figure also suggests that, for fixed $`\gamma `$, the relaxation effects are more pronounced in the inner regions of NFW haloes, compared to the M99 case. Nevertheless, as we will see in Section 4.2, the introduction of variation in $`\gamma `$ reverses this situation, since $`\gamma `$ is significantly larger in the inner regions of NFW systems (this follows from the variations of the ratio of velocity dispersion to circular speed shown in Fig. B1).
### 3.2 Characteristic relaxation time for substructure
We consider a subhalo accreted at redshift $`z_a`$ and remaining a separate dynamical system up to $`z=0`$ — i.e., it stops growing, its outer regions are in fact stripped, but it keeps a dynamically distinct core. For this purpose we exploit the possibility of relating the relaxation and the local Hubble time; and in order to get simple closed form estimates, we will neglect relaxation at times $`t<t(z_a)`$ and focus on relaxation effects since accretion. From equations (15), (10) and (11) we can define a characteristic relaxation time
$$t_{\mathrm{relax}}(r,z_a)=\frac{0.087\pi \gamma ^3\sqrt{2/\nu }}{\mathrm{ln}\mathrm{\Lambda }}\frac{M}{m}H^1(z_a),$$
(21)
where $`M,\nu ,\gamma `$ and $`\mathrm{\Lambda }`$ are all measured inside radius $`r`$ at $`z=z_a`$. This timescale should be compared to the time interval separating the accretion epoch $`t(z_a)`$ and the end of the simulation at $`t_0=t(0)`$. This determines a characteristic substructure relaxation time given by
$$t_{rsub}(r,z_a)=\frac{t_{\mathrm{relax}}(r,z_a)}{t_0t(z_a)}.$$
(22)
Now further assume that $`t(z)=\frac{2}{3}H`$ and $`z=(t_0/t(z))^{2/3}1`$, as is appropriate for an Einstein de-Sitter universe (in the next section we verify our results by undertaking the full calculation, including pre-accretion evolution, in the currently favoured “concordance” cosmology<sup>4</sup><sup>4</sup>4Note that while $`t=t(z)`$ for $`\mathrm{\Lambda }`$CDM is quite different from that in an Einstein-de-sitter model the ratios $`t(z_1)/t(z_2)`$, entering into the equation below, differ by a factor of at most by $`3\%`$ for $`10z1`$ and by $`30\%`$ up to $`z=0`$.). In this case we have
$$t_{rsub}(r,z_a)=\frac{3}{2}\frac{0.087\pi \gamma ^3\sqrt{2/\nu }}{\mathrm{ln}\mathrm{\Lambda }}\frac{N(r,z_a)}{(z_a+1)^{3/2}1}.$$
(23)
If, within some given radius, this quantity is less than unity the implication is that, by $`z=0`$, the dynamics has become dominated by artificial discreteness noise.
In Fig. 3 we reproduce several plots where this dimensionless relaxation time, thus defined, is shown for haloes assumed to have $`N(r_{\mathrm{vir}},z_a)=1000`$ when they are subsumed into their parent structures at different $`z_a`$ (in the case of satellites haloes we only reproduce the results for NFW haloes, those for the M99 haloes are very similar). As can be seen, especially at relatively small radii, the relaxation time can be significantly smaller than the subhalo lifetime. Furthermore, it is to be noted that it will be this inner core of the subhalo that will survive stripping (see also Fig. 6 for more detail concerning the radial distribution of the relaxation time and Section 5.2 for a discussion of some possible consequences).
## 4 Direct Calculations
### 4.1 General considerations
From equations (12) and (13), the expected relative mean square perturbation due to discreetness imposed on to a test particle moving within an evolving cosmological halo is
$$v_p^2/\sigma ^2=\frac{1}{2K}\sqrt{\frac{3G}{\pi }}_{t_f}^{t_0}\mathrm{ln}\mathrm{\Lambda }(r,t)\frac{\sqrt{\nu (r,t)}}{\gamma ^3(r,t)}\frac{\sqrt{\rho _c(t)}}{N(r,t)}𝑑t.$$
(24)
The integral is evaluated within a given fixed fraction of the final virial radius $`r_{vir}(t_0)`$ — with the consequence that the radial coordinate contains an implicit time dependence (i.e, $`r=r(t)`$). The relations in Section 2.1 can be used to determine this dependence, as well as that of the other quantities entering into (24), provided the evolution of redshift as a function of time is known. For a flat universe with cosmological constant, the required transformations are given in Appendix A (in the calculations below we use $`\mathrm{\Omega }_V=0.7`$ and $`h=0.7`$). Since the virial radius increases with time, a given fraction of the virial radius at $`t=t_0`$ will correspond to the whole virial radius at some earlier time. This naturally determines the lower limit of integration, the formation time $`t_f`$ — for some fixed fraction of $`r_{\mathrm{vir}}(z=0)`$, it corresponds to a redshift $`z_f`$ where this fraction is equal to whole virial radius $`r_{\mathrm{vir}}(z_f)`$.
### 4.2 Relaxation of the most massive progenitors
In Fig. 4 we show the RMS perturbation due to particle noise, as a function of radius, expected for a halo with mass evolving according to Eq. 5 (with $`\alpha =1`$) for several final values of the final total particle number $`N_0`$. They are obtained by solving (24) using an adaptive integrator (NAG D01AJF) with a tolerance of $`10^4`$.
An interesting characteristic of the plots in Fig. 4 is the perfect power law behaviour of $`v_p^2/\sigma ^2^{1/2}`$ over most of the radial interval. This is a consequence of a curious property of cosmological haloes, already noted by ?; namely the power law form of the phase space density. It is this phase space density, $`\rho /\sigma ^3`$, that determines the rate of relaxation. This property enables one to deduce a particularity simple fit for the variation of the relative RMS perturbation:
$$v_p^2/\sigma ^2^{1/2}10^3\sqrt{\frac{10^8}{N_0}}\frac{r_{\mathrm{vir}}}{r},$$
(25)
where $`r/r_{\mathrm{vir}}`$ refers to the fraction of the virial radius at $`z=0`$. The number of relaxation times inside radius $`r`$ can be counted as follows:
$$n(t_{\mathrm{relax}})=v_p^2/\sigma ^210^6\frac{10^8}{N_0}\left(\frac{r_{\mathrm{vir}}}{r}\right)^2,$$
(26)
which of course needs to be much smaller than one if artificial relaxation is to be negligible within radius $`r`$.
For $`N_010^6`$ then, particle motion is expected to be entirely dominated by noise in the inner percent of the virial radius; regions bounded by radii an order of magnitude smaller still having undergone $`100`$ relaxation times. These results are in agreement with the simple estimates presented in Section 3.1 (with the difference that the effect on M99 haloes are larger here, because variations in $`\gamma `$ is taken into account).
The inner $`0.1\%`$ is one relaxation time old even for $`N_0=10^8`$. Indeed, a final particle number inside the halo’s virial radius of $`N_0=10^{12}`$ is required for noise to be reduced to the level of a few percent there. For a million particles this noise level is only reached at almost a third of the final virial radius; and this is larger than the average halo scale length $`r_s=r_{\mathrm{vir}}/c`$ (c.f. Eq. 7).
Fig. 5 shows, for the inner percent of the final virial radius, the $`z`$-evolution of the relaxation rate (the integrand in 24) along with the principal components determining it. It shows that the relaxation rate decreases with redshift, for both the NFW and M99 profiles — the effect being more pronounced in the former case because the ratio of velocity dispersion to circular velocities that determines $`\gamma `$ decreases faster with radius (Appendix B); and at smaller $`z`$ because, for a given fixed fraction of the final virial radius, smaller radii are probed (recall that at $`z_f`$ the whole virial radius corresponds to the inner one percent at $`z=0`$).
Finally, note that, since new particles are continuously being introduced into the region, the way in which particles are affected by relaxation will differ, with the probable consequence that relaxation driven evolution can be expected to be different than in a static NFW type system.
### 4.3 Relaxation of subhaloes
In Fig. 6 we show the RMS perturbation in velocities expected as a result of particle noise for a subhalo composed of a thousand particles when it stops growing (after being incorporated into a larger structure). Before this, we assume that its mass grows at the exponential rate described by Eq. (5) with (as always) $`\alpha =1`$. These results are in close agreement with the simple estimates derived in Section 3.2; they reinforce the conclusion that relaxation can be a significant factor determining the structure of subhaloes in simulations, their distribution and their effect on the parent halo structure (Section 5.2).
## 5 Conclusion
This paper presented an attempt at assessing the importance of particle noise in the evolution of collapsed cold dark matter structures, by generalising the diffusion formulation first proposed by Chandrasekhar to the case when the mass distribution and radial extent of the system under consideration are variable. Although a most massive progenitor of a $`z=0`$ halo will typically have increased its collapse mass by more than an an order of magnitude by accreting subhaloes, at any given instant $`80\%`$ of its mass is in a smooth component, most material in subhaloes having been stripped (e.g., Gao et al. 2004). This enables one to divide the analysis in two parts. In this context, ‘parent’ haloes, i.e ones characterised by continuously increasing (virial) mass and radius throughout their evolution, were grown according to the empirical formula discovered by Wechsler et al. (2002). Subhaloes were treated differently, by assuming that their growth is arrested at accretion. Stripping was not explicitly taken into account. But since it is in the inner region of subhaloes that relaxation is most significant, this does not appear to represent a major idealisation — as high resolution simulations, specially designed for the purpose of examining the issue, show the inner regions of stripped haloes to be largely unaffected (Kazantzidis et al. 2004).
We have derived simple closed form formulas for the characteristic relaxation times (Section 3), and also integrated the relevant equation characterising the expected RMS perturbation to the velocity directly (Section 4), The predictions of the two approaches agree quite closely. In what follows we summarise our results and sketch some possible consequences.
### 5.1 Relaxation in parent haloes
Our results were not found to be significantly affected by the exact form chosen for the evolution of the concentration parameter $`c`$ or the parameter $`\alpha `$ characterising the mass growth rate (cf. Eq. 5), both quantities that are dependent on the final halo mass (in physical units), so this latter quantity does not enter directly into our presentation here, which is concerned only with the relaxation of particles after they have been accreted onto the most massive progenitor and are subsequently part of its smooth component (in the next two subsection we consider relaxation inside substructure, and briefly comment on the expected mass dependence of noise-induced relaxation in terms of merger history).
The power law (with index $`2`$) radial variation of the phase space density of cosmological haloes ensures a steep dependence of the relaxation rate, a linear function of this quantity, on radius. The perturbation due to discreetness noise is adequately quantified by the empirical fits given in Section 4.2.
A test particle that is present inside the inner percent of the final virial radius, from the point in time when the virial radius of the halo was equal to this fraction, will experience a relative RMS perturbation to its velocity $`1`$ when the final halo is resolved with $`10^6`$ particles — implying that its dynamics is completely corrupted by discreteness noise. In the very inner few thousandth of the final radius the RMS perturbation due to discreteness noise is an order of magnitude larger, meaning that these regions are of the order of a hundred relaxation times old.
In an evolving cosmological halo, particles are continually added within any given radius; it is also likely that some can gain energy (e.g. by interaction with sinking satellites; cf. El-Zant et al. 2004, El-Zant 2005) and move into trajectories with larger average radii. Particles that are accreted at a later stage are more mildly affected by relaxation. At any given fraction of the final radius therefore, there will be a range in the severity of artificial relaxation, depending on the accretion epoch of individual particles. This somewhat complex situation implies that effects of relaxation may be rather different than the familiar situation of fixed-mass systems, and therefore difficult to detect.
Our calculations thus represent a quantification of the magnitude of the perturbation of the discreteness noise, rather than its effect on the detailed dynamics and macroscopic structure. While it is quite probable that the significant role for artificial relaxation predicted by our analysis is reflected in the contradictory claims concerning the convergence and shape of the inner density profile of CDM haloes (Navarro et al. 2004; Diemand et al. 2005), with the predicted effect for static systems being a weakening of the cusp, further work is needed to determine the consequence in the case of evolving cosmological haloes.
The RMS perturbation to the velocities decreases as $`1/\sqrt{N}`$, implying that a parent halo needs to be resolved with $`N10^{10}`$ particles at $`z=0`$ if discreteness noise is not to dominate the dynamics of the inner $`0.1\%`$ of the virial radius; and if it is to be negligible ($`v_p^2/\sigma ^2^{1/2}1\%`$) in the whole region where $`r0.1r_s0.01r_{\mathrm{vir}}`$. While perturbations of order unity are likely required for significant energy relaxation and restructuring of the azimuthally averaged density profile, much smaller perturbations may affect particle trajectories, with accompanying modification (isotropisation) of the velocity distribution and loss of triaxiality in spatially asymmetric systems (e.g., ?; see also Section 2.2).<sup>5</sup><sup>5</sup>5The reason this happens is due to a transition from regular to chaotic motion, which is very sensitive to perturbation. A very simple example of this phenomenon is the case of a pendulum on a ‘separatrix’ trajectory passing near the unstable (upper) equilibrium. Small perturbations can turn (regular) trajectories rotating in one direction into ones rotating into the opposite direction, or oscillating without a definite sense of rotation, or even transit between all these different possibilities (chaotic trajectories).
### 5.2 The relaxation of substructure
The results presented in this paper suggest that the situation with substructure is still more severe. A parent halo that is identified with $`N_010^7`$ has (from Eq. 5) $`10^6`$ particles at $`z2`$. Suppose a subhalo with $`N_{\mathrm{sub}}=0.001N`$ is then accreted. Assume that between redshifts two and one this $`1000`$-particle halo will be stripped of $`90\%`$ of its mass (as would be expected from Gao et al. 2004 Fig. 14). For an NFW halo this also corresponds to a truncation of about $`90\%`$ in radius. But according to Fig. 3, this inner region will be significantly affected by relaxation. In fact, by $`z=0`$, it is expected to have undergone enough relaxation times to have approached core collapse! But this fate is only likely if it is not, in the meantime, dissolved by stripping.<sup>6</sup><sup>6</sup>6Note that, because relaxation is quite a sensitive function of radius ($`t_{\mathrm{relax}}r^2`$ from Fig. 6 and Eq. 26) this same argument can be made for more central regions of more massive subhaloes (e.g., for the inner $`1\%`$ of a subhalo consisting of $`10\%`$ the mass of mass of the parent halo, accreted at $`z2`$ when the parent halo had $`N10^6`$), or more extended regions of less massive ones.
What effect is the significant relaxation expected to have on stripping? Cosmological haloes have a ‘temperature inversion’ in their inner velocity distribution (cf. Appendix B). Relaxation therefore causes initial expansion of the inner region, forming an isothermal core, before it recontracts like a relaxing globular cluster (e.g., ?). In the first phase, the core becomes less dense and more easily stripped; in the second the situation is reversed. Stripping is most efficient for subhaloes that venture near the centre of the parent, experiencing strong tidal forces. If, while the outer regions of a subhalo are stripped during successive passages, the inner region relaxes to a more diffuse density distribution, further stripping will be accelerated and the core may completely dissolve. Conversely, if the stripping is slow, there may be sufficient time for core contraction to take place. Further stripping is subdued and core dissolution becomes less likely.
If, as is suggested in several studies (e.g., ???), the structure of the halo profiles found in numerical simulations is dependent on the interaction between the subhaloes and their parents, then it is clear that such significant re-engineering of the subhaloes can be of major importance in determining the parent halo profile. It obviously also has implications for the number and spatial and orbital distribution of substructure. Given the discussion above, one would expect artificial relaxation to enhance the number of haloes that do not venture into the inner regions (i.e., those accreted on less eccentric orbits) and decrease the number of those that do venture there.
### 5.3 The N-scaling of the substructure-dominated relaxation time
The expectation from the results just outlined is that, beyond the inner percent or so of the final virial radius, relaxation will be dominated by particles inside substructure, or which have spent considerable time inside subhaloes before being stripped.
Simulations suggest that substructure has a mass function such that $`dn/dmm^{9/5}`$, quite independently of redshift (e.g., Gao et al 2004). If one ignores the logarithmic dependence, the relaxation time when expressed in units of dynamical time is, for any given subhalo, a linear function of the number of particles in it. Consider then a mean relaxation time averaged over subhaloes and normalised over an averaged dynamical time,
$$t_{ra}\frac{N_0_{m_{\mathrm{min}}}^{m_{\mathrm{max}}}m\times m^{9/5}𝑑m}{_{m_{\mathrm{min}}}^{m_{\mathrm{max}}}m^{9/5}𝑑m},$$
(27)
where the integration limits correspond to the least and most massive subhalo present. When the latter is orders of magnitudes more massive than the former, the above relation gives $`t_{ra}N_0m_{\mathrm{max}}m_{\mathrm{min}}^{4/5}`$. If both $`m_{\mathrm{min}}`$ and $`m_{\mathrm{max}}`$ are independent of $`N_0`$ then the relaxation time scales linearly in with $`N_0`$ as in the standard case.
If, on the other hand, one supposes that, because of increasing resolution, $`m_{\mathrm{min}}1/N_0`$ (proportionately smaller haloes appear in the resolved field as $`N_0`$ increases) then $`t_{ra}N_0^{1/5}`$, which is roughly the scaling claimed by Diemand et al. (2004) on the basis of direct calculations of the relaxation time in the context of cosmological simulations. Their results also show that the N-scaling is closer to the canonical linear relation as one moves towards the centre of the main halo. This is also to be expected; since, near the centre, stripping causes the fraction of mass in subhaloes to decrease, while the relaxation in the smooth (parent halo) component becomes more efficient; and this scales roughly linearly with particle number.
Finally, note that although we have not followed the detailed merger history, the expectation derived from the work presented here is that relaxation is more enhanced in cluster haloes — since these form relatively recently, from particles that spend most of their existence inside poorly resolved structures. In contrast, a galaxy sized halo acquires most of its mass (i.e., particles) at larger redshift. It follows that, given the same resolution at $`z=0`$, cluster halo particles would have, on average, been more affected by discreteness noise. This is in line with the conclusion reached by Diemand et al. (2004).
## ACKNOWLEDGMENTS
I am indebted to Miloš Milosavljević
for many valuable discussions and suggestions, as well as critical comments on an earlier version of this paper that led to significant improvement.
## Appendix A Time dependence of $`\rho _c`$ and $`z`$ in flat universes with cosmological constant
According to Kolb & Turner (P. 55), for a flat universe with cosmological constant one has
$$t=\frac{2}{3}H^1\mathrm{\Omega }_V^{1/2}\mathrm{ln}[\frac{1+\mathrm{\Omega }_V^{1/2}}{(1\mathrm{\Omega }_V)^{1/2}}],$$
(28)
which, using $`\mathrm{\Omega }_V=\frac{\mathrm{\Lambda }c^2}{3H^2}`$ and $`C_V=\sqrt{3\mathrm{\Lambda }}c`$ gives
$$\mathrm{\Omega }_V^{1/2}=\frac{e^{C_Vt}1}{e^{C_Vt}+1}$$
(29)
or equivalently
$$\frac{C_V}{3}\frac{e^{C_Vt}+1}{e^{C_Vt}1}dt=\frac{dz}{1+z},$$
(30)
which by integration (with $`A`$ the associated constant) gives
$$1+z=\frac{e^{\frac{C_V}{3}t}}{A[e^{C_Vt}1]^{2/3}}.$$
(31)
In a similar manner we find the critical density to vary as
$$\sqrt{\rho _C}=\frac{C_V}{2\sqrt{6\pi G}}\frac{e^{C_Vt}+1}{e^{C_Vt}1}.$$
(32)
Using these formulas the integral in Eq. (24) can be evaluated with respect to the variable $`C_Vt`$. When the constant in Eq. (31) is evaluated by assuming $`z=0`$ at $`t=t_0`$, the result does not depend on the value of $`t_0`$. This is analogous to the familiar case whereas two body relaxation in an isolated system does not depend on the time elapsed in physical units, but instead on the number of dynamical times the system has been through.
## Appendix B Velocity dispersions
The Jeans equation for spherical system with isotropic velocity dispersion can be written as (e.g., Binney & Tremaine 1987)
$$\frac{d(\rho \overline{v_r^2})}{dr}=\rho \frac{d\mathrm{\Phi }}{dr},$$
(33)
where $`\overline{v_r^2}=\frac{1}{3}\sigma ^2`$ is the radial velocity dispersion. The general solution is
$$\rho \overline{v_r^2}=\rho \frac{d\mathrm{\Phi }}{dr}𝑑r+C.$$
(34)
For any $`\rho 0`$, as $`r\mathrm{}`$, one must have $`C=0`$, if the velocity dispersion is to be bound at large radii. The ratio of this unique4 physical solution to the local circular velocity is
$$\frac{\overline{v_r^2}}{v_c^2}=\frac{_r^{\mathrm{}}\rho \frac{d\mathrm{\Phi }}{dr}𝑑r}{\rho \frac{d\mathrm{\Phi }}{dr}r},$$
(35)
the form reflecting the fact that, as opposed to circular motion, the velocity dispersion at any $`r`$ is caused by particles that venture in and out of that radius.
For power law density distributions $`\rho r^n`$ solutions of Eq. (34), with $`C=0`$, imply that
$$\frac{\overline{v_r^2}}{v_c^2}=\frac{1}{2n2},$$
(36)
for $`n1`$ and
$$\frac{\overline{v_r^2}}{v_c^2}=\mathrm{ln}r,$$
(37)
($`r<1`$). Note that, in the central regions, $`\overline{v_r^2}r\mathrm{ln}r`$ for the NFW profile and $`r^{1/2}`$ for the M99 profile. Spherical cosmological haloes thus, necessarily, have a ‘temperature inversion’, at least as long as the velocity dispersion can be considered isotropic.
In Fig. 7 we reproduce the ratio of the three dimensional velocity dispersion to circular velocity for both profiles dealt with in this paper. It determines the value of $`gamma`$ (corresponding to the average of this quantity within any given radius) that enter into the equations calculating the effects of relaxation.
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# The first eigenvalue of the Laplacian, isoperimetric constants, and the Max Flow Min Cut Theorem
## 1. Introduction
For a domain $`\mathrm{\Omega }^n`$ the *fundamental frequency* is defined by
(1)
$$\lambda _\mathrm{\Omega }=\underset{uC_0^{\mathrm{}}(\mathrm{\Omega })}{inf}R(u),R(u)=\frac{_\mathrm{\Omega }|u|^2}{_\mathrm{\Omega }u^2}.$$
If $`\mathrm{\Omega }`$ is bounded and has a Lipschitz boundary then this is the smallest eigenvalue of the Laplacian $`_i^2/x_i^2`$ on $`\mathrm{\Omega }`$ with Dirichlet boundary conditions. The minimum is attained by the corresponding eigenfunction, which lies in $`H_0^1(\mathrm{\Omega })C^{\mathrm{}}(\mathrm{\Omega })`$.
For most domains it is impossible to determine $`\lambda _\mathrm{\Omega }`$ precisely, so it is a fundamental problem to give estimates in terms of the geometry of $`\mathrm{\Omega }`$. Upper estimates can be obtained by choice of any ’test function’ $`u`$ in (1).
It is less clear how to obtain lower estimates for $`\lambda _\mathrm{\Omega }`$. One such estimate was given by Cheeger in \[Che:LBSEL\]: Define the *Cheeger constant* by
(2)
$$h_\mathrm{\Omega }=\underset{S\mathrm{\Omega }}{inf}\frac{|S|}{|S|}.$$
The infimum is taken over open subsets $`S`$. The absolute value signs denote $`(n1)`$-dimensional Hausdorff measure in the numerator and $`n`$-dimensional Hausdorff (= Lebesgue) measure in the denominator. Then *Cheeger’s inequality* says that
(3)
$$\lambda _\mathrm{\Omega }h_\mathrm{\Omega }^2/4.$$
The Cheeger constant is sometimes called an isoperimetric constant since it resembles the classical isoperimetric constant $`h_{\text{class}}=inf_{S\mathrm{\Omega }}\frac{|S|^{n/(n1)}}{|S|}`$. $`h_{\text{class}}`$ is scale invariant and in fact independent of $`\mathrm{\Omega }`$, and the solution of the classical isoperimetric problem is that the mimimizers are precisely the balls.
In contrast, $`h_\mathrm{\Omega }`$ depends strongly on $`\mathrm{\Omega }`$ and clearly scales as $`h_{r\mathrm{\Omega }}=r^1h_\mathrm{\Omega }`$. In general it is difficult to determine $`h_\mathrm{\Omega }`$ precisely. It is known that minimizers $`S`$ exist if $`\mathrm{\Omega }`$ is Lipschitz, and that $`S\mathrm{\Omega }`$ is smooth (if $`n7`$) and has constant mean curvature (see \[FriKaw:IEFEPLOCC\], Theorem 8 and Remark 9). This may be used to determine $`h_\mathrm{\Omega }`$ explicitly in some cases, for example for polygons (see \[KawLac:CCSCSP\]), and to give conditions when $`\mathrm{\Omega }`$ is itself a minimizer (see \[AltCasCha:CCCSR\], \[Giu:ESPMCEUWBC\]; such sets are called calibrable in the image processing literature). In general, one may hope for estimates on $`h_\mathrm{\Omega }`$ in terms of geometric data.
Again, upper estimates for $`h_\mathrm{\Omega }`$ are obtained by using a suitable ’test domain’ $`S`$, while it is less obvious how to obtain lower estimates.
One purpose of the present note is to point out a very simple idea how to obtain a lower estimate for $`h_\mathrm{\Omega }`$:
###### Proposition 1.
Let $`V:\mathrm{\Omega }^n`$ be a smooth vector field on $`\mathrm{\Omega }`$, $`h`$, and assume
(4) $`|V|`$ $``$ $`1`$
(5) $`\mathrm{div}V`$ $``$ $`h,`$
both pointwise in $`\mathrm{\Omega }`$.
Then $`h_\mathrm{\Omega }h`$.
###### Proof.
Clearly, one may restrict to sets $`S`$ with smooth boundary in (2). For such $`S`$ we have, by Green’s formula and (4), (5),
(6)
$$|S|_SV𝑑n=_S\mathrm{div}Vh|S|.$$
Proposition 1 seems to be little known in the geometric analysis community, although it is implicit in McKean’s proof of lower bounds for $`\lambda _\mathrm{\Omega }`$ in case $`\mathrm{\Omega }`$ is a complete, simply connected Riemannian manifold of strictly negative curvature \[McK:UBSDMNC\] (here, $`V`$ is taken as gradient of the distance to a fixed point). See \[CheLeu:EESWBMCHS\] and \[BesMon:EESWLBMC\] for other applications of the same idea, and also \[BesMon:EBTGA\] for another lower bound on $`\lambda _\mathrm{\Omega }`$ in terms of vector fields.
###### Example 2.
Let $`\mathrm{\Omega }=\{x^n:|x|<1\}`$. Since $`|\mathrm{\Omega }|/|\mathrm{\Omega }|=n`$, we have $`h_\mathrm{\Omega }n`$. The vector field $`V(x)=x`$ has $`|V|1`$ and $`\mathrm{div}Vn`$, so that, in fact, $`h_\mathrm{\Omega }=n`$.
It is important for the sequel to allow non-smooth vector fields. Consider the following classes:
$`𝒳_{\mathrm{div}}(\mathrm{\Omega })`$ $`=`$ $`\{VL^{\mathrm{}}(\mathrm{\Omega },^n):\mathrm{div}VL^2(\mathrm{\Omega })\}`$
$`𝒳_{\mathrm{BV}}(\mathrm{\Omega })`$ $`=`$ $`\{VL^{\mathrm{}}(\mathrm{\Omega },^n):V\text{ has bounded variation}\}.`$
$`\mathrm{div}V`$ is understood in the sense of distributions. Recall that, by definition, $`VL^{\mathrm{}}(\mathrm{\Omega },^n)`$ has *bounded variation* if all of its first derivatives $`V_i/x_j`$ (in the sense of distributions) are (signed) measures. For a vector field $`VL^{\mathrm{}}(\mathrm{\Omega },^n)`$, (4) is meant to hold almost everywhere and (5) in the sense of distributions. If $`V𝒳_{\mathrm{div}}𝒳_{\mathrm{BV}}`$ then $`\mathrm{div}V`$ is a measure, so (5) then holds also in the sense of measures.
Below, the class $`𝒳_{\mathrm{div}}`$ will occur in the context of the Max Flow Min Cut Theorem, and the class $`𝒳_{\mathrm{BV}}`$ will arise for vector fields defined via the distance function.
Addendum to Proposition 1. *Proposition 1 holds for vector fields $`V𝒳_{\mathrm{div}}𝒳_{\mathrm{BV}}`$.*
###### Proof.
The proof (6) still works, since for such $`V`$ one may define a ’restriction to the boundary’ $`V_{|S}`$ (for open $`S\mathrm{\Omega }`$ with Lipschitz boundary), which satisfies $`V_{|S}_{L^{\mathrm{}}(S)}V_{L^{\mathrm{}}(S)}`$ and Green’s formula. For $`V𝒳_{\mathrm{div}}`$ this is shown in \[Anz:PBMBFCC\]. For $`V𝒳_{\mathrm{BV}}`$ this follows from results in \[EvaGar:MTFPF\], Section 5.3. Theorem 1 there states that a function $`fL^1(S)`$ of bounded variation on a Lipschitz domain $`S^n`$ has a well-defined restriction to the boundary $`f_{|S}L^1(S)`$ satisfying, for any $`WC^1(^n,^n)`$, $`_Sf_{|S}W𝑑n=_Sf\mathrm{div}W+_SWf.`$ Applying this to $`f=V_i`$ and $`We_i`$ (the $`i`$th standard unit vector) for each $`i=1,\mathrm{},n`$ and summing over $`i`$ yields Green’s formula for $`V`$. Also, by Theorem 2 loc.cit. one has $`V_{|S}_{L^{\mathrm{}}(S)}V_{L^{\mathrm{}}(S)}`$.
It is a remarkable fact that the estimate in Proposition 1 is sharp. That is, $`h_\mathrm{\Omega }`$ may be characterized using vector fields:
###### Theorem 3.
We have
$$h_\mathrm{\Omega }=sup\{h:V\text{ satisfying }(\text{4}),(\text{5})\},$$
where the supremum is taken over smooth vector fields $`V`$ on $`\mathrm{\Omega }`$.
If $`\mathrm{\Omega }`$ is Lipschitz then there is a maximizing $`V𝒳_{\mathrm{div}}`$.
Theorem 3 may be regarded as a continuous version of the classical Max Flow Min Cut Theorem for networks. It was first proved by Strang \[Str:MFTD\] in two dimensions and by Nozawa \[Noz:MFMCTAN\] in general (see Theorem 4.4 there; Nozawa actually establishes a maximizing $`V`$ with $`\mathrm{div}VL^n(\mathrm{\Omega })`$). We explain the relation to the Max Flow Min Cut Theorem and sketch the proof of Theorem 3 in Section 2.
Given the Theorem, one can prove Cheeger’s inequality easily:
###### Proof.
(of Cheeger’s inequality (3) using Theorem 3.) If $`uC_0^{\mathrm{}}(\mathrm{\Omega })`$ and $`V`$ is a smooth vector field satisfying (4), (5) then, using Green’s formula, (with all integrals over $`\mathrm{\Omega }`$)
$$\begin{array}{c}hu^2(\mathrm{div}V)u^2=V(u^2)\hfill \\ \hfill 2|u||u|2\sqrt{u^2}\sqrt{|u|^2},\end{array}$$
so $`R(u)h^2/4`$. Taking a sequence $`V_k`$ with $`h_k=inf\mathrm{div}V_k`$ approaching $`h_\mathrm{\Omega }`$ one obtains $`R(u)h_\mathrm{\Omega }^2/4`$ for all $`uC_0^{\mathrm{}}(\mathrm{\Omega })`$ and therefore (3).
This is not a substantially new proof of Cheeger’s inequality: Cheeger’s original proof is essentially a similar estimate, plus a clever use of the coarea formula applied to $`u^2`$. But the proof of Theorem 3 also relies on the coarea formula (see Section 2)!
We remark that Proposition 1 and Theorem 3 extend directly to Riemannian manifolds with boundary, although (for Theorem 3) this is not stated explicitly in \[Noz:MFMCTAN\]. The relationship of Cheeger’s inequality and Max Flow Min Cut Theorems was first noted by Alon \[Alo:EE\] in the context of graphs (see also \[Chu:SGT\]).
In this paper we consider Cheeger’s inequality for the Dirichlet problem only since the case of closed manifolds or the Neumann problem is reduced to this by consideration of nodal domains.
In Section 2 we explain the Max Flow Min Cut Theorem, and in Section 3 we show how a classical inequality bounding the Cheeger constant of a plane domain in terms of its inradius may be understood in terms of vector fields.
## 2. Max Flow Min Cut Theorems
The classical Max Flow Min Cut Theorem deals with a discrete network, consisting of a finite set $`V`$ and a function $`c:V\times V[0,\mathrm{})`$. $`c(v,w)`$ may be considered as the capacity of a pipe connecting the ’nodes’ $`v,wV`$. Two nodes are distinguished, the source $`s`$ and the sink $`t`$.
We want to transport some liquid from $`s`$ to $`t`$. A (stationary, i.e. time independent) ’transport plan’ is modelled by a flow, i.e. a function $`f:V\times V[0,\mathrm{})`$ satisfying the capacity constraint $`f(v,w)c(v,w)v,wV`$ and the ’Kirchhoff law’
(7)
$$N_f(v):=\underset{w}{}[f(v,w)f(w,v)]=0vV\{s,t\},$$
that is, the total flow out of $`v`$ equals the total flow into $`v`$, except at the source and the sink. We let
$$\mathrm{value}(f):=N_f(s),$$
the net flow out of the source. (7) implies that this equals $`N_f(t)`$.
The question is how big the value of a flow can be, given the capacity constraint. A simple upper bound can be given by any cut, i.e. subset $`SV`$ containing $`s`$ but not $`t`$. Clearly, if we define
$$\mathrm{cap}(S):=\underset{vS,wVS}{}c(v,w),$$
the total capacity of pipes leaving $`S`$, then
(8)
$$\mathrm{value}(f)\mathrm{cap}(S)$$
for any flow $`f`$ and any cut $`S`$: Any net flow out ouf the source must reach the sink, so it must leave $`S`$ at some point.
###### Theorem 4 (Max Flow Min Cut Theorem, \[EliFeiSha:NMFTN\], \[ForFul:MFTN\]).
In any network there is a flow $`f_{\text{max}}`$ and a cut $`S_{\text{min}}`$ satisfying $`\mathrm{value}(f_{\text{max}})=\mathrm{cap}(S_{\text{min}})`$.
Various generalizations of Theorem 4 have been proposed. First, one may allow several sources and several sinks of prescribed relative ’strengths’. As an example analogous to our continuous setup below, all but one nodes could be sources, of equal strength, the remaining node being the sink $`t`$. Then condition (7) is empty. If we let $`\mathrm{value}(f)`$ be the minimum net flow out of any source node then (8) becomes $`\mathrm{value}(f)\mathrm{cap}(S)/|S|`$ for any set $`S`$ not containing $`t`$, and the corresponding Max Flow Min Cut theorem holds again.
More challenging are generalizations to infinite sets $`V`$, modelling continuously distributed sources or sinks, for example. Measure theoretic versions were proved in \[ChaHar:DTCAFFFN\] and \[FucLus:CC\]. Here we are more interested in a geometric model. Several, slightly different, models were proposed in \[Iri:TFCAFN\], \[Str:MFTD\], \[TagIri:CADNAAURN\] and later unified and generalized by Nozawa \[Noz:MFMCTAN\]. (I recommend \[Str:MFTD\] for enjoyable reading.)
We explain the relation of network flows to Theorem 3, which is a special case of Nozawa’s general Max Flow Min Cut Theorem: The domain $`\mathrm{\Omega }^n`$ is the network. A (stationary) flow is modelled by a vector field $`V`$, as is common in continuum fluid dynamics. The capacity constraint is (4). Sources are distributed uniformly over $`\mathrm{\Omega }`$, and (5) states that they produce liquid at a rate $`h`$, at least. The complement (or boundary) of $`\mathrm{\Omega }`$ should be considered as the sink (one should think of a single sink, i.e. collapse the boundary to a point; in this way one does not need to prescribe the relative strengths of the sinks along the boundary). A cut is a subset $`S\mathrm{\Omega }`$, and (6) states the obvious fact that anything that is produced within $`S`$ must leave $`S`$ through its boundary, which yields the bound in Proposition 1.
The discrete Max Flow Min Cut Theorem is proved, in most texts on discrete optimization (see \[KorVyg:CO\], for example), by inductively constructing $`f_{\text{max}}`$. However, as is already remarked in \[ForFul:FN\], the theorem is also an instance of the very general duality principle in convex optimization, and this approach also yields the generalizations mentioned above.
The duality principle associates to our optimization problem (maximize $`h=inf_{x\mathrm{\Omega }}\mathrm{div}V(x)`$, subject to the constraint (4)) a dual problem, which turns out to be:
(9)
$$\text{Minimize }Q(\varphi ):=\frac{\varphi _{BV}}{\varphi _{L^1}},\text{ subject to }\varphi 0.$$
Here, $`\varphi `$ is a function of bounded variation on $`^n`$ which vanishes outside $`\mathrm{\Omega }`$, $`\varphi _{BV}`$ is its total variation (in $`^n`$), which equals $`_\mathrm{\Omega }|\varphi |`$ in case $`\varphi `$ is smooth, and $`\varphi _{L^1}=_\mathrm{\Omega }\varphi `$. The general duality theorem says that $`sup_Vinf_x(\mathrm{div}V(x))=inf_\varphi Q(\varphi )`$. To see the relation to Cheeger’s constant, first recall (see \[EvaGar:MTFPF\], Chapter 5) that the perimeter of $`S^n`$ is $`|S|=\chi _S_{BV}`$, if this is finite, where $`\chi _S`$ is the characteristic function of $`S`$. Therefore, $`Q(\chi _S)=|S|/|S|`$ if $`S\mathrm{\Omega }`$. Next, the coarea formula (loc. cit.) states that for $`\varphi 0`$ of bounded variation, the sets $`\{\varphi >t\}`$ have finite perimeter for almost all $`t`$, and
(10)
$$\varphi _{BV}=_0^{\mathrm{}}\chi _{\{\varphi >t\}}_{BV}𝑑t.$$
Since $`\varphi =_0^{\mathrm{}}\chi _{\{\varphi >t\}}𝑑t`$, one also has $`\varphi _{L^1}=_0^{\mathrm{}}\chi _{\{\varphi >t\}}_{L^1}𝑑t`$, and therefore $`Q(\varphi )inf_tQ(\chi _{\{\varphi >t\}})`$. This shows that in (9) one may restrict $`\varphi `$ to characteristic functions, so the infimum is precisely Cheeger’s constant.
## 3. Cheeger’s constant and inradius
In this section let $`\mathrm{\Omega }`$ be a simply connected plane domain, and let $`\rho _\mathrm{\Omega }`$ denote its inradius. Also, define the ’reduced inradius’
(11)
$$\stackrel{~}{\rho }_\mathrm{\Omega }:=\frac{\rho _\mathrm{\Omega }}{1+\pi \rho _\mathrm{\Omega }^2/|\mathrm{\Omega }|}.$$
Clearly, $`\rho _\mathrm{\Omega }/2<\stackrel{~}{\rho }_\mathrm{\Omega }<\rho _\mathrm{\Omega }`$. A well-known lower bound for $`\lambda _\mathrm{\Omega }`$ is
(12)
$$\lambda _\mathrm{\Omega }\frac{1}{4\stackrel{~}{\rho }_\mathrm{\Omega }^2}.$$
The weaker estimate $`\lambda _\mathrm{\Omega }>1/4\rho _\mathrm{\Omega }^2`$ is sometimes called Hayman’s inequality or Osserman’s inequality, since it was proved by Hayman (with 4 replaced by 900) \[Hay:SBPF\] and Osserman \[Oss:NHTBND\], but in fact it was first proved by E. Makai \[Mak:LEPFSCM\].<sup>1</sup><sup>1</sup>1The constant 1/4 in $`\lambda _\mathrm{\Omega }>1/4\rho _\mathrm{\Omega }^2`$ has since been improved using ideas from probability and conformal mapping, the currently best value is 0.6197, see \[BanCar:BMFFD\]. There are similar estimates in the multiply connected case (to which the considerations below apply as well), but there is no direct higher dimensional generalization. See \[Bog:IISNEP\] for a generalization to a certain pseudo-Laplacian.
(12) follows from
(13)
$$h_\mathrm{\Omega }\frac{1}{\stackrel{~}{\rho }}_\mathrm{\Omega }$$
and Cheeger’s inequality. Note that (13) is sharp for the disk. (13) is implicit in \[Mak:LEPFSCM\] and \[Oss:NHTBND\], but does not seem to be stated explicitly in the literature. Let us give the proof along the lines of \[Mak:LEPFSCM\], \[Oss:NHTBND\]: Let
(14)
$$𝒮=\{S\mathrm{\Omega }:S\text{ open and simply connected},S\text{ smooth}\}.$$
Clearly,
(15)
$$h_\mathrm{\Omega }=\underset{S𝒮}{inf}\frac{|S|}{|S|}$$
since filling in all ’holes’ in an arbitrary $`S\mathrm{\Omega }`$, making it simply connected, increases $`|S|`$, decreases $`|S|`$, and results in a subset of $`\mathrm{\Omega }`$ (since $`\mathrm{\Omega }`$ is simply connected). Also,
(16)
$$S\mathrm{\Omega }\stackrel{~}{\rho }_S\stackrel{~}{\rho }_\mathrm{\Omega }.$$
To see this, first note that, for $`A,\rho >0`$, the function $`f_A(\rho )=\rho \rho /(1+\pi \rho ^2/A)`$ is increasing in $`\rho `$ for $`\pi \rho ^2A`$. Now $`|S||\mathrm{\Omega }|`$ yields $`\stackrel{~}{\rho }_S=f_{|S|}(\rho _S)f_{|\mathrm{\Omega }|}(\rho _S)`$; also $`\pi \rho _S^2\pi \rho _\mathrm{\Omega }|\mathrm{\Omega }|`$, so the monotonicity of $`f_{|\mathrm{\Omega }|}`$ gives (16). Now the main step is ’Bonnesen’s inequality’: For simply connected $`S^2`$
(17)
$$\rho _S|S||S|+\pi \rho _S^2.$$
This, together with (15) and (16) proves (13).<sup>2</sup><sup>2</sup>2Note that (17) is equivalent to $`|S|^24\pi |S|(|S|2\pi \rho _S)^2`$ and therefore a sharper version of the classical isoperimetric inequality $`|S|^24\pi |S|`$. (17) was proved by Bonnesen for convex $`\mathrm{\Omega }`$ \[Bon:UVIUKEKAMUKK\] and by Besicovitch \[Bes:VCIP\] for general simply connected domains, see also Sz.-Nagy \[SzNag:UPNEB\].
The question arises naturally whether one may prove (13) by constructing a vector field $`V`$ on $`\mathrm{\Omega }`$ satisfying (4), (5) with $`h=1/\stackrel{~}{\rho }_\mathrm{\Omega }`$. It is not clear how to do this. It seems more interesting to infer from Theorem 3 and (13):
###### Corollary 5.
Let $`\mathrm{\Omega }`$ be a simply connected plane domain with Lipschitz boundary and reduced inradius $`\stackrel{~}{\rho }_\mathrm{\Omega }`$ defined by (11). Then there is a vector field $`V`$ on $`\mathrm{\Omega }`$ satisfying (4), (5) with $`h=1/\stackrel{~}{\rho }_\mathrm{\Omega }`$.
Although there seems to be no natural, geometrically defined candidate for this vector field, we now proceed to show how certain geometric vector fields for subdomains of $`\mathrm{\Omega }`$ yield (13).
First, we need the following variant of Proposition 1:
###### Proposition 6.
Let $`\mathrm{\Omega }^n`$ be open, and let $`𝒮`$ be a class of Lipschitz subdomains of $`\mathrm{\Omega }`$ satisfying (15). Let $`h`$. Suppose that for each $`S𝒮`$ there is a vector field $`V_S𝒳_{\mathrm{BV}}(S)`$ on $`S`$ satisfying
(18) $`|V_S|`$ $``$ $`1,\text{pointwise on }|S|,`$
(19) $`{\displaystyle _S}\mathrm{div}V_S`$ $``$ $`h|S|.`$
Then $`h_\mathrm{\Omega }h`$.
###### Proof.
This follows from (6) applied to $`V_S`$.
Note that condition (19) is weaker than the pointwise condition (5) (if this is applied to $`V_S`$ on $`S`$). So in order to get effective lower bounds on $`h_\mathrm{\Omega }`$, one has more flexibility in choosing ’test’ vector fields, but one needs to do it for all $`S𝒮`$ simultaneously.
As before, we define $`𝒮`$ by (14). Using (16), we then see that (13) follows from Proposition 6 and the following:
###### Proposition 7.
Let $`S^2`$ be a smooth, simply connected domain, of inradius $`\rho _S`$ and reduced inradius $`\stackrel{~}{\rho }_S`$. Then there is a vector field $`V𝒳_{\mathrm{BV}}(S)`$ on $`S`$ satisfying
(20) $`|V|`$ $``$ $`1,`$
(21) $`{\displaystyle _S}\mathrm{div}V`$ $``$ $`{\displaystyle \frac{|S|}{\stackrel{~}{\rho }_S}}={\displaystyle \frac{|S|}{\rho _S}}+\pi \rho _S.`$
###### Proof.
By scale invariance we may assume $`\rho _S=1`$. Let $`\varphi (x)`$ denote the distance of $`xS`$ to $`S`$, and define
(22)
$$V=\frac{1}{2}(1\varphi )^2=(1\varphi )\varphi .$$
This is motivated by the case of the disk, Example 2.
We will use the following facts about the distance function:
1. $`\varphi L^{\mathrm{}}(S,^n)`$, and $`|\varphi |=1`$ almost everywhere.
2. $`\varphi `$ has bounded variation.
3. For almost all $`t[0,1]`$ the level set $`\varphi ^1(t)`$ consists of a finite union of piecewise smooth, simple closed curves, with non-zero angles, and $`\{\varphi <t\}=S\varphi ^1(t).`$
(a) and (c) are proved, for example, in \[Har:GPCL\], by a detailed analysis of the function $`F:[0,L]\times [0,1]^2`$, with $`L`$ the length of $`S`$, defined by requiring that $`sF(s,0)`$ is an arclength parametrization of $`S`$ and $`tF(s,t)`$ is the unit speed normal to $`S`$ starting inward at $`F(s,0)`$. For higher dimensional generalizations of (c), stating that $`\varphi ^1(t)`$ is Lipschitz for a.e. $`t`$ and using the notion of Clarke gradient of $`\varphi `$, see \[ItoTan:STDF\], \[Rif:MSTDFRM\]. (b) is folklore<sup>3</sup><sup>3</sup>3Proof: Let $`x_0S`$ and $`a=\varphi (x_0)>0`$. By an easy calculation, $`x\varphi _y(x):=|xy|\frac{2}{a}|xx_0|^2`$ has negative definite Hessian at $`x_0`$ (if $`yS`$) and so is concave near $`x_0`$. Therefore, $`\stackrel{~}{\varphi }(x)=inf_{yS}\varphi _y(x)`$ is concave near $`x_0`$, so $`\stackrel{~}{\varphi }`$ has bounded variation near $`x_0`$ (see \[EvaGar:MTFPF\], Section 6.3, Theorem 3). Finally, $`\varphi (x)=\stackrel{~}{\varphi }(x)+\frac{2}{a}|xx_0|^2`$ shows that $`\varphi `$ has locally bounded variation in $`S`$, and since $`\varphi `$ is smooth near $`S`$ it has bounded variation..
Therefore, to prove Proposition 7 it remains to verify $`\mathrm{div}V|S|+\pi `$. Now $`\mathrm{div}V=|\varphi |^2(1\varphi )\mathrm{\Delta }\varphi `$, with $`\mathrm{\Delta }`$ the Laplace operator. Since $`|\varphi |=1`$ we need to show
(23)
$$_S(1\varphi )\mathrm{\Delta }\varphi \pi .$$
Now we have $`_S(1\varphi )\mathrm{\Delta }\varphi =_S(_{\varphi (x)}^1𝑑t)\mathrm{\Delta }\varphi =_0^1(_{\{\varphi <t\}}\mathrm{\Delta }\varphi )𝑑t`$, by applying Fubini’s theorem for measures (see \[EvaGar:MTFPF\], for example) to $`_U\mathrm{\Delta }\varphi 𝑑t`$, where $`U=\{(t,x):\varphi (x)<t\}[0,1]\times S`$. If $`t`$ is as in (c) above then, by the divergence theorem, $`_{\varphi <t}\mathrm{\Delta }\varphi =_{\{\varphi <t\}}(\varphi )𝑑n=L_tL_0`$, where $`L_t`$ denotes the length of $`\varphi ^1(t)`$. Finally, $`_0^1L_t𝑑t=|S|`$ (by the coarea formula (10), for example, using (a)), and therefore
$$_0^1(_{\{\varphi <t\}}\mathrm{\Delta }\varphi )𝑑t=_0^1(L_tL_0)𝑑t=|S||S|\pi $$
by (17), and this proves (23).
Note that the vector field (22) cannot be used directly in Proposition 1, since $`\mathrm{\Delta }\varphi >0`$ near concave parts of $`\mathrm{\Omega }`$, so that the pointwise estimate $`\mathrm{div}V1`$ is false.
Remark: If $`\varphi `$ was smooth everywhere, one could prove (23) without appealing to Bonnesen’s inequality (and, in effect, reprove this inequality), as follows: The general coarea formula gives
$$_S(1\varphi )\mathrm{\Delta }\varphi =_0^1(_{\varphi ^1(t)}\frac{1\varphi }{|\varphi |}\mathrm{\Delta }\varphi )𝑑t=_0^1(1t)(_{\varphi ^1(t)}\mathrm{\Delta }\varphi )𝑑t,$$
where the line integrals are with respect to arclength measure. It is elementary to see that $`\mathrm{\Delta }\varphi `$ equals minus the curvature of $`\varphi ^1(t)`$ (whereever $`\varphi `$ is smooth). Also, the integral of the curvature along a smooth simple closed curve equals $`2\pi `$. So if almost all level sets were smooth (instead of piecewise smooth) we would obtain (23).
The problem with this ’proof’ is that, typically, $`\mathrm{\Delta }\varphi `$ is not a function but a measure (so that the coarea formula is not applicable), and a positive measure of level sets may be non-smooth. Consider, for example, a rectangle: $`\varphi `$ has a jump at its ’center line’, leading to a $`\delta `$ type singularity of $`\mathrm{\Delta }\varphi `$ there.
It should be possible and would be interesting to find generalizations of the coarea formula (using suitable transversality hypotheses) and of the curvature argument that make this proof work.
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# Generalizations of Quantum Mechanics
## 1 Introduction
According to the so-called “Copenhagen Interpretation,” standard quantum theory is limited to describing experimental situations. It is at once remarkably successful in its predictions, and remarkably ill-defined in its conceptual structure: what is an experiment? what physical objects do or do not require quantization? how are the states realized in nature to be characterized? how and when is the wave-function “collapse postulate” to be invoked? Because of its success, one may suspect that quantum theory can be promoted from a theory of measurement to a theory of reality. But, that requires there to be an unambiguous specification (S) of the possible real states of Nature and their probabilities of being realized.
There are several approaches that attempt to achieve S. The more conservative approaches (consistent histories, environmental decoherence, many worlds) do not produce any predictions that differ from the standard ones because they do not tamper with the usual basic mathematical formalism. Rather, they utilize structures compatible with standard quantum theory to elucidate S. These approaches, which will not be discussed in this article, have arguably been less successful so far at achieving S than approaches that introduce significant alterations to quantum theory.
This article will largely deal with the two most well-developed realistic models that reproduce quantum theory in some limit and yield potentially new and testable physics outside that limit. First, the pilot-wave model, which will be discussed in the broader context of “hidden-variables theories.” Second, the continuous spontaneous localization (CSL) model, which describes wave function collapse as a physical process. Other related models will also be discussed briefly.
Due to bibliographic space limitations, this article contains a number of uncited references, of the form “(author) in (year).” Those in the next section can be found in Valentini (2002b, 2004a,b) or at www.arxiv.org. Those in subsequent sections can be found in Adler (2004), Bassi and Ghirardi (2003), Pearle (1999) (or in subsequent papers by these authors, or directly, at www.arxiv.org), and in Wallstrom (1994).
## 2 Hidden Variables and Quantum Nonequilibrium
A deterministic hidden-variables theory defines a mapping $`\omega =\omega (M,\lambda )`$ from initial hidden parameters $`\lambda `$ (defined, e.g., at the time of preparation of a quantum state) to final outcomes $`\omega `$ of quantum measurements. The mapping depends on macroscopic experimental settings $`M`$, and fixes the outcome for each run of the experiment. Bell’s theorem of 1964 shows that, for entangled quantum states of widely separated systems, the mapping must be nonlocal: some outcomes for (at least) one system must depend on the setting for another distant system.
In a viable theory, the statistics of quantum measurement outcomes – over an ensemble of experimental trials with fixed settings $`M`$ – will agree with quantum theory for some special distribution $`\rho _{\mathrm{QT}}(\lambda )`$ of hidden variables. For example, expectation values will coincide with the predictions of the Born rule
$$\omega _{\mathrm{QT}}𝑑\lambda \rho _{\mathrm{QT}}(\lambda )\omega (M,\lambda )=\mathrm{Tr}(\widehat{\rho }\widehat{\mathrm{\Omega }})$$
for an appropriate density operator $`\widehat{\rho }`$ and Hermitian observable $`\widehat{\mathrm{\Omega }}`$. (As is customary in this context, $`𝑑\lambda `$ is to be understood as a generalised sum.) However, given the mapping $`\omega =\omega (M,\lambda )`$ for individual trials, one may, in principle, consider nonstandard distributions $`\rho (\lambda )\rho _{\mathrm{QT}}(\lambda )`$ that yield statistics outside the domain of ordinary quantum theory (Valentini 1991, 2002a). We may say that such distributions correspond to a state of quantum nonequilibrium.
Quantum nonequilibrium is characterised by the breakdown of a number of basic quantum constraints. In particular, nonlocal signals appear at the statistical level. We shall first illustrate this for the hidden-variables model of de Broglie and Bohm. Then we shall generalize the discussion to all (deterministic) hidden-variables theories.
At present there is no experimental evidence for quantum nonequilibrium in nature. However, from a hidden-variables perspective, it is natural to explore the theoretical properties of nonequilibrium distributions, and to search experimentally for the statistical anomalies associated with them.
From this point of view, quantum theory is a special case of a wider physics, much as thermal physics is a special case of a wider (nonequilibrium) physics. (The special distribution $`\rho _{\mathrm{QT}}(\lambda )`$ is analogous to, say, Maxwell’s distribution of molecular speeds.) Quantum physics may be compared with the physics of global thermal equilibrium, which is characterised by constraints – such as the impossibility of converting heat into work (in the absence of temperature differences) – that are not fundamental but contingent on the state. Similarly, quantum constraints such as statistical locality (the impossibility of converting entanglement into a practical signal) are seen as contingencies of $`\rho _{\mathrm{QT}}(\lambda )`$.
### 2.1 Pilot-Wave Theory
The de Broglie-Bohm “pilot-wave theory” – as it was originally called by de Broglie, who first presented it at the Fifth Solvay Congress in 1927 – is the classic example of a deterministic hidden-variables theory of broad scope (Bohm 1952; Bell 1987; Holland 1993). We shall use it to illustrate the above ideas. Later, the discussion will be generalised to arbitrary theories.
In pilot-wave dynamics, an individual closed system with (configuration-space) wave function $`\mathrm{\Psi }(X,t)`$ satisfying the Schrödinger equation
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=\widehat{H}\mathrm{\Psi }$$
(1)
has an actual configuration $`X(t)`$ with velocity
$$\dot{X}(t)=\frac{J(X,t)}{\left|\mathrm{\Psi }(X,t)\right|^2}$$
(2)
where $`J=J\left[\mathrm{\Psi }\right]=J(X,t)`$ satisfies the continuity equation
$$\frac{\left|\mathrm{\Psi }\right|^2}{t}+J=0$$
(3)
(which follows from (1)). In quantum theory, $`J`$ is the “probability current”. In pilot-wave theory, $`\mathrm{\Psi }`$ is an objective physical field (on configuration space) guiding the motion of an individual system.
Here, the objective state (or ontology) for a closed system is given by $`\mathrm{\Psi }`$ and $`X`$. A probability distribution for $`X`$ – discussed below – completes an unambiguous specification S (as mentioned in the introduction).
Pilot-wave dynamics may be applied to any quantum system with a locally conserved current in configuration space. Thus, $`X`$ may represent a many-body system, or the configuration of a continuous field, or perhaps some other entity.
For example, at low energies, for a system of $`N`$ particles with positions $`𝐱_i(t)`$ and masses $`m_i`$ ($`i=1,2,\mathrm{}.,N)`$, with an external potential $`V`$, (1) (with $`X(𝐱_1,𝐱_2,\mathrm{}.,𝐱_N)`$) reads
$$i\mathrm{}\frac{\mathrm{\Psi }}{t}=\underset{i=1}{\overset{N}{}}\frac{\mathrm{}^2}{2m_i}_i^2\mathrm{\Psi }+V\mathrm{\Psi }$$
(4)
while (2) has components
$$\frac{d𝐱_i}{dt}=\frac{\mathrm{}}{m_i}\text{Im}\left(\frac{_i\mathrm{\Psi }}{\mathrm{\Psi }}\right)=\frac{_iS}{m_i}$$
(5)
(where $`\mathrm{\Psi }=\left|\mathrm{\Psi }\right|e^{(i/\mathrm{})S}`$).
In general, (1) and (2) determine $`X(t)`$ for an individual system, given the initial condtions $`X(0)`$, $`\mathrm{\Psi }(X,0)`$ at $`t=0`$. For an arbitrary initial distribution $`P(X,0)`$, over an ensemble with the same wave function $`\mathrm{\Psi }(X,0)`$, the evolution $`P(X,t)`$ of the distribution is given by the continuity equation
$$\frac{P}{t}+(P\dot{X})=0$$
(6)
The outcome of an experiment is determined by $`X(0)`$, $`\mathrm{\Psi }(X,0)`$, which may be identified with $`\lambda `$. For an ensemble with the same $`\mathrm{\Psi }(X,0)`$, we have $`\lambda =`$ $`X(0)`$.
#### 2.1.1 Quantum Equilibrium
From (3) and (6), if we assume $`P(X,0)=\left|\mathrm{\Psi }(X,0)\right|^2`$ at $`t=0`$, we obtain $`P(X,t)=\left|\mathrm{\Psi }(X,t)\right|^2`$ – the Born-rule distribution of configurations – at all times $`t`$.
Quantum measurements are, like any other process, described and explained in terms of evolving configurations. For measurement devices whose pointer readings reduce to configurations, the distribution of outcomes of quantum measurements will match the statistical predictions of quantum theory (Bohm 1952; Bell 1987; Dürr, Goldstein and Zanghì 2003). Thus quantum theory emerges phenomenologically for a “quantum equilibrium” ensemble with distribution $`P(X,t)=\left|\mathrm{\Psi }(X,t)\right|^2`$ (or $`\rho (\lambda )=\rho _{\mathrm{QT}}(\lambda )`$).
#### 2.1.2 Quantum Nonequilibrium
In principle, as we saw for general hidden-variables theories, we may consider a nonequilibrium distribution $`P(X,0)\left|\mathrm{\Psi }(X,0)\right|^2`$ of initial configurations while retaining the same deterministic dynamics (1), (2) for individual systems (Valentini 1991). The time evolution of $`P(X,t)`$ will be determined by (6).
As we shall see, in appropriate circumstances (with a sufficiently complicated velocity field $`\dot{X}`$), (6) generates relaxation $`P\left|\mathrm{\Psi }\right|^2`$ on a coarse-grained level, much as the analogous classical evolution on phase space generates thermal relaxation. But for as long as the ensemble is in nonequilibrium, the statistics of outcomes of quantum measurements will disagree with quantum theory.
Quantum nonequilibrium may have existed in the very early universe, with relaxation to equilibrium occurring soon after the big bang. Thus, a hidden-variables analogue of the classical thermodynamic “heat death of the universe” may have actually taken place (Valentini 1991). Even so, relic cosmological particles that decoupled sufficiently early could still be in nonequilibrium today, as suggested by Valentini in 1996 and 2001. It has also been speculated that nonequilibrium could be generated in systems entangled with degrees of freedom behind a black-hole event horizon (Valentini 2004a).
Experimental searches for nonequilibrium have been proposed. Nonequilibrium could be detected by the statistical analysis of random samples of particles taken from a parent population of (for example) relics from the early universe. Once the parent distribution is known, the rest of the population could be used as a resource, to perform tasks that are currently impossible (Valentini 2002b).
### 2.2 H-Theorem: Relaxation to Equilibrium
Before discussing the potential uses of nonequilibrium, we should first explain why all systems probed so far have been found in the equilibrium state $`P=\left|\mathrm{\Psi }\right|^2`$. This distribution may be accounted for along the lines of classical statistical mechanics, noting that all currently-accessible systems have had a long and violent astrophysical history.
Dividing configuration space into small cells, and introducing coarse-grained quantities $`\overline{P}`$, $`\overline{\left|\mathrm{\Psi }\right|^2}`$, a general argument for relaxation $`\overline{P}\overline{\left|\mathrm{\Psi }\right|^2}`$ is based on an analog of the classical coarse-graining H-theorem. The coarse-grained H-function
$$\overline{H}=𝑑X\overline{P}\mathrm{ln}(\overline{P}/\overline{\left|\mathrm{\Psi }\right|^2})$$
(7)
(minus the relative entropy of $`\overline{P}`$ with respect to $`\overline{\left|\mathrm{\Psi }\right|^2}`$) obeys the H-theorem (Valentini 1991)
$$\overline{H}(t)\overline{H}(0)$$
(assuming no initial fine-grained microstructure in $`P`$ and $`\left|\mathrm{\Psi }\right|^2`$). Here, $`\overline{H}0`$ for all $`\overline{P}`$, $`\overline{\left|\mathrm{\Psi }\right|^2}`$ and $`\overline{H}=0`$ if and only if $`\overline{P}=\overline{\left|\mathrm{\Psi }\right|^2}`$ everywhere.
The H-theorem expresses the fact that $`P`$ and $`\left|\mathrm{\Psi }\right|^2`$ behave like two “fluids” that are “stirred” by the same velocity field $`\dot{X}`$, so that $`P`$ and $`\left|\mathrm{\Psi }\right|^2`$ tend to become indistinguishable on a coarse-grained level. Like its classical analog, the theorem provides a general understanding of how equilibrium is approached, while not proving that equilibrium is actually reached. (And of course, for some simple systems – such as a particle in the ground state of a box, for which the velocity field $`S/m`$ vanishes – there is no relaxation at all.) A strict decrease of $`\overline{H}(t)`$ immediately after $`t=0`$ is guaranteed if $`\dot{X}_0(P_0/\left|\mathrm{\Psi }_0\right|^2)`$ has non-zero spatial variance over a coarse-graining cell, as shown by Valentini in 1992 and 2001.
A relaxation timescale $`\tau `$ may be defined by $`1/\tau ^2\left(d^2\overline{H}/dt^2\right)_0/\overline{H}_0`$. For a single particle with quantum energy spread $`\mathrm{\Delta }E`$, a crude estimate given by Valentini in 2001 yields $`\tau (1/\epsilon )\mathrm{}^2/m^{1/2}(\mathrm{\Delta }E)^{3/2}`$, where $`\epsilon `$ is the coarse-graining length. For wave functions that are superpositions of many energy eigenfunctions, the velocity field (generally) varies rapidly, and detailed numerical simulations (in two dimensions) show that relaxation occurs with an approximately exponential decay $`\overline{H}(t)\overline{H}_0e^{t/t_c}`$, with a time constant $`t_c`$ of order $`\tau `$ (Valentini and Westman 2005).
Equilibrium is then to be expected for particles emerging from the violence of the big bang. The possibility is still open that relics from very early times may not have reached equilibrium before decoupling.
### 2.3 Nonlocal Signaling
We now show how nonequilibrium, if it were ever discovered, could be used for nonlocal signaling.
Pilot-wave dynamics is nonlocal. For a pair of particles $`A`$, $`B`$ with entangled wave function $`\mathrm{\Psi }(𝐱_A,𝐱_B,t)`$, the velocity $`\dot{𝐱}_A(t)=_AS(𝐱_A,𝐱_B,t)/m_A`$ of $`A`$ depends instantaneously on $`𝐱_B`$, and local operations at $`B`$ – such as switching on a potential – instantaneously affect the motion of $`A`$. For an ensemble $`P(𝐱_A,𝐱_B,t)=|\mathrm{\Psi }(𝐱_A,𝐱_B,t)|^2`$, local operations at $`B`$ have no statistical effect at $`A`$: the individual nonlocal effects vanish upon averaging over an equilibrium ensemble.
Nonlocality is (generally) hidden by statistical noise only in quantum equilibrium. If instead $`P(𝐱_A,𝐱_B,0)|\mathrm{\Psi }(𝐱_A,𝐱_B,0)|^2`$, a local change in the Hamiltonian at $`B`$ generally induces an instantaneous change in the marginal $`p_A(𝐱_A,t)d^3𝐱_BP(𝐱_A,𝐱_B,t)`$ at $`A`$. For example, in one dimension a sudden change $`\widehat{H}_B\widehat{H}_B^{}`$ in the Hamiltonian at $`B`$ induces a change $`\mathrm{\Delta }p_Ap_A(x_A,t)p_A(x_A,0)`$ (for small $`t`$) (Valentini 1991),
$$\mathrm{\Delta }p_A=\frac{t^2}{4m}\frac{}{x_A}\left(a(x_A)𝑑x_Bb(x_B)\frac{P(x_A,x_B,0)|\mathrm{\Psi }(x_A,x_B,0)|^2}{|\mathrm{\Psi }(x_A,x_B,0)|^2}\right)$$
(8)
(Here $`m_A=m_B=m`$, $`a(x_A)`$ depends on $`\mathrm{\Psi }(x_A,x_B,0)`$, while $`b(x_B)`$ also depends on $`\widehat{H}_B^{}`$ and vanishes if $`\widehat{H}_B^{}=\widehat{H}_B`$.) The signal is generally nonzero if $`P_0|\mathrm{\Psi }_0|^2`$.
Nonlocal signals do not lead to causal paradoxes if, at the hidden-variable level, there is a preferred foliation of spacetime with a time parameter that defines a fundamental causal sequence. Such signals, if they were observed, would define an absolute simultaneity as discussed by Valentini in 1992 and 2005. Note that in pilot-wave field theory, Lorentz invariance emerges as a phenomenological symmetry of the equilibrium state, conditional on the structure of the field-theoretical Hamiltonian (as discussed by Bohm and Hiley in 1984, Bohm, Hiley and Kaloyerou in 1987, and Valentini in 1992 and 1996).
### 2.4 Subquantum Measurement
In principle, nonequilibrium particles could also be used to perform “subquantum measurements” on ordinary, equilibrium systems. We illustrate this with an exactly solvable one-dimensional model (Valentini 2002b).
Consider an apparatus “pointer” coordinate $`y`$, with known wave function $`g_0(y)`$ and known (ensemble) distribution $`\pi _0(y)\left|g_0(y)\right|^2`$, where $`\pi _0(y)`$ has been deduced by statistical analysis of random samples from a parent population with known wave function $`g_0(y)`$. (We assume that relaxation may be neglected: for example, if $`g_0`$ is a box ground state, $`\dot{y}=0`$ and $`\pi _0(y)`$ is static.) Consider also a “system” coordinate $`x`$ with known wave function $`\psi _0(x)`$ and known distribution $`\rho _0(x)=\left|\psi _0(x)\right|^2`$. If $`\pi _0(y)`$ is arbitrarily narrow, $`x_0`$ can be measured without disturbing $`\psi _0(x)`$, to arbitrary accuracy (violating the uncertainty principle).
To do this, at $`t=0`$ we switch on an interaction Hamiltonian $`\widehat{H}=a\widehat{x}\widehat{p}_y`$, where $`a`$ is a constant and $`p_y`$ is canonically conjugate to $`y`$. For relatively large $`a`$, we may neglect the Hamiltonians of $`x`$ and $`y`$. For $`\mathrm{\Psi }=\mathrm{\Psi }(x,y,t)`$, we then have $`\mathrm{\Psi }/t=ax\mathrm{\Psi }/y`$. For $`\left|\mathrm{\Psi }\right|^2`$ we have the continuity equation $`\left|\mathrm{\Psi }\right|^2/t=ax\left|\mathrm{\Psi }\right|^2/y`$, which implies the hidden-variable velocity fields $`\dot{x}=0,\dot{y}=ax`$ and trajectories $`x(t)=x_0,y(t)=y_0+ax_0t`$.
The initial product $`\mathrm{\Psi }_0(x,y)=\psi _0(x)g_0(y)`$ evolves into $`\mathrm{\Psi }(x,y,t)=\psi _0(x)g_0(yaxt)`$. For $`at0`$ (with $`a`$ large but fixed), $`\mathrm{\Psi }(x,y,t)\psi _0(x)g_0(y)`$ and $`\psi _0(x)`$ is undisturbed: for small $`at`$, a standard quantum pointer with the coordinate $`y`$ would yield negligible information about $`x_0`$. Yet, for arbitrarily small $`at`$, the hidden-variable pointer coordinate $`y(t)=y_0+ax_0t`$ does contain complete information about $`x_0`$ (and $`x(t)=x_0`$). This “subquantum” information will be visible to us if $`\pi _0(y)`$ is sufficiently narrow.
For, over an ensemble of similar experiments, with initial joint distribution $`P_0(x,y)=\left|\psi _0(x)\right|^2\pi _0(y)`$ (equilibrium for $`x`$ and nonequilibrium for $`y`$), the continuity equation $`P/t=axP/y`$ implies that $`P(x,y,t)=\left|\psi _0(x)\right|^2\pi _0(yaxt)`$. If $`\pi _0(y)`$ is localised around $`y=0`$ ($`\pi _0(y)=0`$ for $`\left|y\right|>w/2`$), then a standard (faithful) measurement of $`y`$ with result $`y_{\mathrm{meas}}`$ will imply that $`x`$ lies in the interval $`(y_{\mathrm{meas}}/atw/2at,y_{\mathrm{meas}}/at+w/2at)`$ (so that $`P(x,y,t)0`$). Taking the simultaneous limits $`at0`$, $`w0`$, with $`w/at0`$, the midpoint $`y_{\mathrm{meas}}/atx_0`$ (since $`y_{\mathrm{meas}}=y_0+ax_0t`$ and $`\left|y_0\right|w/2`$), while the error $`w/2at0`$.
If $`w`$ is arbitrarily small, a sequence of such measurements will determine the hidden trajectory $`x(t)`$ without disturbing $`\psi (x,t)`$, to arbitrary accuracy.
### 2.5 Subquantum Information and Computation
From a hidden-variables perspective, immense physical resources are hidden from us by equilibrium statistical noise. Quantum nonequilibrium would probably be as useful technologically as thermal or chemical nonequilibrium.
#### 2.5.1 Distinguishing Nonorthogonal States
In quantum theory, nonorthogonal states $`|\psi _1`$, $`|\psi _2`$ ($`\psi _1|\psi _20`$) cannot be distinguished without disturbing them. This theorem breaks down in quantum nonequilibrium (Valentini 2002b). For example, if $`|\psi _1`$, $`|\psi _2`$ are distinct states of a single spinless particle, then the associated de Broglie-Bohm velocity fields will in general be different, even if $`\psi _1|\psi _20`$, and so will the hidden-variable trajectories. Subquantum measurement of the trajectories could then distinguish the states $`|\psi _1`$, $`|\psi _2`$.
#### 2.5.2 Breaking Quantum Cryptography
The security of standard protocols for quantum key distribution depends on the validity of the laws of quantum theory. These protocols would become insecure given the availability of nonequilibrium systems (Valentini 2002b).
The protocols known as BB84 and B92 depend on the impossibility of distinguishing nonorthogonal quantum states without disturbing them. An eavesdropper in possession of nonequilibrium particles could distinguish the nonorthogonal states being transmitted between two parties, and so read the supposedly secret key. Further, if subquantum measurements allow an eavesdropper to predict quantum measurement outcomes at each “wing” of a (bipartite) entangled state, then the EPR (Einstein-Podolsky-Rosen) protocol also becomes insecure.
#### 2.5.3 Subquantum Computation
It has been suggested that nonequilibrium physics would be computationally more powerful than quantum theory, because of the ability to distinguish nonorthogonal states (Valentini 2002b). However, this ability depends on the (less-than-quantum) dispersion $`w`$ of the nonequilibrium ensemble. A well-defined model of computational complexity requires that the resources be quantified in some way. Here, a key question is how the required $`w`$ scales with the size of the computational task. So far, no rigorous results are known.
### 2.6 Extension to All Deterministic Hidden-Variables Theories
Let us now discuss arbitrary (deterministic) theories.
#### 2.6.1 Nonlocal Signaling
Consider a pair of two-state quantum systems $`A`$ and $`B`$, which are widely separated and in the singlet state. Quantum measurements of observables $`\widehat{\sigma }_A𝐦_A\widehat{𝝈}_A`$, $`\widehat{\sigma }_B𝐦_B\widehat{𝝈}_B`$ (where $`𝐦_A`$, $`𝐦_B`$ are unit vectors in Bloch space and $`\widehat{𝝈}_A`$, $`\widehat{𝝈}_B`$ are Pauli spin operators) yield outcomes $`\sigma _A`$, $`\sigma _B=\pm 1`$, in the ratio $`1:1`$ at each wing, with a correlation $`\widehat{\sigma }_A\widehat{\sigma }_B=𝐦_A𝐦_B`$. Bell’s theorem shows that for a hidden-variables theory to reproduce this correlation – upon averaging over an equilibrium ensemble with distribution $`\rho _{\mathrm{QT}}(\lambda )`$ – it must take the nonlocal form
$$\sigma _A=\sigma _A(𝐦_A,𝐦_B,\lambda ),\sigma _B=\sigma _B(𝐦_A,𝐦_B,\lambda )$$
(9)
More precisely, to obtain $`\sigma _A\sigma _B_{\mathrm{QT}}=𝐦_A𝐦_B`$ (where $`\sigma _A\sigma _B_{\mathrm{QT}}𝑑\lambda \rho _{\mathrm{QT}}(\lambda )\sigma _A\sigma _B`$), at least one of $`\sigma _A`$, $`\sigma _B`$ must depend on the measurement setting at the distant wing. Without loss of generality, we assume that $`\sigma _A`$ depends on $`𝐦_B`$.
For an arbitrary nonequilibrium ensemble with distribution $`\rho (\lambda )\rho _{\mathrm{QT}}(\lambda )`$, in general $`\sigma _A\sigma _B𝑑\lambda \rho (\lambda )\sigma _A\sigma _B`$ differs from $`𝐦_A𝐦_B`$, and the outcomes $`\sigma _A`$, $`\sigma _B=\pm 1`$ occur in a ratio different from $`1:1`$. Further, a change of setting $`𝐦_B𝐦_B^{}`$ at $`B`$ will generally induce a change in the outcome statistics at $`A`$, yielding a nonlocal signal at the statistical level. To see this, note that, in a nonlocal theory, the “transition sets”
$`T_A(,+)`$ $`\left\{\lambda \right|\sigma _A(𝐦_A,𝐦_B,\lambda )=1,\sigma _A(𝐦_A,𝐦_B^{},\lambda )=+1\}`$
$`T_A(+,)`$ $`\left\{\lambda \right|\sigma _A(𝐦_A,𝐦_B,\lambda )=+1,\sigma _A(𝐦_A,𝐦_B^{},\lambda )=1\}`$
cannot be empty for arbitrary settings. Yet, in quantum equilibrium, the outcomes $`\sigma _A=\pm 1`$ occur in the ratio $`1:1`$ for all settings, so the transition sets must have equal equilibrium measure, $`\mu _{\mathrm{QT}}[T_A(,+)]=\mu _{\mathrm{QT}}[T_A(+,)]`$ ($`d\mu _{\mathrm{QT}}\rho _{\mathrm{QT}}(\lambda )d\lambda `$). That is, the fraction of the equilibrium ensemble making the transition $`\sigma _A=1\sigma _A=+1`$ under $`𝐦_B𝐦_B^{}`$ must equal the fraction making the reverse transition $`\sigma _A=+1\sigma _A=1`$. (This “detailed balancing” is analogous to the principle of detailed balance in statistical mechanics.) Since $`T_A(,+)`$, $`T_A(+,)`$ are fixed by the deterministic mapping, they are independent of the ensemble distribution $`\rho (\lambda )`$. Thus, for $`\rho (\lambda )\rho _{\mathrm{QT}}(\lambda )`$, in general $`\mu [T_A(,+)]\mu [T_A(+,)]`$ ($`d\mu \rho (\lambda )d\lambda `$): the fraction of the nonequilibrium ensemble making the transition $`\sigma _A=1\sigma _A=+1`$ will not in general balance the fraction making the reverse transition. The outcome ratio at $`A`$ will then change under $`𝐦_B𝐦_B^{}`$ and there will be an instantaneous signal at the statistical level from $`B`$ to $`A`$ (Valentini 2002a).
Thus, in any deterministic hidden-variables theory, nonequilibrium distributions $`\rho (\lambda )\rho _{\mathrm{QT}}(\lambda )`$ generally allow entanglement to be used for nonlocal signaling (just as, in ordinary statistical physics, differences of temperature make it possible to convert heat into work).
#### 2.6.2 Experimental Signature of Nonequilibrium
Quantum expectations are additive, $`c_1\widehat{\mathrm{\Omega }}_1+c_2\widehat{\mathrm{\Omega }}_2=c_1\widehat{\mathrm{\Omega }}_1+c_2\widehat{\mathrm{\Omega }}_2`$, even for noncommuting observables ($`[\widehat{\mathrm{\Omega }}_1,\widehat{\mathrm{\Omega }}_2]0`$, with $`c_1`$, $`c_2`$ real). As emphasised by Bell in 1966, this seemingly trivial consequence of the (linearity of the) Born rule $`\widehat{\mathrm{\Omega }}=\mathrm{Tr}(\widehat{\rho }\widehat{\mathrm{\Omega }})`$ is remarkable because it relates statistics from distinct, “incompatible” experiments. In nonequilibrium, such additivity generically breaks down (Valentini 2004b).
Further, for a two-state system with observables $`𝐦\widehat{𝝈}`$, the “dot-product” structure of the quantum expectation $`𝐦\widehat{𝝈}=\mathrm{Tr}\left(\widehat{\rho }𝐦\widehat{𝝈}\right)=𝐦𝐏`$ (for some Bloch vector $`𝐏`$) is equivalent to expectation additivity (Valentini 2004b). Nonadditive expectations then provide a convenient signature of nonequilibrium for any two-state system. For example, the sinusoidal modulation of the quantum transmission probability for a single photon through a polariser
$$p_{\mathrm{QT}}^+(\mathrm{\Theta })=\frac{1}{2}\left(1+𝐦\widehat{𝝈}\right)=\frac{1}{2}\left(1+P\mathrm{cos}2\mathrm{\Theta }\right)$$
(10)
(where an angle $`\theta `$ on the Bloch sphere corresponds to a physical angle $`\mathrm{\Theta }=\theta /2`$) will generically break down in nonequilibrium. Deviations from (10) would provide an unambiguous violation of quantum theory (Valentini 2004b).
Such deviations were searched for by Papaliolios in 1967, using laboratory photons and successive polarisation measurements over very short times, to test a hidden-variables theory (distinct from pilot-wave theory) due to Bohm and Bub (1966), in which quantum measurements generate nonequilibrium for short times. Experimentally, successive measurements over timescales $`10^{13}`$ $`\mathrm{sec}`$ agreed with the (quantum) sinusoidal modulation $`\mathrm{cos}^2\mathrm{\Theta }`$ to $`1\%`$. Similar tests might be performed with photons of a more exotic origin.
## 3 Continuous Spontaneous Localization Model (CSL)
The basic postulate of CSL is that the state vector $`|\psi ,t`$ represents reality. Since, for example, in describing a measurement, the usual Schrödinger evolution readily takes a real state into a non-real state, that is, into a superposition of real states (such as apparatus states describing different experimental outcomes), CSL requires a modification of Schrödinger’s evolution. To the Hamiltonian is added a term which depends upon a classical randomly fluctuating field $`w(𝐱,t)`$ and a mass-density operator $`\widehat{A}(𝐱,t)`$. This term acts to collapse a superposition of states, which differ in their spatial distribution of mass density, to one of these states. The rate of collapse is very slow for a superposition involving a few particles, but very fast for a superposition of macroscopically different states. Thus, very rapidly, what you see (in nature) is what you get (from the theory). Each state vector evolving under each $`w(𝐱,t)`$ corresponds to a realizable state, and a rule is given for how to associate a probability with each. In this way, an unambiguous specification S as mentioned in the introduction is achieved.
### 3.1 Requirements for Stochastic Collapse Dynamics
Consider a normalized state vector $`|\psi ,t=_n\alpha _n(t)|a_n`$ ($`a_n|a_n^{}=\delta _{nn^{}}`$) which undergoes a stochastic dynamical collapse process. This means that, starting from the initial superposition at $`t=0`$, for each run of the process, the squared amplitudes $`x_n(t)|\alpha _n(t)|^2`$ fluctuate until all but one vanish, that is, $`x_m(\mathrm{})=1`$, ($`x_m(\mathrm{})=0)`$ with probability $`x_m(0)`$.
This may be achieved simply, assuming negligible effect of the usual Schrödinger evolution, if the stochastic process enjoys the following properties (Pearle 1979):
$`{\displaystyle \underset{n}{}}x_n(t)`$ $`=1`$ (11a)
$`\overline{x_n(t)}`$ $`=x_n(0)`$ (11b)
$`\overline{x_n(\mathrm{})x_m(\mathrm{})}`$ $`=0\text{for}mn,`$ (11c)
where the overbar indicates the ensemble average at the indicated time. The only way that a sum of products of non-negative terms can vanish is for at least one term in each product to vanish. Thus, according to (11c), for each run, at least one of each pair {$`x_n(\mathrm{})`$, $`x_m(\mathrm{})`$} ($`nm`$) must vanish. This means that at most one $`x_n(\mathrm{})`$ might not vanish and, by (11a), applied at $`t=\mathrm{}`$, one $`x_n(\mathrm{})`$ must not vanish and, in fact, must equal 1: hence, each run produces collapse. Now, let the probability of the outcome {$`x_n(\mathrm{})=1`$, $`x_n(\mathrm{})=0`$} be denoted $`P_n`$. Since $`\overline{x_n(\mathrm{})}=1P_n+_{mn}0P_m=P_n`$ then, according to the Martingale property (11b), applied at $`t=\mathrm{}`$, $`P_n=x_n(0)`$: hence the ensemble of runs produces the probability postulated by the usual “collapse rule” of standard quantum theory.
A (nonquantum) stochastic process which obeys these equations is the gambler’s ruin game. Suppose one gambler initially possesses the fraction $`x_1(0)`$ of their joint wealth, and the other has the fraction $`x_2(0)`$. They toss a coin: heads, a dollar goes from gambler 1 to gambler 2, tails the dollar goes the other way. (11a) is satisfied since the sum of money in the game remains constant, (11b) holds because it is a fair game, and (11c) holds because each game eventually ends. Thus, gambler $`i`$ wins all the money with probability $`x_i(0)`$.
### 3.2 CSL in Essence
Consider the (non-unitary) Schrödinger picture evolution equation
$$|\psi ,t_w=𝒯e^{_0^t𝑑t^{}\{i\widehat{H}+(4\lambda )^1[w(t^{})2\lambda \widehat{A}]^2\}}|\psi ,0,$$
(12)
where $`\widehat{H}`$ is the usual Hamiltonian, $`w(t^{})`$ is an arbitrary function of white noise class, $`\widehat{A}`$ is a Hermitian operator ($`\widehat{A}|a_n=a_n|a_n`$), $`\lambda `$ is a collapse rate parameter, $`𝒯`$ is the time-ordering operator and $`\mathrm{}=1`$. Associated with this, the probability rule
$$P_t(w)Dw_w\psi ,t|\psi ,t_w\underset{j=0}{\overset{t/dt}{}}dw(t_j)/(2\pi \lambda /dt)^{1/2}$$
(13)
is defined, which gives the probability that nature chooses a noise which lies in the range $`\{w(t^{}),w(t^{})+dw(t^{})\}`$ (for calculational purposes, time is discretized, with $`t_0=0`$).
Equations (12) and (13) contain the essential features of CSL, and are all that is needed to discuss the simplest collapse behavior. Set $`\widehat{H}=0`$, so there is no competition between collapse and the usual Schrödinger evolution, and let the initial statevector be $`|\psi ,0=_n\alpha _n|a_n`$. Equations (12) and (13) become
$`|\psi ,t_w`$ $`=`$ $`{\displaystyle \underset{n}{}}\alpha _n|a_ne^{(4\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda a_n]^2}`$ (14a)
$`P_t(w)`$ $`=`$ $`{\displaystyle \underset{n}{}}|\alpha _n|^2e^{(2\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda a_n]^2}.`$ (14b)
When the un-normalized state vector in (14a) is divided by $`P_t^{1/2}(w)`$ and so normalized, the squared amplitudes are
$$x_n(t)=|\alpha _n|^2e^{(2\lambda )^1_0^t𝑑t^{}[w(t^{})2\lambda a_n]^2}/P_t(w),$$
which are readily shown to satisfy (11a), (11b), and (11c) in the form $`\overline{x_n^{1/2}(\mathrm{})x_m^{1/2}(\mathrm{})}=0`$ ($`mn`$) (which does not change the argument in section 3.1, but makes for an easier calculation). Thus, (14a), (14b) describe collapse dynamics.
To describe collapse to a joint eigenstate of a set of mutually commuting operators $`\widehat{A}^r`$, replace $`(4\lambda )^1[w(t^{})2\lambda \widehat{A}]^2`$ in the exponent of (12) by $`_r(4\lambda )^1[w^r(t^{})2\lambda \widehat{A}^r]^2`$. The interaction picture state vector in this case is (12) multiplied by $`\mathrm{exp}(i\widehat{H}t)`$:
$$|\psi ,t_w=𝒯e^{(4\lambda )^1_0^t𝑑t^{}_r[w^r(t^{})2\lambda \widehat{A}^r(t^{})]^2}|\psi ,0,$$
(15)
where $`\widehat{A}^r(t^{})\mathrm{exp}(i\widehat{H}t^{})\widehat{A}^r\mathrm{exp}(i\widehat{H}t^{})`$. The density matrix follows from (15), (13):
$$\widehat{\rho }(t)P_t(w)Dw|\psi ,t_w_w\psi ,t|/P_t(w)=𝒯e^{\lambda /2_0^t𝑑t^{}_r[\widehat{A}_L^r(t^{})\widehat{A}_R^r(t^{})]^2}\widehat{\rho }(0)$$
(16)
where $`\widehat{A}_L^r(t^{})`$ ($`\widehat{A}_R^r(t^{})`$) appears to the left (right) of $`\widehat{\rho }(0)`$, and is time-ordered (time reverse-ordered). In the example described by (14), the density matrix (16) is
$$\widehat{\rho }(t)=\underset{n,m}{}e^{(\lambda t/2)(a_na_m)^2}\alpha _n\alpha _m^{}|a_na_m|,$$
which encapsulates the ensemble’s collapse behavior.
### 3.3 CSL
The CSL proposal (Pearle 1989) is that collapse is engendered by distinctions between states at each point of space, so the index $`r`$ of $`\widehat{A}^r`$ in (15) becomes $`𝐱`$,
$$|\psi ,t_w=𝒯e^{(4\lambda )^1_0^t{\scriptscriptstyle 𝑑t^{}𝑑𝐱^{}[w(𝐱^{},t^{})2\lambda \widehat{A}(𝐱^{},t^{})]^2}}|\psi ,0,$$
(17)
and the distinction looked at is mass density. However, one cannot make the choice $`\widehat{A}(𝐱,0)=\widehat{M}(𝐱)`$, where $`\widehat{M}(𝐱)=_im_i\widehat{\xi }_i^{}(𝐱)\widehat{\xi }_i(𝐱)`$ is the mass density operator ($`m_i`$ is the mass of the $`i`$th type of particle, so $`m_e`$, $`m_p`$, $`m_n`$… are the masses of electrons, protons, neutrons…, and $`\widehat{\xi }_i^{}(𝐱)`$ is the creation operator for such a particle at location $`𝐱`$), because this entails an infinite rate of energy increase of particles ((23) with $`a=0`$). Instead, adapting a “gaussian smearing” idea from the Ghirardi et al. (1986) spontaneous localization (SL) model (see section 3.6), choose $`\widehat{A}^𝐱`$ as, essentially, proportional to the mass in a sphere of radius $`a`$ about $`𝐱`$:
$$\widehat{A}(𝐱,𝐭)e^{i\widehat{H}t}\frac{1}{(\pi a^2)^{3/4}}𝑑𝐳\frac{\widehat{M}(𝐳)}{m_p}e^{(2a^2)^1(𝐱𝐳)^2}e^{i\widehat{H}t}$$
(18)
The parameter value choices of SL, $`\lambda 10^{16}`$sec<sup>-1</sup> (according to (17), the collapse rate for protons) and $`a10^5`$cm are, so far, consistent with experiment (see section 3.4), and will be adopted here.
The density matrix associated with (17) is, as in (16),
$$\widehat{\rho }(t)=𝒯e^{(\lambda /2)_0^t𝑑t^{}𝑑𝐱^{}[\widehat{A}_L(𝐱^{},t^{})\widehat{A}_R(𝐱^{},t^{})]^2}\widehat{\rho }(0),$$
(19)
which satisfies the differential equation
$$\frac{d\widehat{\rho }(t)}{dt}=\frac{\lambda }{2}𝑑𝐱^{}[\widehat{A}(𝐱^{},t),[\widehat{A}(𝐱^{},t),\widehat{\rho }(t)]]$$
(20)
of Lindblad-Kossakowski form.
### 3.4 Consequences of CSL
Since the state vector dynamics of CSL is different from that of standard quantum theory, there are phenomena for which the two make different predictions, allowing for experimental tests. Consider an $`N`$-particle system with position operators $`\widehat{X}_i`$ ($`\widehat{X}_i|𝐱=x_i|𝐱`$). Substitution of $`\widehat{A}(𝐱^{})`$ from (18) in the Schrödinger picture version of (20), integration over $`𝐱^{}`$, and utilization of
$$f(𝐳)\widehat{M}(𝐳)|𝐱=\underset{i=1}{\overset{N}{}}m_if(\widehat{𝐗}_i)\delta (𝐳\widehat{𝐗}_i)|𝐱$$
results in
$`{\displaystyle \frac{d\widehat{\rho }(t)}{dt}}`$ $`=`$ $`i[\widehat{\rho }(t),\widehat{H}]{\displaystyle \frac{\lambda }{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{m_i}{m_p}}{\displaystyle \frac{m_j}{m_p}}`$ (21)
$`\times \left[e^{(4a^2)^1(\widehat{𝐗}_{Li}\widehat{𝐗}_{Lj})^2}+e^{(4a^2)^1(\widehat{𝐗}_{Ri}\widehat{𝐗}_{Rj})^2}2e^{(4a^2)^1(\widehat{𝐗}_{Li}\widehat{𝐗}_{Rj})^2}\right]\widehat{\rho }(t)`$
which is a useful form for calculations first suggested by Pearle and Squires in 1994.
#### 3.4.1 Interference
Consider the collapse rate of an initial state $`|\varphi =\alpha _1|1+\alpha _2|2`$, where $`|1`$, $`|2`$ describe a clump of matter, of size $`<<a`$, at different locations with separation $`>>a`$. Electrons may be neglected because of their small collapse rate compared to the much more massive nucleons, and the nucleon mass difference may be neglected. In using (21) to calculate $`d1|\widehat{\rho }(t)|2/dt`$, since $`\mathrm{exp}[(4a^2)^1(\widehat{𝐗}_i\widehat{𝐗}_j)^2]1`$ when acting on state $`|1`$ or $`|2`$, and $`0`$ when $`\widehat{𝐗}_i`$ acts on $`|1`$ and $`\widehat{𝐗}_j`$ acts on $`|2`$, (21) yields, for $`N`$ nucleons, the collapse rate $`\lambda N^2`$:
$$\frac{d1|\widehat{\rho }(t)|2}{dt}=i1|[\widehat{\rho }(t),\widehat{H}]|2\lambda N^21|\widehat{\rho }(t)|2.$$
(22)
If the clump undergoes a two-slit interference experiment, where the size and separation conditions above are satisfied for $`\mathrm{\Delta }T`$sec, and if the result agrees with the standard quantum theory prediction to 1%, it also agrees with CSL provided $`\lambda ^1>100N^2\mathrm{\Delta }T`$. So far, interference experiments with $`N`$ as large as $`10^3`$ have been performed by Nairz, Arndt and Zeilinger in 2000. The SL value of $`\lambda `$ would be testable, that is, the quantum predicted interference pattern would be “washed out” to 1% accuracy, if the clump were an $`10^6`$cm radius sphere of mercury, which contains $`N10^8`$ nucleons, interfered for $`\mathrm{\Delta }T=.01`$sec. Currently envisioned but not yet performed experiments (e.g., by Marshall, Simon, Penrose and Bouwmester in 2003) have been analyzed (e.g., by Bassi, Ippoliti and Adler in 2004 and by Adler in 2005) involving a superposition of a larger clump of matter in slightly displaced positions, entangled with a photon whose interference pattern is measured: these proposed experiments are still too crude to detect the SL value of $`\lambda `$, or the gravitationally-based collapse rate proposed by Penrose in 1996 (see section 4 and papers by Christian in 1999 and 2005).
#### 3.4.2 Bound State Excitation
Collapse narrows wave packets, thereby imparting energy to particles. If $`\widehat{H}=_{i=1}^N\widehat{𝐏}_i^2/2m_i+\widehat{V}(𝐱_1,\mathrm{}𝐱_N)`$, it is straightforward to calculate from (21) that
$$\frac{d}{dt}\widehat{H}\frac{d}{dt}\text{Tr}[\widehat{H}\widehat{\rho }(t)]=\underset{i=1}{\overset{N}{}}\frac{3\lambda \mathrm{}^2}{4m_ia^2}.$$
(23)
For a nucleon, the mean rate of energy increase is quite small, $`3\times 10^{25}`$eV/sec. However, deviations from the mean can be significantly greater.
For, (21) predicts excitation of atoms and nucleii. Let $`|E_0`$ be an initial bound energy eigenstate. Expanding (21) in a power series in (bound state size/$`a`$)<sup>2</sup>, the excitation rate of state $`|E`$ is
$$\mathrm{\Gamma }\frac{dE|\widehat{\rho }(t)|E}{dt}|_{t=0}=\frac{\lambda }{2a^2}E|\underset{i=1}{\overset{N}{}}\frac{m_i}{m_p}\widehat{𝐗}_i|E_0E_0|\underset{i=1}{\overset{N}{}}\frac{m_i}{m_p}\widehat{𝐗}_i|E+\text{O(size}/a)^4.$$
(24)
Since $`|E_0`$, $`|E`$ are eigenstates of the center of mass operator $`_{i=1}^Nm_i\widehat{𝐗}_i/_{i=1}^Nm_i`$ with eigenvalue 0, the dipole contribution explicitly given in (24) vanishes identically. This leaves the quadrupole contribution as the leading term, which is too small to be measured at present.
However, the choice of $`\widehat{A}(𝐱)`$ as mass density operator was made only after experimental indication. Let $`g_i`$ replace $`m_i/m_p`$ in (21), (24), so that $`\lambda g_i^2`$ is the collapse rate for the $`i`$th particle. Then, experiments looking for the radiation expected from “spontaneously” excited atoms and nucleii, in large amounts of matter for a long time, as shown by Collett, Pearle, Avignone and Nussinov in 1995, Pearle, Ring, Collar and Avignone in 1999 and Jones, Pearle and Ring in 2004, have placed the following limits:
$$|g_e/g_pm_e/m_p|<12m_e/m_p,|g_n/g_pm_n/m_p|<3(m_nm_p)/m_p.$$
#### 3.4.3 Random Walk
According to (17), (13), the center of mass wave packet, of a piece of matter of size $`a`$ or smaller, containing $`N`$ nucleons, achieves equilibrium size $`s`$ in a characteristic time $`\tau _s`$, and undergoes a random walk through a root mean square distance $`\mathrm{\Delta }Q`$:
$$s\left[\frac{a^2\mathrm{}}{\lambda m_pN^3}\right]^{1/4},\tau _s\frac{Nm_ps^2}{\mathrm{}},\mathrm{\Delta }Q\frac{\mathrm{}\lambda ^{1/2}t^{3/2}}{m_pa}.$$
(25)
The results in (25) were obtained by Collett and Pearle in 2003. These quantitative results can be qualitatively understood as follows.
In time $`\mathrm{\Delta }t`$, the usual Schrödinger equation expands a wave packet of size $`s`$ to $`s+(\mathrm{}/Nm_ps)\mathrm{\Delta }t`$. CSL collapse, by itself, narrows the wave packet to $`s[1\lambda N^2(s/a)^2\mathrm{\Delta }t]`$. The condition of no change in $`s`$ is the result quoted above. $`\tau _s`$ is the time it takes the Schrödinger evolution to expand a wavepacket near size $`s`$ to size $`s`$: $`(\mathrm{}/Nm_ps)\tau _ss`$. The $`t^{3/2}`$ dependence of $`\mathrm{\Delta }Q`$ arises because this is a random walk without damping (unlike Brownian motion, where $`\mathrm{\Delta }Qt^{1/2}`$). The mean energy increase $`\lambda N\mathrm{}^2m_p^1a^2t`$ of (23) implies the root-mean-square velocity increase $`[\lambda \mathrm{}^2m_p^2a^2t]^{1/2}`$, whose product with $`t`$ is $`\mathrm{\Delta }Q`$.
For example, a sphere of density 1gm/cc and radius $`10^5`$cm has $`s4\times 10^7`$cm, $`\tau _s0.6`$sec and $`\mathrm{\Delta }Q5[t\text{ in days}]^{3/2}`$cm. At the reported achieved low pressure of $`5\times 10^{17}`$Torr at $`4.2^{}`$K reported by Gabrielse’s group in 1990, the mean collision time with gas molecules is $`80`$min, over which $`\mathrm{\Delta }Q0.7`$mm. Thus, observation of this effect should be feasible.
### 3.5 Further Remarks
It is possible to define energy for the $`w(𝐱,t)`$ field so that total energy is conserved: as the particles gain energy, the $`w`$-field loses energy, as shown by Pearle in 2005.
Attempts to construct a special relativistic CSL-type model have not yet succeeded although Pearle in 1990, 1992, 1999, Ghirardi, Grassi and Pearle in 1990 and Nicrosini and Rimini in 2003 have made valiant attempts. The problem is that the white noise field $`w(𝐱,t)`$ contains all wavelengths and frequencies, exciting the vacuum in lowest order in $`\lambda `$ to produce particles at the unacceptable rate of infinite energy/sec-cc. Collapse models which utilize a colored noise field $`w`$ have a similar problem in higher order. In 2005, Pearle suggested a “quasi-relativistic” model which reduces to CSL in the low speed limit.
CSL is a phenomenological model which describes dynamical collapse so as to achieve S. Besides needing decisive experimental verification, it needs identification of the $`w(𝐱,t)`$ field with a physical entity.
Other collapse models which have been investigated are briefly described below.
### 3.6 Spontaneous Localization Model (SL)
The SL model of Ghirardi et al. (1986), although superseded by CSL, is historically important and conceptually valuable. Let $`\widehat{H}=0`$ for simplicity, and consider a single particle whose wave function at time $`t`$ is $`\psi (𝐱,t)`$. Over the next interval $`dt`$, with probability $`1\lambda dt`$, it does not change. With probability $`\lambda dt`$ it does change, by being “spontaneously localized” or “hit.” A hit means that the new (unnormalized) wavefunction suddenly becomes
$$\psi (𝐱,t+dt)=\psi (𝐱,t)(\pi a^2)^{3/4}e^{(2a^2)^1(𝐱𝐳)^2}\text{with probability}\lambda dtd𝐳𝑑𝐱|\psi (𝐱,t+dt)|^2.$$
Thus $`𝐳`$, the “center” of the hit, is most likely to be located where the wavefunction is large. For a single particle in the superposition described in section 3.4.1, a single hit is overwhelmingly likely to reduce the wave function to one or the other location, with total probability $`|\alpha _i|^2`$, at the rate $`\lambda `$.
For an $`N`$ particle clump, it is considered that each particle has the same independent probability, $`\lambda dt`$, of being hit. But, for the example in section 3.4.1, a single hit on any particle in one location of the clump has the effect of multiplying the wave function part describing the clump in the other location by the tail of the gaussian, thereby collapsing the wave function at the rate $`\lambda N`$.
By use of the gaussian hit rather than a delta function hit, SL solves the problem of giving too much energy to particles as mentioned in section 3.3. It also solves the problem of achieving a slow collapse rate for a superposition of small objects and a fast collapse rate for a superposition of large objects. However, the SL hits on individual particles destroys the (anti-) symmetry of wave functions. The CSL collapse toward mass density eigenstates removes that problem. Also, while SL modifies the Schrödinger evolution of a wave function, it involves discontinuous dynamics and so is not described by a modified Schrödinger equation as is CSL.
## 4 Other Models
For a single (low-energy) particle, the polar decomposition $`\mathrm{\Psi }=Re^{(i/\mathrm{})S}`$ of the Schrödinger equation implies two real equations,
$$\frac{R^2}{t}+(R^2\frac{S}{m})=0$$
(26)
(the continuity equation for $`R^2=\left|\mathrm{\Psi }\right|^2`$) and
$$\frac{S}{t}+\frac{(S)^2}{2m}+V+Q=0$$
(27)
where $`Q(\mathrm{}^2/2m)^2R/R`$ is the “quantum potential”. (These equations have an obvious generalisation to higher-dimensional configuration space.) In 1926, Madelung proposed that one should start from (26) and (27) – regarded as hydrodynamical equations for a classical charged fluid with mass density $`mR^2`$ and fluid velocity $`S/m`$ – and construct $`\mathrm{\Psi }=Re^{(i/\mathrm{})S}`$ from the solutions.
This “hydrodynamical” interpretation suffers from many difficulties, especially for many-body systems. In any case, a criticism by Wallstrom (1994) seems decisive: (26) and (27) (and their higher-dimensional analogs) are not, in fact, equivalent to the Schrödinger equation. For, as usually understood, the quantum wave function $`\mathrm{\Psi }`$ is a single-valued and continuous complex field, which typically possesses nodes ($`\mathrm{\Psi }=0`$), in the neighborhood of which the phase $`S`$ is multivalued, with values differing by integral multiples of $`2\pi \mathrm{}`$. If one allows $`S`$ in (26), (27) to be multi-valued, there is no reason why the allowed values should differ by integral multiples of $`2\pi \mathrm{}`$, and in general $`\mathrm{\Psi }`$ will not be single-valued. On the other hand, if one restricts $`S`$ in (26), (27) to be single-valued, one will exclude wave functions – such as those of nonzero angular momentum – with a multivalued phase. (This problem does not exist in pilot-wave theory as we have presented it here, where $`\mathrm{\Psi }`$ is regarded as a basic entity.)
Stochastic mechanics, introduced by Fényes in 1952 and Nelson (1966), has particle trajectories $`𝐱(t)`$ obeying a “forward” stochastic differential equation $`d𝐱(t)=𝐛(𝐱(t),t)dt+d𝐰(t)`$, where $`𝐛`$ is a drift (equal to the mean forward velocity) and $`𝐰`$ a Wiener process, and also a similar “backward” equation. Defining the “current velocity” $`𝐯=\frac{1}{2}(𝐛+𝐛_{})`$, where $`𝐛_{}`$ is the mean backward velocity, and using an appropriate time-symmetric definition of mean acceleration, one may impose a stochastic version of Newton’s second law. If one assumes, in addition, that $`𝐯`$ is a gradient ($`𝐯=S/m`$ for some $`S`$), then one obtains (26), (27) with $`R\sqrt{\rho }`$, where $`\rho `$ is the particle density. Defining $`\mathrm{\Psi }\sqrt{\rho }e^{(i/\mathrm{})S}`$, it appears that one recovers the Schrödinger equation for the derived quantity $`\mathrm{\Psi }`$. However, again, there is no reason why $`S`$ should have the specific multivalued structure required for the phase of a single-valued complex field. It then seems that, despite appearances, quantum theory cannot in fact be recovered from stochastic mechanics (Wallstrom 1994). The same problem occurs in models that use stochastic mechanics as an intermediate step (e.g., Markopoulou and Smolin in 2004): the Schrödinger equation is obtained only for exceptional, nodeless wave functions.
Bohm and Bub (1966) first proposed dynamical wave function collapse through deterministic evolution. Their collapse outcome is determined by the value of a Wiener-Siegel hidden variable (a variable distributed uniformly over the unit hypersphere in a Hilbert space identical to that of the statevector). In 1976, Pearle proposed dynamical wave function collapse equations where the collapse outcome is determined by a random variable, and suggested (Pearle 1979) that the modified Schrödinger equation be formulated as an Itô stochastic differential equation, a suggestion which has been widely followed. (The equation for the state vector given here, which is physically more transparent, has its time derivative equivalent to a Stratonovich stochastic differential equation, which is readily converted to the Itô form.) The importance of having the density matrix describing collapse be of the Lindblad-Kossakowski form was emphasized by Gisin in 1984 and Diosi in 1988. The stochastic differential Schrödinger equation which achieves this was found independently by Diosi in 1988 and by Belavkin, Gisin and Pearle in separate papers in 1989 (see also Ghirardi et al. 1990).
A gravitationally motivated stochastic collapse dynamics was proposed by Diosi in 1989 (and somewhat corrected by Ghirardi et al. in 1990). Penrose emphasized in 1996 that a quantum state, such as that describing a mass in a superposition of two places, puts the associated space-time geometry also in a superposition, and has argued that this should lead to wave function collapse. He suggests that the collapse time should be $`\mathrm{}/\mathrm{\Delta }E`$, where $`\mathrm{\Delta }E`$ is the gravitational potential energy change obtained by actually displacing two such masses: for example, the collapse time $`\mathrm{}/(Gm^2/R)`$, where the mass is $`m`$, its size is $`R`$, and the displacement is $`R`$ or larger. No specific dynamics is offered, just the vision that this will be a property of a correct future quantum theory of gravity.
Collapse to energy eigenstates was first proposed by Bedford and Wang in 1975 and 1977 and, in the context of stochastic collapse (e.g., (11) with $`\widehat{A}=\widehat{H}`$), by Milburn in 1991 and Hughston in 1996, but it has been argued by Finkelstein in 1993 and Pearle in 2004 that such energy-driven collapse cannot give a satisfactory picture of the macroscopic world. Percival in 1995 and in a 1998 book, and Fivel in 1997 have discussed energy-driven collapse for microscopic situations.
Adler (2004) has presented a classical theory (a hidden-variables theory) from which it is argued that quantum theory “emerges” at the ensemble level. The classical variables are $`N\times N`$ matrix field amplitudes at points of space. They obey appropriate classical Hamiltonian dynamical equations which he calls “trace dynamics,” since the Hamiltonian, Lagrangian, Poisson bracket, etc. expressions have the form of the trace of products of matrices and their sums with constant coefficients. Using classical statistical mechanics, canonical ensemble averages of (suitably projected) products of fields are analyzed and it is argued that they obey all the properties associated with Wightman functions, from which quantum field theory, and its non-relativistic limit quantum mechanics, may be derived. As well as obtaining the algebra of quantum theory in this way, it is argued that statistical fluctuations around the canonical ensemble can give rise to wave function collapse behavior, of the kind discussed here, both energy-driven and CSL-type mass-density-driven collapse. The Hamiltonian needed for this theory to work is not provided but, as the argument progresses, its necessary features are delimited.
BIBLIOGRAPHY
Adler SL (2004) Quantum Theory as an Emergent Phenomenon. Cambridge University Press, Cambridge.
Bassi A and Ghirardi GC (2003) Dynamical reduction models. Physics Reports 379: 257–426. ArXiv: quant-ph/0302164.
Bell JS (1987) Speakable and Unspeakable in Quantum Mechanics. Cambridge University Press, Cambridge.
Bohm D (1952) A suggested interpretation of the quantum theory in terms of ‘hidden’ variables. I and II. Physical Review 85: 166–179; 180–193.
Bohm D and Bub J (1966) A proposed solution of the measurement problem in quantum mechanics by a hidden variable theory. Reviews of Modern Physics 38: 453–469.
Dürr D, Goldstein S and Zanghì N (2003) Quantum equilibrium and the role of operators as observables in quantum theory. ArXiv: quant-ph/0308038.
Ghirardi G, Rimini A and Weber T (1986) Unified dynamics for microscopic and macroscopic systems. Physical Review D 34: 470–491.
Ghirardi G, Pearle P and Rimini A (1990) Markov processes in Hilbert space and continuous spontaneous localization of systems of identical particles. Physical Review A 42: 78–89.
Holland PR (1993) The Quantum Theory of Motion: an Account of the de Broglie-Bohm Causal Interpretation of Quantum Mechanics. Cambridge University Press, Cambridge.
Nelson E (1966) Derivation of the Schrödinger equation from Newtonian mechanics. Physical Review 150: 1079–1085.
Pearle P (1979) Toward explaining why events occur. International Journal of Theoretical Physics 18: 489–518.
Pearle P (1989) Combining stochastic dynamical state-vector reduction with spontaneous localization Physical Review A 39: 2277–2289.
Pearle P (1999) Collapse models. In: Petruccione F and Breuer HP (eds.) Open Systems and Measurement in Relativistic Quantum Theory, pp 195–234. Springer-Verlag, Heidelberg. ArXiv: quant-ph/9901077.
Valentini A (1991) Signal-locality, uncertainty, and the subquantum H-theorem. I and II. Physics Letters A156: 5–11; A158: 1–8.
Valentini A (2002a) Signal-locality in hidden-variables theories. Physics Letters A297: 273–278.
Valentini A (2002b) Subquantum information and computation. Pramana – Journal of Physics 59: 269–277. ArXiv: quant-ph/0203049.
Valentini A (2004a) Black holes, information loss, and hidden variables. ArXiv: hep-th/0407032.
Valentini A (2004b) Universal signature of non-quantum systems. Physics Letters A332: 187–193. ArXiv: quant-ph/0309107.
Valentini A and Westman H (2005) Dynamical origin of quantum probabilities. Proceedings of the Royal Society of London A461: 253–272.
Wallstrom TC (1994) Inequivalence between the Schrödinger equation and the Madelung hydrodynamic equations. Physical Review A49: 1613–1617.
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# Homotopy groups of moduli spaces of representations
## 1. Introduction
Given a closed oriented surface, $`X`$, and a Lie group $`G`$, moduli spaces of surface group representations in $`G`$ have rich geometric and topological structure which reflects properties of both $`X`$ and $`G`$. In this paper we consider the cases where $`G`$ is $`\mathrm{GL}(n,)`$ or $`\mathrm{U}(p,q)`$.
Our main tools rely on an interpretation of the moduli spaces in terms of holomorphic bundles. Such an interpretation starts from the basic correspondence between representations of the fundamental group and flat principal bundles. Holomorphic bundles enter the picture if we fix a complex structure on the surface $`X`$ — thereby turning it into a Riemann surface. By results of Hitchin , Donaldson , Simpson and Corlette , if $`G`$ is complex semisimple then the flat principal $`G`$-bundles corresponding to semisimple representations of $`\pi _1X`$ in $`G`$ are equivalent to polystable $`G`$-Higgs bundles over the Riemann surface. More generally, such Higgs bundles exist if $`G`$ is complex reductive, in which case the polystable $`G`$-Higgs bundles correspond to semisimple representations not of $`\pi _1X`$ but of a central extension of the fundamental group.
Referring to $`\pi _1X`$ and its central extensions as surface groups, we can thus identify the moduli spaces of surface group representations with moduli spaces of polystable Higgs bundles. This identification puts a natural Kähler structure on the moduli spaces and also reveals a compatible $`^{}`$-action. The restriction of this action to $`S^1`$ leads to a symplectic moment map whose squared norm serves as a proper Morse function. In a striking example of the interplay between geometry and topology, these geometric features on the moduli space of Higgs bundles provide powerful tools for studying the topology of the underlying moduli spaces of surface group representations.
Holomorphic bundle techniques can also be adapted to the case in which $`G`$ is a real reductive Lie group, in particular when $`G`$ is a real form of a complex reductive group. If $`G`$ is the compact real form of a complex reductive group $`G_{}`$, then the theorem of Narasimhan and Seshadri and its generalization by Ramanathan identify representations into $`G`$ with polystable principal $`G_{}`$-bundles. For non-compact real forms the basic ideas were first introduced by Nigel Hitchin. In he outlined how to define the appropriate Higgs bundles and applied his methods to the case $`G=\mathrm{SL}(n,)`$ and also to other split real forms. Other special cases have been considered in a similar way<sup>1</sup><sup>1</sup>1Notably $`\mathrm{Sp}(4,)`$ and $`\mathrm{SU}(2,2)`$ , $`\mathrm{U}(2,1)`$ , $`\mathrm{U}(p,q)`$ and $`\mathrm{PU}(p,q)`$ , $`\mathrm{GL}(n,)`$ . Higgs bundle methods have also been applied, albeit in a more algebraic way in the cases $`\mathrm{U}(p,1)`$ , $`\mathrm{PU}(2,1)`$ , and $`\mathrm{PU}(p,p)`$ .. In we began an in-depth study of the groups $`\mathrm{U}(p,q)`$ and their adjoint forms $`\mathrm{PU}(p,q)`$ (for any $`p`$ and $`q`$) from this point of view. This paper is a continuation of that work.
The most primitive topological feature of the moduli spaces is their number of connected components, i.e. $`\pi _0`$. The above methods have been effective in addressing this question, mainly by exploiting the properness of the above mentioned Morse function. This transfers questions about $`\pi _0`$ for the moduli spaces into questions about the connected components of the minimal submanifolds for the Morse function.
In good cases, there is additional useful Morse theoretic information which has thus far gone unexploited. Our goal is to correct this oversight. In particular, using information about the Morse indices of non-minimal critical points, we can relate higher homotopy groups for the full moduli spaces to those of their minimal submanifolds. For the latter we rely on the calculations by Daskalopoulos and Uhlenbeck for higher homotopy groups of the moduli space of stable vector bundles. Our main results for the moduli spaces of $`\mathrm{GL}(n,)`$ and $`\mathrm{U}(p,q)`$ Higgs bundles, and hence for the corresponding moduli spaces of representations, are given in Theorems 4.4 and 4.22 respectively.
Acknowledgments. We thank Tamás Hausel for enlightening discussions. We are greatly indebted to an anonymous referee for pointing out that our initial estimate in (2) of Proposition 3.12 could be improved and for providing the argument outlined in Remark 3.13.
## 2. Surface group representations and Higgs bundles
For a more thorough account of the material in this section see .
### 2.1. Surface group representations
Let $`X`$ be a smooth closed oriented surface of genus $`g2`$. The fundamental group, $`\pi _1X`$, of $`X`$ is a finitely generated group generated by $`2g`$ generators, say $`A_1,B_1,\mathrm{},A_g,B_g`$, subject to the single relation $`_{i=1}^g[A_i,B_i]=1`$. It has a universal central extension
(2.1)
$$0\mathrm{\Gamma }\pi _1X1$$
generated by the same generators as $`\pi _1X`$, together with a central element $`J`$ subject to the relation $`_{i=1}^g[A_i,B_i]=J`$.
By a representation of $`\mathrm{\Gamma }`$ in $`\mathrm{GL}(n,)`$ we mean a homomorphism $`\rho :\mathrm{\Gamma }\mathrm{GL}(n,)`$. We say that a representation of $`\mathrm{\Gamma }`$ in $`\mathrm{GL}(n,)`$ is *semisimple* if the $`^n`$-representation of $`\mathrm{\Gamma }`$ induced by the fundamental representation of $`\mathrm{GL}(n,)`$ is semisimple<sup>2</sup><sup>2</sup>2In general a representation of $`\mathrm{\Gamma }`$ in a reductive Lie group $`G`$ is said to be semisimple if the induced (adjoint) representation on the Lie algebra of $`G`$ is semisimple. For $`G=\mathrm{GL}(n,)`$ this is equivalent to the definition given here.. The group $`\mathrm{GL}(n,)`$ acts on the set of representations via conjugation. Restricting to the semisimple representations, denoted by $`\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{GL}(n,))`$, we get the *moduli space* of representations of $`\mathrm{\Gamma }`$ in $`\mathrm{GL}(n,)`$,
(2.2)
$$(\mathrm{\Gamma },\mathrm{GL}(n,))=\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{GL}(n,))/\mathrm{GL}(n,).$$
The set $`\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{GL}(n,))`$ can be embedded in $`\mathrm{GL}(n,)^{2g+1}`$ via the map
$`\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{GL}(n,))`$ $`\mathrm{GL}(n,)^{2g+1}`$
$`\rho `$ $`(\rho (A_1),\mathrm{}\rho (B_g),\rho (J)).`$
We can then give $`\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{GL}(n,))`$ the subspace topology and $`(\mathrm{\Gamma },\mathrm{GL}(n,))`$ the quotient topology. This topology is Hausdorff because we have restricted attention to semisimple representations.
There is a topological invariant of a representation $`\rho (\mathrm{\Gamma },\mathrm{GL}(n,))`$ given by $`\rho (J)`$, which coincides with the first Chern class of the vector bundle with central curvature associated to $`\rho `$. Fixing this invariant, we define
$$(n,d):=\{\rho (\mathrm{\Gamma },\mathrm{GL}(n,))|\rho (J)=[d]_nZ(\mathrm{GL}(n,)\}.$$
In particular the representations with vanishing degree correspond to representations of the fundamental group of $`X`$, that is,
(2.3)
$$(n,0)=(\pi _1X,\mathrm{GL}(n,)):=\mathrm{Hom}^+(\pi _1X,\mathrm{GL}(n,))/\mathrm{GL}(n,).$$
Similarly to the case of $`\mathrm{GL}(n,)`$ we consider the moduli space
(2.4)
$$(\mathrm{\Gamma },\mathrm{U}(p,q))=\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{U}(p,q))/\mathrm{U}(p,q).$$
The moduli space $`(\mathrm{\Gamma },\mathrm{U}(p,q))`$ can be identified with the moduli space of $`\mathrm{U}(p,q)`$-bundles on $`X`$ with projectively flat connections. Taking a reduction to the maximal compact $`\mathrm{U}(p)\times \mathrm{U}(q)`$, we thus associate to each class $`\rho (\mathrm{\Gamma },\mathrm{U}(p,q))`$ a vector bundle of the form $`VW`$, where $`V`$ and $`W`$ are rank $`p`$ and $`q`$ respectively, and thus a pair of integers $`(a,b)=(\mathrm{deg}(V),\mathrm{deg}(W))`$. There is thus a map
$$c:(\mathrm{\Gamma },\mathrm{U}(p,q))$$
given by $`c(\rho )=(a,b)`$. The corresponding map on $`\mathrm{Hom}^+(\mathrm{\Gamma },\mathrm{U}(p,q))`$ is clearly continuous and thus locally constant. Since $`\mathrm{U}(p,q)`$ is connected, the map $`c`$ is likewise continuous and thus constant on connected components. The subspace of $`(\mathrm{\Gamma },(\mathrm{U}(p,q))`$ corresponding to representations with invariants $`(a,b)`$ is denoted by
(2.5)
$$(p,q,a,b)=c^1(a,b)=\{\rho (\mathrm{\Gamma },\mathrm{U}(p,q))|c(\rho )=(a,b)\}.$$
The representations for which $`a+b=0`$ correspond to representations of the fundamental group of $`X`$, that is,
(2.6)
$$(p,q,a,a)=c^1(a,a)=\{\rho (\pi _1X,\mathrm{U}(p,q))|c(\rho )=(a,a)\}.$$
### 2.2. $`\mathrm{GL}(n,)`$-Higgs bundles
A $`\mathrm{GL}(n,)`$-Higgs bundle on a compact Riemann surface $`X`$ is a pair $`(E,\mathrm{\Phi })`$, where $`E`$ is a rank $`n`$ holomorphic vector bundle over $`X`$ and $`\mathrm{\Phi }H^0(\mathrm{End}(E)K)`$ is a holomorphic endomorphism of $`E`$ twisted by the canonical bundle $`K`$ of $`X`$. The $`\mathrm{GL}(n,)`$-Higgs bundle $`(E,\mathrm{\Phi })`$ is *stable* if the slope stability condition
(2.7)
$$\mu (E^{})<\mu (E)$$
holds for all proper $`\mathrm{\Phi }`$-invariant subbundles $`E^{}`$ of $`E`$. Here the *slope* is defined by $`\mu (E)=\mathrm{deg}(E)/\mathrm{rk}(E)`$ and *$`\mathrm{\Phi }`$-invariance* means that $`\mathrm{\Phi }(E^{})E^{}K`$. *Semistability* is defined by replacing the above strict inequality with a weak inequality. A Higgs bundle is called *polystable* if it is the direct sum of stable Higgs bundles with the same slope.
Given a hermitian metric on $`E`$, let $`A`$ denote the unique unitary connection compatible with the holomorphic structure, and let $`F_A`$ be its curvature. *Hitchin’s equations* on $`(E,\mathrm{\Phi })`$ are
(2.8) $`F_A+[\mathrm{\Phi },\mathrm{\Phi }^{}]`$ $`=\sqrt{1}\mu \text{Id}_E\omega ,`$
$`\overline{}_A\mathrm{\Phi }`$ $`=0,`$
where $`\omega `$ is the Kähler form on $`X`$, $`\text{Id}_E`$ is the identity on $`E`$, $`\mu =\mu (E)`$ and $`\overline{}_A`$ is the anti-holomorphic part of the covariant derivative $`d_A`$. A solution to Hitchin’s equations is *irreducible* if there is no proper subbundle of $`E`$ preserved by $`A`$ and $`\mathrm{\Phi }`$.
If we define a Higgs connection (as in ) by
(2.9)
$$D=d_A+\theta $$
where $`\theta =\mathrm{\Phi }+\mathrm{\Phi }^{}`$, then Hitchin’s equations are equivalent to the conditions
(2.10) $`F_D`$ $`=\sqrt{1}\mu \text{Id}_E\omega ,`$
$`d_A^{}\theta `$ $`=0.`$
In particular, the first equation says that $`D`$ is a projectively flat connection<sup>3</sup><sup>3</sup>3 The other equation is an harmonicity constraint.. If $`\mathrm{deg}(E)=0`$ then $`D`$ is actually flat. It follows that in this case the pair $`(E,D)`$ defines a representation of $`\pi _1X`$ in $`\mathrm{GL}(n,)`$. If $`\mathrm{deg}(E)0`$, then the pair $`(E,D)`$ defines a representation of $`\pi _1X`$ in $`\mathrm{PGL}(n,)`$, or equivalently, a representation of $`\mathrm{\Gamma }`$ in $`\mathrm{GL}(n,)`$. By the theorem of Corlette (), every semisimple representation of $`\mathrm{\Gamma }`$ (and therefore every semisimple representation of $`\pi _1X`$) arises in this way.
If we fix the rank and degree (say $`n`$ and $`d`$ respectively) of the bundle $`E`$, i.e. on bundles of fixed topological type, the isomorphism classes of polystable Higgs bundles are parameterized by a quasi-projective variety of dimension $`2+2n^2(g1)`$. We denote this moduli space of rank $`n`$ degree $`d`$ polystable Higgs bundles by $`(n,d)`$.
If we fix a hermitian metric on a smooth rank $`n`$ degree $`d`$ complex vector bundle on $`X`$, then there is a gauge theoretic moduli space of pairs $`(A,\mathrm{\Phi })`$, consisting of a unitary connection $`A`$ and an endomorphism valued $`(1,0)`$-form $`\mathrm{\Phi }`$, which are solutions to Hitchin’s equations (2.8), modulo $`\mathrm{U}(n)`$-gauge equivalence.
The gauge theory moduli space and $`(n,d)`$ are related by virtue of the Hitchin-Kobayashi correspondence: a $`\mathrm{GL}(n,)`$-Higgs bundle $`(E,\mathrm{\Phi })`$ is polystable if and only if it admits a hermitian metric such that Hitchin’s equations (2.8) are satisfied, and $`(E,\mathrm{\Phi })`$ is stable if and only if the corresponding solution is irreducible. There is, moreover, a map from the gauge theoretic moduli space to this moduli space given by taking a solution $`(A,\mathrm{\Phi })`$ to Hitchin’s equations to the Higgs bundle $`(E,\mathrm{\Phi })`$, where the holomorphic structure on $`E`$ is given by $`\overline{}_A`$. This map is a homeomorphism, and a diffeomorphism on the smooth locus.
In view of the relation between Hitchin’s equations and projectively flat connections, this correspondence gives rise to a homeomorphism between $`(n,d)`$ and the component $`(n,d)`$ of the moduli space of semisimple representations of $`\mathrm{\Gamma }`$ in $`\mathrm{GL}(n,)`$. If the degree of the Higgs bundle is zero, then the moduli space $`(n,0)`$ is homeomorphic to the moduli space of representations of $`\pi _1X`$ in $`\mathrm{GL}(n,)`$.
###### Theorem 2.1.
If $`(n,d)`$ is such that $`\mathrm{GCD}(n,d)=1`$ then the moduli space $`(n,d)`$ is a non-empty connected smooth hyperkähler manifold.
### 2.3. $`\mathrm{U}(p,q)`$-Higgs bundles
There is a special class of $`\mathrm{GL}(n,)`$-Higgs bundles, related to representations in $`\mathrm{U}(p,q)`$ given by the requirements that
(2.11) $`E`$ $`=VW`$
$`\mathrm{\Phi }`$ $`=\left(\begin{array}{cc}0& \beta \\ \gamma & 0\end{array}\right)`$
where $`V`$ and $`W`$ are holomorphic vector bundles of rank $`p`$ and $`q`$ respectively and the non-zero components in the Higgs field are $`\beta H^0(\mathrm{Hom}(W,V)K)`$, and $`\gamma H^0(\mathrm{Hom}(V,W)K)`$. We say $`(E,\mathrm{\Phi })`$ is a *stable* $`\mathrm{U}(p,q)`$-Higgs bundle if the slope stability condition $`\mu (E^{})<\mu (E)`$, is satisfied for all $`\mathrm{\Phi }`$-invariant subbundles $`E^{}=V^{}W^{}`$, i.e. for all subbundles $`V^{}V`$ and $`W^{}W`$ such that
(2.12) $`\beta `$ $`:W^{}V^{}K`$
(2.13) $`\gamma `$ $`:V^{}W^{}K.`$
Semistability and polystability are defined analogously to the way they are defined for $`\mathrm{GL}(n,)`$-Higgs bundles.
Hitchin’s equations make sense for $`\mathrm{U}(p,q)`$-Higgs bundles, with a $`\mathrm{U}(p,q)`$ solution being a metric with respect to which $`E=VW`$ is an orthogonal decomposition. With $`\mathrm{\Phi }`$ as in (2.11) and $`\theta =\mathrm{\Phi }+\mathrm{\Phi }^{}`$, the corresponding $`\mathrm{U}(p,q)`$-Higgs connection $`D=d_A+\theta `$ is not only projectively flat but has $`\mathrm{U}(p,q)`$ holonomy. This provides the link between $`\mathrm{U}(p,q)`$-Higgs bundles and surface group representations in $`\mathrm{U}(p,q)`$, leading to:
###### Theorem 2.2.
Let $`(p,q,a,b)`$ be the moduli space of polystable $`\mathrm{U}(p,q)`$-Higgs bundles with $`\mathrm{deg}(V)=a`$ and $`\mathrm{deg}W=b`$. Then with $`(p,q,a,b)`$ as in (2.5) there is a homeomorphism $`(p,q,a,b)(p,q,a,b)`$.
The *Toledo invariant* of the representation corresponding to $`(E=VW,\mathrm{\Phi })`$ is defined by
(2.14)
$$\tau =\tau (p,q,a,b)=2\frac{qapb}{p+q}$$
where $`a=\mathrm{deg}(V)`$ and $`b=\mathrm{deg}(W)`$. This invariant satisfies the following Milnor-Wood-type inequality (proved by Domic and Toledo )
(2.15)
$$|\tau (p,q,a,b)|\mathrm{min}\{p,q\}(2g2).$$
###### Theorem 2.3.
Let $`(p,q,a,b)`$ such that $`\mathrm{GCD}(p+q,a+b)=1`$. Then $`(p,q,a,b)`$ (and hence $`(p,q,a,b)`$) is a connected smooth Kähler manifold which is non-empty if and only if $`|\tau (p,q,a,b)|\mathrm{min}\{p,q\}(2g2)`$.
## 3. Morse theory on the moduli space
### 3.1. The Morse function
Let $``$ be either $`(n,d)`$ or $`(p,q,a,b)`$. We will assume that $`\mathrm{GCD}(n,d)=1`$ and $`\mathrm{GCD}(p+q,a+b)=1`$. Under this coprimality condition, there are no strictly semistable Higgs bundles and the moduli space $``$ is smooth. The non-zero complex numbers $`^{}`$ act on $``$ via the map $`\lambda (E,\mathrm{\Phi })=(E,\lambda \mathrm{\Phi })`$. However, to have an action on the gauge theory moduli space (i.e. on the set of solutions to Hitchin’s equations (2.8), cf. Section 2), one must restrict to the action of $`S^1^{}`$. This is a Hamiltonian action and the associated moment map is
$$[(A,\mathrm{\Phi })]\frac{1}{2}\mathrm{\Phi }^2=i_X\mathrm{tr}(\mathrm{\Phi }\mathrm{\Phi }^{})$$
where the adjoint $`\mathrm{\Phi }^{}`$ is taken with respect to the hermitian metric on $`E`$. We shall, however, prefer to consider the positive function
(3.1)
$$f([A,\mathrm{\Phi }])=\frac{1}{2}\mathrm{\Phi }^2.$$
Next we recall a general result of Frankel , which was first used in the context of moduli spaces of Higgs bundles by Hitchin .
###### Theorem 3.1.
Let $`\stackrel{~}{f}:M`$ be a proper moment map for a Hamiltonian circle action on a Kähler manifold $`M`$. Then $`\stackrel{~}{f}`$ is a perfect Bott–Morse function.
### 3.2. Morse theory and homotopy groups
In this Section we recall some basic facts of Bott–Morse theory. Let $`_l`$ be the critical submanifolds of $`f`$ and $`\nu (_l)`$ be the normal bundle of $`_l`$ in $``$. The Hessian of $`f`$ is non-degenerate on $`\nu (_l)`$ and we have the decomposition in positive and negative eigenspace bundles
$$\nu (_l)=\nu ^+(_l)\nu ^{}(_l).$$
The index of $`_l`$ is defined as
$$\mathrm{index}(_l):=\mathrm{rk}\nu ^{}(_l).$$
Let $`_l^+`$ be the stable set of $`_l`$, i.e., the subset of $``$ defined by the points of $``$ which flow to $`_l`$. It follows from Bott–Morse theory that $`_l^+`$ is a submanifold of $``$ of codimension
(3.2)
$$\mathrm{codim}_{}(_l^+)=\mathrm{index}(_l),$$
and that there is a stratification
(3.3)
$$=\underset{l}{}_l^+.$$
###### Proposition 3.2.
Let $`𝒩=_0`$ be the submanifold of local minima of $`f`$. If $`\mathrm{index}(_l)m2`$ for every $`l0`$ then
$$\pi _i()\pi _i(𝒩)\text{for}im2.$$
###### Proof.
The stratification (3.3) shows that
$$_0^+=\underset{l0}{}_l^+$$
and the Morse flow defines a retraction from $`_0^+`$ to $`𝒩=_0`$. Thus the result is an immediate consequence of standard homotopy theory, using (3.2). ∎
### 3.3. Deformation theory of Higgs bundles
In the following we recall some standard facts about the deformation theory of Higgs bundles (this has been treated in many places, a convenient reference is Biswas–Ramanan ). In order to describe the results in a uniform way for a $`G`$-Higgs bundle $`(E,\mathrm{\Phi })`$ when $`G=\mathrm{GL}(n,)`$ or $`\mathrm{U}(p,q)`$, we introduce bundles $`U_G^+`$, $`U_G^{}`$ and $`U_G`$ defined by
$`U_{\mathrm{GL}(n,)}^+`$ $`=U_{\mathrm{GL}(n,)}^{}=U_{\mathrm{GL}(n,)}=\mathrm{End}(E),`$
$`U_{\mathrm{U}(p,q)}^+`$ $`=\mathrm{End}(V)\mathrm{End}(W),`$
$`U_{\mathrm{U}(p,q)}^{}`$ $`=\mathrm{Hom}(W,V)\mathrm{Hom}(V,W),`$
$`U_{\mathrm{U}(p,q)}`$ $`=U_{\mathrm{U}(p,q)}^+U_{\mathrm{U}(p,q)}^{}=\mathrm{End}(VW),`$
where the bundles $`V`$ and $`W`$ are as in Section 2.3. Note that, with this notation, $`\mathrm{\Phi }H^0(U_G^{}K)`$.
###### Remark 3.3.
Both for $`G=\mathrm{GL}(n,)`$ and for $`G=\mathrm{U}(p,q)`$, there is an inner product on $`U_G`$ which is invariant under the adjoint action of $`U_G`$, i.e.,
(3.4)
$$\mathrm{ad}(\psi )x,y+x,\mathrm{ad}(\psi )y=0$$
for local sections $`x`$, $`y`$ and $`\psi `$ of $`U_G`$. This inner product restricts to an inner product on $`U_G^{}`$ and $`U_G^+`$, giving rise to an isomorphism
(3.5)
$$U_G^\pm \stackrel{}{}(U_G^\pm )^{}.$$
Note that under this duality
$$\mathrm{ad}(\mathrm{\Phi })^t=\mathrm{ad}(\mathrm{\Phi })1_{K^1}.$$
###### Proposition 3.4.
Let $`(E,\mathrm{\Phi })`$ be a $`G`$-Higgs bundle for $`G=\mathrm{GL}(n,)`$ or $`G=\mathrm{U}(p,q)`$ and define the following complex of sheaves
$$C_G^{}(E,\mathrm{\Phi }):U_G^+\stackrel{\mathrm{ad}(\mathrm{\Phi })}{}U_G^{}K.$$
Then the following holds:
1. The space of endomorphisms of $`(E,\mathrm{\Phi })`$ is naturally isomorphic to $`^0(C_G^{})`$.
2. The infinitesimal deformation space of $`(E,\mathrm{\Phi })`$ is naturally isomorphic to $`^1(C_G^{})`$.
The following proposition is simply a statement of the fact that a stable Higgs bundle is simple.
###### Proposition 3.5.
Let $`(E,\mathrm{\Phi })`$ be a stable $`G`$-Higgs bundle for $`G=\mathrm{GL}(n,)`$ or $`G=\mathrm{U}(p,q)`$. Then
$$^0(C_G^{}(E,\mathrm{\Phi })),$$
generated by the scalar multiples of the identity morphism.
### 3.4. Critical points and Morse indices
In the following $`(E,\mathrm{\Phi })`$ continues to denote a $`G`$-Higgs bundle for $`G=\mathrm{GL}(n,)`$ or $`G=\mathrm{U}(p,q)`$ and for ease of notation we omit the subscript $`G`$ on the bundles $`U_G^\pm `$ and the complex $`C_G^{}`$. The critical points of the function $`f`$ are exactly the fixed points of the $`S^1`$-action on $``$. This allows one to describe the corresponding Higgs bundles as “complex variations of Hodge structure”, as follows (cf. Hitchin and also Simpson ).
###### Proposition 3.6.
If $`(E,\mathrm{\Phi })`$ corresponds to a critical point of $`f`$, then there is a semisimple element $`\psi H^0(U^+)`$ and a corresponding decomposition in eigenspace bundles
(3.6)
$$U_G^\pm =\underset{k}{}U_k^\pm $$
for the adjoint action of $`\psi `$, such that $`\mathrm{ad}(\psi )`$ has eigenvalue $`ik`$ on $`U_k^\pm `$. Furthermore, $`[\psi ,\mathrm{\Phi }]=i\mathrm{\Phi }`$, i.e.,
$$\mathrm{\Phi }H^0(U_1^{}K).$$
In particular, this means that the deformation complex of $`(E,\mathrm{\Phi })`$ decomposes as
(3.7)
$$C^{}(E,\mathrm{\Phi })=\underset{k}{}C_k^{}(E,\mathrm{\Phi }),$$
where we have defined for each $`k`$ the complex
$$C_k^{}(E,\mathrm{\Phi }):U_k^+\stackrel{\mathrm{ad}(\mathrm{\Phi })}{}U_{k+1}^{}K.$$
Thus the tangent space to $``$ at $`(E,\mathrm{\Phi })`$ has a decomposition
(3.8)
$$^1(C^{}(E,\mathrm{\Phi }))=\underset{k}{}^1(C_k^{}(E,\mathrm{\Phi })).$$
###### Remark 3.7.
Using the definition of the $`U_k`$ and (3.4), we have that
$$U_k^\pm U_k^{\pm ,}$$
under the duality (3.5). Moreover, writing
$$\mathrm{ad}(\mathrm{\Phi })_k^\pm =\mathrm{ad}(\mathrm{\Phi })_{|U_k^\pm }:U_k^\pm U_{k+1}^{}K,$$
we have
$$\mathrm{ad}(\mathrm{\Phi })_{k,t}^\pm =(\mathrm{ad}(\mathrm{\Phi })_{k1}^{})1_{K^1}.$$
The calculations of Hitchin \[19, §8\] show that eigenvalues of the Hessian of $`f`$ at a critical point can be calculated as follows.
###### Proposition 3.8.
Let $`(E,\mathrm{\Phi })`$ be a stable $`G`$-Higgs bundle which corresponds to a critical point of $`f`$, for $`G=\mathrm{GL}(n,)`$ or $`G=\mathrm{U}(p,q)`$. In the decomposition (3.8) the eigenvalue $`k`$ subspace for the Hessian of $`f`$ is isomorphic to $`^1(C_k^{}(E,\mathrm{\Phi }))`$. In particular, the negative eigenspace at $`(E,\mathrm{\Phi })`$ for the Hessian is given by
$$\nu _{(E,\mathrm{\Phi })}^{}(_l)\underset{k>0}{}^1(C_k^{}).$$
###### Lemma 3.9.
Let $`(E,\mathrm{\Phi })`$ be a stable $`G`$-Higgs bundle which corresponds to a critical point of $`f`$. Then
$$^0(C_k^{}(E,\mathrm{\Phi }))=0\text{and}^2(C_k^{}(E,\mathrm{\Phi }))=0$$
for $`k>0`$.
###### Proof.
From (3.7) we have a decomposition
$$^0(C^{}(E,\mathrm{\Phi }))=\underset{k}{}^0(C_k^{}(E,\mathrm{\Phi })).$$
and we know from Proposition 3.5 that the only trivial endomorphisms of $`(E,\mathrm{\Phi })`$ are the scalars, which have weight zero in this decomposition. This gives the vanishing of $`^0`$.
For the vanishing of $`^2`$, consider first the case $`G=\mathrm{GL}(n,)`$. Then $`U_k^+=U_k^{}`$ and, using Remark 3.7, we see that the dual complex of $`C_k^{}(E,\mathrm{\Phi })`$ is isomorphic to the complex
$$C_{k1}^{}(E,\mathrm{\Phi })K^1:U_{k1}^+K^1\stackrel{\mathrm{ad}(\mathrm{\Phi })}{}U_k^{}.$$
The change in sign of $`\mathrm{ad}(\mathrm{\Phi })`$ does not influence the cohomology and hence Serre duality for hypercohomology gives
$$^2(C_k^{}(E,\mathrm{\Phi }))^0(C_{k1}^{}(E,\mathrm{\Phi }))^{}.$$
It follows that $`^2(C_k^{}(E,\mathrm{\Phi }))`$ vanishes for $`k1`$. The case $`G=\mathrm{U}(p,q)`$ follows essentially from this, by using the fact that stability as a $`\mathrm{U}(p,q)`$-Higgs bundle implies stability as a $`\mathrm{GL}(n,)`$-Higgs bundle (see \[5, Proposition 3.19\] for a detailed argument). ∎
###### Proposition 3.10.
Let $`(E,\mathrm{\Phi })`$ be a stable $`G`$-Higgs bundle which corresponds to a critical point of $`f`$. Then the Morse index of the corresponding critical submanifold $`_l`$ is
$$\mathrm{index}(_l)=2\underset{k>0}{}dim^1(C_k^{}(E,\mathrm{\Phi })),$$
where
$$dim^1(C_k^{}(E,\mathrm{\Phi }))=\chi (C_k^{}(E,\mathrm{\Phi })).$$
###### Proof.
This is immediate from Proposition 3.8 and the vanishing of Lemma 3.9 (note that the Morse index is the real dimension of the $`^1`$, hence the factor of $`2`$). ∎
The following lemma is essentially Proposition 4.14 of . We provide a complete proof, taking this opportunity to correct some inaccuracies in the argument given in .
###### Lemma 3.11.
Let $`(E,\mathrm{\Phi })`$ be a stable $`G`$-Higgs bundle which corresponds to a critical point of $`f`$, for $`G=\mathrm{GL}(n,)`$ or $`G=\mathrm{U}(p,q)`$. Then
$$\chi (C_k^{}(E,\mathrm{\Phi }))(g1)(2\mathrm{rk}(\mathrm{ad}(\mathrm{\Phi })_k^+)\mathrm{rk}(U_k^+)\mathrm{rk}(U_{k+1}^{})).$$
Furthermore, the vanishing $`\chi (C_k^{}(E,\mathrm{\Phi }))=0`$ holds if and only if $`\mathrm{ad}(\mathrm{\Phi })_k^+:U_k^+U_{k+1}^{}K`$ is an isomorphism.
###### Proof.
In the following we shall use the abbreviated notations $`C_k^{}=C_k^{}(E,\mathrm{\Phi })`$ and
$$\mathrm{\Phi }_k^\pm =\mathrm{ad}(\mathrm{\Phi })_k^\pm :U_k^+U_{k1}^{}K.$$
By the Riemann–Roch formula we have
(3.9)
$$\chi (C_k^{})=(1g)\left(\mathrm{rk}(U_k^+)+\mathrm{rk}(U_{k+1}^{})\right)+\mathrm{deg}(U_k^+)\mathrm{deg}(U_{k+1}^{}),$$
thus we can prove the inequality stated in the Lemma by estimating the difference $`\mathrm{deg}(U_k^+)\mathrm{deg}(U_k^{})`$. In order to do this, we note first that there are short exact sequences of sheaves
$$0\mathrm{ker}(\mathrm{\Phi }_k^+)U_k^+\mathrm{im}(\mathrm{\Phi }_k^+)0$$
and
$$0\mathrm{im}(\mathrm{\Phi }_k^+)U_{k+1}^{}K\mathrm{coker}(\mathrm{\Phi }_k^+)0.$$
It follows that
(3.10)
$$\mathrm{deg}(U_k^+)\mathrm{deg}(U_{k+1}^{})=\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^+))+(2g2)\mathrm{rk}(U_{k+1}^{})\mathrm{deg}(\mathrm{coker}(\mathrm{\Phi }_k^+)).$$
We shall prove the following inequalities below.
(3.11) $`\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^+))`$ $`0,`$
(3.12) $`\mathrm{deg}(\mathrm{coker}(\mathrm{\Phi }_k^+))`$ $`(2g2)(\mathrm{rk}(U_{k+1}^{})+\mathrm{rk}(\mathrm{\Phi }_k^+)).`$
Combining these with (3.10) we obtain
$$\mathrm{deg}(U_k^+)\mathrm{deg}(U_{k+1}^{})(2g2)\mathrm{rk}(\mathrm{\Phi }_k^+),$$
which, together with (3.9), proves the inequality stated in the Lemma.
It remains to prove (3.11) and (3.12). For this we use the fact that the adjoint Higgs bundle $`(U_G,\mathrm{ad}(\mathrm{\Phi }))`$ is semistable (one way of seeing this is to note that it supports a solution to Hitchin’s equations). Clearly, the subbundle $`\mathrm{ker}(\mathrm{\Phi }_k^+)U_G`$ is $`\mathrm{ad}(\mathrm{\Phi })`$-invariant and hence
$$\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^+))\mathrm{deg}(U_G)=0,$$
thus proving (3.11).
In order to prove (3.12) a bit more work needs to be done. Consider the dual of $`\mathrm{\Phi }_k^+`$,
$$\mathrm{\Phi }_k^{+,t}:U_{k+1}^,K^1U_k^{+,},$$
and note that the image of $`\mathrm{\Phi }_k^+`$ goes to zero under the restriction map
$$U_{k+1}^{}K\mathrm{ker}(\mathrm{\Phi }_k^{+,t})^{}.$$
Hence there is an induced map
$$\mathrm{coker}(\mathrm{\Phi }_k^+)\mathrm{ker}(\mathrm{\Phi }_k^{+,t})^{}$$
which is generically an isomorphism — in fact, its kernel is the torsion subsheaf of $`\mathrm{coker}(\mathrm{\Phi }_k^+)`$. It follows that
$$\mathrm{deg}(\mathrm{coker}(\mathrm{\Phi }_k^+))\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^{+,t})^{}).$$
Since $`\mathrm{ker}(\mathrm{\Phi }_k^{+,t})`$ is locally free (in fact a subbundle) this shows that
(3.13)
$$\mathrm{deg}(\mathrm{coker}(\mathrm{\Phi }_k^+))\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^{+,t})),$$
the difference being the degree of the torsion subsheaf of $`\mathrm{coker}(\mathrm{\Phi }_k^+)`$. Now Remark 3.7 tells us that we have a commutative diagram
$$\begin{array}{ccc}U_{k+1}^,K^1& \stackrel{\mathrm{\Phi }_k^{+,t}}{}& U_k^{+,}\\ & & & & \\ U_{k1}^{}K^1& \stackrel{\mathrm{\Phi }_{k1}^{}1_{K^1}}{}& U_k^+,\end{array}$$
and thus
$$\mathrm{ker}(\mathrm{\Phi }_k^{+,t})\mathrm{ker}(\mathrm{\Phi }_{k1}^{})K^1$$
from which we conclude that
$$\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^{+,t}))=\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_{k1}^{}))(2g2)\mathrm{rk}(\mathrm{ker}(\mathrm{\Phi }_{k1}^{})).$$
Again we apply semistability of $`(U_G,\mathrm{ad}(\mathrm{\Phi }))`$ to the $`\mathrm{ad}(\mathrm{\Phi })`$-invariant subbundle $`\mathrm{ker}(\mathrm{\Phi }_{k1}^{})`$ to obtain
(3.14)
$$\mathrm{deg}(\mathrm{ker}(\mathrm{\Phi }_k^{+,t}))(2g2)\mathrm{rk}(\mathrm{ker}(\mathrm{\Phi }_{k1}^{})).$$
But clearly, $`\mathrm{rk}(\mathrm{\Phi }_k^+)=\mathrm{rk}(\mathrm{\Phi }_k^{+,t})=\mathrm{rk}(\mathrm{\Phi }_{k1}^{})`$ and $`\mathrm{rk}(U_{k+1}^{})=\mathrm{rk}(U_{k1}^,)=\mathrm{rk}(U_{k1}^{})`$ so
$$\mathrm{rk}(\mathrm{ker}(\mathrm{\Phi }_{k1}^{}))=\mathrm{rk}(U_{k+1}^{})\mathrm{rk}(\mathrm{\Phi }_k^+).$$
Combining this fact with (3.13) and (3.14) concludes the proof of (3.12).
Finally, to prove the last statement of the Lemma we note that if $`\chi (C_k^{})=0`$ then $`\mathrm{rk}(\mathrm{\Phi }_k^+)=\mathrm{rk}(U_k^+)=\mathrm{rk}(U_{k+1}^{}K)`$ and equality holds in (3.13), thus showing that $`\mathrm{\Phi }_k^+`$ is an isomorphism. ∎
###### Proposition 3.12.
1. For $`=(p,q,a,b)`$
$$\mathrm{index}(_l)2g2$$
for every non-minimal critical submanifold $`_l`$.
2. For $`=(n,d)`$
$$\mathrm{index}(_l)(n1)(2g2)$$
for every non-minimal critical submanifold $`_l`$.
###### Proof.
(1) Let $`k_0`$ be the largest $`k`$ such that $`\chi (C_k^{}(E,\mathrm{\Phi }))0`$. Since $`_l`$ is non-minimal, by Proposition 3.10 we have $`k_0>0`$. The proof of \[5, Proposition 4.17\] shows that the restriction of $`\mathrm{ad}(\mathrm{\Phi })_k^+:U_k^+U_{k+1}^{}K`$ to a fibre is never an isomorphism (in the notation of that proof, $`k_0=m1`$). Hence the right hand side of the inequality of Lemma 3.9 is strictly negative. Now the result follows from this inequality and Proposition 3.10.
(2) We recall (cf. ) that the decomposition $`U=U_k`$ comes from a decomposition $`E=E_1\mathrm{}E_m`$ with $`U_k=_{k=ji}\mathrm{Hom}(E_i,E_j)`$. In particular, the weights $`k`$ are consecutive integers. Thus Proposition 3.10, together with Riemann–Roch and the fact that $`U_G^+=U_G^{}`$ for $`G=\mathrm{GL}(n,)`$, gives
$`\frac{1}{2}\mathrm{index}(_l)`$ $`=(g1){\displaystyle \underset{k1}{}}(\mathrm{rk}(U_k)+\mathrm{rk}(U_{k+1})2\mathrm{rk}(\mathrm{ad}(\mathrm{\Phi })_k)`$
$`=(g1)(\mathrm{rk}(U_1)+2\mathrm{rk}(U_{k2})2\mathrm{rk}(\mathrm{ad}(\mathrm{\Phi })_{k1})).`$
But clearly the rank of $`\mathrm{ad}(\mathrm{\Phi })_{k1}:U_{k1}U_{k2}K`$ is less than or equal to the rank of $`U_{k2}`$ and hence
$$\frac{1}{2}\mathrm{index}(_l)(g1)\mathrm{rk}(U_1).$$
Let $`\nu _i=\mathrm{rk}(E_i)`$. Then $`\nu _i=n`$ and $`\mathrm{rk}(U_1)=\nu _1\nu _2+\mathrm{}+\nu _{m1}\nu _m`$. One easily shows that $`\nu _1\nu _2+\mathrm{}+\nu _{m1}\nu _mn1`$. This finishes the proof of (2). ∎
###### Remark 3.13.
Our initial estimate in (2) of Proposition 3.12 was $`\mathrm{index}(_l)2g2`$. It was pointed out to us by an anonymous referee that this could be improved, and also that an alternative way of proving this estimate is as follows. The absolute minimum of $`f`$ on $`(n,d)`$ is $`M(n,d)`$, so $`H^{}(M(n,d))`$ injects into $`H^{}((n,d))`$ because $`f`$ is perfect. For the same reason, any critical submanifold of index $`l`$ gives a non-trivial contribution to the cohomology of $`(n,d)`$ in dimension $`l`$, which is not in the image of $`H^{}(M(n,d))`$. Now Markman shows that $`H^{}(B\overline{𝒢})`$ (the cohomology of the classifying space of the reduced gauge group) surjects onto $`H^{}((n,d))`$. On the other hand, Uhlenbeck–Daskalopoulos prove that $`H^r(B\overline{𝒢})`$ is isomorphic to $`H^r(M(n,d))`$ for $`r<(2g2)(n1)`$. Hence no critical submanifold can have index $`l<(2g2)(n1)`$.
### 3.5. Local Minima
The minima of the Morse function on $`(n,d)`$ is given by the following .
###### Proposition 3.14.
Let $`𝒩(n,d)(n,d)`$ be the set of local minima. Then
$$𝒩(n,d)=\{(E,\mathrm{\Phi })(n,d)|\mathrm{\Phi }=0\}.$$
Hence $`𝒩(n,d)`$ coincides with $`M(n,d)`$, the moduli space of semistable vector bundles of rank $`n`$ and degree $`d`$, which equals the subvariety $`M^s(n,d)M(n,d)`$ corresponding to stable bundles if $`\mathrm{GCD}(n,d)=1`$.
The minima of the Morse function on $`(p,q,a,b)`$ have been characterized in . One has the following.
###### Proposition 3.15.
Let $`𝒩(p,q,a,b)(p,q,a,b)`$ be the set of local minima. Then
$$𝒩(p,q,a,b)=\{(E,\mathrm{\Phi })(p,q,a,b)|\beta =0\text{or}\gamma =0\}.$$
More precisely, let $`(E,\mathrm{\Phi })𝒩(p,q,a,b)`$. Then
1. $`\beta =0`$ if and only if $`a/p>b/q`$ (i.e. $`\tau >0`$).
2. $`\gamma =0`$ if and only if $`a/p<b/q`$ (i.e. $`\tau <0`$).
###### Remark 3.16.
Since we are assuming $`\mathrm{GCD}(p+q,a+b)=1`$ then $`\tau 0`$.
## 4. Homotopy groups
### 4.1. Homotopy groups of $`(n,d)`$
Combining Propositions 3.2, 3.12 and 3.14 we have the following.
###### Theorem 4.1.
Let $`\mathrm{GCD}(n,d)=1`$. Then
$$\pi _i((n,d))\pi _i(M(n,d)),fori2(g1)(n1)2.$$
Now, the homotopy groups of $`M(n,d)`$ have been computed by Daskalopoulos and Uhlenbeck (here $`n`$ and $`d`$ are not assumed to be coprime). Their result is the following.
###### Theorem 4.2.
Let $`M^s(n,d)`$ be the moduli space of stable vector bundles of rank $`n`$ and degree $`d`$. Assume that $`n>1`$ and $`(n,g)(2,2)`$. Then
1. $`\pi _1(M^s(n,d))H_1(X,)`$;
2. $`\pi _2(M^s(n,d))_{\mathrm{GCD}}(n,d)`$;
3. $`\pi _i(M^s(n,d))\pi _{i1}(𝒢),\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}<i2(g1)(n1)2`$, where $`𝒢`$ is the unitary gauge group.
###### Remark 4.3.
The proof of (1) when $`n`$ and $`d`$ are coprime is given by Atiyah–Bott .
As a corollary of Theorems 4.1 and 4.2 we have the following.
###### Theorem 4.4.
Assume that $`n>1`$ and $`\mathrm{GCD}(n,d)=1`$ and let $`g3`$. Then
1. $`\pi _1((n,d))H_1(X,)`$;
2. $`\pi _2((n,d))`$;
3. $`\pi _i((n,d))\pi _{i1}(𝒢),\text{for}\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}<i2(g1)(n1)2`$.
###### Remark 4.5.
As a consequence of Theorem 4.1 and the connectedness of $`M(n,d)`$ one obtains that $`(n,d)`$ is also connected .
A proof of (1) when $`n=2`$ is given by Hitchin .
###### Remark 4.6.
When $`n=2`$ Hausel \[17, Theorem 7.5.7\] proved that the isomorphism (3) holds for $`i4g8`$, which is twice as good as our estimate. It would be very interesting to see if this result can be generalized to higher $`n`$.
### 4.2. Moduli space of triples
The next step is to identify the spaces $`𝒩(p,q,a,b)`$ as moduli spaces in their own right. They turn out to be examples of the moduli spaces of triples studied in and . We briefly recall the relevant definitions and results. See for details.
A *holomorphic triple* on $`X`$, $`T=(E_1,E_2,\varphi )`$, consists of two holomorphic vector bundles $`E_1`$ and $`E_2`$ on $`X`$ and a holomorphic map $`\varphi :E_2E_1`$. Denoting the ranks $`E_1`$ and $`E_2`$ by $`n_1`$ and $`n_2`$, and their degrees by $`d_1`$ and $`d_2`$, we refer to $`(n_1,n_2,d_1,d_2)`$ as the type of the triple.
A homomorphism from $`T^{}=(E_1^{},E_2^{},\varphi ^{})`$ to $`T=(E_1,E_2,\varphi )`$ is a commutative diagram
$$\begin{array}{ccc}E_2^{}& \stackrel{\varphi ^{}}{}& E_1^{}\\ & & & & \\ E_2& \stackrel{\varphi }{}& E_1.\end{array}$$
$`T^{}=(E_1^{},E_2^{},\varphi ^{})`$ is a subtriple of $`T=(E_1,E_2,\varphi )`$ if the homomorphisms of sheaves $`E_1^{}E_1`$ and $`E_2^{}E_2`$ are injective.
For any $`\alpha `$ the *$`\alpha `$-degree* and *$`\alpha `$-slope* of $`T`$ are defined to be
$`\mathrm{deg}_\alpha (T)`$ $`=\mathrm{deg}(E_1)+\mathrm{deg}(E_2)+\alpha \mathrm{rk}(E_2),`$
$`\mu _\alpha (T)`$ $`={\displaystyle \frac{\mathrm{deg}_\alpha (T)}{\mathrm{rk}(E_1)+\mathrm{rk}(E_2)}}`$
$`=\mu (E_1E_2)+\alpha {\displaystyle \frac{\mathrm{rk}(E_2)}{\mathrm{rk}(E_1)+\mathrm{rk}(E_2)}}.`$
The triple $`T=(E_1,E_2,\varphi )`$ is *$`\alpha `$-stable* if
(4.1)
$$\mu _\alpha (T^{})<\mu _\alpha (T)$$
for any proper sub-triple $`T^{}=(E_1^{},E_2^{},\varphi ^{})`$. Define *$`\alpha `$-semistability* by replacing (4.1) with a weak inequality. A triple is called *$`\alpha `$-polystable* if it is the direct sum of $`\alpha `$-stable triples of the same $`\alpha `$-slope. It is strictly $`\alpha `$-semistable (polystable) if it is $`\alpha `$-semistable (polystable) but not $`\alpha `$-stable.
We denote the moduli space of isomorphism classes of $`\alpha `$-polystable triples of type $`(n_1,n_2,d_1,d_2)`$ by
(4.2)
$$𝒩_\alpha =𝒩_\alpha (n_1,n_2,d_1,d_2).$$
Using Jordan–Hölder filtrations of $`\alpha `$-semistable triples one can define $`S`$-equivalence, and view $`𝒩_\alpha `$ as the moduli space of $`S`$-equivalence classes of $`\alpha `$-semistable triples. The isomorphism classes of $`\alpha `$-stable triples form a subspace which we denoted by $`𝒩_\alpha ^s`$.
###### Proposition 4.7 ().
The moduli space $`𝒩_\alpha (n_1,n_2,d_1,d_2)`$ is a complex projective variety. A necessary condition for the moduli space $`𝒩_\alpha (n_1,n_2,d_1,d_2)`$ to be non-empty is
(4.3)
$$\{\begin{array}{cc}0\alpha _m\alpha \alpha _M\hfill & \text{if }n_1n_2\hfill \\ 0\alpha _m\alpha \hfill & \text{if }n_1=n_2\hfill \end{array}$$
where
(4.4) $`\alpha _m`$ $`=\mu _1\mu _2,`$
(4.5) $`\alpha _M`$ $`=(1+{\displaystyle \frac{n_1+n_2}{|n_1n_2|}})(\mu _1\mu _2)`$
and $`\mu _1=\frac{d_1}{n_1}`$, $`\mu _2=\frac{d_2}{n_2}`$.
Whenever necessary we shall indicate the dependence of $`\alpha _m`$ and $`\alpha _M`$ on $`(n_1,n_2,d_1,d_2)`$ by writing $`\alpha _m=\alpha _m(n_1,n_2,d_1,d_2)`$, and similarly for $`\alpha _M`$.
Within the allowed range for $`\alpha `$ there is a discrete set of critical values. These are the values of $`\alpha `$ for which it is numerically possible to have a subtriple $`T^{}=(E_1^{},E_2^{},\varphi ^{})`$ such that $`\mu (E_1^{}E_2^{})\mu (E_1E_2)`$ but $`\mu _\alpha (T^{})=\mu _\alpha (T^{})`$. All other values of $`\alpha `$ are called generic. The critical values of $`\alpha `$ are precisely the values for $`\alpha `$ at which the stability properties of a triple can change, i.e. there can be triples which are strictly $`\alpha `$-semistable, but either $`\alpha ^{}`$-stable or $`\alpha ^{}`$-unstable for $`\alpha ^{}\alpha `$.
The following result relates the stability conditions for holomorphic triples and that for $`\mathrm{U}(p,q)`$-Higgs bundles.
###### Proposition 4.8.
A $`\mathrm{U}(p,q)`$-Higgs bundle $`(E,\mathrm{\Phi })`$ with $`\beta =0`$ or $`\gamma =0`$ is (semi)stable if and only if the corresponding holomorphic triple is $`\alpha `$-(semi)stable for $`\alpha =2g2`$.
Combining Propositions 3.15 and 4.8, we have the following characterization of the subspace of local minima $`𝒩(p,q,a,b)`$.
###### Theorem 4.9.
Let $`𝒩(p,q,a,b)`$ be the subspace of local minima of $`f`$ on $`(p,q,a,b)`$ and let $`\tau `$ be the Toledo invariant.
If $`a/p<b/q`$, or equivalently if $`\tau <0`$, then $`𝒩(p,q,a,b)`$ can be identified with the moduli space of $`\alpha `$-polystable triples of type $`(p,q,a+p(2g2),b)`$, with $`\alpha =2g2`$.
If $`a/p>b/q`$, or equivalently if $`\tau >0`$, then $`𝒩(p,q,a,b)`$ can be identified with the moduli space of $`\alpha `$-polystable triples of type $`(q,p,b+q(2g2),a)`$, with $`\alpha =2g2`$.
That is,
$$𝒩(p,q,a,b)\{\begin{array}{cc}𝒩_{2g2}(p,q,a+p(2g2),b)\hfill & \text{if }a/p<b/q\text{ (equivalently }\tau <0\text{)}\hfill \\ 𝒩_{2g2}(q,p,b+q(2g2),a)\hfill & \text{if }a/p>b/q\text{ (equivalently }\tau >0\text{)}\hfill \end{array}$$
In view of Theorem 4.9 it is important to understand where $`2g2`$ lies in relation to the range (given by Proposition 4.7) for the stability parameter $`\alpha `$. One has the following.
###### Proposition 4.10.
Fix $`(p,q,a,b)`$. Then
(4.6)
$$0|\tau |\mathrm{min}\{p,q\}(2g2)0<\alpha _m(p,q,a,b)2g2\alpha _M(p,q,a,b)\text{if }pq$$
Proposition 4.10 shows that in order to study $`𝒩(p,q,a,b)`$ for different values of the Toledo invariant, we need to understand the moduli spaces of triples for values of $`\alpha `$ that may lie anywhere (including at the extremes $`\alpha _m`$ and $`\alpha _M`$) in the $`\alpha `$-range given in Proposition 4.7.
We can assume $`n_1>n_2`$, since by triples duality one has the following.
###### Proposition 4.11.
$`𝒩_\alpha (n_1,n_2,d_1,d_2)𝒩_\alpha (n_2,n_1,d_2,d_1)`$.
Recall that the allowed range for the stability parameter is $`\alpha _m\alpha \alpha _M`$, where $`\alpha _m=\mu _1\mu _2`$ and $`\alpha _M=\frac{2n_1}{n_1n_2}\alpha _m`$, and we assume that $`\mu _1\mu _2>0`$. We describe the moduli space $`𝒩_\alpha `$ for $`2g2\alpha <\alpha _M`$.
Let $`\alpha _L`$ be the largest critical value in $`(\alpha _m,\alpha _M)`$, and let $`𝒩_L`$ (respectively $`𝒩_L^s`$) denote the moduli space of $`\alpha `$-polystable (respectively $`\alpha `$-stable) triples for $`\alpha _L<\alpha <\alpha _M`$. We refer to $`𝒩_L`$ as the ‘large $`\alpha `$’ moduli space.
###### Proposition 4.12.
Let $`T=(E_1,E_2,\varphi )`$ be an $`\alpha `$-semistable triple for some $`\alpha `$ in the range $`\alpha _L<\alpha <\alpha _M`$. Then $`T`$ is of the form
$$0E_2\stackrel{\mathit{\varphi }}{}E_1F0,$$
with $`F`$ locally free, and $`E_2`$ and $`F`$ are semistable.
In the converse direction we have<sup>4</sup><sup>4</sup>4This Proposition replaces Proposition 7.6 of . We thank Stefano Pasotti, Francesco Prantil and Carlos Tejero for pointing out errors in this Proposition.:
###### Proposition 4.13.
Let $`T=(E_1,E_2,\varphi )`$ be a triple of the form
(4.7)
$$0E_2\stackrel{\mathit{\varphi }}{}E_1F0,$$
with $`F`$ locally free and such that the extension is non-trivial. If $`E_2`$ and $`F`$ are stable then $`T`$ is $`\alpha `$-stable for $`\alpha `$ in the range $`\alpha _L<\alpha <\alpha _M`$.
###### Proof.
Regarding the top line of the diagram
(4.8)
$$\begin{array}{ccc}E_2& & \varphi (E_2)\\ & & & & \\ E_2& \stackrel{\varphi }{}& E_1\\ & & & & \\ 0& & F\end{array}$$
as a subtriple $`T^{}`$ of $`T`$, and the bottom line as a quotient triple $`T^{\prime \prime }`$, we can consider $`T`$ as an extension of triples
$$0T^{}TT^{\prime \prime }0.$$
It follows from stability of $`E_2`$ that the subtriple $`T^{}`$ is $`\alpha `$-stable for any $`\alpha >0`$. It follows similarly from stability of $`F`$, that the quotient triple $`T^{\prime \prime }`$ is $`\alpha `$-stable for any $`\alpha `$. In particular $`T^{}`$ and $`T^{\prime \prime }`$ are both $`\alpha _M`$-stable.
A simple calculation shows that
(4.9)
$$\mu _{\alpha _M}(T^{})=\mu (E_2)+\frac{1}{2}\alpha _M=\mu (F)=\mu _{\alpha _M}(T^{\prime \prime }).$$
It is a general fact that an extension of $`\alpha `$-semistable triples of the same $`\alpha `$-slope is itself $`\alpha `$-semistable. Thus we deduce from (4.9) that the triple $`T`$ is $`\alpha _M`$-semistable.
It remains to show that $`T`$ is $`\alpha `$-stable for $`\alpha _L<\alpha <\alpha _M`$. We do this by showing that there is no $`\alpha `$-destabilizing subtriple, i.e., a subtriple $`S`$ of $`T`$ such that $`\mu _\alpha (S)\mu _\alpha (T)`$.
To do this, we first observe that the extension (4.8) is a Jordan–Hölder filtration of $`T`$ considered as an $`\alpha _M`$-semistable object. This follows since $`T^{}`$ and $`T^{\prime \prime }`$ are $`\alpha _M`$-stable and have the same $`\alpha _M`$-slope. Hence the associated graded object for $`T`$ in the category of $`\alpha _M`$-semistable triples is
(4.10)
$$\mathrm{Gr}(T)=T^{}T^{\prime \prime }.$$
Now assume that $`ST`$ is $`\alpha `$-destabilizing for $`\alpha `$ in the range $`\alpha _L<\alpha <\alpha _M`$. By continuity of $`\mu _\alpha (S)`$ in $`\alpha `$, we have $`\mu _{\alpha _M}(S)\mu _{\alpha _M}(T)`$ and hence $`\alpha _M`$-semistability of $`T`$ implies that $`\mu _{\alpha _M}(S)=\mu _{\alpha _M}(T)`$. It follows that in $`\mathrm{Gr}(T)=T^{}T^{\prime \prime }`$ the triple $`S`$ must be isomorphic to either $`T^{}`$ or $`T^{\prime \prime }`$.
We first show that $`S`$ cannot be isomorphic to $`T^{}`$, i.e, that the subtriple $`T^{}`$ is not destabilizing for $`\alpha <\alpha _M`$. The key piece of information is that
$`{\displaystyle \frac{n_2(T^{})}{n(T^{})}}`$ $`={\displaystyle \frac{n_2}{2n_2}}={\displaystyle \frac{1}{2}},`$
$`{\displaystyle \frac{n_2(T)}{n(T)}}`$ $`={\displaystyle \frac{n_2}{n_1+n_2}}<{\displaystyle \frac{1}{2}},`$
where, for any triple $`T=(E_2,E_1,\varphi )`$, we write $`n_i(T)=\mathrm{rk}(E_i)`$ and $`n(T)=\mathrm{rk}(E_2E_1)`$. Hence
(4.11)
$$\frac{n_2(T^{})}{n(T^{})}>\frac{n_2(T)}{n(T)}>\frac{n_2(T^{\prime \prime })}{n(T^{\prime \prime })}.$$
But, since $`\mu _{\alpha _M}(T^{})=\mu _{\alpha _M}(T)`$,
$$\mu _\alpha (T^{})\mu _\alpha (T)=(\alpha \alpha _M)\left(\frac{n_2(T^{})}{n(T^{})}\frac{n_2}{n}\right)<0$$
for $`\alpha <\alpha _M`$.
Finally we show that $`S`$ cannot be isomorphic to $`T^{\prime \prime }`$. In fact, if $`T`$ has a subtriple isomorphic to $`T^{\prime \prime }`$, then $`E_1`$ has a subbundle, $`\stackrel{~}{F}`$, isomorphic to $`F`$. The composition of the isomorphism from $`F`$ to $`\stackrel{~}{F}`$ with the projection from $`E_1`$ to $`F`$ produces a homomorphism
$$\psi :FF.$$
Since $`F`$ is stable, $`\psi `$ is either zero or a multiple of the identity. If it is zero, then there must be a non-trivial homomorphism from $`F`$ to $`E_2`$. This is impossible since $`\mu (\stackrel{~}{F})>\mu (E_2)`$, and both are stable bundles. Hence the the isomorphism from $`F`$ to $`\stackrel{~}{F}`$ splits the extension (4.7). But this contradicts our assumptions. ∎
###### Theorem 4.14.
Assume that $`n_1>n_2`$ and $`d_1/n_1>d_2/n_2`$. Then the moduli space $`𝒩_L^s=𝒩_L^s(n_1,n_2,d_1,d_2)`$ is smooth of dimension
$$(g1)(n_1^2+n_2^2n_1n_2)n_1d_2+n_2d_1+1,$$
and includes a $`^N`$-fibration $`𝒫`$ over $`M^s(n_1n_2,d_1d_2)\times M^s(n_2,d_2)`$, where $`M^s(n,d)`$ is the moduli space of stable bundles of rank $`n`$ and degree $`d`$, and $`N=n_2d_1n_1d_2+n_1(n_1n_2)(g1)1`$. Moreover, the complex codimension of $`𝒩_L^s𝒫`$ is at least $`g1`$. In particular, $`𝒩_L^s(n_1,n_2,d_1,d_2)`$ is non-empty and irreducible.
If $`\mathrm{GCD}(n_1n_2,d_1d_2)=1`$ and $`\mathrm{GCD}(n_2,d_2)=1`$, then $`𝒩_L^s(n_1,n_2,d_1,d_2)`$ is isomorphic to $`𝒫`$.
###### Proof.
The birational equivalence between $`𝒫`$ and $`𝒩_L^s`$ is proved in . To obtain the precise estimate of the codimension of $`𝒩_L^s𝒫`$ in $`𝒩_L^s`$ we see that, by Proposition 4.12, it suffices to estimate the dimension of stable triples like (4.7) with $`E_2`$ and $`F`$ semistable.
Now, for any family of semistable bundles the complex codimension of the set of strictly semistable bundles is at least $`g1`$. A computation of the precise estimate can be found in . The proof is finished by observing that for a stable triple of the form (4.7) $`H^0(X,E_2F^{})=0`$ (see ). ∎
The following is proved in .
###### Theorem 4.15.
Let $`\alpha `$ be any value in the range $`\alpha _m<2g2\alpha <\alpha _M`$. Then $`𝒩_\alpha ^s`$ is birationally equivalent to $`𝒩_L^s`$. Moreover, they are isomorphic outside of a set of complex codimension greater or equal than $`g1`$. In particular, $`𝒩_\alpha ^s`$ is non-empty and irreducible.
### 4.3. Homotopy groups of moduli spaces of triples
The strategy to compute the homotopy groups of $`𝒩(p,q,a,b)`$ is to compute first those of the moduli space of $`\alpha `$-stable triples $`𝒩_\alpha ^s`$ for large $`\alpha `$.
Let $`n_1>n_2`$ and let $`𝒫𝒩_L`$ be the $`^N`$-fibration over $`M^s(n_1n_2,d_1d_2)\times M^s(n_2,d_2)`$ given in Theorem 4.14. As a consequence of Theorems 4.14 and 4.15 we have the following.
###### Proposition 4.16.
Let $`2g2\alpha <\alpha _M`$. Then
$$\pi _i(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))\pi _i(𝒩_L^s(n_1,n_2,d_1,d_2))\pi _i(𝒫)\text{for}i2g4.$$
Associated to the $`^N`$-fibration $`𝒫`$ over $`M^s(n_1n_2,d_1d_2)\times M^s(n_2,d_2)`$ there is a homotopy sequence
(4.12)
$$\begin{array}{c}\mathrm{}\pi _i(^N)\pi _i(𝒫)\pi _i(M^s(n_2,d_2))\times \pi _i(M^s(n_1n_2,d_1d_2))\hfill \\ \hfill \pi _{i1}(^N)\mathrm{}\end{array}$$
###### Proposition 4.17.
Let $`n_1>n_2`$ and $`n_2d_1>n_1d_2`$. Assume that $`(n_2,g)(2,2)`$ and $`(n_1n_2,g)(2,2)`$ (for our applications we will actually assume $`g3`$). Then
1. $`\pi _1(𝒫)\pi _1(M^s(n_2,d_2))\times \pi _1(M^s(n_1n_2,d_1d_2))H_1(X,)H_1(X,)`$;
2. $`\pi _2(𝒫)`$ is the middle term in an exact sequence
$$0\pi _2(𝒫)Q(n_1,n_2,d_1,d_2)0$$
where
$$Q(n_1,n_2,d_1,d_2)=\{\begin{array}{cc}_{\mathrm{GCD}(n_2,d_2)}_{\mathrm{GCD}(n_1n_2,d_1d_2)}\hfill & \text{if }n_2>1\text{ and }(n_1n_2)>1\hfill \\ _{\mathrm{GCD}(n_2,d_2)}\hfill & \text{if }n_2>1\text{ and }(n_1n_2)=1\hfill \\ _{\mathrm{GCD}(n_1n_2,d_1d_2)}\hfill & \text{if }n_2=1\text{ and }(n_1n_2)>1\hfill \\ 0\hfill & \text{if }n_2=1\text{ and }n_1=2\hfill \end{array}$$
###### Remark 4.18.
It follows immediately from the exact sequence in (2) of Proposition 4.17 that the free part of the finitely generated abelian group $`\pi _2(𝒫)`$ equals the direct sum of $``$ and the free part of $`Q(n_1,n_2,d_1,d_2)`$. In particular, if the co-primality conditions $`\mathrm{GCD}(n_2,d_2)=1`$ and $`\mathrm{GCD}(n_1n_2,n_2d_2)=1`$ hold, then we have a complete description of $`\pi _2(𝒫)`$ as the direct sum $`Q(n_1,n_2,d_1,d_2)`$. Also, under any circumstances, it follows that $`\pi _2(𝒫)Q(n_1,n_2,d_1,d_2)`$ so, for rational homotopy, our results are complete.
###### Proof of Proposition 4.17.
From the homotopy sequence (4.12), since $`\pi _0(^N)=\pi _1(^N)=0`$, we deduce that $`\pi _1(𝒫)\pi _1(M^s(n_2,d_2))\times \pi _1(M^s(n_1n_2,d_1d_2))`$. Statement (1) follows from Theorem 4.2.
Since $`\pi _1(^N)=0`$, (4.12) gives
(4.13)
$$\mathrm{}\pi _2(^N)\pi _2(𝒫)\pi _2(M^s(n_2,d_2))\times \pi _2(M^s(n_1n_2,d_1d_2))0.$$
On the other hand, by Hurewicz’ theorem $`\pi _2(^N))H_2(^N,)`$. Now, the map $`f:\pi _2(^N)\pi _2(𝒫)`$ in (4.13) is injective since one has the commutative diagram
$$\begin{array}{ccc}\pi _2(^N)& \stackrel{f}{}& \pi _2(𝒫)\\ & & & & \\ H_2(^N,)& & H_2(𝒫,).\end{array}$$
and $`H_2(^N,)H_2(𝒫,)`$ must be injective because the restriction of an ample line bundle over $`𝒫𝒩_L`$ to $`^N`$ must give an ample line bundle. Note that the natural map $`H_2(^N)H_2(𝒩_L)`$ is injective and factors through $`H_2(^N)H_2(𝒫)H_2(𝒩_L)`$. Now, we obtain (2) from Theorem 4.2 and the fact that if $`n=1`$ then the moduli space $`M^s(n,d)`$ is the Jacobian of degree $`d`$ line bundles. ∎
As a corollary of Proposition 4.16 and Proposition 4.17 we have the following.
###### Theorem 4.19.
Assume $`n_1>n_2`$, $`n_2d_1>n_1d_2`$, $`g3`$, and $`2g2\alpha <\alpha _M`$. Then
1. $`\pi _1(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))H_1(X,)H_1(X,)`$;
2. $`\pi _2(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))`$ is the middle term in an exact sequence
$$0\pi _2(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))Q(n_1,n_2,d_1,d_2)0$$
where
$$Q(n_1,n_2,d_1,d_2)=\{\begin{array}{cc}_{\mathrm{GCD}(n_2,d_2)}_{\mathrm{GCD}(n_1n_2,d_1d_2)}\hfill & \text{if }n_2>1\text{ and }(n_1n_2)>1\hfill \\ _{\mathrm{GCD}(n_2,d_2)}\hfill & \text{if }n_2>1\text{ and }(n_1n_2)=1\hfill \\ _{\mathrm{GCD}(n_1n_2,d_1d_2)}\hfill & \text{if }n_2=1\text{ and }(n_1n_2)>1\hfill \\ 0\hfill & \text{if }n_2=1\text{ and }n_1=2\hfill \end{array}$$
###### Remark 4.20.
Theorem 4.19 gives a complete description of $`\pi _2(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))`$ when the co-primality conditions $`\mathrm{GCD}(n_2,d_2)=1`$ and $`\mathrm{GCD}(n_1n_2,n_2d_2)=1`$ hold, and of $`\pi _2(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))`$ under all circumstances (cf. Remark 4.18).
Using the results of , a complete description of $`\pi _2(𝒩_\alpha ^s(n_1,n_2,d_1,d_2))`$ can also be given in the case when $`n_2=1`$ and $`n_1n_2>1`$, as we now explain. In that paper, the moduli space of *stable pairs* $`(V,\varphi )`$ was studied. Here $`V`$ is a vector bundle and $`\varphi H^0(X,V)`$ is a holomorphic section of $`V`$. Viewing the section $`\varphi `$ as a map $`\varphi :𝒪V`$, a pair $`(V,\varphi )`$ gives rise to a triple $`(E_1,E_2,\varphi )=(V,𝒪,\varphi )`$. Through this correspondence it can be seen that the moduli space of triples $`𝒩_\alpha ^s`$ of triples with $`n_2=1`$ fibres over the Jacobian variety of the curve, with fibres isomorphic to the moduli space of pairs. Among other things, in the second homotopy group of the moduli space of pairs was calculated to be $``$ for $`\alpha `$ between $`\alpha _m`$ and the first critical value of $`\alpha `$ larger than $`\alpha _m`$. It follows from these results that, when $`d_2=1`$, one has $`\pi _2(𝒩_\alpha ^s)=`$ for such $`\alpha `$. Combining this fact with Proposition 4.16 it follows that $`\pi _2(𝒩_\alpha ^s)=`$ for $`2g2\alpha <\alpha _M`$ in the case $`n_2=1`$ and $`n_1n_2>1`$.
### 4.4. Homotopy groups of $`(p,q,a,b)`$
Combining Propositions 3.2, 3.12 and 3.15 we have the following.
###### Theorem 4.21.
Let $`\mathrm{GCD}(p+q,a+b)=1`$. Then
$$\pi _i((p,q,a,b))\pi _i(𝒩(p,q,a,b)),fori2g4.$$
As a corollary of Theorems 4.21, 4.9 and 4.19 and Proposition 4.11, we conclude the following.
###### Theorem 4.22.
Let $`pq`$ and $`\mathrm{GCD}(p+q,a+b)=1`$ and let $`g3`$. Then
1. $`\pi _1((p,q,a,b))H_1(X,)H_1(X,)`$;
2. $`\pi _2((p,q,a,b))`$ is the middle term in an exact sequence
$$0\pi _2((p,q,a,b))Q(n_1,n_2,d_1,d_2)0,$$
where
$$(n_1,n_2,d_1,d_2)=\{\begin{array}{cc}(p,q,a+p(2g2),b)\hfill & \text{if }\tau <0\text{ and }p>q\hfill \\ (q,p,b,ap(2g2))\hfill & \text{if }\tau <0\text{ and }p<q\hfill \\ (p,q,a,bq(2g2))\hfill & \text{if }\tau >0\text{ and }p>q\hfill \\ (q,p,b+q(2g2),a)\hfill & \text{if }\tau >0\text{ and }p<q\hfill \end{array}$$
and where $`Q(n_1,n_2,d_1,d_2)`$ is as in Theorem 4.19.
###### Remark 4.23.
Theorem 4.22 gives a complete description of $`\pi _2((p,q,a,b))`$ when the co-primality conditions $`\mathrm{GCD}(n_2,d_2)=1`$ and $`\mathrm{GCD}(n_1n_2,n_2d_2)=1`$ hold, and of $`\pi _2((p,q,a,b))`$ $``$ under all circumstances (cf. Remarks 4.18 and 4.20).
###### Remark 4.24.
As a consequence of Theorem 4.21 and the connectedness of $`𝒩(p,q,a,b)`$ we have that $`(p,q,a,b)`$ is also connected, as proved in .
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# The Second Quantized Quantum Turing Machine and Kolmogorov Complexity
## Abstract
The Kolmogorov complexity of a physical state is the minimal physical resources required to reproduce that state. We define a second quantized quantum Turing machine and use it to define second quantized Kolmogorov complexity. There are two advantages to our approach – our measure of second quantized Kolmogorov complexity is closer to physical reality and unlike other quantum Kolmogorov complexities it is continuous. We give examples where second quantized and quantum Kolmogorov complexity differ.
Quantum physics, as far as we know, is the most accurate and universal description of all physical phenomena in the universe. If we therefore wish to speak about the complexity of some processes or some physical states in nature, we need to use quantum physics as our best available theory. An important question is (in very simple terms), given a physical system in some state, how difficult is it for us to reproduce it. If we wish to have a universal measure of this difficulty (which applies to all systems and states) a way to proceed is to follow the prescription of Kolmogorov. First, we define a universal computer (which is capable of simulating all other computers) and then we look for the shortest input (another physical state) to this computer that reproduces as the output the desired physical state. This way, all the complexity is defined with respect to the same universal computer, and we can thus speak about universal complexity. The universal computer (such as a universal Turing machine) needs to be fully quantum mechanical in order to capture all the possibilities available in nature.
In this letter, we define a fully quantized Turing machine which we use to define a fully quantum Kolmogorov complexity and apply it to a number of different problems. Our approach is different from others in that we consider indeterminate length input programs whose expected length is our measure of complexity. Others, whose work is discussed in detail below, have perhaps avoided our approach, as allowing programs in superposition can lead to a superposition of halting times. Since traditionally computation is viewed as giving a deterministic output after a fixed amount of time, having a superposition of halting times seems to contradict the very concept of computation. We, on the other hand, view the superposition of different length inputs which can lead to a superposition of halting times as a necessity dictated by a fully quantum mechanical description of nature. Moreover allowing superpositions of different programs can lead to a program which has on average a smaller Kolmogorov complexity than when programs are restricted to having a variable but determinate length (we give examples of this in this letter). Since we want to know what is physically the shortest input which produces a given output, we need allow the computation with superpositions of different length programs.
This letter is organized as follows. First we discuss the concept of a Turing machine and review previous work. Then we define the fully quantum Turing machine (we call it second quantized for reasons that will become apparent below). We use this second quantized Turing machine (SQTM) to define the concept of second quantized Kolmogorov complexity (SQKC). This notion is then tested by applying it to various simple examples such as the average length of two programs of different sizes and programs which halt at a superposition of different times. We consider conditional complexity and show that multiple copies of a quantum state can be compressed asymptotically further using SQTM than on a standard quantum Turing machine (QTM) given the number of copies $`n`$.
Quantum information is usually studied in $`n`$-particle systems described by an $`n`$-fold tensor product Hilbert space $`H^n`$ in which $`n`$ is fixed. In the second quantization $`n`$, the number of particles, can exist in superposition. The corresponding space which describes a system in the second quantization is the Fock space $`H^{}=_{n=0}^{\mathrm{}}H^n`$.
A physical example of such a system is the quantized electromagnetic field. The second quantized electromagnetic field consists of a number of modes $`\omega _i`$ (labelled by frequency and polarization) which can each be populated by a number of photons $`m_i+n_i`$, $`m_i`$ in the horizontal polarization $`H`$ and $`n_i`$ in the vertical polarization $`V`$. The state of the system can be written as:
$$|\psi =|m_1_{\omega _1,H}|n_1_{\omega _1,V}|m_2_{\omega _2,H}|n_2_{\omega _2,V}\mathrm{}$$
(1)
For the purposes of this letter, we will encode qubits into the polarization degree of freedom, which then restricts the number $`m_i+n_i`$ of quanta in the modes to either $`0`$ (the vacuum state) or $`1`$. The modes are ordered in some predefined way, and the initial contiguous sequence which are occupied by a single photon before the first vacuum state represents a quantum string. For example, the state
$`|\psi `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|1_{\omega _1,H}|0_{\omega _1,V}|0_{\omega _2,H}|1_{\omega _2,V}`$ (2)
$``$ $`\left((|1_{\omega _3,H}|0_{\omega _3,V})+(|0_{\omega _3,H}|1_{\omega _3,V})\right)`$
$``$ $`|0_{\omega _4,H}|0_{\omega _4,V}\mathrm{}`$ (3)
can be written as a quantum string $`|\psi =|01(|0+|1)/\sqrt{2}`$, where a photon in the horizontal or vertical polarizations represents a $`|0`$ or $`|1`$ respectively. If a photon exists at a frequency in superposition, the quantum string may be in a superposition of two different lengths. For example,
$`|\psi `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|1_{\omega _1,H}|0_{\omega _1,V}|0_{\omega _2,H}`$ (4)
$``$ $`(|0_{\omega _2,V}+|1_{\omega _2,V})|0_{\omega _3,H}|0_{\omega _3,V}`$
$``$ $`\mathrm{}`$ (5)
which can be written as the quantum string $`|\psi =(|01+|0)/\sqrt{2}`$. Note that the superposition of different program lengths immediately implies a superposition of different times of computation. If a head moves along the string from left to right, it will reach the end of the string at different times directly proportional to the length. Since variable length encoding naturally leads to the concept of creation and annihilation of particles Rallan and Vedral (2003), we call the Turing machine described in this letter the Second Quantized Turing machine (SQTM).
Now we briefly review similar work. Svozil Svozil (1996) originally defined quantum Kolmogorov complexity on a circuit based model and then Berthiaume, van Dam and Laplante Berthiaume et al. (2001) defined quantum Kolmogorov complexity Li and Vitanyi (1997) on a quantum Turing machine. Svozil (1996) and Berthiaume et al. (2001) both restricted inputs to be of variable but determined length (i.e. quantum but not second quantized inputs). Vitanyi Vitanyi (2001) also provided a definition of quantum Kolmogorov complexity based on classical descriptions, where a penalty factor was added depending on the accuracy of the classical description. Mora and Briegel Mora and Briegel (2006, 2005) defined the algorithmic complexity of a quantum state to be the classical Kolmogorov complexity of describing a circuit that produces the quantum state (up to an error parameter $`ϵ`$).
Gacs Gacs (2001) and Tadaki Tadaki defined quantum Kolmogorov complexity without reference to a computation device by generalizing classical universal semi-measures to quantum universal semi-POVM’s. Tadaki Tadaki (2002) went on to derive $`\mathrm{\Omega }_Q`$, a quantum generalization of Chaitin’s halting probability Chaitin (1987). More recently, Benatti et al Benatti et al. (2006) have given a quantum Brudno theorem and Müller Mueller (2006) has given a detailed proof of the invariance theorem for Berthiaume et al’s complexity Berthiaume et al. (2001).
The standard models of quantum Turing machines were defined by Deutsch Deutsch (1985) and Berstein and Vazirani Bernstein and Vazirani (1997). Nielsen Nielsen and Chuang (1997) also defined a programmable quantum gate array. Nishimura and Ozawa Nishimura and Ozawa (2002) compared various models of quantum computation. Various papers Myers (1997); Ozawa (1998); Shi (2002); Miyadera and Ohya (2002) have discussed how these models can halt coherently (all computation paths halting simultaneously) when two programs which halt at different times are used as input in superposition. Miyadera and Ohya (2002) showed that it is undecidable whether a quantum Turing machine halts coherently for a general input.
Schumacher and Westmoreland Schumacher and Westmoreland (2001), Boström and Felbinger Bostroem and Felbinger (2002) and Rallan and Vedral Rallan and Vedral (2003) studied lossless quantum compression with variable length quantum strings. One of the authors Rogers and Nagarajan studied exact lossless quantum compression and universal lossless quantum compression with indeterminate length quantum strings.
Now we discuss the SQTM. The aim of developing an SQTM is to fully quantize every aspect of the quantum Turing machine (QTM) to produce a fully quantum computational model of nature. Unlike a QTM, the SQTM’s definition makes it explicit that it can hold second quantized states on the tape(s) and so that it can halt at a superposition of different times.
The contents of the tape including the input and output is allowed to exist in the second quantization. States such as $`|\psi =\frac{1}{\sqrt{2}}(|0+|10)`$ are forbidden on the QTM but allowed on the SQTM.
Several papers Myers (1997); Ozawa (1998); Shi (2002); Miyadera and Ohya (2002) have pointed out issues in how quantum computations halt coherently (how all the computational paths halt simultaneously) on the quantum Turing machine Bernstein and Vazirani (1997); Deutsch (1985) and Nielsen’s parallel gate array Nielsen and Chuang (1997). The SQTM allows programs of different sizes to be input in superposition. These programs may halt at different times, however as natural processes may also halt at a superposition of different times, the SQTM allows this. Informally, the halting condition says that the SQTM comes to a halt when all the computation paths of the working tape stop evolving.
The SQTM is made up of a two-way infinite working tape, an environment tape, a transition function, a head which points to the current cell being processed (perhaps in superposition) and a current instruction being executed (perhaps also in superposition). Each cell in the tape contains either a $`|0`$, $`|1`$ or a vacuum state. The string held on the tape is made up of the $`|0`$’s and $`|1`$’s before the first vacuum state which can exist in a superposition of positions.
The head of the Turing machine contains two quantum states. The first $`|C`$, contains the current superposition of cells being processed, the second $`|I`$, the current instruction being executed. The transition of the Turing machine from one state to the next is controlled by the transition function $`\delta `$ which can act on the current and neighboring cells. At each step, the head can move left, right or the current cell being scanned can be modified and the internal state of the head can be changed with the restriction that the whole operation of the Turing machine is restricted to be a unitary operation. This unitary operation is computable and can be computably written down classically. Using this classical description, there exists a universal SQTM capable of simulating any other SQTM with a constant length description.
The environment tape is defined in the same way as the working tape, except that it receives no input and gives no output. The SQTM can continue to carry out computations on the environment tape (for example, counting up to infinite) to ensure that its operation is unitary. The proposed halting scheme for the SQTM (though others may be possible) is as follows. On input $`|\psi _I`$, the SQTM is said to halt in the state $`|\psi _F`$ if the contents of the working tape converges to $`|\psi _F`$ computably, i.e. if there exists a computable sequence of integers $`t_1`$, $`t_2`$, $`\mathrm{}`$ such that for all $`tt_i`$ steps, the SQTM has executed $`t`$ steps and the working tape of the SQTM (with the state of the other parts of the SQTM traced out) is in state $`T_{SQ}^t(|\psi _I)`$ where $`|T_{SQ}^t(|\psi _I)|\psi _F|1/2^i`$. The restriction of the sequence $`t_i`$ to be computable allows the SQTM to be simulated by some deterministic classical Turing machine to compute a classical description of $`|\psi _F`$ with arbitrary accuracy, ensuring that the SQTM halts in a computable state when given computable input. The halting condition also guarantees that there exists some SQTM which on input $`|i`$ will halt in the state $`|\psi _F`$ with a fidelity of $`1/2^i`$.
The classical Kolmogorov complexity of a string is defined as the shortest classical description of that string with respect to a universal Turing machine. Kolmogorov complexity gives a universal measure of the complexity of a string and has many applications such as the incompressibility method, data compression, universal induction and absolute information. Following from classical Kolmogorov complexity, there are several possible ways to define second quantized Kolmogorov complexity (SQKC) Kolmogorov (1965); Li and Vitanyi (1997).
The combinatorial approach is to study the number of dimensions used to described a quantum state $`|\psi `$. The combinatorial approach is similar to the restriction of input to be variable but determined length, which is used by Svozil Svozil (1996) and then Berthiaume et al Berthiaume et al. (2001) in their definitions of quantum Kolmogorov complexity.
This approach is to define Kolmogorov complexity in terms of probability theory avoiding the need to refer to a Turing machine. Tadaki Tadaki ; Tadaki (2002) and Gacs Gacs (2001) took this approach for quantum Kolmogorov complexity.
The algorithmic approach is to consider the complexity of a string as the length of the shortest algorithm or description which describes that string, the idea being that simpler strings can be described by shorter descriptions. For example, $`x=1111111111`$ is very simple and regular and has a very short description as “ten ones” whereas $`y=1000111010`$ is apparently much more random and it seems difficult to find a short description of $`y`$. A second quantized description may be in a superposition of many different lengths. It is the average length of the description which is used to measure the second quantized Kolmogorov complexity of a string giving the average physical resources used to describe the state.
We believe the average length is the appropriate measure of complexity because it in some sense corresponds to the average energy of the input Rallan and Vedral (2003) and energy is frequently the most crucial physical resource. Unlike the combinatorial approach, it leads to a continuous measure of complexity which is a much more natural physical quantity. We now study the algorithmic approach and by SQKC we mean the average length complexity used in the algorithmic approach.
Let $`|\psi =_{i=0}^{\mathrm{}}\alpha _i|i`$ be a second quantized state. The average length Bostroem and Felbinger (2002) $`\overline{l}`$ of $`|\psi `$ is
$`\overline{l}(|\psi )`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}|\alpha _i|^2l(i)`$ (6)
The SQKC of a state $`|\psi `$ is
$`C_{SQ}(|\psi )`$ $`=`$ $`\underset{U_{SQ}(|\varphi )=|\psi }{inf}\overline{l}(|\varphi )`$ (7)
where $`U_{SQ}`$ is a universal SQTM. The SQKC $`C_{SQ}`$ is invariant up to an additive constant factor (this invariance theorem can be proved in the same way as the classical invariance theoremLi and Vitanyi (1997)).
A simple example of SQKC, which is different from QKC, is to try to find the shortest string which can describe the state $`|\psi _n=(\sqrt{2}|0+|1^n)/\sqrt{3}`$ for $`n>0`$. A shortest program to describe $`|\psi _n`$ is $`|p_{\psi _n}=(\sqrt{2}|0+|n)/\sqrt{3}`$ for $`n>0`$. $`|p_{\psi _n}`$ can be input into an SQTM $`T_{SQ}`$ defined by $`T_{SQ}|n=|1^n`$ and $`T_{SQ}|0=|0`$ to produce $`|\psi _n`$ as output. The SQKC of $`|\psi _n`$ is then at most the average length of $`|p_{\psi _n}`$ plus a constant factor $`O(1)`$ which comes from describing $`T_{SQ}`$ on the universal machine $`U_{SQ}`$. The SQKC of $`|\psi _n`$ is at most
$`C_{SQ}(|\psi _n)`$ $``$ $`\overline{l}(|\varphi _{\psi _n})+O(1)`$ (8)
$`=`$ $`{\displaystyle \frac{2}{3}}O(1)+{\displaystyle \frac{1}{3}}(\mathrm{log}(n)+O(1))`$ (9)
$`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{log}(n)+O(1)`$ (10)
On the other hand, the quantum Kolmogorov complexity of $`1^n`$ is $`\mathrm{log}(n)+O(1)`$ which is the number of qubits used in a fixed length description of $`|\psi _n`$. Thus the SQKC of $`|\psi _n`$ is multiplicatively smaller than the quantum Kolmogorov complexity of $`|\psi _n`$. Notice also that if the amplitude $`1/\sqrt{3}`$ is changed to $`\alpha `$ and $`|\alpha |0`$ then $`C_{SQ}(\psi _n)C_{SQ}(0)`$, which is not true of the quantum Kolmogorov complexity defined by Berthiame et al Berthiaume et al. (2001).
The $`|\psi _n`$ defined above can also be used to show how halting at a superposition of different times is used in SQKC. As above, let $`|\psi _n=(\sqrt{2}|0+|1^n)/\sqrt{3}`$. Suppose that an SQTM can compute $`|0`$ in $`\tau `$ steps. For sufficiently large $`n`$, $`|1^n`$ takes more than $`\tau `$ steps to compute. Thus, for large $`n`$, the SQTM must be allowed to halt at a superposition of different times to describe $`|\psi _n`$. Thus it is by allowing the SQTM to halt at a the number of steps in superposition that we can describe states which are in a superposition of different lengths.
The conditional SQKC of a string $`|\psi `$ given $`|\varphi `$ is the complexity of $`|\psi `$ given $`|\varphi `$ as input to the universal SQTM.
$$C_{SQ}(|\psi ||\varphi )=inf_{U(|\varphi ,p)}\overline{l}(p)$$
(11)
The SQKC of $`n`$ copies of a state $`|\psi `$ assuming that $`n`$ is known in advance is also a non-trivial example of how SQKC differs from QKC. Taking $`|\psi =\alpha |0+\beta |1`$, $`|\psi ^n`$ can be expanded out into the symmetric basis as
$`|\psi ^n`$ $`=`$ $`(\alpha |0+\beta |1)^n`$ (12)
$`=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}\alpha ^i\beta ^{ni}S(i,n)`$ (13)
where $`S(i,n)`$ is the sum of strings containing $`i`$ $`0`$’s and $`ni`$ $`1`$’s.
Berthiaume et al Berthiaume et al. (2001) showed that $`|\psi ^n`$ can be described using $`\mathrm{log}(n)`$ qubits using a fixed length code to encode each $`S(i,n)`$. On the other hand, using a variable length code to encode $`S(i,n)`$ as a string of $`\mathrm{log}\left(\genfrac{}{}{0pt}{}{n}{i}\right)`$ qubits, $`|\psi ^n`$ can be encoded using $`S(\rho _\psi )`$ qubits asymptotically where $`S`$ is the von Neumann entropy and $`\rho _\psi =_i|\alpha ^i\beta ^{ni}|^2|S(i,n)S(i,n)|`$ is the state of $`|\psi `$ after a measurement has been carried out in the $`\{S(i,n)\}_i`$ basis. Thus as $`n`$ grows, the number of bits used to describe $`|\psi ^n`$ with an SQTM can be far fewer than are used by a QTM.
We have defined a second quantized Turing machine which formally models information processing in second quantized physical systems and addresses how computational processes halt (the halting of the SQTM). We have used this SQTM to define second quantized Kolmogorov complexity which improves on the previous measure, quantum Kolmogorov complexity, in that it is continuous and therefore a much more natural physical quantity.
Other measures of information such as von Neumann entropy and its variants are widely used for studying various properties of quantum systems such as entanglement, distances between systems, etc. These other information measures can be seen as computable approximations of (second quantized) Kolmogorov complexity. Kolmogorov complexity is much more powerful. For instance, if there is some underlying mechanical description of a set of experiments (e.g. the Bohm interpretation) then the measurement outcomes of the experiment may have a lower Kolmogorov complexity than if the measurement outcomes are completely random while the entropy of the outcomes are the same in cases.
There are many other possible applications of the SQTM. One example is Maxwell’s Demon. An attempt was made (and subsequently resolved) to show that the second law can be violated in the classical setting Lloyd (1997); Caves (1990); Zurek (1999), using the fact that a sequence of particles might have a very compact description while their joint entropy is high. Maxwell’s demon can be considered in its fully quantum setting with the help of the SQTM. As we have argued SQKC can be smaller than QKC, which may help the demon in trying to violate the second law in the quantum case. This question is left for future research.
In this paper, we have studied a few examples of SQKC. Classical Kolmogorov complexity is a well-studied area with a plethora of applications in many areas of physics. We hope that the future development of SQKC will lead to a plethora of tools for studying quantum systems and that the differences in classical and second quantized Kolmogorov complexity will yield further insights into the differences between classical and quantum systems.
The authors thank Charles Bennett, Johann Summhammer, Karl Svozil and Andreas Winter for useful discussion. The authors thank the EPSRC for financial support. V.V. also thanks the European Union and the British Council in Austria.
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# PRE-MULTISYMPLECTIC CONSTRAINT ALGORITHM FOR FIELD THEORIES
## 1 Introduction
Systems of singular differential equations have been a matter of increasing interest, especially during the last 30 years, and they have been studied separately in theoretical physics and in some technical areas such as engineering of electric networks or control theory. The fundamental characteristic of these kinds of systems is that the existence and uniqueness of solutions are not assured.
In particular, this situation arises in mechanics when dynamical systems described by singular Lagrangians are considered. Furthermore, these systems do not have a nice Hamiltonian description, since not all the momenta are available, and there is a submanifold of the momentum phase space where, in general, the dynamical equations have no solution everywhere. The same problems arise when considering systems of PDE’s associated with field theories described by singular Lagrangians (indeed, many field theories are singular, for instance electromagnetism), as well as in some other applications related with optimal control theories.
Dirac was pioneering in solving the problem for the Hamiltonian formalism of singular mechanical systems, by developing a constraint algorithm which gives, in the favourable cases, a final constraint submanifold where admissible solutions to the dynamics exist (in the sense that the dynamical evolution remains on this manifold). Dirac’s main aim was to apply this procedure to field theories. After Dirac, a lot of work was done in order to geometrize his algorithm. The first important step was the work by Gotay et al , and its application to the Lagrangian formalism . Other algorithms were given later, in order to find consistent solutions of the dynamical equations in the Lagrangian formalism of singular systems (including the problem of finding holonomic solutions) , and afterwards, new geometric algorithms were developed to be applied both in the Hamiltonian and the Lagrangian formalisms .
The Lagrangian and Hamiltonian descriptions of field theories, termed the multisymplectic approach, is the natural extension of time-dependent mechanics. Therefore, in order to understand the constraint algorithm for field theories in a covariant formalism, the first step was to develop the algorithmic procedures for time-dependent systems. This work was provided in . A basic geometric study of these systems can be found in . Furthermore, a qualitative description of constraint algorithms for field theories was made in .
Working within the framework of the multisymplectic description for these theories, we present in this paper a geometric algorithm for finding the maximal submanifold where there are consistent solutions to the field equations of singular theories. This algorithm gives an intrinsic description of all the constraint submanifolds. The problem is stated in a generic pre-multisymplectic fibre bundle, in order to give a solution to both Lagrangian and Hamiltonian field theories, as well as other possible kinds of systems of partial differential equations. In this framework, the solutions to these equations are given geometrically by integrable connections or, what is equivalent, by integrable locally decomposable $`m`$-vector fields which are transverse to the fibre projection. The key point consists in using an auxiliar connection for constructing different geometrical structures needed to develop the algorithm, by following the same methods introduced in for time-dependent singular systems. This technique (the use of a connection) was used for the first time in , in order to obtain (global) Hamiltonian functions, and afterwards applied both in the Lagrangian and Hamiltonian formalisms for this and other purposes (see ). An exhaustive use of this technique in mechanics and field theory can be found in .
First, the problem is reduced to another in the realm of linear algebra, and solved in this context, and then the results are applied to the general pre-multisymplectic framework. In this way, a constraint algorithm can be developed giving a sequence of submanifolds which, in the best case, ends in some final constraint submanifold where field equations have consistent solutions (connections or $`m`$-vector fields), although not necessarily integrable. The problem of integrability is considered and solved separately. Finally, Lagrangian and Hamiltonian field theories are particular cases where the above results are applied straightforwardly, although in the Lagrangian case the problem of finding holonomic solutions must be also analized.
The paper is organized as follows:
First, in Section 2, we state and solve the algebraic version of the problem. Then, in Section 3, we pose the general problem in the context of a pre-multisymplectic fiber bundle and, applying the results obtained in the previous Section, the solution is achieved after studying the additional problem of integrability. After this, Section 4 is devoted to giving the application to Lagrangian and Hamiltonian field theories, including the problem of finding holonomic solutions in the Lagrangian formalism. Finally, as a classical example, field theories described by affine Lagrangians are analyzed in Section 5. An Appendix about multivector fields and connections is also included, in order to make the paper more self-contained and readable.
Manifolds are real, paracompact, connected and $`C^{\mathrm{}}`$. Maps are $`C^{\mathrm{}}`$. Sum over crossed repeated indices is understood.
## 2 Linear theory
### 2.1 Statement of the problem. Equivalences
The problem we want to solve can be first posed and solved in a linear algebraic way. In fact, let $`𝒲`$ and $``$ be $``$-vector spaces (although, instead of $``$, another field of characteristic different from $`2`$ can be used), with $`dim=m`$, and $`dim𝒲=m+n`$. Let $`\sigma :𝒲`$ be a surjective morphism, and denote $`\mathrm{V}(\sigma )=\mathrm{ker}\sigma `$, and by $`ȷ:\mathrm{V}(\sigma )𝒲`$ the natural injection. Consider the exact sequence
$$0\mathrm{V}(\sigma )\text{}𝒲\text{}0$$
(1)
Suppose that $`\eta \mathrm{\Lambda }^m^{}`$ is a volume element; denote $`\omega =\sigma ^{}\eta `$, and assume that a form $`\mathrm{\Omega }\mathrm{\Lambda }^{m+1}𝒲^{}`$ and a subspace $`𝒞`$ of $`𝒲`$ are given. We denote this collection of data as $`(\sigma ;\eta ,\mathrm{\Omega };𝒞)`$.
Next we consider the following problems in $`(\sigma ;\eta ,\mathrm{\Omega };𝒞)`$:
###### Statement 1
To find a $`m`$-vector $`𝒳\mathrm{\Lambda }^m𝒞`$ satisfying that:
1. $`𝒳`$ is decomposable. 2. $`𝑖(𝒳)\omega =1`$. 3. $`𝑖(𝒳)\mathrm{\Omega }=0`$.
###### Statement 2
To find a subspace $`𝒞`$ satisfying that:
1. $`dim=dim=m`$. 2. $`\sigma |_{}:`$ is an isomorphism. 3. $`[𝑖(w)\mathrm{\Omega }]|_{}=0`$, $`w𝒲`$.
Observe that condition 2 is equivalent to $`𝒲=\mathrm{V}(\sigma )`$.
###### Statement 3
To find a linear map $`𝐡:𝒞𝒲`$ satisfying that:
1. $`\sigma 𝐡=\mathrm{Id}_{}`$. 2. $`[𝑖(w)\mathrm{\Omega }]|_{\mathrm{Im}𝐡}=0`$, for every $`w𝒲`$.
###### Proposition 1
Statements 1, 2, and 3 are equivalent, that is, from every solution to some of these problems we can obtain a solution to the others.
( Proof ) (1 $``$ 2) Let $`𝒳\mathrm{\Lambda }^m𝒞`$ be a solution to the problem 1. As a consequence of the first condition, we have $`𝒳=w_1\mathrm{}w_m`$, with $`w_\alpha 𝒞`$. If $`e_\alpha =\sigma (w_\alpha )`$, for every $`\alpha =1,\mathrm{},m`$, by the second condition we have $`\eta (e_1,\mathrm{},e_m)=1`$, and hence $`\{e_\alpha \}`$ is a basis of $``$.
Consider the subspace $`=w_1,\mathrm{},w_m`$. We have obviously that $`dim=m`$ and that the restriction $`\sigma |_{}:`$ is an isomorphism. Furthermore, $`[𝑖(w)\mathrm{\Omega }](w_1,\mathrm{},w_m)=0`$, for every $`w𝒲`$. Thus $``$ is a solution to problem 2.
(2 $``$ 3) Let $``$ be a solution to problem 2. So $`\sigma |_{}`$ is an isomorphism. If $`ȷ_{}:𝒞`$ is the natural injection, let $`𝐡:𝒞`$ be defined as $`𝐡:=ȷ_{}(\sigma |_{})^1`$. This map is a solution to problem 3 because the first condition holds straightforwardly and, as $`\mathrm{Im}𝐡=`$, the second condition holds.
(3 $``$ 1) Let $`𝐡`$ be a solution to problem 3. If $`\{e_1,\mathrm{},e_m\}`$ is a basis of $``$ satisfying that $`\eta (e_1,\mathrm{},e_m)=1`$, let $`w_\alpha =𝐡(e_\alpha )`$, and $`𝒳=w_1\mathrm{}w_m`$. Then $`𝒳\mathrm{\Lambda }^m𝒞`$ is a solution to problem 1, because it is decomposable, and
$$𝑖(𝒳)\omega =\omega (w_1,\mathrm{},w_m)=\eta (e_1,\mathrm{},e_m)=1.$$
Furthermore, if $`w𝒲`$,
$$[𝑖(𝒳)\mathrm{\Omega }](w)=\mathrm{\Omega }(w_1,\mathrm{},w_m,w)=(1)^m[𝑖(w)\mathrm{\Omega }](w_1,\mathrm{},w_m)=0$$
since $`𝐡`$ is a solution to problem 3, and $`w_1,\mathrm{},w_m\mathrm{Im}𝐡`$.
### 2.2 Maps induced by a section
Consider the exact sequence (1), and let $`:𝒲`$ be a section of $`\sigma `$. Denote $`\mathrm{H}():=\mathrm{Im}`$. We have the splitting
$$𝒲=\mathrm{H}()\mathrm{V}(\sigma )$$
$`\mathrm{H}()`$ is called the horizontal subspace of $``$, and $`\mathrm{V}(\sigma )`$ is the vertical subspace of $`\sigma `$. Note that $`\sigma |_{\mathrm{H}()}`$ is an isomorphism. The above splitting induces the natural projections
$$\sigma _{}^H:𝒲\mathrm{H}()𝒲;\sigma _{}^V:𝒲\mathrm{V}(\sigma )𝒲$$
with $`\sigma _{}^H+\sigma _{}^V=\mathrm{Id}_𝒲`$; and, for every $`w𝒲`$, we write $`w=w_{}^H+w_{}^V`$, where $`w_{}^H\mathrm{H}()`$ and $`w_{}^V\mathrm{V}(\sigma )`$ are called the horizontal and vertical components of $`w`$ induced by $``$. In the same way we have the induced splitting
$$𝒲^{}=\mathrm{H}^{}()\mathrm{V}^{}(\sigma )$$
where $`\mathrm{H}^{}()`$ is identified with the set $`\{\beta 𝒲^{};\beta \sigma _{}^V=0\}`$, and $`\mathrm{V}^{}(\sigma )`$ with $`\{\beta 𝒲^{};\beta \sigma _{}^H=0\}`$, in a natural way. This splitting of $`𝒲^{}`$ induced by $``$ gives rise to a bigradation in $`\mathrm{\Lambda }^k𝒲^{}`$ given by
$$\mathrm{\Lambda }^k𝒲^{}=\underset{p,q=0,\mathrm{},k;p+q=k}{}(\mathrm{\Lambda }^p\mathrm{H}^{}()\mathrm{\Lambda }^q\mathrm{V}^{}(\sigma ))$$
Now, let $`𝒵\mathrm{\Lambda }^m`$ such that $`\eta (𝒵)=1`$. With this condition, $`𝒵`$ is unique and decomposable, since $`dim=m`$. Consider $`𝒴_\eta ^{}=\mathrm{\Lambda }^m(𝒵)\mathrm{\Lambda }^m𝒲`$, which verifies the following properties:
1. $`𝒴_\eta ^{}`$ is decomposable, because if $`𝒵=e_1\mathrm{}e_m`$, then $`𝒴_\eta ^{}=(e_1)\mathrm{}(e_m)`$.
2. $`\omega (𝒴_\eta ^{})=1`$, since
$$\omega (𝒴_\eta ^{})=\sigma ^{}\eta (\mathrm{\Lambda }^m(𝒵))=\eta [(\mathrm{\Lambda }^m\sigma \mathrm{\Lambda }^m)(𝒵)]=\eta [\mathrm{\Lambda }^m(\sigma )(𝒵)]=\eta (𝒵)=1.$$
$`𝒴_\eta ^{}`$ is said to be the $`m`$-vector associated to $``$ and $`\eta `$, and it generates $`\mathrm{\Lambda }^m\mathrm{H}()`$.
The bigradation in $`\mathrm{\Lambda }^k𝒲^{}`$ induces a splitting of $`\mathrm{\Omega }`$ as follows: $`\mathrm{\Omega }=\mathrm{\Omega }^{(m,1)}+\mathrm{\Omega }^{}`$, $`\mathrm{\Omega }^{(m,1)}`$ being a $`(m+1)`$-form of bidegree $`(m,1)`$, and $`\mathrm{\Omega }^{}`$ a $`(m+1)`$-form that includes the rest of components. Moreover, we have:
###### Proposition 2
$`\mathrm{\Omega }^{(m,1)}=\omega \gamma _\eta ^{}`$, where $`\gamma _\eta ^{}:=𝑖(𝒴_\eta ^{})\mathrm{\Omega }`$. Then $`\mathrm{\Omega }=\mathrm{\Omega }^{}+\omega \gamma _\eta ^{}`$.
( Proof ) As $`𝒴_\eta ^{}`$ generates $`\mathrm{\Lambda }^m\mathrm{H}()`$, it suffices to prove that $`\mathrm{\Omega }^{(m,1)}`$ and $`\omega \gamma _\eta ^{}`$ coincide when acting on $`𝒴_\eta ^{}v`$, for every $`v\mathrm{V}(\sigma )`$. Thus, as $`\gamma _\eta ^{}`$ vanishes on $`\mathrm{H}()`$, we obtain
$`\mathrm{\Omega }^{(m,1)}(𝒴_\eta ^{}v)`$ $`=`$ $`\mathrm{\Omega }(𝒴_\eta ^{}v)=[𝑖(𝒴_\eta ^{})\mathrm{\Omega })](v)=\gamma ^{}_\eta (v)`$
$`(\omega \gamma _\eta ^{})(𝒴_\eta ^{}v)`$ $`=`$ $`\omega (𝒴_\eta ^{})\gamma _\eta ^{}(v)=\gamma _\eta ^{}(v)`$
Finally, if $`𝐡:𝒞`$ is a linear map, $``$ induces a splitting $`𝐡=𝐡_{}^H+𝐡_{}^V`$, where $`𝐡_{}^H=\sigma _{}^H𝐡`$, and $`𝐡_{}^V=\sigma _{}^V𝐡`$. Then, we introduce the map (endomorphism of $`𝒲`$)
$$\stackrel{~}{𝐡_{}^V}=𝐡_{}^V\sigma =\sigma _{}^V𝐡\sigma :𝒲\mathrm{V}(\sigma )𝒲$$
### 2.3 Characterization of solutions
In what follows, we assume that:
###### Assumption 1
The $`(m+1)`$-form $`\mathrm{\Omega }^{}`$ is of bidegree $`(m1,2)`$. Hence
$$\mathrm{\Omega }=\mathrm{\Omega }^{(m,1)}+\mathrm{\Omega }^{(m1,2)}$$
(2)
This is equivalent to demanding that $`𝑖(v_1)𝑖(v_2)𝑖(v_3)\mathrm{\Omega }=0`$, for every $`v_1,v_2,v_3\mathrm{V}(\sigma )`$.
Note that if $`U`$ and $`V`$ are real vector spaces of finite dimension then $`U^{}V\{𝐡:UV|𝐡\text{is linear}\}`$. Thus, the auxiliar section $``$ induces the $``$-bilinear map
$$\begin{array}{ccccc}\mathrm{}_\mathrm{\Omega }^{}& :& ^{}𝒞& & (^{}H())\times \mathrm{V}^{}(\sigma )\\ & & 𝐡& & (𝐡_{}^H,𝑖(𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{}))\mathrm{\Omega }|_{\mathrm{V}(\sigma )})\end{array}$$
(3)
where $`𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{}`$ is the $`m`$-vector on $`𝒲`$ defined as follows: for every $`\beta ^1,\mathrm{},\beta ^m𝒲^{}`$,
$$(𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{})(\beta ^1,\mathrm{},\beta ^m):=\underset{\alpha =1}{\overset{m}{}}𝒴_\eta ^{}(\beta ^1,\mathrm{},[\stackrel{~}{𝐡_{}^V}]^t(\beta ^\alpha ),\mathrm{},\beta ^m)$$
Observe that, if $`𝒴_\eta ^{}=w_1\mathrm{}w_m`$, with $`w_\alpha 𝒲`$, then
$$𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{}=\underset{\alpha =1}{\overset{m}{}}w_1\mathrm{}\stackrel{~}{𝐡_{}^V}(w_\alpha )\mathrm{}w_m$$
###### Theorem 1
The necessary and sufficient condition for a linear map $`𝐡:𝒞`$ to be a solution to the problem posed in Statement 3 is that
$$\mathrm{}_\mathrm{\Omega }^{}(𝐡)=(ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1,(\gamma _\eta ^{})|_{V(\sigma )})$$
(4)
where $`ȷ_{\mathrm{H}()}:\mathrm{H}()𝒲`$ denotes the natural injection, and $`ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1`$ is the horizontal lift associated with $``$.
( Proof ) ($``$) Suppose that the linear map $`𝐡:𝒞`$ is a solution to the problem posed in Statement 3. Consider the linear map $`\phi :𝒲`$ defined by
$$\phi :=𝐡ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1:𝒲.$$
We have that
$$\sigma \phi =\sigma 𝐡\sigma ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1=\mathrm{Id}\mathrm{Id}=0$$
and therefore,
$$𝐡_{}^H=ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1,𝐡_{}^V=\phi $$
Now, suppose that $`e_1,\mathrm{},e_m`$ such that $`\eta (e_1,\mathrm{},e_m)=1`$, and let $`w_\alpha =(e_\alpha )`$, for $`\alpha =1,\mathrm{},m`$; thus $`𝒴_\eta ^{}=w_1\mathrm{}w_m`$. We obtain that
$$𝐡(e_\alpha )=𝐡_{}^H(e_\alpha )+𝐡_{}^V(e_\alpha )=w_\alpha +𝐡_{}^V(e_\alpha )$$
As $`𝐡`$ is a solution to the problem, using the splitting (2), for every $`v𝒲`$, we have
$`0`$ $`=`$ $`\mathrm{\Omega }(𝐡(e_1),\mathrm{},𝐡(e_m),v)=\mathrm{\Omega }(w_1+𝐡_{}^V(e_1),\mathrm{},w_m+𝐡_{}^V(e_m),v)`$
$`=`$ $`\mathrm{\Omega }^{(m,1)}(w_1,\mathrm{},w_m,v)+{\displaystyle \underset{\alpha =1}{\overset{m}{}}}\mathrm{\Omega }^{(m1,2)}(w_1,\mathrm{},𝐡_{}^V(e_\alpha ),\mathrm{},w_m,v)`$
$`=`$ $`\gamma _\eta ^{}(v)+(𝑖(𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{})\mathrm{\Omega })(v)`$
and the result follows.
($``$) Suppose that there exists a linear map $`𝐡:𝒞`$ such that (4) holds; that is,
$`𝐡_{}^H=\sigma _{}^H𝐡`$ $`=`$ $`ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1,`$
$`(𝑖(𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{})\mathrm{\Omega })|_{\mathrm{V}(\sigma )}=\gamma _\eta ^{}|_{\mathrm{V}(\sigma )}`$ $`=`$ $`𝑖(𝒴_\eta ^{})\mathrm{\Omega }|_{V(\sigma )}.`$
First we prove that $`𝐡`$ is a section of $`\sigma `$. In fact,
$$\sigma 𝐡=\sigma (𝐡_{}^H+𝐡_{}^V)=\sigma 𝐡_{}^H=\sigma ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1=\mathrm{Id}_{}.$$
Furthermore, let $`e_1,\mathrm{},e_m`$, with $`\eta (e_1,\mathrm{},e_m)=1`$, and let $`w_\alpha =(e_\alpha )`$, for $`\alpha =1,\mathrm{},m`$. We have $`𝒴_\eta ^{}=w_1\mathrm{}w_m`$. Now
$$𝐡(e_\alpha )=(ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1)(e_\alpha )+𝐡_{}^V(e_\alpha )=w_\alpha +𝐡_{}^V(e_\alpha )$$
and we must prove that, if $`w𝒲`$, then $`\mathrm{\Omega }(𝐡(e_1),\mathrm{},𝐡(e_m),w)=0`$. Note that, as $`𝐡`$ is a section of $`\sigma `$, it induces a splitting $`𝒲=𝐡()\mathrm{V}(\sigma )`$, and hence $`w=w_𝐡^H+w_𝐡^V`$, where $`w_𝐡^H\mathrm{Im}𝐡`$ and $`w_𝐡^VV(\sigma )`$. Then
$$\mathrm{\Omega }(𝐡(e_1),\mathrm{},𝐡(e_m),w_𝐡^H)=0,$$
and it suffices to prove that $`\mathrm{\Omega }(𝐡(e_1),\mathrm{},𝐡(e_m),v)=0`$, for every $`v\mathrm{V}(\sigma )`$. In fact,
$`\mathrm{\Omega }(𝐡(e_1),\mathrm{},𝐡(e_m),v)`$ $`=`$ $`\mathrm{\Omega }(w_1+𝐡_{}^V(e_1),\mathrm{},w_m+𝐡_{}^V(e_m),v)`$
$`=`$ $`\mathrm{\Omega }^{(m,1)}(w_1,\mathrm{},w_m,v)+{\displaystyle \underset{\alpha =1}{\overset{m}{}}}\mathrm{\Omega }^{(m1,2)}(w_1,\mathrm{},𝐡_{}^V(e_\alpha ),\mathrm{},w_m,v)`$
$`=`$ $`(𝑖(𝒴_\eta ^{})\mathrm{\Omega })(v)+{\displaystyle \underset{\alpha =1}{\overset{m}{}}}\mathrm{\Omega }^{(m1,2)}(w_1,\mathrm{},\stackrel{~}{𝐡_{}^V}(w_\alpha ),\mathrm{},w_m,v)`$
$`=`$ $`\gamma _\eta ^{}(v)+(𝑖(𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{})\mathrm{\Omega })(v)=0`$
Now, from Theorem 1, we deduce that:
###### Corollary 1
A linear map $`𝐡:𝒞`$ is a solution to the problem posed in Statement 3 if, and only if
$$𝐡_{}^H=ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1,[𝑖(𝑖([\stackrel{~}{𝐡_{}^V}]^t)𝒴_\eta ^{})\mathrm{\Omega }]|_{\mathrm{V}(\sigma )}=\gamma _\eta ^{}|_{\mathrm{V}(\sigma )}$$
Let $`\mathrm{V}(\sigma )^0𝒲^{}`$ be the annihilator of $`\mathrm{V}(\sigma )`$. It is clear that the vector spaces $`\mathrm{H}^{}()`$ and $`\mathrm{V}(\sigma )^0`$ are isomorphic. The orthogonal complement of $`𝒞`$ with respect to $`\mathrm{\Omega }`$ and $``$ is the subspace $`(𝒞^{})_\mathrm{\Omega }^{}`$ of $`(\mathrm{V}(\sigma )^0)\times \mathrm{V}(\sigma )`$ defined by
$$(𝒞^{})_\mathrm{\Omega }^{}:=(\mathrm{Im}\mathrm{}_\mathrm{\Omega }^{})^0.$$
(5)
Then, from Theorem 1, we obtain:
###### Theorem 2
There exists a solution to the problem posed in Statement 3 if, and only if,
$$𝐡^{}(ȷ_{\mathrm{H}()}(\sigma |_{\mathrm{H}()})^1)\gamma _\eta ^{}(Z)=0,\text{for every }(𝐡^{},Z)(𝒞^{})_\mathrm{\Omega }^{}.$$
(6)
Note that, if $`𝒞=𝒲`$ and $`(𝒲^{})_\mathrm{\Omega }^{}=\{0\}`$, then it is clear that (6) holds. This is the case in the following Proposition:
###### Proposition 3
If the $`(m+1)`$-form $`\mathrm{\Omega }^{}\mathrm{\Lambda }^{m+1}𝒲^{}`$ given by
$$\mathrm{\Omega }^{}=\mathrm{\Omega }\omega \gamma _\eta ^{}$$
(7)
is $`1`$-nondegenerate (that is, the map $`\mathrm{}_\mathrm{\Omega }^{}:𝒲\mathrm{\Lambda }^m𝒲^{}`$, defined by $`\mathrm{}_\mathrm{\Omega }^{}(v)=𝑖(v)\mathrm{\Omega }^{}`$, for every $`v𝒲`$, is injective), then $`(𝒲^{})_\mathrm{\Omega }^{}=\{0\}`$.
( Proof ) Let $`(𝐡^{},Z)(𝒲^{})_\mathrm{\Omega }^{}`$. From the definitions of $`\mathrm{}_\mathrm{\Omega }^{}`$ and $`(𝒲^{})_\mathrm{\Omega }^{}`$ (eqs. (3) and (5)), we obtain that
$$𝐡^{}(𝐡_{}^{}{}_{}{}^{H})+𝑖(Z)𝑖(𝑖([\stackrel{~}{𝐡_{}^{}{}_{}{}^{V}}]^{})𝒴_\eta ^{})\mathrm{\Omega }=0$$
(8)
for every $`𝐡^{}Lin(,𝒲)`$. In particular, this implies that $`𝐡^{}(𝐡^{})=0`$, for every $`𝐡^{}^{}H()`$, and hence $`𝐡^{}=0`$. Therefore, using (8), we deduce that
$$𝑖(Z)𝑖(𝑖([\stackrel{~}{𝐡_{}^V}]^{})𝒴_\eta ^{})\mathrm{\Omega }=0,\text{for every }𝐡^{}𝒲.$$
As a consequence, from (7) and from assumption 1, it follows that $`𝑖(Z)\mathrm{\Omega }^{}=0`$ and, since $`\mathrm{\Omega }^{}`$ is $`1`$-nondegenerate, we have that $`Z=0`$.
## 3 The general multisymplectic case
### 3.1 Statement of the problem
The problem we wish to solve arises from the Lagrangian and Hamiltonian formalisms in field theories, although other kinds of systems can also be stated in this way.
The general geometrical setting for these kinds of systems consists in giving a fibred manifold $`\kappa :FM`$ (which in what follows is assumed to be a fibre bundle), where $`dimM=m>1`$ and $`dimF=n+m`$, and $`M`$ is an orientable manifold with volume form $`\eta \mathrm{\Omega }^m(M)`$. We denote $`\omega =\kappa ^{}\eta `$. We write $`(U;x^\mu ,y^j)`$, $`\mu =1,\mathrm{},m`$, $`j=1,\mathrm{},n`$, for local charts of coordinates in $`F`$ adapted to the fibred structure, and such that $`\omega =\mathrm{d}x^1\mathrm{}\mathrm{d}x^m\mathrm{d}^mx`$. Let $`\mathrm{\Omega }\mathrm{\Omega }^{m+1}(F)`$ be a closed form, and consider the triad $`(F,\mathrm{\Omega },\omega )`$. The form $`\mathrm{\Omega }`$ is said to be a multisymplectic form if it is $`1`$-nondegenerate, that is, if the map $`\mathrm{}_\mathrm{\Omega }:\mathrm{T}F\mathrm{\Lambda }^m\mathrm{T}^{}F`$, defined by $`\mathrm{}_\mathrm{\Omega }(v)=𝑖(v)\mathrm{\Omega }`$, for every $`v𝒲`$, is injective. In this case, the system described by the above triad is called a multisymplectic system. Otherwise, the form is said to be a pre-multisymplectic form, and the system is pre-multisymplectic.
The problem is stated as follows:
###### Statement 4
Given a pre-multisymplectic system $`(F,\mathrm{\Omega },\omega )`$, we want to find a submanifold $`ȷ_C:CF`$, and a $`\kappa `$-transverse, locally decomposable and integrable $`m`$-vector field $`𝒳_C`$ along $`C`$, in the fibration $`\kappa :FM`$, such that
$$𝑖(𝒳_C(y))\mathrm{\Omega }(y)=0,\text{for every }yC.$$
(9)
First we obviate the integrability condition. Hence the problem consists in finding a submanifold $`CF`$ and a locally decomposable $`m`$-vector field $`𝒳_C\text{X}^m(F)`$ along $`C`$ such that
$$𝑖(𝒳_C(y))\omega (y)=1,𝑖(𝒳_C(y))\mathrm{\Omega }(y)=0,\text{for every }yC.$$
(10)
(Note that the first equation implies that $`𝒳_C`$ is $`\kappa `$-transverse).
Taking into account Remark 5 in the Appendix and Proposition 1, we have:
###### Proposition 4
If $`C`$ is a submanifold of $`F`$, then there exists a solution to the problem stated in Statement 4 if, and only if, at every point $`yC`$, there is $`𝐡_y\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yCLin(\mathrm{T}_{\kappa (y)}M,T_yC)`$ such that
1. $`𝐡_y`$ is $`\kappa `$-transverse (that is, it is a connection along $`C`$):
$$\mathrm{T}_y\kappa |_{\mathrm{T}_yC}𝐡_y=Id.$$
(11)
2. For every $`(X_1^{})_{\kappa (y)},\mathrm{},(X_m^{})_{\kappa (y)}\mathrm{T}_{\kappa (y)}M`$, and $`Y_y\mathrm{T}_yF`$,
$$\mathrm{\Omega }(y)(𝐡_y((X_1^{})_{\kappa (y)}),\mathrm{},𝐡_y((X_m^{})_{\kappa (y)}),Y_y)=0.$$
(12)
In order to solve this problem, the use of an arbitrary connection in the fibration $`\kappa :FM`$ is required. Thus, let $``$ be a connection in $`\kappa :FM`$, and $`𝒴_\eta ^{}`$ the corresponding locally decomposable $`m`$-vector field on $`F`$ such that $`𝑖(𝒴_\eta ^{})\omega =1`$. As is well-known (see Appendix and Section 2.2), the connection $``$ induces a splitting
$$\mathrm{\Lambda }^k\mathrm{T}^{}F=\underset{p,q=0,\mathrm{},k;p+q=k}{}(\mathrm{\Lambda }^p\mathrm{H}^{}()\mathrm{\Lambda }^q\mathrm{V}^{}(\kappa ))$$
where $`\mathrm{H}()F`$ is the horizontal subbundle associated with the connection $``$ and $`\mathrm{V}(\kappa )F`$ is the vertical subbundle of the fibration $`\kappa :FM`$. Thus, we have that
$$\mathrm{\Omega }=\mathrm{\Omega }^{(m,1)}+\mathrm{\Omega }^{},$$
$`\mathrm{\Omega }^{(m,1)}`$ being a $`(m+1)`$-form of bidegree $`(m,1)`$ and $`\mathrm{\Omega }^{}`$ a $`(m+1)`$-form. Moreover, as a straightforward consequence of Proposition 2, we have that:
###### Proposition 5
$`\mathrm{\Omega }^{(m,1)}=\omega \gamma _\eta ^{}`$, where $`\gamma _\eta ^{}:=𝑖(𝒴_\eta ^{})\mathrm{\Omega }`$. Hence $`\mathrm{\Omega }=\mathrm{\Omega }^{}+\omega \gamma _\eta ^{}`$.
In what follows, we assume that the following condition holds:
###### Assumption 2
The $`(m+1)`$-form $`\mathrm{\Omega }^{}`$ is of bidegree $`(m1,2)`$.
By Proposition 5, this is equivalent to demanding that
$$𝑖(Z_1)𝑖(Z_2)𝑖(Z_3)\mathrm{\Omega }=0,\text{for every }Z_1,Z_2,Z_3\text{X}^{\mathrm{V}(\kappa )}(F).$$
###### Remark 1
The above assumption is justified because this is the situation in the Lagrangian and Hamiltonian formalism of field theories (see Propositions 7 and 11).
### 3.2 Conditions for the existence of solutions on a submanifold of the total space
Taking into account the above considerations, the necessary and sufficient condition for the existence of solutions to the problem posed in the Statement 10 arises from the results obtained in Sections 2.2 and 2.3. The key consists in working at every point of the manifolds involved in this problem. Thus, if $`yC`$, the following identifications can be made:
$$\mathrm{T}_{\kappa (y)}M,𝒲\mathrm{T}_yF,𝒞\mathrm{T}_yC,\mathrm{V}(\sigma )\mathrm{V}_y(\kappa )$$
Then we may consider the $``$-linear map
$$\mathrm{}_\mathrm{\Omega }^{}(y):\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yC(\mathrm{T}_{\kappa (y)}^{}MH_y())\times V_y^{}(\kappa )$$
defined by
$$\mathrm{}_\mathrm{\Omega }^{}(y)(𝐡_y)=((𝐡_y)_{}^H,𝑖(𝑖([\stackrel{~}{(𝐡_y)_{}^V}]^t)(𝒴_\eta ^{}(y)))(\mathrm{\Omega }(y))|_{\mathrm{V}_y(\kappa )}).$$
(13)
Therefore, Theorem 1 and Corollary 1 lead to the following results:
###### Theorem 3
Let $`yC`$. Then, there exists a linear map $`𝐡_y\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yC`$ such that (11) and (12) hold if, and only if,
$$\mathrm{}_\mathrm{\Omega }^{}(y)(𝐡_y)=((\mathrm{T}_y\kappa _{\mathrm{H}()})^1,\gamma _\eta ^{}(y)|_{\mathrm{V}_y(\kappa )})$$
where $`(\mathrm{T}_y\kappa _{\mathrm{H}()})^1:\mathrm{T}_{\kappa (y)}M\mathrm{H}_y()`$ is the horizontal lift at $`y`$ associated with the connection $``$. (Observe that $`(\mathrm{T}_y\kappa _{\mathrm{H}()})^1\mathrm{T}_{\kappa (y)}^{}M\mathrm{H}_y()`$).
###### Corollary 2
If $`yC`$, and $`𝐡_y\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yC`$, then (11) and (12) hold if, and only if,
$$(𝐡_y)_{}^H=(\mathrm{T}_y\kappa _{\mathrm{H}()})^1,[𝑖(𝑖([\stackrel{~}{(𝐡_y)_{}^V}]^t)(𝒴_\eta ^{}(y)))(\mathrm{\Omega }(y))]|_{\mathrm{V}_y(\kappa )}=(\gamma _\eta ^{}(y))|_{\mathrm{V}_y(\kappa )}.$$
###### Remark 2
If $`yC`$, let $`\mathrm{V}_y(\kappa )^0\mathrm{T}_y^{}𝒲`$ be the annihilator of the vertical subspace $`\mathrm{V}_y(\kappa )`$ at the point $`y`$. Then we have that
$$[Lin(\mathrm{T}_{\kappa (y)}M,\mathrm{H}_y())]^{}\mathrm{T}_{\kappa (y)}M\mathrm{H}_y^{}()\mathrm{T}_{\kappa (y)}M\mathrm{V}_y(\kappa )^0Lin(\mathrm{T}_{\kappa (y)}^{}M,\mathrm{V}_y(\kappa )^0).$$
If $`𝐡_y:\mathrm{T}_{\kappa (y)}M\mathrm{H}_y()`$ and $`𝐡_y^{}:\mathrm{T}_{\kappa (y)}^{}M\mathrm{V}_y(\kappa )^0`$ are linear maps, $`\{(X_1)_{\kappa (y)},\mathrm{},(X_m)_{\kappa (y)}\}`$ is a basis of $`\mathrm{T}_{\kappa (y)}M`$ such that $`\{\alpha _{\kappa (y)}^1,\mathrm{},\alpha _{\kappa (y)}^m\}`$ is the dual basis of $`\mathrm{T}_{\kappa (y)}^{}M`$, and
$$\eta (\kappa (y))=(\alpha ^1)_{\kappa (y)}\mathrm{}(\alpha ^m)_{\kappa (y)},𝒳_\eta (\kappa (y))=(X_1)_{\kappa (y)}\mathrm{}(X_m)_{\kappa (y)},$$
then, taking $`(Y_i)_y=(\mathrm{T}_y\kappa _{\mathrm{H}()})((X_i)_{\kappa (y)})`$ and $`\beta _y^i=\alpha _{\kappa (y)}^i\mathrm{T}_y\kappa `$, for all $`i\{1,\mathrm{},m\}`$, we deduce that $`\{(Y_1)_y,\mathrm{},(Y_m)_y\}`$ and $`\{\beta _y^1,\mathrm{},\beta _y^m\}`$ are a basis of $`\mathrm{H}_y()`$ and $`\mathrm{V}_y(\kappa )^0`$, respectively. Moreover, if
$$𝐡_y((X_i)_{\kappa (y)})=(𝐡_y)_i^j(Y_j)_y,𝐡_y^{}(\alpha _{\kappa (y)}^i)=(𝐡_y^{})_j^i\beta _y^j,\text{for every }i\{1,\mathrm{},m\}$$
it follows that $`<𝐡_y^{},𝐡_y>=𝐡_y^{}(𝐡_y)=(𝐡_y)_i^j(𝐡_y^{})_j^i`$.
Now, if $`yC`$, the orthogonal complement $`(\mathrm{T}_y^{}C)_\mathrm{\Omega }^{}`$ with respect to $`\mathrm{\Omega }`$ and $``$ is the subspace of $`(\mathrm{T}_{\kappa (y)}M\mathrm{V}_y(\kappa )^0)\times \mathrm{V}_y(\kappa )`$ defined by
$$(\mathrm{T}_y^{}C)_\mathrm{\Omega }^{}=(\mathrm{Im}\mathrm{}_\mathrm{\Omega }^{}(y))^0.$$
(14)
As in Theorem 2, from Theorem 3 we obtain
###### Theorem 4
Let $`yC`$. Then, there exists a linear map $`𝐡_y\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yC`$ such that
$$(𝐡_y)_{}^H=(\mathrm{T}_y\kappa _{\mathrm{H}()})^1,[𝑖(𝑖([\stackrel{~}{(𝐡_y)_{}^V}]^t)(𝒴_\eta ^{}(y)))(\mathrm{\Omega }(y))]|_{\mathrm{V}_y(\kappa )}=(\gamma _\eta ^{}(y))|_{\mathrm{V}_y(\kappa )}$$
if, and only if,
$$𝐡_y^{}(\mathrm{T}_y\kappa _{\mathrm{H}()})^1\gamma _\eta ^{}(y)(Z_y)=0,\text{for every}(𝐡_y^{},Z_y)(\mathrm{T}_y^{}C)_\mathrm{\Omega }^{}.$$
(15)
Note that if $`(\mathrm{T}_y^{}C)_\mathrm{\Omega }^{}=\{0\}`$ then it is clear that (15) holds. Thus, from Proposition 3, we have:
###### Proposition 6
If the $`(m+1)`$-form $`\mathrm{\Omega }^{}`$ on $`F`$ given by $`\mathrm{\Omega }^{}=\mathrm{\Omega }\omega \gamma _\eta ^{}`$ is $`1`$-nondegenerate, that is, the map $`\mathrm{}_\mathrm{\Omega }^{}:\mathrm{T}F\mathrm{\Lambda }^m\mathrm{T}^{}F`$ is injective, then
$$(\mathrm{T}_y^{}F)_\mathrm{\Omega }^{}=\{0\},\text{for every}yF.$$
### 3.3 The pre-multisymplectic constraint algorithm
Now we apply the above results in order to solve the problem stated in Section 3.1. The procedure is algorithmic, and gives a sequence of subsets $`\{C_i\}`$ of $`F`$. Then, we assume that:
###### Assumption 3
Every subset $`C_i`$ of this sequence is a regular submanifold of $`F`$, and its natural injection is an embedding.
Thus, we consider the submanifold $`C_1F`$ where a solution exists, that is,
$$\begin{array}{ccc}C_1\hfill & =& \hfill \{yF𝐡_yLin(\mathrm{T}_{\kappa (y)}M,\mathrm{T}_yF)\text{such that}(𝐡_y)_{}^H=(\mathrm{T}_y\kappa _{\mathrm{H}()})^1,\\ & & \hfill [𝑖(𝑖([\stackrel{~}{(𝐡_y)_{}^V}]^{})(𝒴_\eta ^{}(y)))(\mathrm{\Omega }(y))]|_{\mathrm{V}_y(\kappa )}=(\gamma _\eta ^{}(y))|_{\mathrm{V}_y(\kappa )}\}.\end{array}$$
Then, using the results of Section 3.2, we deduce that there is a locally decomposable section $`𝒳_1`$ of the vector bundle $`\mathrm{\Lambda }^m\mathrm{T}_{C_1}FC_1`$ such that $`(𝑖(𝒳_1)\omega )|_{C_1}=1`$ and $`(𝑖(𝒳_1)\mathrm{\Omega })|_{C_1}=0`$. However, in general, $`𝐡_y(\mathrm{T}_{\kappa (y)}M)`$ is not a subspace of $`\mathrm{T}_yC_1`$ and then $`𝒳_1`$ is not tangent to $`C_1`$ or, in other words, in general, $`𝒳_1`$ is not a connection in the fibration $`\kappa :FM`$ along $`C_1`$. Therefore, we consider the submanifold
$$\begin{array}{ccc}C_2\hfill & =& \hfill \{y_1C_1𝐡_{y_1}Lin(\mathrm{T}_{\kappa (y_1)}M,\mathrm{T}_{y_1}C_1)\text{such that}(𝐡_{y_1})_{}^H=(\mathrm{T}_{y_1}\kappa _{\mathrm{H}()})^1,\\ & & \hfill [𝑖(𝑖([\stackrel{~}{(𝐡_{y_1})_{}^V}]^{})(𝒴_\eta ^{}(y_1)))(\mathrm{\Omega }(y_1))]|_{\mathrm{V}_{y_1}(\kappa )}=(\gamma _\eta ^{}(y_1))|_{\mathrm{V}_{y_1}(\kappa )}\}.\end{array}$$
Then, there is a locally decomposable section $`𝒳_2`$ of the vector bundle $`\mathrm{\Lambda }^m\mathrm{T}_{C_2}C_1C_2`$ such that $`(𝑖(𝒳_2)\omega )|_{C_2}=1`$ and $`(𝑖(𝒳_2)\mathrm{\Omega })|_{C_2}=0`$. However, in general, $`𝒳_2`$ is not a connection in the fibration $`\kappa :FM`$ along $`C_2`$. Following this process, we obtain a sequence of constraint submanifolds
$$\mathrm{}\stackrel{j_{i+1}^i}{}C_i\stackrel{j_i^{i1}}{}\mathrm{}\stackrel{j_2^1}{}C_1\stackrel{j_1}{}C_0F.$$
(16)
For every $`i1`$, $`C_i`$ is called the $`i`$th constraint submanifold.
This procedure is called the pe-multisymplectic constraint algorithm. We have two possibilities:
* There exists an integer $`k>0`$ such that $`dimC_km1`$. This means that the equations have no solution on a submanifold of $`F`$.
* There exists an integer $`k>0`$ such that $`C_{k+1}=C_kC_f`$. In such a case, there exists a connection $`𝒳_f`$ in the fibration $`\kappa :FM`$ along $`C_f`$ such that
$$𝑖(𝒳_f(y_f))(\mathrm{\Omega }(y_f))=0,\text{for every}y_fC_f.$$
In this case, $`C_f`$ is called the final constraint submanifold. This is the situation which is interesting to us. Note that the existence of a connection in the fibration $`\kappa :FM`$ along $`C_f`$ implies that $`\kappa (C_f)`$ is an open subset of $`M`$ and that $`\kappa |_{C_f}:C_f\kappa (C_f)`$ is a fibration (see Remark 5 in the Appendix). In particular, $`dimC_fm`$.
Next we give an intrinsic characterization of the constraints which define the constraint submanifolds $`C_i`$. For this purpose, we consider the vector bundle over $`F`$,
$$W(\kappa ,)=(\kappa ^{}(\mathrm{T}^{}M)\mathrm{H}())_F\mathrm{V}^{}(\kappa )$$
whose fiber over the point $`yF`$ is
$$W_y(\kappa ,)=(\mathrm{T}_{\kappa (y)}^{}M\mathrm{H}_y())\times \mathrm{V}_y^{}(\kappa )Lin(\mathrm{T}_{\kappa (y)}M,\mathrm{H}_y())\times \mathrm{V}_y^{}(\kappa ).$$
The horizontal lift associated with the connection $``$ and the $`1`$-form $`\gamma _\eta ^{}`$ induce a section $`((\mathrm{T}\kappa _{\mathrm{H}()})^1,(\gamma _\eta ^{})|_{\mathrm{V}(\kappa )})`$ of this vector bundle given by
$$((\mathrm{T}\kappa _{\mathrm{H}()})^1,(\gamma _\eta ^{})|_{\mathrm{V}(\kappa )})(y)=((\mathrm{T}_y\kappa _{\mathrm{H}()})^1,(\gamma _\eta ^{}(y))|_{\mathrm{V}_y(\kappa )}),\text{for every}yF.$$
Furthermore, let $`W_{C_i}(\kappa ,)`$ be the vector bundle over the submanifold $`C_i`$ whose fiber at the point $`y_iC_i`$ is $`W_{y_i}(\kappa ,)`$. Moreover, we may consider the orthogonal complement $`(\mathrm{T}_{y_i}^{}C_i)_\mathrm{\Omega }^{}`$ of $`\mathrm{T}_{y_i}C_i`$ with respect to $`\mathrm{\Omega }`$ and $``$ given by (see (14))
$`(\mathrm{T}_{y_i}^{}C_i)_\mathrm{\Omega }^{}=\{(𝐡_{y_i}^{},Z_{y_i})Lin(\mathrm{T}_{\kappa (y_i)}^{}M,\mathrm{V}_{y_i}(\kappa )^0)\times \mathrm{V}_{y_i}(\kappa )𝐡_{y_i}^{}((𝐡_{y_i})_{}^H)+`$
$`𝑖(Z_{y_i})𝑖(𝑖([\stackrel{~}{(𝐡_{y_i})_{}^V}]^t)(𝒴_\eta ^{}(y)))(\mathrm{\Omega }(y_i))=0,\text{for every}𝐡_{y_i}Lin(\mathrm{T}_{\kappa (y_i)}M,\mathrm{T}_{y_i}C_i)\}.`$
Note that $`(\mathrm{T}_{y_i}^{}C_i)_\mathrm{\Omega }^{}W_{y_i}^{}(\kappa ,)`$. Furthermore, if $`(\mathrm{T}^{}C_i)_\mathrm{\Omega }^{}`$ is a vector subbundle of rank $`r`$ of $`W_{C_i}^{}(\kappa ,)`$ (that is, the dimension of $`(\mathrm{T}_{y_i}^{}C_i)_\mathrm{\Omega }^{}`$ is $`r`$, for every $`y_iC_i`$) then one may choose a set of $`r`$ local sections $`\{(𝐡_1^{},Z_1),\mathrm{},(𝐡_r^{},Z_r)\}`$ of the vector bundle $`W^{}(\kappa ,)F`$ such that $`\{(𝐡_1^{},Z_1)|_{C_i},\mathrm{},(𝐡_r^{},Z_r)|_{C_i}\}`$ is a local basis of the space $`\mathrm{\Gamma }((\mathrm{T}^{}C_i)_\mathrm{\Omega }^{})`$ of sections of the vector subbundle $`(\mathrm{T}^{}C_i)_\mathrm{\Omega }^{}C_i`$. In addition, using Theorem 4, we deduce
###### Theorem 5
Every submanifold $`C_i`$ $`(i1)`$ in the sequence (16) may be defined as
$$C_i=\{y_{i1}C_{i1}((\mathrm{T}\kappa _{\mathrm{H}()})^1,(\gamma _\eta ^{})|_{\mathrm{V}(\kappa )})(y_{i1}),(\mathrm{T}_{y_{i1}}^{}C_{i1})_\mathrm{\Omega }^{}=0\}.$$
Therefore, if $`(\mathrm{T}^{}C_{i1})_\mathrm{\Omega }^{}`$ is a vector subbundle of rank $`r`$ of $`W_{C_{i1}}^{}(\kappa ,)`$ and $`\{(𝐡_1^{},Z_1)^{(i1)},\mathrm{},(𝐡_r^{},Z_r)^{(i1)}\}`$ is a set of sections of the vector bundle $`W^{}(\kappa ,)F`$ spanning locally the space $`\mathrm{\Gamma }((\mathrm{T}^{}C_{i1})_\mathrm{\Omega }^{})`$, then $`C_i`$, is defined locally, as a submanifold of $`C_{i1}`$, as the zero set of the functions $`\xi _j^{(i)}C^{\mathrm{}}(F)`$ given by
$$\xi _j^{(i)}=((\mathrm{T}\kappa _{\mathrm{H}()})^1,(\gamma _\eta ^{})|_{\mathrm{V}(\kappa )})((𝐡_j^{},Z_j)^{(i1)}).$$
These functions are called $`i`$th-generation constraints.
### 3.4 The integrability algorithm
Suppose that after applying the premultisymplectic constraint algorithm we have a final constraint submanifold $`C_fF`$ and a connection defined by the multivector field $`𝒳_f`$ in the fibration $`\kappa :FM`$ along $`C_f`$ such that (9) holds on $`C_f`$, that is,
$$𝑖(𝒳_f(y))\mathrm{\Omega }(y)=0,\text{for every}yC_f.$$
(17)
However, $`𝒳_f`$ is not, in general, a flat connection. Nevertheless, in many cases, one may find a submanifold $`_f`$ of $`C_f`$ such that $`(𝒳_f)|__f`$ is a flat connection in the fibration $`\kappa :FM`$ along $`_f`$ and (9) holds for $`(𝒳_f)|__f`$.
Next we present an algorithm which enables us to find this submanifold (which is an adapted version of that given in ). This is a local algorithm, that is, we are in fact working on suitable open sets in $`C_f`$. Hence, let $`𝒳_f_{\mu =1}^mX_\mu `$ be a solution to (17).
* Integrability condition: The condition that $`𝒳_f`$ is flat is equivalent to demanding that the distribution spanned by $`X_1,\mathrm{},X_m`$ is involutive. Then, if $`c_f=dimC_f`$, let $`Z_1,\mathrm{},Z_{nm}\text{X}(F)`$, such that $`\{X_1,\mathrm{},X_m,Z_1,\mathrm{},Z_{c_fm}\}`$ is a local basis of the module of vector fields on $`C_f`$. Therefore, for every pair $`X_\mu ,X_\nu `$ ($`1\mu ,\nu m`$) we have
$$[X_\mu ,X_\nu ]=f_{\mu \nu }^\rho X_\rho +\zeta _{\mu \nu }^lZ_l$$
for some functions $`f_{\mu \nu }^\rho ,\zeta _{\mu \nu }^l`$. Consider the system $`\zeta _{\mu \nu }^l=0`$ and let
$$_1=\{yC_f;\zeta _{\mu \nu }^l(y)=0,\mu ,\nu ,l\}.$$
We have three options:
1. $`_1=C_f`$. Then the distribution spanned by $`X_1,\mathrm{},X_m`$ is involutive, and $`(𝒳_f)|_{C_f}`$ is a flat connection in the fibration $`\kappa :FM`$ along $`C_f`$.
2. $`_1=\mathrm{}`$. Then the distribution spanned by $`X_1,\mathrm{},X_m`$ is not involutive at any point in $`C_f`$, and hence the $`m`$-vector field $`𝒳_f`$ is not integrable.
3. $`_1`$ is a proper subset of $`C_f`$. In this case we assume that $`_1`$ is a closed submanifold of $`C_f`$ and the functions $`\zeta _{\mu \nu }^l`$ are the constraints locally defining $`_1`$. The distribution spanned by $`X_1,\mathrm{},X_m`$ is involutive on $`_1`$; that is, the $`m`$-vector field $`𝒳_f`$ is integrable on $`_1`$.
If $`𝒳_f`$ is tangent to $`_1`$, then $`(𝒳_f)|__1`$ defines a flat connection in $`\kappa :FM`$ along $`_1`$ and (9) holds on $`_1`$ which implies that the problem is solved. Nevertheless, this is not the case in general, so we need the following:
* Tangency condition: Consider the set
$$_2:=\{y_1;𝒳_f(y)\mathrm{\Lambda }^m\mathrm{T}_y_1\}$$
For $`_2`$ we have the same problem, so we define inductively, for $`i>1`$,
$$_i:=\{y_{i1};𝒳_f(y)\mathrm{\Lambda }^m\mathrm{T}_y_{i1}\}$$
and assume that we obtain a sequence $`\mathrm{}_i\mathrm{}_1C_f`$ such that $`_i`$ is a non-empty (closed) submanifold of $`F`$, for all $`i`$, or $`_i=\mathrm{}`$, for some $`i`$.
Observe that the locally decomposable $`m`$-vector field $`𝒳_f=X_1\mathrm{}X_m`$ is tangent to $`_i`$ (with $`_i\mathrm{}`$) if, and only if, $`X_\mu `$ is tangent to $`_i`$, for every $`\mu `$.
Thus, using the constraints, we have that, if $`\{\zeta _{\alpha _i}^{(i)}\}`$ is a basis of constraints defining locally $`_i`$ in $`_{i1}`$, the tangency condition is $`0\underset{_i}{=}X_\mu (\zeta _{\alpha _i}^{(i)})`$ (for every $`\mu ,\alpha _i`$), that is, we have
$$_{i+1}:=\{y_i;X_\mu (\zeta _{\alpha _i}^{(i)})(y)=0,\mu ,\alpha _i\},\text{for every}i1.$$
The above algorithm ends at step $`f`$ in one of the following two options:
1. $`dim_fm1`$. In such a case, we deduce that it is not possible to find a submanifold $``$ of $`C_f`$ such that $`(𝒳_f)|_{}`$ is a flat connection in the fibration $`\kappa :FM`$ along $``$. Therefore, we must consider (if it exists) another connection $`𝒳_f^{}`$ along $`C_f`$ such that $`𝑖(𝒳_f^{}(y))\mathrm{\Omega }(y)=0`$, for every $`yC_f`$, and then we must repeat the above procedure.
2. $`_{f+1}=_f`$. In this case $`_f`$ is a submanifold of $`F`$ and we deduce that $`(𝒳_f)|__f`$ is a flat connection in the fibration $`\kappa :FM`$ along $`_f`$ such that $`𝑖(𝒳_f(y))(\mathrm{\Omega }(y))=0`$, for every $`y_f`$. Thus, the problem is solved. As in Section 3.3, we remark that the existence of a connection in the fibration $`\kappa :FM`$ along $`_f`$ implies that $`\kappa (_f)`$ is an open subset of $`M`$ and that $`\kappa |__f:_f\kappa (_f)`$ is a fibration. In particular, $`dim_fm`$.
We will call this procedure the integrability algorithm for decomposable $`m`$-vector fields.
## 4 Application to Lagrangian and Hamiltonian field theories
### 4.1 Lagrangian and Hamiltonian field theories
(For details on the construction of the Lagrangian and Hamiltonian formalisms of field theories, see for instance, , , , , , , , , , , , .)
A first-order classical field theory is described by its configuration fibre bundle $`\pi :EM`$ and a Lagrangian density which is a $`\overline{\pi }^1`$-semibasic $`m`$-form, $``$, on $`J^1\pi `$ (the first-order jet bundle of $`\pi :EM`$). $``$ is usually written as $`=L(\overline{\pi }^1\eta )L\omega `$, where $`L\mathrm{C}^{\mathrm{}}(J^1\pi )`$ is the Lagrangian function associated with $``$ and $`\omega `$, and $`\pi ^1:J^1\pi E`$ and $`\overline{\pi }^1:=\pi \overline{\pi }:J^1\pi M`$ are the natural projections. The Poincaré-Cartan $`m`$ and $`(m+1)`$-forms associated with the Lagrangian density $``$ are defined using the vertical endomorphism $`𝒱`$ of the bundle $`J^1\pi `$
$$\mathrm{\Theta }_{}:=𝑖(𝒱)+\mathrm{\Omega }^m(J^1\pi );\mathrm{\Omega }_{}:=\mathrm{d}\mathrm{\Theta }_{}\mathrm{\Omega }^{m+1}(J^1\pi )$$
Then a Lagrangian system is a couple $`(J^1\pi ,\mathrm{\Omega }_{})`$. The Lagrangian system is regular if $`\mathrm{\Omega }_{}`$ is $`1`$-nondegenerate. Elsewhere it is called singular. In a natural chart of coordinates $`(x^\alpha ,y^A,v_\alpha ^A)`$ in $`J^1\pi `$ (adapted to the bundle structure, and such that $`\omega =\mathrm{d}x^1\mathrm{}\mathrm{d}x^m\mathrm{d}x^m`$) we have
$`\mathrm{\Omega }_{}`$ $`=`$ $`{\displaystyle \frac{^2L}{v_\nu ^Bv_\alpha ^A}}\mathrm{d}v_\nu ^B\mathrm{d}y^A\mathrm{d}^{m1}x_\alpha {\displaystyle \frac{^2L}{y^Bv_\alpha ^A}}\mathrm{d}y^B\mathrm{d}y^A\mathrm{d}^{m1}x_\alpha `$ (18)
$`+{\displaystyle \frac{^2L}{v_\nu ^Bv_\alpha ^A}}v_\alpha ^A\mathrm{d}v_\nu ^B\mathrm{d}^mx+\left({\displaystyle \frac{^2L}{y^Bv_\alpha ^A}}v_\alpha ^A{\displaystyle \frac{L}{y^B}}+{\displaystyle \frac{^2L}{x^\alpha v_\alpha ^B}}\right)\mathrm{d}y^B\mathrm{d}^mx`$
(where $`\mathrm{d}^mx=\mathrm{d}x^1\mathrm{}\mathrm{d}x^m`$ and $`\mathrm{d}^{m1}x^\alpha 𝑖\left({\displaystyle \frac{}{x^\alpha }}\right)\mathrm{d}^mx`$). Locally, the regularity condition is equivalent to $`det\left({\displaystyle \frac{^2L}{v_\alpha ^Av_\nu ^B}}(\overline{y})\right)0`$, for every $`\overline{y}J^1\pi `$.
The Lagrangian problem associated with a Lagrangian system $`(J^1\pi ,\mathrm{\Omega }_{})`$ consists in finding sections $`\varphi \mathrm{\Gamma }(M,E)`$ (where $`\mathrm{\Gamma }(M,E)`$ denotes the set of sections of $`\pi `$), such that
$$(j^1\varphi )^{}𝑖(X)\mathrm{\Omega }_{}=0,\text{for every }X\text{X}(J^1\pi )$$
In natural coordinates this is equivalent to demanding that $`\varphi `$ satisfies the Euler-Lagrange equations. The problem of finding these sections can be formulated equivalently as follows: to find the integral sections of a class of holonomic $`m`$-vector fields $`\{𝒳_{}\}\text{X}^m(J^1\pi )`$, such that
$$𝑖(𝒳_{})\mathrm{\Omega }_{}=0,\text{for every }𝒳_{}\{𝒳_{}\}$$
(Holonomic means that $`𝒳_{}`$ is integrable and its integral sections are holonomic. This is equivalent to demanding that $`𝒳_{}`$ is integrable and semi-holonomic, that is, it satisfies the condition $`𝑖(𝒳_{})𝒱=0`$. Semi-holonomic (not necessarily integrable) locally decomposable $`m`$-vector fields which are solution to these equations are called Euler-Lagrange $`m`$-vector fields for $`(J^1\pi ,\mathrm{\Omega }_{})`$.
For the Hamiltonian formalism of field theories, we take as the multimomentum bundle the manifold $`J^1\pi ^{}\mathrm{\Lambda }_2^m\mathrm{T}^{}E/\pi ^{}\mathrm{\Lambda }^m\mathrm{T}^{}M`$, where $`\mathrm{\Lambda }_2^m\mathrm{T}^{}E\pi `$ is the bundle of $`m`$-forms on $`E`$ vanishing by the action of two $`\pi `$-vertical vector fields. It is a bundle $`\overline{\tau }^1=\pi \tau ^1:J^1\pi ^{}M`$, where $`\tau ^1:J^1\pi ^{}E`$ is the natural projection. Natural charts of coordinates in $`\pi `$ and $`J^1\pi `$ (adapted to the bundle structure, and such that $`\omega ^{}\overline{\tau }^1\eta =\mathrm{d}x^1\mathrm{}\mathrm{d}x^m\mathrm{d}x^m`$) are denoted by $`(x^\alpha ,y^A,p_A^\alpha ,p)`$ and $`(x^\alpha ,y^A,p_A^\alpha )`$, respectively.
As $`\pi `$ is a subbundle of $`\mathrm{\Lambda }^m\mathrm{T}^{}E`$ (the multicotangent bundle of $`E`$ of order $`m`$), then $`\pi `$ is endowed with canonical forms: the “tautological form” $`\mathrm{\Theta }\mathrm{\Omega }^m(\pi )`$, and the multisymplectic form $`\mathrm{\Omega }:=\mathrm{d}\mathrm{\Theta }\mathrm{\Omega }^{m+1}(\pi )`$. They are known as the multimomentum Liouville $`m`$ and $`(m+1)`$-forms. Their local expressions are
$$\mathrm{\Theta }=p_A^\alpha \mathrm{d}y^A\mathrm{d}^{m1}x_\alpha +p\mathrm{d}^mx,\mathrm{\Omega }=\mathrm{d}p_A^\alpha \mathrm{d}y^A\mathrm{d}^{m1}x_\alpha \mathrm{d}p\mathrm{d}^mx$$
Now, if $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a Lagrangian system, the extended Legendre map associated with $``$, $`\stackrel{~}{}:J^1\pi \pi `$, is defined by: $`(\stackrel{~}{}\overline{y}))(Z_1,\mathrm{},Z_m):=(\mathrm{\Theta }_{})_{\overline{y}}(\overline{Z}_1,\mathrm{},\overline{Z}_m)`$, for $`\overline{y}J^1\pi `$, where $`Z_1,\mathrm{},Z_m\mathrm{T}_{\pi ^1(\overline{y})}E`$, and $`\overline{Z}_1,\mathrm{},\overline{Z}_m\mathrm{T}_{\overline{y}}J^1\pi `$ are such that $`\mathrm{T}_{\overline{y}}\pi ^1\overline{Z}_\alpha =Z_\alpha `$. Then, using the natural projection $`\mu :\pi J^1\pi ^{}`$, we define the restricted Legendre map associated with $``$ as $`:=\mu \stackrel{~}{}`$. Their local expressions are
$$\begin{array}{ccccccc}\stackrel{~}{}^{}x^\nu =x^\nu & ,& \stackrel{~}{}^{}y^A=y^A& ,& \stackrel{~}{}^{}p_A^\nu =\frac{L}{v_\nu ^A}& ,& \stackrel{~}{}^{}p=Lv_\nu ^A\frac{L}{v_\nu ^A}\\ ^{}x^\nu =x^\nu & ,& ^{}y^A=y^A& ,& ^{}p_A^\nu =\frac{L}{v_\nu ^A}& & \end{array}$$
We have that $`\stackrel{~}{}^{}\mathrm{\Theta }=\mathrm{\Theta }_{}`$, and $`\stackrel{~}{}^{}\mathrm{\Omega }=\mathrm{\Omega }_{}`$.
$`(J^1\pi ,\mathrm{\Omega }_{})`$ is a regular Lagrangian system if $``$ is a local diffeomorphism (this definition is equivalent to that given above). Elsewhere $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a singular Lagrangian system. As a particular case, $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a hyper-regular Lagrangian system if $``$ is a global diffeomorphism. A singular Lagrangian system $`(J^1\pi ,\mathrm{\Omega }_{})`$ is almost-regular if $`𝒫:=(J^1\pi )`$ is a closed submanifold of $`J^1\pi ^{}`$, $``$ is a submersion onto its image, and for every $`\overline{y}J^1\pi `$, the fibres $`^1((\overline{y}))`$ are connected submanifolds of $`J^1\pi `$.
If $`(J^1\pi ,\mathrm{\Omega }_{})`$ is an almost-regular Lagrangian system then $`𝒫`$ is a fibre bundle over $`E`$ and $`M`$ (the natural projections are denoted by $`\tau _0^1:𝒫E`$ and $`\overline{\tau }_0^1:=\pi \tau _0^1:𝒫M`$) and the $`\mu `$-transverse submanifold $`\stackrel{~}{𝒫}=\stackrel{~}{}(J^1\pi )\pi `$ is diffeomorphic to $`𝒫`$ (and we denote by $`\stackrel{~}{ȷ}_0:\stackrel{~}{𝒫}\pi `$ the natural imbedding). This diffeomorphism is denoted $`\stackrel{~}{\mu }:\stackrel{~}{𝒫}𝒫`$, and it is just the restriction of the projection $`\mu `$ to $`\stackrel{~}{𝒫}`$. Then, taking $`\stackrel{~}{h}:=\stackrel{~}{\mu }^1`$, we define the Hamilton-Cartan $`(m+1)`$-form $`\mathrm{\Omega }_h^0=(\stackrel{~}{ȷ}_0\stackrel{~}{h})^{}\mathrm{\Omega }`$, which verifies that $`_0^{}\mathrm{\Omega }_h^0=\mathrm{\Omega }_{}`$ (where $`_0`$ is the restriction map of $``$ onto $`𝒫`$). Then $`\stackrel{~}{h}`$ is called a Hamiltonian section, and $`(𝒫,\mathrm{\Omega }_\mathrm{h}^0)`$ is the Hamiltonian system associated with the almost-regular Lagrangian system $`(J^1\pi ,\mathrm{\Omega }_{})`$ (see ).
If $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a hyper-regular Lagrangian system, then $`𝒫=J^1\pi ^{}`$, and the construction is the same. In addition, $`\stackrel{~}{}(J^1\pi )`$ is a 1-codimensional embedded submanifold of $`\pi `$, which is transverse to the projection $`\mu `$, and is diffeomorphic to $`J^1\pi ^{}`$. This diffeomorphism is $`\mu ^1`$, when $`\mu `$ is restricted to $`\stackrel{~}{}(J^1\pi )`$, and coincides with the map $`h:=\stackrel{~}{}^1`$, when it is restricted onto its image. $`h`$ is the Hamiltonian section in this case, and the associated Hamiltonian system is denoted by $`(J^1\pi ^{},\mathrm{\Omega }_h)`$, where $`\mathrm{\Omega }_h=h^{}\mathrm{\Omega }`$. In a local chart of natural coordinates, the Hamiltonian section is specified by a local Hamiltonian function $`H\mathrm{C}^{\mathrm{}}(U)`$, $`UJ^1\pi ^{}`$, such that $`h(x^\alpha ,y^A,p_A^\alpha )(x^\alpha ,y^A,p_A^\alpha ,p=H)`$, where
$$H(x^\alpha ,y^A,p_A^\alpha )=(^1)^{}\left(v_\alpha ^A\frac{L}{v_\alpha ^A}L\right)=p_A^\alpha (^1)^{}v_\alpha ^A(^1)^{}L$$
and $`\mathrm{\Omega }_h=\mathrm{d}p_A^\alpha \mathrm{d}y^A\mathrm{d}^{m1}x_\alpha +\mathrm{d}H\mathrm{d}^mx`$.
The Hamiltonian problem associated with the Hamiltonian system $`(𝒫,\mathrm{\Omega }_\mathrm{h}^0)`$ (for $`(J^1\pi ^{},\mathrm{\Omega }_h)`$ is analogous), consists in finding sections $`\psi _o\mathrm{\Gamma }(M,𝒫)`$ such that
$$\psi _o^{}𝑖(X_0)\mathrm{\Omega }_h^0=0,\text{for every }X_0\text{X}(𝒫)$$
As in the Lagrangian case, these sections are the integral sections of a class of integrable and $`\overline{\tau }_0^1`$-transverse $`m`$-vector fields $`\{𝒳__o\}\text{X}^m(𝒫)`$ satisfying that
$$𝑖(𝒳__o)\mathrm{\Omega }_h^0=0,\text{for every }𝒳__o\{𝒳__o\}.$$
$`m`$-vector fields satisfying these conditions (but not necessarily integrable) are called Hamilton-De Donder-Weyl $`m`$-vector fields for $`(𝒫,\mathrm{\Omega }_\mathrm{h}^0)`$.
### 4.2 Lagrangian and Hamiltonian algorithms
Let $`(J^1\pi ,\mathrm{\Omega }_{})`$ be a Lagrangian system. If $``$ is an Ehresmann connection in the fibration $`\overline{\pi }^1:J^1\pi M`$, let $`𝒴_\eta ^{}`$ be the corresponding $`m`$-vector field on $`J^1\pi `$. Then, we have:
###### Proposition 7
The Poincaré-Cartan $`(m+1)`$-form may be written as
$$\mathrm{\Omega }_{}=\omega (\gamma _{})_\eta ^{}+\mathrm{\Omega }_{}^{},$$
where $`(\gamma _{})_\eta ^{}=𝑖(𝒴_\eta ^{})\mathrm{\Omega }_{}\Omega ^1(J^1\pi )`$, and $`\mathrm{\Omega }_{}^{}`$ is a $`(m+1)`$-form on $`J^1\pi `$ of bidegree $`(m1,2)`$ with respect to the connection $``$.
( Proof ) If $`yJ^1\pi `$ and $`v_1,v_2,v_3V_y(\overline{\pi }^1)`$ ($`V(\overline{\pi }^1)`$ being the vertical bundle of $`\overline{\pi }^1`$) then, from (18) we have that $`i(v_1v_2v_3)\mathrm{\Omega }_{}(y)=0`$. Thus, the result follows from Propositions 2 and 5.
If $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a hyperregular Lagrangian system (the regular case is analogous) $``$ is a global diffeomorphism. Moreover, if $`\mathrm{\Omega }_h`$ is the Hamilton-Cartan $`(m+1)`$-form on $`J^1\pi ^{}`$, then
$$^{}\mathrm{\Omega }_h=\mathrm{\Omega }_{}.$$
(19)
Furthermore, as $``$ is a global diffeomorphism, the connection $``$ induces a connection $`^{}`$ in the fibration $`\overline{\tau }^1:J^1\pi ^{}M`$ in such a way that
$$_{}𝒴_\eta ^{}=𝒴_\eta ^{^{}},$$
(20)
where $`𝒴_\eta ^{^{}}`$ is the $`m`$-vector field on $`J^1\pi ^{}`$ associated with $`^{}`$ and the volume form $`\eta `$. Thus, from (19), (20) and Proposition 7, we obtain:
###### Proposition 8
The Hamilton-Cartan $`(m+1)`$-form may be written as
$$\mathrm{\Omega }_h=\mathrm{\Omega }_h^{^{}}+\omega ^{}(\gamma _h)_\eta ^{^{}},$$
where $`(\gamma _h)_\eta ^{^{}}=𝑖(𝒴_\eta ^{^{}})\mathrm{\Omega }_h`$, and $`\mathrm{\Omega }_h^{^{}}`$ is a $`(m+1)`$-form on $`J^1\pi ^{}`$ of bidegree $`(m1,2)`$ with respect to the connection $`^{}`$.
Furthermore, we may prove the following result:
###### Proposition 9
If $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a regular Lagrangian system, then the $`(m+1)`$-forms $`\mathrm{\Omega }_{}^{}`$ and $`\mathrm{\Omega }_h^{^{}}`$ are $`1`$-nondegenerate.
( Proof ) As $``$ is a diffeomorphism and $`^{}\mathrm{\Omega }_h^{^{}}=\mathrm{\Omega }_{}^{}`$, it suffices to prove that $`\mathrm{\Omega }_h^{^{}}`$ is $`1`$-nondegenerate. The local expression of $`\mathrm{\Omega }_h^{^{}}`$ is
$$\mathrm{\Omega }_h^{^{}}=dp_A^\alpha dy^Ad^{m1}x_\alpha +\theta d^mx,$$
$`\theta `$ being a $`1`$-form such that $`\theta \left({\displaystyle \frac{}{x_\alpha }}\right)=0`$, for every $`\alpha `$. As a consequence,
$$𝑖(\frac{}{x_\beta })\mathrm{\Omega }_h^{^{}}=\underset{A,\alpha ;\alpha \beta }{}dp_A^\alpha dy^Ad^{m2}x_{\alpha \beta }\theta \left(\frac{}{y^A}\right)dy^Ad^{m1}x_\beta \theta \left(\frac{}{p_A^\alpha }\right)dp_A^\alpha d^{m1}x_\beta $$
(21)
$$𝑖\left(\frac{}{y^A}\right)\mathrm{\Omega }_h^{^{}}=\underset{\alpha }{}dp_A^\alpha d^{m1}x_\alpha +\theta \left(\frac{}{y^A}\right)d^mx,$$
(22)
$$𝑖\left(\frac{}{p_A^\alpha }\right)\mathrm{\Omega }_h^{^{}}=dy^Ad^{m1}x_\alpha +\theta \left(\frac{}{p_A^\alpha }\right)d^mx.$$
(23)
Thus, if $`X=\lambda _\beta {\displaystyle \frac{}{x_\beta }}+\mu ^A{\displaystyle \frac{}{y^A}}+\nu _A^\alpha {\displaystyle \frac{}{p_A^\alpha }}`$ is a local vector field such that $`𝑖(X)\mathrm{\Omega }_h^{^{}}=0`$ then, from (21), it follows that $`\lambda _\beta =0`$, for every $`\beta `$, which implies that (see (22) and (23))
$$\mu ^Adp_A^\alpha d^{m1}x_\alpha +\mu ^A\theta (\frac{}{y^A})d^mx\nu _A^\alpha dy^Ad^{m1}x_\alpha +\nu _\alpha ^A\theta (\frac{}{p_A^\alpha })d^mx=0.$$
Therefore, $`\mu ^A=0`$ and $`\nu _A^\alpha =0`$, for every $`A`$ and $`\alpha `$, that is, $`X=0`$.
If the Lagrangian is regular, then from Propositions 6 and 9, we obtain that $`(T_y^{}J^1\pi )_\mathrm{\Omega }_{}^{}=\{0\}`$, for every $`yJ^1\pi `$. Thus, there exist locally decomposable $`m`$-vector fields $`𝒳_{}`$ on $`J^1\pi `$ such that
$$i(𝒳_{})\omega =1,i(𝒳_{})\mathrm{\Omega }_{}=0.$$
Moreover, we have
###### Proposition 10
If $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a regular Lagrangian system and $`𝒳_{}`$ is a locally decomposable $`m`$-vector field on $`J^1\pi `$ such that $`i(𝒳_{})\omega =1`$ and $`i(𝒳_{})\mathrm{\Omega }_{}=0`$ then $`𝒳_{}`$ is an Euler-Lagrange $`m`$-vector field for $``$.
( Proof ) We must prove that $`𝒳_{}`$ is semi-holonomic, that is, $`i(𝒳_{})𝒱=0.`$ For this purpose, we consider local fibred coordinates $`(x^\alpha ,y^A,v_\alpha ^A)`$ on $`J^1\pi `$. Then, since $`i(𝒳_{})\omega =1`$, it follows that
$$𝒳_{}=\mathrm{\Lambda }_{\alpha =1}^m\left(\frac{}{x^\alpha }+\mathrm{\Gamma }_\alpha ^A\frac{}{y^A}+\mathrm{\Gamma }_{\alpha \beta }^A\frac{}{v_\beta ^A}\right)$$
with $`\mathrm{\Gamma }_\alpha ^A`$ and $`\mathrm{\Gamma }_{\alpha \beta }^A`$ local real functions on $`J^1\pi `$. Furthermore, from (18), we deduce that
$$0=(i(𝒳_{})\mathrm{\Omega }_{})(\frac{}{v_\nu ^B})=(1)^mi(𝒳_{})(i(\frac{}{v_\nu ^B})\mathrm{\Omega }_{})$$
$$=(1)^{m+1}(\mathrm{\Gamma }_\alpha ^Av_\alpha ^A)\frac{^2L}{v_\alpha ^Av_\nu ^B},\text{ for all }B\text{ and }\nu .$$
Therefore, using the fact that $``$ is regular, we conclude that
$$\mathrm{\Gamma }_\alpha ^A=v_\alpha ^A,\text{ for all }A\text{ and }\alpha ,$$
which implies that $`𝒳_{}`$ is semi-holonomic.
Hence, if $`(J^1\pi ,\mathrm{\Omega }_{})`$ is a regular Lagrangian system, then the existence of classes of Euler-Lagrange $`m`$-vector fields for $``$ is assured in $`J^1\pi `$. In the same way, for the Hamiltonian formalism, the existence of Hamilton-De Donder-Weyl $`m`$-vector fields is assured everywhere in $`J^1\pi ^{}`$ (note that if $`𝒳_{}`$ is an Euler-Lagrange $`m`$-vector field for $``$ then $`()_{}𝒳_{}`$ is a Hamilton-De Donder-Weyl $`m`$-vector field on $`J^1\pi ^{}`$). In both cases, the solution is not unique.
For singular (almost-regular) Lagrangian systems, the existence of Euler-Lagrange $`m`$-vector fields is not assured except perhaps on some submanifold $`S_fJ^1\pi `$, where the solution is not unique. In order to find this submanifold we apply the algorithm developed in Section 3.3 to the system $`(J^1\pi ,\mathrm{\Omega }_{})`$, by doing the identifications $`\kappa :FM`$ with $`\overline{\pi }^1:J^1\pi M`$, and $`\mathrm{\Omega }`$ with $`\mathrm{\Omega }_{}`$. Thus we obtain obtain a sequence
$$\mathrm{}\stackrel{j_{i+1}^i}{}N_i\stackrel{j_i^{i1}}{}\mathrm{}\stackrel{j_2^1}{}N_1\stackrel{j_1}{}N_0J^1\pi .$$
(24)
which, in the best of cases stabilizes in the final constraint submanifold $`N_f`$ where there exist $`m`$-vector fields $`𝒳^{N_f}`$ on $`N_f`$, solution to the equations
$$(𝑖(𝒳^{N_f})\mathrm{\Omega }_{})|_{N_f}=0,(𝑖(𝒳^{N_f})\omega )|_{N_f}=1.$$
(25)
But $`𝒳^{N_f}`$ will not be, in general, an Euler-Lagrange $`m`$-vector field on $`N_f`$ (that is, it is not semi-holonomic), and, in addition, $`𝒳^{N_f}`$ will not in general be an integrable $`m`$-vector field. The problem of finding integrable Euler-Lagrange $`m`$-vector fields (i.e., holonomic) is discussed and solved in the next Section.
Now, we consider the Hamiltonian system $`(𝒫,\mathrm{\Omega }_\mathrm{h}^0)`$. Let $`_0^{}`$ be a connection in the bundle $`\overline{\tau }_0^1:𝒫M`$ and denote by $`𝒴_\eta ^_0^{}`$ the corresponding $`m`$-vector field on $`𝒫`$ associated with $`_0^{}`$ and $`\eta `$. Then, we have:
###### Proposition 11
The Hamilton-Cartan $`(m+1)`$-form may be written as
$$\mathrm{\Omega }_h^0=\mathrm{\Omega }_h^_0^{}+\omega _0^{}(\gamma _h)_\eta ^_0^{},$$
where $`\omega _0^{}=\overline{\tau }_0^1\eta `$, $`(\gamma _h)_\eta ^_0^{}=𝑖(𝒴_\eta ^_0^{})\mathrm{\Omega }_h`$, and $`\mathrm{\Omega }_h^_0^{}`$ is a $`(m+1)`$-form on $`𝒫`$ of bidegree $`(m1,2)`$ with respect to the connection $`_0^{}`$.
( Proof ) If $`\overline{y}=_0(y)𝒫`$, with $`yJ^1\pi `$, and $`\overline{v}_1,\overline{v}_2,\overline{v}_3V_y(\overline{\tau }_0^1)`$ then, since $`_0:J^1\pi 𝒫`$ is a submersion and $`\overline{\tau }_0^1_0=\overline{\pi }^1`$, it follows that there exist $`v_1,v_2,v_3V_y(\overline{\pi }^1)`$ such that
$$(T_y_0)(v_i)=\overline{v}_i,\text{ for }i\{1,2,3\}.$$
Thus, using that $`(_0)^{}\mathrm{\Omega }_h^0=\mathrm{\Omega }_{}`$, we deduce that
$$i(\overline{v}_1\overline{v}_2\overline{v}_3)\mathrm{\Omega }_h^0(\overline{y})=0.$$
This proves the result.
Hamilton-De Donder-Weyl $`m`$-vector fields do not exist, in general, in $`𝒫`$, and then we must apply the algorithmic procedure developed in Section 3.3 to the system $`(𝒫,\mathrm{\Omega }_\mathrm{h}^0)`$, by doing the identifications $`\kappa :FM`$ with $`\overline{\tau }^1|_𝒫:𝒫M`$, and $`\mathrm{\Omega }`$ with $`\mathrm{\Omega }_h^0`$. Thus we obtain a sequence
$$\mathrm{}\stackrel{j_{i+1}^i}{}P_i\stackrel{j_i^{i1}}{}\mathrm{}\stackrel{j_2^1}{}P_1\stackrel{j_1}{}P_0𝒫.$$
(26)
which, in the best of cases stabilizes in the final constraint submanifold $`P_f`$ of $`𝒫`$ where there exist $`m`$-vector fields $`𝒳^{P_f}`$ on $`P_f`$, solution to the equations
$$(𝑖(𝒳^{P_f})\mathrm{\Omega }_h^0)|_{P_f}=0,(𝑖(𝒳^{P_f})\omega _0^{})|_{P_f}=1.$$
(27)
Of course the solution $`𝒳^{P_f}`$ is not unique.
###### Remark 3
The Lagrangian and Hamiltonian pre-multisymplectic algorithms are equivalent in the following sense: at every level $`j`$ of the Lagrangian and Hamiltonian algorithms, the submanifolds of the sequences (24) and (26) are $``$-related, that is, $`(N_j)=P_j`$ and $`_j=_{|N_j}:N_jP_j`$ is a submersion such that $`_j^1(_j(x_j))=_0^1(_0(x_j))`$, for $`x_jN_j`$. Moreover, if $`N_f`$ is the final constraint submanifold (in the Lagrangian level) and $`𝒳^{N_f}`$ is a locally decomposable $`m`$-vector field on $`N_f`$ such that equations (25) hold and, in addition, $`𝒳^{N_f}`$ is $`_f`$-projectable to an $`m`$-vector field $`𝒳^{P_f}`$ on $`P_f`$ then $`𝒳^{P_f}`$ is locally decomposable and equations (27) hold. Conversely, if $`𝒳^{P_f}`$ is a locally decomposable $`m`$-vector field on $`P_f`$ satisfying equations (27) and $`𝒳^{N_f}`$ is a locally decomposable $`m`$-vector field on $`N_f`$ which is $`_f`$-projectable on $`𝒳^{P_f}`$ then $`𝒳^{N_f}`$ satisfies equations (25) (see for a detailed discussion on this topic).
Finally, the Hamilton-De Donder-Weyl $`m`$-vector fields $`𝒳^{P_f}`$ are not integrable, in general. In fact, if we have that $`𝒳^{P_f}=X_1^{P_f}\mathrm{}X_m^{P_f}`$, where $`X_\alpha ^{P_f}`$ are (local) vector fields on $`P_f`$, for all $`\alpha `$, and $`\{X_1^{P_f},\mathrm{},X_m^{P_f},\overline{Z}_1,\mathrm{},\overline{Z}_p\}`$ is a local basis of the vector bundle $`TP_fP_f`$ then
$$[X_\alpha ^{P_f},X_\beta ^{P_f}]=\overline{f}_{\alpha \beta }^\gamma X_\gamma ^{P_f}+\overline{\zeta }_{\alpha \beta }^l\overline{Z}_l$$
for some functions $`\overline{f}_{\alpha \beta }^\gamma `$ and $`\overline{\zeta }_{\alpha \beta }^l`$ on $`P_f`$. Therefore, we must apply the integrability algorithm of Section 3.4, and we obtain a sequence $`\mathrm{}𝒥_i\mathrm{}𝒥_1P_f`$, such that $`𝒥_i`$ is a non-empty (closed) submanifold of $`S_f`$, with
$`𝒥_1`$ $`=`$ $`\{yP_f|\overline{\zeta }_{\alpha \beta }^l(y)=0\}`$
$`𝒥_i`$ $`=`$ $`\{y𝒥_{i1}|𝒳^{P_f}(y)\mathrm{\Lambda }^m\mathrm{T}_y𝒥_{i1}\},\text{for}i2.`$
In the best cases, there exists an integer $`i`$ such that $`𝒥_{i+1}=𝒥_i`$. Then, $`𝒥_f=𝒥_{i+1}=𝒥_i`$ is a submanifold of $`P_f`$, and $`𝒳^{𝒥_f}=(𝒳^{P_f})|_{𝒥_f}`$ is an integrable Hamilton-De Donder-Weyl $`m`$-vector field in $`𝒥_f`$.
### 4.3 Almost-regular Lagrangians and integrable Euler-Lagrange $`m`$-vector fields
Let $`(J^1\pi ,\mathrm{\Omega }_{})`$ be an almost-regular Lagrangian system, and $`N_f`$ the final constraint submanifold (in the Lagrangian setting). Then, there exists a locally decomposable $`m`$-vector field $`𝒳^{N_f}`$ on $`N_f`$ such that
$$(𝑖(𝒳^{N_f})\mathrm{\Omega }_{})|_{N_f}=0,(𝑖(𝒳^{N_f})\omega )|_{N_f}=1.$$
But, in general, $`𝒳^{N_f}`$ is not an Euler-Lagrange $`m`$-vector field on $`N_f`$ and, in addition, $`𝒳^{N_f}`$ will not in general be an integrable $`m`$-vector field.
In order to solve these problems, first we construct a submanifold $`S_f`$ of $`N_f`$ where there exists a locally decomposable $`m`$-vector field $`𝒳^{S_f}`$ such that
$$(𝑖(𝒳^{S_f})\mathrm{\Omega }_{})|_{S_f}=0,(𝑖(𝒳^{S_f})\omega )|_{S_f}=1,(𝑖(𝒳^{S_f})𝒱)|_{S_f}=0.$$
In fact, from the above discussion we know that we can choose the $`m`$-vector field $`𝒳^{N_f}`$ on $`N_f`$ such that it projects via $`_f`$ (the restriction of $``$ to $`N_f`$) onto an $`m`$-vector field $`𝒳^{P_f}`$ on $`P_f`$. Then, we consider the subset $`S_f`$ of $`N_f`$ defined by
$$S_f=\{xN_f/(𝑖(𝒳^{N_f})𝒱)(x)=0\}.$$
(28)
In (see also ), it was proved that
$$(𝑖(𝒳^{N_f})𝒱)(x)Ker\mathrm{T}_x(_f)=Ker\mathrm{T}_x(_0)$$
(29)
and that for every $`xN_f`$, $`S_f_f^1(_f(x))=S_f_0^1(_0(x))`$ is a single point in $`S_f`$.
The above result allows us to introduce a well-defined map $`s_f:P_fN_f`$ such that
$$S_f=s_f(P_f),_fs_f=Id.$$
Thus, $`s_f:P_fN_f`$ is a global section of the submersion $`_f:N_fP_f`$ and, therefore, $`S_f`$ is an embedded submanifold of $`N_f`$ and the map $`s_f:P_fS_f`$ is a diffeomorphism (for more details, see ).
Now, defining the $`m`$-vector field $`𝒳^{S_f}`$ on $`S_f`$ by $`𝒳^{S_f}=(\mathrm{\Lambda }^m\mathrm{T}s_f)𝒳^{P_f}`$, then we have :
###### Theorem 6
$`𝒳^{S_f}`$ is an Euler-Lagrange $`m`$-vector field on $`S_f`$ for the Lagrangian $``$, that is, $`𝒳^{S_f}`$ is a locally decomposable $`m`$-vector field on $`S_f`$ and
$$(𝑖(𝒳^{S_f})\mathrm{\Omega }_{})|_{S_f}=0,(𝑖(𝒳^{S_f})\omega )|_{S_f}=1,(𝑖(𝒳^{S_f})𝒱)|_{S_f}=0.$$
Next, we give a local description of the submanifold $`S_f`$ and of the Euler-Lagrange $`m`$-vector field $`𝒳^{S_f}`$ on $`S_f`$. Since $``$ is almost-regular, it follows that the rank of the partial Hessian matrix $`\left({\displaystyle \frac{^2L}{v_\alpha ^Av_\beta ^B}}\right)`$ is constant. Let $`rank\left({\displaystyle \frac{^2L}{v_\alpha ^Av_\beta ^B}}\right)=pm+q`$, with $`0pn1`$ and $`0qm`$, and assume that the first $`pm+q`$ rows of this matrix are independent. Denote by $`\stackrel{~}{V(\pi ^1)}J^1\pi `$ the vector subbundle of the vertical bundle $`V(\pi ^1)J^1\pi `$ of $`\pi ^1:J^1\pi M`$ generated by the local vector fields
$$\{\frac{}{v_\alpha ^A},\frac{}{v_1^{p+1}},\mathrm{},\frac{}{v_q^{p+1}}\},\text{for}\mathrm{\hspace{0.33em}1}Ap\text{and}\mathrm{\hspace{0.33em}\hspace{0.33em}1}\alpha m.$$
Then, there exist sections $`\{X_{q+1}^{p+1},\mathrm{},X_m^{p+1},X_\alpha ^A\}`$, with $`p+2An`$ and $`1\alpha m`$, of the vector bundle $`\stackrel{~}{V(\pi ^1)}J^1\pi `$ such that $`\{W_{q+1}^{p+1},\mathrm{},W_m^{p+1},W_\alpha ^A\}`$, with $`p+2An`$ and $`1\alpha m`$, is a local basis of $`Ker(\mathrm{T}(_0))`$, where
$$W_\beta ^{p+1}=\frac{}{v_\beta ^{p+1}}+X_\beta ^{p+1},W_\alpha ^A=\frac{}{v_\alpha ^A}+X_\alpha ^A,(p+2An,1\alpha m,q+1\beta m).$$
(30)
Now, suppose that $`𝒳^{N_f}=𝒳_1^{N_f}\mathrm{}𝒳_m^{N_f}`$, with
$$𝒳_\alpha ^{N_f}=\left(\frac{}{x^\alpha }+\mathrm{\Gamma }_\alpha ^A\frac{}{y^A}+\mathrm{\Gamma }_{\alpha \beta }^A\frac{}{v_\beta ^A}\right)|_{N_f},\text{for}\alpha \{1,\mathrm{},m\}.$$
Then, using that $`𝒳^{N_f}`$ is $`_f`$-projectable, it follows that the functions $`\mathrm{\Gamma }_\alpha ^A`$ are constant on the fibers of $`_f:N_fP_f`$. But, as $`_f^1(_f(x))=_0^1(_0(x))`$, for every $`xN_f`$ (see Remark 3), we obtain that
$$\begin{array}{ccc}W_\gamma ^{p+1}(\mathrm{\Gamma }_\alpha ^A)=0,\hfill & & \hfill \gamma \{q+1,\mathrm{},m\}\\ W_\gamma ^{p+1+i}(\mathrm{\Gamma }_\alpha ^A)=0,\hfill & & \hfill i\{1,\mathrm{},np1\},\gamma \{1,\mathrm{},m\}.\end{array}$$
(31)
Furthermore
$$𝑖(𝒳^{N_f})𝒱=(\mathrm{\Gamma }_\alpha ^Av_\alpha ^A)\frac{}{v_\alpha ^A}.$$
Thus, from (29) and (30), we have that
$$𝑖(𝒳^{N_f})𝒱=(\mathrm{\Gamma }_\gamma ^{p+1}v_\gamma ^{p+1})W_\gamma ^{p+1}+(\mathrm{\Gamma }_\gamma ^{p+1+i}v_\gamma ^{p+1+i})W_\gamma ^{p+1+i}.$$
(32)
Note that the functions
$$\begin{array}{ccc}\zeta _\gamma ^{p+1}=\mathrm{\Gamma }_\gamma ^{p+1}v_\gamma ^{p+1},\hfill & & \hfill \gamma \{q+1,\mathrm{},m\}\\ \zeta _{\overline{\gamma }}^{p+1+i}=\mathrm{\Gamma }_{\overline{\gamma }}^{p+1+i}v_{\overline{\gamma }}^{p+1+i},\hfill & & \hfill i\{1,\mathrm{},np1\},\overline{\gamma }\{1,\mathrm{},m\}.\end{array}$$
(33)
are independent on $`N_f`$. In fact (see (30), (31) and (33))
$$\begin{array}{ccc}W_\gamma ^{p+1}(\zeta _\gamma ^{}^{p+1})=\delta _{\gamma \gamma ^{}},\hfill & & \hfill W_\gamma ^{p+1}(\zeta _{\overline{\gamma }}^{p+1+i})=0,\\ W_{\overline{\gamma }}^{p+1+i}(\zeta _\gamma ^{p+1})=0,\hfill & & \hfill W_{\overline{\alpha }}^{p+1+i}(\zeta _{\overline{\gamma }}^{p+1+j})=\delta _{ij}\delta _{\overline{\alpha }\overline{\gamma }}.\end{array}$$
Moreover, using (28) and (32), we conclude that $`\{\zeta _\gamma ^{p+1},\zeta _{\overline{\gamma }}^{p+1+i}\}`$, with $`\gamma \{1,\mathrm{},m\}`$, $`i\{1,\mathrm{},np1\}`$ and $`\overline{\gamma }\{1,\mathrm{},m\}`$, is a set of local independent constraint functions defining $`S_f`$ as a submanifold of $`N_f`$, that is,
$$S_f=\{xN_f/(\mathrm{\Gamma }_\gamma ^{p+1}v_\gamma ^{p+1})(x)=0,(\mathrm{\Gamma }_{\overline{\gamma }}^{p+1+i}v_{\overline{\gamma }}^{p+1+i})(x)=0\}.$$
Finally, a direct calculation proves that the Euler-Lagrange $`m`$-vector field $`𝒳^{S_f}`$ on $`S_f`$ is given by $`𝒳^{S_f}=𝒳_1^{S_f}\mathrm{}𝒳_m^{S_f}`$, with
$$𝒳_\alpha ^{S_f}=(𝒳_\alpha ^{N_f}+𝒳_\alpha ^{N_f}(\zeta _\gamma ^{p+1})W_\gamma ^{p+1}+𝒳_\alpha ^{N_f}(\zeta _{\overline{\gamma }}^{p+1+i})W_{\overline{\gamma }}^{p+1+i})|_{S_f},\text{for every}\alpha .$$
$`𝒳^{S_f}`$ is not, in general, integrable. In fact, if $`\{𝒳_1^{S_f},\mathrm{},𝒳_m^{S_f},Z_1,\mathrm{},Z_s\}`$ is a local basis of the vector bundle $`TS_fS_f`$ then we have that
$$[𝒳_\alpha ^{S_f},𝒳_\beta ^{S_f}]=f_{\alpha \beta }^\gamma 𝒳_\gamma ^{S_f}+\zeta _{\alpha \beta }^lZ_l,$$
for some functions $`f_{\alpha \beta }^\gamma `$ and $`\zeta _{\alpha \beta }^l`$.
Therefore, we must apply the integrability algorithm of Section 3.4. Then, we obtain a sequence $`\mathrm{}_i\mathrm{}_1S_f`$, such that $`_i`$ is a non-empty (closed) submanifold of $`S_f`$, with
$`_1`$ $`=`$ $`\{xS_f|\zeta _{\alpha \beta }^\gamma (x)=0\}`$
$`_i`$ $`=`$ $`\{x_{i1}|𝒳^{S_f}(x)\mathrm{\Lambda }^m\mathrm{T}_x_{i1}\},\text{for}i2`$
In the best cases, there exists an integer $`i`$ such that $`_{i+1}=_i`$. Then, $`_f=_{i+1}=_i`$ is a submanifold of $`S_f`$ and $`𝒳^_f=(𝒳^{S_f})|__f`$ is an integrable Euler-Lagrange $`m`$-vector field on $`_f`$, and hence it is holonomic. In fact:
###### Theorem 7
If $`U`$ is an open subset of $`M`$ and $`s:UM_f`$ is an integral section of $`𝒳^_f`$ then there exists a section $`\varphi :UME`$ of the projection $`\pi :EM`$ such that $`s=j^1\varphi `$ and $`\varphi `$ is a solution to the Euler-Lagrange equations for $``$.
( Proof ) We have that
$$(𝑖(𝒳^f)\mathrm{\Omega }_{})|__f=0,$$
(34)
$$(𝑖(𝒳^f)𝒱)|__f=0.$$
(35)
We can assume, without loss of generality, that $`s(U)\stackrel{~}{U}`$, with $`\stackrel{~}{U}`$ an open subset of $`J^1\pi `$ and $`(x^\alpha ,y^A,v_\alpha ^A)`$ a system of local coordinates on $`\stackrel{~}{U}`$. Then, since $`𝒳^_f`$ is locally decomposable and $`𝑖(𝒳^f)\omega |__f=1,`$ we deduce that
$$𝒳^_f|_{\stackrel{~}{U}_f}=𝒳_1^_f\mathrm{}𝒳_m^_f,$$
(36)
with $`𝒳_\alpha ^_f𝔛(\stackrel{~}{U}_f)`$ given by
$$𝒳_\alpha ^_f=\left(\frac{}{x^\alpha }+\mathrm{\Gamma }_\alpha ^A\frac{}{y^A}+\mathrm{\Gamma }_{\alpha \beta }^A\frac{}{v_\beta ^A}\right)|_{\stackrel{~}{U}_f},$$
(37)
for all $`\alpha `$, where $`\mathrm{\Gamma }_\alpha ^A`$ and $`\mathrm{\Gamma }_{\alpha \beta }^A`$ are local functions on $`\stackrel{~}{U}`$. Now, using that
$$(𝑖(𝒳^_f)𝒱)|_{\stackrel{~}{U}_f}=\left((\mathrm{\Gamma }_\alpha ^Av_\alpha ^A)\frac{}{v_\alpha ^A}\right)|_{\stackrel{~}{U}_f},$$
it follows that (see (35))
$$𝒳_\alpha ^_f=\left(\frac{}{x^\alpha }+v_\alpha ^A\frac{}{y^A}+\mathrm{\Gamma }_{\alpha \beta }^A\frac{}{v_\beta ^A}\right)|_{\stackrel{~}{U}_f}.$$
(38)
Furthermore, from (18), we obtain that
$$\begin{array}{ccc}𝑖(\frac{}{y^A})\mathrm{\Omega }_{}\hfill & =\hfill & \frac{^2L}{v_\nu ^Bv_\alpha ^A}\mathrm{d}v_\nu ^B\mathrm{d}^{m1}x^\alpha +\left(\frac{^2L}{y^Bv_\alpha ^A}\frac{^2L}{y^Av_\alpha ^B}\right)\mathrm{d}y^B\mathrm{d}^{m1}x^\alpha \hfill \\ & & +\left(\frac{^2L}{y^Av_\alpha ^B}v_\alpha ^B\frac{L}{y^A}+\frac{^2L}{x^\alpha v_\alpha ^A}\right)\mathrm{d}^mx\hfill \end{array}$$
(39)
Therefore, using (34), (36), (37) and (39), we conclude that
$$\frac{^2L}{x^\alpha v_\alpha ^A}+\frac{^2L}{y^Bv_\alpha ^A}v_\alpha ^B+\frac{^2L}{v_\nu ^Bv_\alpha ^A}\mathrm{\Gamma }_{\alpha \nu }^B\frac{L}{y^A}=0,\text{for every}A.$$
(40)
Next, suppose that $`U`$ is an open subset of $`M`$ and that $`s:UM\stackrel{~}{U}_fJ^1\pi `$ is an integral section of $`𝒳^_f|_{\stackrel{~}{U}_f}`$ such that the local expression of $`s`$ is $`s(x^\beta )=(x^\beta ,s^A(x^\beta ),s_\alpha ^A(x^\beta ))`$. Using (38) and the fact that
$$(\mathrm{T}s)\left(\frac{}{x^\beta }\right)=\left(\frac{}{x^\beta }+\frac{s^A}{x^\beta }\frac{}{y^A}+\frac{s_\alpha ^A}{x^\beta }\frac{}{v_\beta ^A}\right)|_{\stackrel{~}{U}_f},\text{for every}\beta .$$
we deduce that
$$s_\alpha ^A=\frac{s^A}{x^\alpha },\mathrm{\Gamma }_{\alpha \beta }^As=\frac{^2s^A}{x^\alpha x^\beta },\text{for every}A,\alpha ,\beta .$$
(41)
From (41), it follows that there exists $`\varphi :UME`$ a local section of $`\pi :EM`$ such that $`s=j^1\varphi `$. Moreover, using (40) and (41), we obtain that
$$\frac{^2L}{x^\alpha v_\alpha ^A}+\frac{^2L}{y^Bv_\alpha ^A}\frac{s^B}{x^\alpha }+\frac{^2L}{v_\nu ^Bv_\alpha ^A}\frac{^2s^B}{x^\alpha x^\nu }\frac{L}{y^A}=0,\text{for every}A.$$
This implies that
$$(j^i\varphi )^{}\left(\frac{L}{y^A}\frac{d}{dx^\alpha }\frac{L}{v_\alpha ^A}\right)=0,\text{for every}A.$$
In other words, $`\varphi `$ is a solution to the Euler-Lagrange equations associated with $``$.
###### Remark 4
The behaviour of the integrability algorithm in the Lagrangian and Hamiltonian levels is the same. Indeed, it is easy to prove that $`(_f)(_i)=𝒥_i`$, and that the map $`(_f)|__i:_i𝒥_i`$ is a diffeomorphism, for every $`i`$. Thus, if the integrability algorithm in the Lagrangian level stabilizes at step $`i`$ then the integrability algorithm in the Hamiltonian level also stabilizes at step $`i`$ and, conversely, if the integrability algorithm in the Hamiltonian level stabilizes at step $`i`$ then the integrability algorithm in the Lagrangian level also stabilizes at step $`i`$.
## 5 An example: affine Lagrangian densities
Consider the configuration bundle $`\pi :EM`$, and $`\alpha \mathrm{\Lambda }_1^m\mathrm{T}^{}E`$. Then, $`\alpha `$ induces a function $`L=\widehat{\alpha }\mathrm{C}^{\mathrm{}}(J^1\pi )`$ as follows: given $`xM`$ and a section $`\varphi :ME`$, we define $`L(j_x^1\varphi )`$ by
$$L(j_x^1\varphi )\eta (x)=\left[\varphi ^{}\alpha \right](x).$$
Note that $`L(j_x^1\varphi )`$ is well-defined: if $`\varphi ,\psi `$ are sections such that $`j_x^1\varphi =j_x^1\psi `$, then $`L(j_x^1\varphi )=L(j_x^1\psi )`$.
Taking fibered coordinates $`(x^\alpha ,y^A,v_\alpha ^A)`$ in $`J^1E`$, if $`\alpha =a(x^\alpha ,y^A)d^mx+f_B^\mu (x^\alpha ,y^A)dy^Bd^{m1}x_\mu `$, then
$$L(x^\alpha ,y^A,v_\alpha ^A)=a(x^\alpha ,y^A)+f_B^\mu (x^\alpha ,y^A)v_\mu ^B.$$
Thus, the Lagrangian density $`=L\omega `$ is affine.
A direct computation in local coordinates shows that $`\mathrm{\Theta }_L=(\pi ^1)\alpha `$ and, hence, $`\mathrm{\Omega }_L=(\pi )^1(d\alpha )`$. We also obtain $`\stackrel{~}{}=\alpha \pi ^1`$, and $`=\mu \alpha \pi ^1`$. Therefore, $`\stackrel{~}{𝒫}=\stackrel{~}{}(J^1\pi )=\alpha (E)`$ is an embedded submanifold of $`\pi `$, which is diffeomorphic to $`E`$ by means of the mapping $`\alpha :E\stackrel{~}{𝒫}\text{Im}\alpha `$. Since $`\pi ^1`$ is a surjective submersion with connected fibers, then so is $`\stackrel{~}{}_0:J^1\pi 𝒫`$ (recall that $`\stackrel{~}{}_0`$ is the restriction of $`\stackrel{~}{}`$ onto its image $`𝒫`$). Moreover, since $`\stackrel{~}{}^1(\stackrel{~}{})(\overline{y})=(\pi ^1)^1(\pi ^1(\overline{y}))`$, for all $`\overline{y}J^1\pi `$, and $`\stackrel{~}{}^1(\stackrel{~}{})(\overline{y})^1()(\overline{y})(\pi ^1)^1(\pi ^1(\overline{y}))`$, we obtain $`^1()(\overline{y})=\stackrel{~}{}^1(\stackrel{~}{})(\overline{y})=(\pi ^1)^1(\pi ^1(\overline{y}))`$, and hence the fibers of $``$ are connected submanifolds of $`J^1\pi `$. In conclusion, affine Lagrangian systems are almost regular.
Note that the manifold $`𝒫`$ can be identified with $`E`$, and the mapping $`_0:J^1\pi 𝒫`$ can be identified with the mapping $`\pi ^1:J^1\pi E`$. Hence, the $`(m+1)`$-form $`\mathrm{\Omega }_h^0=(\stackrel{~}{ȷ}_0\stackrel{~}{h})^{}\mathrm{\Omega }`$ (resp. the $`m`$-form $`\omega _0^{}`$) on $`𝒫`$ can be identified with the $`(m+1)`$-form $`d\alpha `$ (resp. $`\pi ^{}(\eta )`$) on $`E`$. Taking these identifications into account, the constrained Hamilton equations on $`E`$ are
$$𝑖(𝒳^𝒫)(d\alpha )=0,𝑖(𝒳^𝒫)(\pi ^{}(\eta ))=1.$$
(42)
Let $`_0^{}`$ be a connection in the bundle $`\tau _0^1:𝒫M`$, and $`𝒴_\eta ^_0^{}={\displaystyle \underset{\mu =1}{\overset{m}{}}}\left({\displaystyle \frac{}{x^\mu }}+\mathrm{\Gamma }_\mu ^A{\displaystyle \frac{}{y^A}}\right)`$ the corresponding $`m`$-vector field on $`𝒫`$ associated with $`_0^{}`$ and $`\eta `$. A direct computation shows that
$$(\gamma _h)_\eta ^_0^{}=𝑖(𝒴_\eta ^_0^{})\mathrm{\Omega }_h^0=(1^m)\left[\frac{f_A^\nu }{x^\nu }\frac{a}{y^A}+\mathrm{\Gamma }_\nu ^B\left(\frac{f_A^\nu }{y^B}\frac{f_B^\nu }{y^A}\right)\right]\left(dy^A\mathrm{\Gamma }_\mu ^Adx^\mu \right),$$
$$\mathrm{\Omega }_h^_0^{}=\mathrm{\Gamma }_\nu ^B\left(\frac{f_B^\nu }{y^A}\frac{f_A^\nu }{y^B}\right)dy^Ad^mx\frac{f_B^\mu }{y^A}dy^Ady^Bd^{m1}x_\mu .$$
It is easy to show that $`\mathrm{\Omega }_h^_0^{}`$ is 1-nondegenerate if, and only if, the matrix $`(f_{AB}^\mu )=\left({\displaystyle \frac{f_B^\mu }{y^A}}{\displaystyle \frac{f_A^\mu }{y^B}}\right)`$ is regular, for every $`\mu \{1,\mathrm{},m\}`$. Then, the Hamiltonian constrained system (42) has solution. A $`m`$-vector field $`𝒳^𝒫={\displaystyle \underset{\mu =1}{\overset{m}{}}}\left({\displaystyle \frac{}{x^\mu }}+F_\mu ^A{\displaystyle \frac{}{y^A}}\right)`$ is a solution to (42) on $`𝒫E`$ if, and only if,
$$\left(\frac{f_A^\mu }{y^B}\frac{f_B^\mu }{y^A}\right)F_\mu ^B=\frac{a}{y^A}\frac{f_A^\nu }{x^\nu }.$$
Note that there are $`n`$ equations and $`mn`$ variables, and that the rank of the matrix $`\left({\displaystyle \frac{f_A^\mu }{y^B}}{\displaystyle \frac{f_B^\mu }{y^A}}\right)`$ of type $`n\times nm`$ is maximum, that is, $`n`$. Thus, the set of solutions of the system is an affine space of dimension $`n(m1)`$ (the solution is not unique if $`m>1`$).
With respect to the integrability of the solutions, a direct computation shows that a $`m`$-vector field $`𝒳^𝒫`$ solution to (42) is integrable if
$$\frac{F_\nu ^A}{x^\mu }\frac{F_\mu ^A}{x^\nu }+F_\mu ^B\frac{F_\nu ^A}{y^B}F_\nu ^B\frac{F_\mu ^A}{y^B}=0,\text{ for all }A\text{ and }\mu ,\nu .$$
Otherwise, the integrability algorithm should be applied.
Taking into account the identification $`𝒫E`$, as $`\mathrm{\Omega }_{}=_0^{}(d\alpha )`$, if $`𝒳^𝒫`$ is a solution to the constrained Hamiltonian equations on $`𝒫`$, then every locally decomposable $`m`$-vector field $`𝒳^{J^1\pi }`$ which projects via $`_0`$ onto $`𝒳^𝒫`$ is a solution to the equations
$$i(𝒳^{J^1\pi })(\mathrm{\Omega }_{})=0,i(𝒳^{J^1\pi })\omega =1.$$
Let $`\mathrm{\Psi }`$ be the first-order jet field with respect to the fibration $`\pi :EM`$ associated to the Ehresmann connection defined by the $`m`$-vector field $`𝒳^𝒫`$. Then, the submanifold $`S`$ of $`J^1\pi `$ where a semi-holonomic $`m`$ vector field satisfying the Lagrangian equations exists is $`\mathrm{\Psi }(E)`$. In fact, if $`𝒳^𝒮=(\mathrm{\Lambda }^m\mathrm{T}\mathrm{\Psi })𝒳^𝒫`$ then $`𝒳^S`$ is an Euler-Lagrange $`m`$-vector field on $`S`$ for $``$, that is, $`𝒳^S`$ is a locally decomposable $`m`$-vector field on $`S`$ and
$$(𝑖(𝒳^S)\mathrm{\Omega }_{})|_S=0,(𝑖(𝒳^S)\omega )|_S=1,(𝑖(𝒳^S)𝒱)|_S=0.$$
If the matrix $`(f_{AB}^\mu )`$ is singular but there are no higher-order constraints, the previous results remain true. Otherwise, we will have to apply the premultisymplectic constraint algorithm. Suppose that we have obtained the final constraint submanifolds $`N_f`$ and $`P_f`$, with the submersion $`(\pi ^1)|_{N_f}:N_fP_f`$. Let $`𝒳^{P_f}`$ be a $`m`$-vector field solution of the constrained Hamiltonian equations. We have that $`\pi (P_f)`$ is an open subset of $`M`$ and that $`\pi _f=\pi _{|P_f}:P_f\pi (P_f)M`$ is a fibration. Moreover, $`J^1\pi _f`$ is a submanifold of $`J^1\pi `$ (see Appendix). Now, let $`\mathrm{\Psi }`$ be the first-order jet field with respect to the fibration $`\pi _f:P_f\pi (P_f)`$ associated to the $`m`$-vector field $`𝒳^{P_f}`$. Then, the submanifold $`S_f`$ of $`J^1\pi `$ where an Euler-Lagrange $`m`$-vector field for $``$ along $`S_f`$ exists is $`\mathrm{\Psi }(P_f)`$, and $`𝒳^{S_f}=(\mathrm{\Lambda }^mT\mathrm{\Psi })𝒳^{P_f}`$ is such an Euler-Lagrange $`m`$-vector field (see Theorem 6).
Example: Let $`\pi :^4^2`$ be the configuration bundle, and $`L=x^2(y^1v_2^1+y^2v_2^2)+y^1y^2`$. In this case, $`\alpha =y^1y^2dx^1dx^2x^2y^1dy^1dx^1x^2y^2dy^2dx^1`$. If $``$ is the trivial connection, $`𝒴_\eta ^{}={\displaystyle \frac{}{x^1}}{\displaystyle \frac{}{x^2}}`$, then $`\gamma _\eta ^{}=(y^1y^2)(dy^1dy^2)`$ and $`\mathrm{\Omega }_{}^{}=0`$. A simple computation shows that, in this case, $`\mathrm{}_\mathrm{\Omega }^{}(𝐡)=((𝐡)_{}^H,0)`$. Therefore, the vector fields $`Z_i`$ in Theorem 5 are all the vertical vector fields in $`\mathrm{V}(\overline{\pi }^1)`$. Hence, the submanifold $`N_1`$ is characterized by the constraint $`y^1y^2=0`$. In fact, every semi-holonomic $`2`$-vector field in $`N_1`$ is an Euler-Lagrange $`2`$-vector field for this problem.
## Appendix: $`m`$-vector fields and Ehresmann connections in fibre bundles
(See for the proofs and other details about the results in this section).
Let $`F`$ be a $`N`$-dimensional differentiable manifold. Sections of $`\mathrm{\Lambda }^m(\mathrm{T}F)`$ are called multivector fields in $`F`$, or more precisely, $`m`$-vector fields in $`F`$ (they are contravariant skew-symmetric tensors of order $`m`$ in $`F`$). The space of $`m`$-vector fields is denoted by $`𝒱^m(F)`$. $`𝒳𝒱^m(F)`$ is locally decomposable if, for every $`pF`$, there exists an open neighbourhood $`U_pF`$ and $`Y_1,\mathrm{},Y_m\text{X}(U_p)`$ such that $`𝒳\underset{U_p}{=}Y_1\mathrm{}Y_m`$. We denote by $`\text{X}^m(F)`$ the set of locally decomposable $`m`$-vector fields in $`F`$. Contraction of $`m`$-vector fields and tensor fields in $`F`$ is the usual one.
We can define an equivalence relation: if $`𝒳,𝒳^{}\text{X}^m(F)`$ are non-vanishing $`m`$-vector fields, and $`UF`$ is a connected open set, then $`𝒳\stackrel{U}{}𝒳^{}`$ if there exists a non-vanishing function $`f\mathrm{C}^{\mathrm{}}(U)`$ such that $`𝒳^{}\underset{U}{=}f𝒳`$. Equivalence classes are denoted by $`\{𝒳\}_U`$. There is a one-to-one correspondence between the set of $`m`$-dimensional orientable distributions $`D`$ in $`F`$ and the set of the equivalence classes $`\{𝒳\}_F`$ of non-vanishing, locally decomposable $`m`$-vector fields in $`F`$. If $`𝒳\text{X}^m(F)`$ is non-vanishing and locally decomposable, the distribution associated with the class $`\{𝒳\}_U`$ is denoted $`𝒟_U(𝒳)`$ (If $`U=F`$ we write $`𝒟(𝒳)`$). A non-vanishing, locally decomposable $`m`$-vector field $`𝒳\text{X}^m(F)`$ is said to be integrable if its associated distribution $`𝒟_U(𝒳)`$ is integrable. Of course, if $`𝒳\text{X}^m(F)`$ is integrable, then so is every $`m`$-vector field in its equivalence class $`\{𝒳\}`$, and all of them have the same integral manifolds. Moreover, from Frobenius’ theorem, a non-vanishing and locally decomposable $`m`$-vector field is integrable if, and only if, $`𝒟(𝒳)`$ is involutive.
Now, let $`\kappa :FM`$ be a fibre bundle ($`dimM=m`$). We are concerned with the case where the integral manifolds of integrable $`m`$-vector fields in $`F`$ are sections of $`\kappa `$. Thus, $`𝒳\text{X}^m(F)`$ is said to be $`\kappa `$-transverse if, at every point $`yF`$, $`(𝑖(𝒳)(\kappa ^{}\eta ))_y0`$, for every $`\eta \mathrm{\Omega }^m(M)`$ such that $`\eta (\kappa (y))0`$. Then, if $`𝒳\text{X}^m(F)`$ is integrable, it is $`\kappa `$-transverse if, and only if, its integral manifolds are local sections of $`\kappa :FM`$. In this case, if $`\varphi :UMF`$ is a local section with $`\varphi (x)=y`$ and $`\varphi (U)`$ is the integral manifold of $`𝒳`$ through $`y`$, then $`\mathrm{T}_y(\mathrm{Im}\varphi )`$ is $`𝒟_y(𝒳)`$. Integral sections $`\varphi `$ of $`𝒳`$ can be characterized by the condition
$$\mathrm{\Lambda }^m\mathrm{T}\varphi =f𝒳\varphi \varrho _M$$
where $`\mathrm{\Lambda }^m\mathrm{T}\varphi :\mathrm{\Lambda }^m\mathrm{T}M\mathrm{\Lambda }^m\mathrm{T}F`$ is the natural lifting of $`\varphi `$, $`\varrho _M:\mathrm{\Lambda }^m\mathrm{T}MM`$ is the natural projection, and$`f\mathrm{C}^{\mathrm{}}(F)`$ is a non-vanishing function (observe that this characterizes the entire class $`\{𝒳\}`$ of integrable $`m`$-vector fields).
Let $``$ be an Ehresmann connection in the fibration $`\kappa :FM`$. As is known, it defines a horizontal subbundle $`\mathrm{H}()\mathrm{T}F`$, such that $`\mathrm{T}F=\mathrm{H}()\mathrm{V}(\kappa )`$, where $`\mathrm{V}(\kappa )`$ is the $`\kappa `$-vertical subbundle. If $`yF`$, then $`\mathrm{H}_y()=Im(y)`$. Thus, we have the horizontal distribution associated with the connection $``$. The connection $``$ is said to be flat (respectively, orientable) if the horizontal distribution is completely integrable (respectively, orientable).
Classes of locally decomposable and $`\kappa `$-transverse $`m`$-vector fields $`\{𝒳\}\text{X}^m(F)`$ are in one-to-one correspondence with orientable Ehresmann connections $``$ in $`\kappa :FM`$. This correspondence is given by the fact that the horizontal subbundle associated with $``$ is $`𝒟(𝒳)`$. Thus, classes of integrable locally decomposable and $`\kappa `$-transverse $`m`$-vector fields correspond to flat orientable Ehresmann connections.
A connection $``$ in the fibration $`\kappa :FM`$ induces a splitting $`\mathrm{T}^{}F=\mathrm{H}^{}()\mathrm{V}^{}(\kappa )`$, where
$$\mathrm{H}_y^{}()=\mathrm{V}_y(\kappa )^0,\mathrm{V}_y^{}(\kappa )=\mathrm{H}_y()^0.$$
Here, $`\mathrm{V}_y(\kappa )^0\mathrm{T}_y^{}F`$ (respectively, $`\mathrm{H}_y()^0\mathrm{T}_y^{}F`$) denotes the annihilator of the subspace $`\mathrm{V}_y(\kappa )\mathrm{T}_yF`$ (respectively, $`\mathrm{H}_y()\mathrm{T}_yF`$). The splittings $`\mathrm{T}F=\mathrm{H}()\mathrm{V}(\kappa )`$ and $`\mathrm{T}^{}F=\mathrm{H}^{}()\mathrm{V}^{}(\kappa )`$ may be extended to the tensor bundles
$$\mathrm{\Lambda }^l\mathrm{T}F=\underset{r,s=0,\mathrm{},l;r+s=l}{}(\mathrm{\Lambda }^r\mathrm{H}()\mathrm{\Lambda }^s\mathrm{V}(\kappa ))$$
(43)
$$\mathrm{\Lambda }^k\mathrm{T}^{}F=\underset{p,q=0,\mathrm{},k;p+q=k}{}(\mathrm{\Lambda }^p\mathrm{H}^{}()\mathrm{\Lambda }^q\mathrm{V}(\kappa )^{})$$
(44)
Thus, for every $`X\text{X}(F)`$, we obtain that $`𝑖(X)X_{}^H`$ is an horizontal vector field, that is, a section of $`\mathrm{H}()F`$. $`X_{}^H`$ is the horizontal component of $`X`$, and we write $`X=X_{}^H+X_{}^V`$, where $`X_{}^V=XX_{}^H`$ is a $`\kappa `$-vertical vector field. Moreover, if $`\alpha \mathrm{\Omega }^1(F)`$, then we have that $`𝑖(\alpha )\alpha _{}^H\mathrm{\Omega }^1(F)`$ is an horizontal $`1`$-form, that is, a section of $`\mathrm{H}()^{}F`$. $`\alpha _{}^H`$ is the horizontal component of $`\alpha `$, and we write $`\alpha =\alpha _{}^H+\alpha _{}^V`$, where $`\alpha _{}^V=\alpha \alpha _{}^H`$ is a $`\kappa `$-vertical $`1`$-form with respect to the connection $``$, that is, it vanishes under the action of every horizontal vector field associated with the connection $``$. Furthermore, if $`X\text{X}(F)`$ is a $`\kappa `$-vertical vector field, then $`𝑖(X)\alpha _{}^H=0`$. In addition, if $`𝒳\text{X}^k(F)`$ and $`\beta \mathrm{\Omega }^l(F)`$, the splittings (43) and (44) allow us to make the following decomposition
$$𝒳=\underset{r,s=0;r+s=k}{}𝒳_{}^{(r,s)},\beta =\underset{p,q=0;p+q=l}{}\beta _{}^{(p,q)},$$
where the superscripts $`(i,j)`$ denote the horizontal and vertical parts respectively, of the $`k`$-vector field $`𝒳`$ and the $`l`$-form $`\beta `$.
Finally, if $``$ is an Ehresmann connection in the fibration $`\kappa :FM`$ and $`yF`$ then the map
$$\mathrm{\Lambda }^k\mathrm{T}_y\kappa _{\mathrm{H}()}:\mathrm{\Lambda }^k\mathrm{H}_y()\mathrm{\Lambda }^k\mathrm{T}_{\kappa (y)}M,1kdimM=m,$$
is a linear isomorphism and the inverse morphism $`(\mathrm{\Lambda }^m\mathrm{T}_y\kappa _{\mathrm{H}()})^1:\mathrm{\Lambda }^m\mathrm{T}_{\kappa (y)}M\mathrm{\Lambda }^m\mathrm{H}_y()`$ is just the horizontal lift at $`y`$ induced by $``$. Denoting by $`\mathrm{\Lambda }^m(\kappa _{\mathrm{H}()})_{}^1`$ the natural extension of this map to $`m`$-vector fields on $`M`$, one may consider $`\mathrm{\Lambda }^m(\kappa _{\mathrm{H}()})_{}^1(𝒳)`$, the horizontal lift of $`𝒳\text{X}^m(M)`$, as the $`m`$-vector field on $`F`$ given by
$$[\mathrm{\Lambda }^m(\kappa _{\mathrm{H}()})_{}^1(𝒳)](y)=(\mathrm{\Lambda }^m\mathrm{T}_y\kappa _{\mathrm{H}()})^1(𝒳(\kappa (y))),\text{for every}yF.$$
In particular, if $`𝒳_\eta `$ is the $`m`$-vector field on $`M`$ characterized by the condition
$$\alpha _1\mathrm{}\alpha _m=𝒳_\eta (\alpha _1,\mathrm{},\alpha _m)\eta ,\text{for every }\alpha _1,\mathrm{},\alpha _m\mathrm{\Omega }^1(M)$$
one may define the $`m`$-vector field $`𝒴_\eta ^{}\text{X}^m(F)`$ by $`𝒴_\eta ^{}=\mathrm{\Lambda }^m(\kappa _{\mathrm{H}()})_{}^1(𝒳_\eta )`$. Note that $`𝒴_\eta ^{}`$ is a locally decomposable and $`\kappa `$-transverse $`m`$-vector field on $`F`$, verifying that $`𝑖(𝒴_\eta ^{})\omega =1`$, and that the distribution $`𝒟(𝒴_\eta ^{})`$ is just the horizontal distribution associated with the connection $``$.
If $`C`$ is a submanifold of $`F`$, and $`𝒳_C`$ is a locally decomposable $`m`$-vector field on $`C`$ such that
$$𝑖(𝒳_C(y))\omega (y)=1,\text{for every}yC$$
then $`\kappa |_C\kappa _C:CM`$ is a submersion. In fact, if $`yC`$ and $`𝒳_C(y)=X_C^1(y)\mathrm{}X_C^m(y)`$, with $`X_C^i(y)\mathrm{T}_yC`$, then
$$\eta (\kappa (y))(\mathrm{T}_y\kappa _C(X_C^1(y)),\mathrm{},\mathrm{T}_y\kappa _C(X_C^m(y)))=1$$
This implies that $`\{\mathrm{T}_y\kappa _C(X_C^1(y)),\mathrm{},\mathrm{T}_y\kappa _C(X_C^m(y))\}`$ is a basis of $`\mathrm{T}_{\kappa (y)}M`$, and thus, $`\mathrm{T}_y\kappa _C:\mathrm{T}_yC\mathrm{T}_{\kappa (y)}M`$ is an epimorphism. Therefore, $`\kappa (C)`$ is an open subset of $`M`$ and $`\kappa _C:C\kappa (C)`$ is a fibre bundle. Consequently, $`𝒳_C`$ defines an oriented Ehresmann connection in the fibration $`\kappa _C:C\kappa (C)`$ which, in the terminology of , is said to be an (oriented) Ehresmann connection in the fibration $`\kappa :FM`$ along the submanifold $`C`$. Note that the canonical inclusion $`\iota :J^1\kappa _CJ^1\kappa `$ is an embedding and, thus, $`J^1\kappa _C`$ is a submanifold of $`J^1\kappa `$.
###### Remark 5
It is well-known that there exists a one-to-one correspondence between Ehresmann connections in the fibration $`\kappa :FM`$ and first-order jet fields with respect to $`\kappa `$, that is, sections of the fibration $`\kappa ^1:J^1FF`$. In fact, let $``$ be a connection in the fibration $`\kappa :FM`$, (that is, an element of $`\mathrm{\Gamma }(E,\kappa ^{}\mathrm{T}^{}M)\mathrm{\Gamma }(F,\mathrm{T}F)`$), such that $`^{}\alpha =\alpha `$, for every $`\kappa `$-semibasic form $`\alpha \mathrm{\Omega }^1(F)`$), and $`\mathrm{H}()`$ the associated horizontal subbundle. If $`(\mathrm{T}_y\kappa _{\mathrm{H}()})^1`$ denotes the horizontal lift at $`y`$; for every $`yF`$, let $`\varphi :MF`$ be a section of $`\kappa `$ passing through $`y`$, such that
$$\mathrm{T}_{\kappa (y)}\varphi =\mathrm{T}_y\kappa _{\mathrm{H}()})^1:\mathrm{T}_{\kappa (y)}M\mathrm{H}_y()\mathrm{T}_yF$$
then we define the map
$$\begin{array}{ccccc}\psi ^{}& :& F& & J^1F\\ & & y& & (j^1\varphi )(\kappa (y))\end{array}$$
which is a section of the fibration $`\kappa ^1:J^1FF`$. Conversely, given a section $`\psi ^{}:FJ^1F`$, for every $`\overline{y}J^1F`$ with $`\kappa ^1(\overline{y})=y`$, and a representative $`\varphi :MF`$ of $`\overline{y}`$, we define the horizontal subspace $`\mathrm{H}_y():=\mathrm{ImT}_y\varphi `$, and $`\mathrm{H}():=_y\mathrm{H}_y()`$. Thus we have identified the fibre $`J_y^1F=(\kappa ^1)^1(y)`$ with the set
$$\{𝐡_y\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yF\mathrm{T}_y\kappa 𝐡_y=Id\}$$
In particular, if we have a connection or, what is equivalent, a class of $`\kappa `$-transverse, locally decomposable $`m`$-vector fields in the fibration $`\kappa :FM`$, along a submanifold $`C`$ of $`F`$, and a representative $`𝒳_C`$ of this class, then $`\kappa (C)`$ is an open subset of $`M`$, $`\kappa _C=\kappa |_C:C\kappa (C)`$ is a fibration, and $`𝒳_C`$ may be seen as a section $`\psi _C^{}`$ of the fibration $`\kappa _C^1:J^1\kappa _CC`$. Thus, $`\psi _C^{}(y)`$ is identified with a linear map from $`\mathrm{T}_{\kappa (y)}M`$ onto $`\mathrm{T}_yC`$, that is, an element $`𝐡_y\mathrm{T}_{\kappa (y)}^{}M\mathrm{T}_yC`$, and
$$(\mathrm{T}_y\kappa _C\psi _C)(y)=(\mathrm{T}_y\kappa _C|_{\mathrm{T}_yC}\psi _C^{})(y)=Id,\text{for every}yC.$$
### Acknowledgments
We acknowledge the financial support of Ministerio de Educación y Ciencia, projects BFM2002-03493, BFM2003-01319 and MTM2004-7832. We thank Mr. Jeff Palmer for his assistance in preparing the English version of the manuscript.
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# Boundary blow-up in nonlinear elliptic equations of Bieberbach–Rademacher type
## 1. Introduction
Let $`\mathrm{\Omega }^N`$ $`(N3)`$ be a smooth bounded domain. We consider semilinear elliptic problems under the following form
(1.1)
$$\mathrm{\Delta }u=g(x,u)\text{in}\mathrm{\Omega },$$
subject to the singular boundary condition
(1.2)
$$u(x)\mathrm{}\text{as}d(x):=\mathrm{dist}(x,\mathrm{\Omega })0(\text{in short, }u=\mathrm{}\text{on}\mathrm{\Omega }).$$
The nonnegative solutions of (1.1)+(1.2) are called large (or blow-up) solutions.
The study of large solutions has been initiated in 1916 by Bieberbach for the particular case $`g(x,u)=\mathrm{exp}(u)`$ and $`N=2`$. He showed that there exists a unique solution of (1.1) such that $`u(x)\mathrm{log}(d(x)^2)`$ is bounded as $`x\mathrm{\Omega }`$. Problems of this type arise in Riemannian geometry; if a Riemannian metric of the form $`|ds|^2=\mathrm{exp}(2u(x))|dx|^2`$ has constant Gaussian curvature $`c^2`$ then $`\mathrm{\Delta }u=c^2\mathrm{exp}(2u)`$. Motivated by a problem in mathematical physics, Rademacher continued the study of Bieberbach on smooth bounded domains in $`^3`$. Lazer–McKenna extended the results of Bieberbach and Rademacher for bounded domains in $`^N`$ satisfying a uniformal external sphere condition and for nonlinearities $`g(x,u)=b(x)\mathrm{exp}(u)`$, where $`b`$ is continuous and strictly positive on $`\overline{\mathrm{\Omega }}`$.
The interest in large solutions extended to $`N`$-dimensional domains and for other classes of nonlinearities (see e.g., , , , , , , , , ).
Let $`g(x,u)=f(u)`$ where $`f`$ satisfies
($`A`$)
$$fC^1[0,\mathrm{}),f^{}(s)0\text{for }s0,f(0)=0\text{and}f(s)>0\text{for}s>0.$$
In this case, Keller and Osserman proved that large solutions of (1.1) exist if and only if
($`A_0`$)
$$_1^{\mathrm{}}\frac{dt}{\sqrt{F(t)}}<\mathrm{},\text{where}F(t)=_0^tf(s)𝑑s.$$
In a celebrated paper, Loewner and Nirenberg linked the uniqueness of the blow-up solution to the growth rate at the boundary. Motivated by certain geometric problems, they established the uniqueness for the case $`f(u)=u^{\frac{N+2}{N2}}`$ $`(N>2)`$. Bandle and Marcus give results on asymptotic behaviour and uniqueness of the large solution for more general nonlinearities including $`f(u)=u^p`$ for any $`p>1`$. Theorem 2.3 in proves that when $`(A)`$ holds and
($`B`$)
$$\mu >0\text{and}s_01\text{such that}f(\tau s)\tau ^{\mu +1}f(s)\tau (0,1)ss_0/\tau $$
then for any large solution of $`\mathrm{\Delta }u=f(u)`$ we have
(1.3)
$$\underset{d(x)0}{lim}\frac{u(x)}{Z(d(x))}=1$$
where $`Z`$ is a chosen solution of
(1.4)
$$\{\begin{array}{cc}& Z^{\prime \prime }(r)=f(Z(r)),r(0,\delta )\text{for some}\delta >0\hfill \\ & Z(r)\mathrm{}\text{as}r0^+.\hfill \end{array}$$
If, in addition, $`f(\tau s)\tau f(s)`$, for all $`\tau (0,1)`$ and $`s>0`$, then the uniqueness of large solutions takes place. Lazer and McKenna consider the case when the $`C^1`$-function $`f`$ is either defined and positive on $``$ or is defined on $`[a_0,\mathrm{})`$ with $`f(a_0)=0`$ and $`f(s)>0`$ for $`s>a_0`$. They prove the uniqueness of large solutions to $`\mathrm{\Delta }u=f(u)`$ in $`\mathrm{\Omega }^N`$, $`N>1`$, under the assumptions (see \[24, Theorem 3.1\]):
(1.5) $`\begin{array}{cc}& \mathrm{\Omega }\text{satisfies both a uniform internal sphere condition and a uniform}\hfill \\ & \text{external sphere condition with the same constant}R_1>0\hfill \end{array}`$
(1.6) $`f^{}(s)0\text{for }s\text{ in the domain of }f;`$
(1.7) $`\text{there exists }a_1\text{such that }f^{}(s)\text{is nondecreasing for }sa_1;`$
(1.8) $`\underset{s\mathrm{}}{lim}f^{}(s)/\sqrt{F(s)}=\mathrm{}.`$
Moreover, the asymptotics of the large solution is found in terms of a difference
$$\underset{d(x)0}{lim}[u(x)Z(d(x))]=0,\text{for any}Z\text{satisfying}(\text{1.4}).$$
We are interested in large solutions of (1.1) when $`g(x,u)=b(x)f(u)au`$, i.e.,
($`P`$)
$$\mathrm{\Delta }u=aub(x)f(u)\text{in}\mathrm{\Omega },$$
where $`fC^1[0,\mathrm{})`$, $`a`$ and $`bC^{0,\mu }(\overline{\mathrm{\Omega }})`$ ($`0<\mu <1`$) satisfies $`b0`$, $`b0`$ in $`\mathrm{\Omega }`$.
Many papers (see e.g., , ) have been written about Eq. ($`P`$), on a bounded domain or $`^N`$, when $`f(u)=u^p`$ ($`p>1`$). For this case of nonlinearity and $`b>0`$ on $`\overline{\mathrm{\Omega }}`$, Eq. ($`P`$) subject to $`u=0`$ on $`\mathrm{\Omega }`$ is referred to as the logistic equation. It is known that it has a unique positive solution if and only if $`a>\lambda _1(\mathrm{\Omega })`$, where $`\lambda _1(\mathrm{\Omega })`$ is the first Dirichlet eigenvalue of $`(\mathrm{\Delta })`$ in $`\mathrm{\Omega }`$. We mention that the logistic equation has been proposed as a model for population density of a steady-state single species $`u(x)`$ when $`\mathrm{\Omega }`$ is fully surrounded by inhospitable areas. However, not until recently was the case of a degenerate logistic type equation considered, which allows $`b`$ to vanish on $`\overline{\mathrm{\Omega }}`$ (see , and ). The understanding of the asymptotics for positive solutions of the degenerate logistic equation leads to the study of large solutions (we refer to and ).
Let $`\mathrm{\Omega }_0`$ denote the interior of the zero set of $`b`$ in $`\mathrm{\Omega }`$, i.e.,
$$\mathrm{\Omega }_0=\mathrm{int}\{x\mathrm{\Omega }:b(x)=0\}.$$
We assume throughout that $`\mathrm{\Omega }_0`$ is connected, $`\mathrm{\Omega }_0`$ satisfies the exterior cone condition (possibly, $`\mathrm{\Omega }_0=\mathrm{}`$), $`\overline{\mathrm{\Omega }}_0\mathrm{\Omega }`$ and $`b>0`$ on $`\mathrm{\Omega }\overline{\mathrm{\Omega }}_0`$. Note that $`b0`$ on $`\mathrm{\Omega }`$.
Let $`\lambda _{\mathrm{},1}`$ be the first Dirichlet eigenvalue of $`(\mathrm{\Delta })`$ in $`\mathrm{\Omega }_0`$. Set $`\lambda _{\mathrm{},1}=\mathrm{}`$ if $`\mathrm{\Omega }_0=\mathrm{}`$.
Alama and Tarantello find the maximal interval $`I`$ for the parameter $`a`$ such that ($`P`$), subject to $`u=0`$ on $`\mathrm{\Omega }`$, has a positive solution $`u_a`$, provided that
($`A_1`$)
$$f0\text{and}f(u)/u\text{is increasing on}(0,\mathrm{}).$$
Moreover, for each $`aI`$, the solution $`u_a`$ is unique (see \[1, Theorem A (bis)\]).
Theorem 1.1 in proves that if ($`A_0`$) and ($`A_1`$) are fulfilled, then Eq. ($`P`$) has large solutions if and only if $`a(\mathrm{},\lambda _{\mathrm{},1})`$. The uniqueness and asymptotic behaviour near $`\mathrm{\Omega }`$ prove to be very challenging in the above generality.
In we advance for the first time the idea of using the regular variation theory arising in applied probability to study the uniqueness of large solutions. There we consider the case when $`f^{}`$ varies regularly at infinity (see Definition 2.1).
Note that there are many nonlinearities $`f(u)`$, such as $`\mathrm{exp}(u)1`$, $`\mathrm{sinh}(u)`$, $`\mathrm{exp}(\mathrm{exp}(u))e`$, $`\mathrm{exp}(u)\mathrm{log}(u+1)`$, which do not fall in the category treated by Theorem 1 in . Although some examples might fit into the framework of \[3, Theorem 2.3\] or \[24, Theorem 3.1\], the uniqueness and growth rate at the boundary for large solutions of ($`P`$) have not yet been studied when $`a0`$ and $`b`$ vanishes in $`\mathrm{\Omega }`$ with $`b0`$ on $`\mathrm{\Omega }`$.
Our purpose is to fill in this gap by analysing a wide range of functions $`f`$ and $`b`$. We develop the research line opened up in to treat here the case when $`f`$ does not vary regularly at infinity. Thus our approach for the uniqueness is different from that of Bandle–Marcus and Lazer–McKenna, being based on Karamata’s theory.
## 2. Framework and main result
We first recall some results from the Karamata regular variation theory (see ).
###### Definition 2.1.
A measurable function $`R:[A,\mathrm{})(0,\mathrm{})`$, for some $`A>0`$, is called regularly varying at infinity of index $`\rho `$, in short $`RRV_\rho `$, provided that
$$\underset{u\mathrm{}}{lim}\frac{R(\xi u)}{R(u)}=\xi ^\rho ,\xi >0.$$
When the index $`\rho `$ is zero, we say that the function is slowly varying.
From now on, we do not write at infinity when the regular variation occurs there. Notice that the transformation $`R(u)=u^\rho L(u)`$ reduces regular variation to slow variation. Examples of slowly varying functions are given by:
* Every measurable function on $`[A,\mathrm{})`$ which has a positive limit at $`\mathrm{}`$.
* The logarithm $`\mathrm{log}u`$, its iterates $`\mathrm{log}_mu`$, and powers of $`\mathrm{log}_mu`$.
* $`\mathrm{exp}\{(\mathrm{log}u)^{\alpha _1}(\mathrm{log}_2u)^{\alpha _2}\mathrm{}(\mathrm{log}_mu)^{\alpha _m}\}`$ where $`\alpha _i(0,1)`$ and $`\mathrm{exp}\{\frac{\mathrm{log}u}{\mathrm{log}\mathrm{log}u}\}`$.
###### Proposition 2.2 (Representation Theorem).
The function $`L(u)`$ is slowly varying if and only if it can be written in the form
$$L(u)=M(u)\mathrm{exp}\left\{_B^u\frac{\varphi (t)}{t}𝑑t\right\},uB$$
for some $`B>0`$, where $`\varphi C[B,\mathrm{})`$ satisfies $`lim_u\mathrm{}\varphi (u)=0`$ and $`M(u)`$ is measurable on $`[B,\mathrm{})`$ such that $`lim_u\mathrm{}M(u)=\widehat{M}(0,\mathrm{})`$.
If $`M(u)`$ is replaced by $`\widehat{M}`$ then the new function, say $`\widehat{L}(u)`$, is referred to as a normalised slowly varying function. We see that $`\varphi (u)=\frac{u\widehat{L}^{}(u)}{\widehat{L}(u)}`$, $`uB`$. Conversely, any function $`\widehat{L}C^1[B,\mathrm{})`$ which is positive and satisfies $`lim_u\mathrm{}\frac{u\widehat{L}^{}(u)}{\widehat{L}(u)}=0`$ is a normalised slowly varying function.
Note that any slowly varying function $`L(u)`$ is asymptotic equivalent to some normalised slowly varying function $`\widehat{L}(u)`$ (i.e., $`lim_u\mathrm{}L(u)/\widehat{L}(u)=1`$).
The notion of regular variation can be extended to any real number. For instance, we say that $`R(u)`$ is regularly varying (on the right) at the origin with index $`\rho `$ (and write $`RRV_\rho (0+)`$) if $`R(1/u)RV_\rho `$. Let $`NRV_\rho (0+)`$ (resp., $`NRV_\rho `$) denote the set of all normalised regularly varying functions at $`0`$ (resp., $`\mathrm{}`$) of index $`\rho `$.
By $`f_1(x)f_2(x)`$ as $`d(x)0`$ we mean that $`lim_{d(x)0}\frac{f_1(x)}{f_2(x)}=1`$.
Our main result is
###### Theorem 2.3.
Let ($`A_1`$) hold and $`f𝔏RV_\rho `$ $`(\rho >0)`$ for some $`𝔏C^2[A,\mathrm{})`$ satisfying $`lim_u\mathrm{}𝔏(u)=\mathrm{}`$ and $`𝔏^{}NRV_1`$. Suppose that
($`H`$) $`b(x)𝔎^2(d(x))\text{as}d(x)0,\text{where}𝔎NRV_\theta (0+),\text{for some }\theta 0\text{and}`$
$`𝔎\text{is nondecreasing near the origin if }\theta =0.`$
Then, for any $`a<\lambda _{\mathrm{},1}`$, Eq. ($`P`$) has a unique large solution $`u_a`$. In addition, the blow-up rate of $`u_a`$ at $`\mathrm{\Omega }`$ can be expressed by
(2.1)
$$u_a(x)(𝔏\mathrm{\Phi })(d(x))\text{as}d(x)0,a<\lambda _{\mathrm{},1}.$$
The function $`\mathrm{\Phi }`$ is defined as follows
(2.2)
$$_{\mathrm{\Phi }(t)}^{\mathrm{}}\frac{[𝔏^{}(y)]^{\frac{1}{2}}}{y^{\frac{\rho +1}{2}}[𝑳_f(y)]^{\frac{1}{2}}}𝑑y=_0^t𝔎(s)𝑑s,t(0,\beta )\text{with}\beta >0\text{small},$$
where $`𝐋_f`$ is a normalised slowly varying function such that $`lim_u\mathrm{}\frac{f(𝔏(u))}{u^\rho 𝐋_f(u)}=1`$.
Note that Theorem 2.3 brings a new insight into the asymptotics of the large solution of ($`P`$) even in the case $`a=0`$ and $`b=1`$. For instance, the function which is used in (2.1) to estimate the blow-up rate of the solution near $`\mathrm{\Omega }`$ is not chosen as a solution of (1.4). This fact will allow us, through Corollary 2.7, to illustrate the explosion pattern followed by the large solution when the nonlinearity $`f`$ is of the form (2.8) at infinity and satisfies $`(A_1)`$. In particular, if $`lim_u\mathrm{}\frac{f(u)}{[\mathrm{exp}_m(u^{\frac{1}{\alpha }})]^\rho }=1`$, ($`\alpha ,\rho >0`$ and $`m1`$ an integer), then the unique large solution of $`\mathrm{\Delta }u=f(u)`$ satisfies $`\frac{u(x)}{\mathrm{\Psi }(d(x))}1`$ as $`d(x)0`$, where
$$\mathrm{\Psi }(d(x))=\{\begin{array}{cc}& [\mathrm{log}(d(x)^{\frac{2}{\rho }})]^\alpha ,\text{if }m=1,\hfill \\ & [\mathrm{log}_m(d(x)^1)]^\alpha ,\text{if }m2.\hfill \end{array}$$
We set $`\mathrm{log}_m()=\underset{m\mathrm{times}}{(\underset{}{\mathrm{log}\mathrm{}\mathrm{log}})}()`$ and $`\mathrm{exp}_m()=\underset{m\mathrm{times}}{(\underset{}{\mathrm{exp}\mathrm{}\mathrm{exp}})}()`$, $`m1`$. If $`f(u)=\mathrm{exp}_2(u)+\mathrm{cos}(\mathrm{exp}_2(u))`$ for $`u`$ large and $`(A_1)`$ holds, the uniqueness of large solutions for $`\mathrm{\Delta }u=f(u)`$ cannot be inferred from the Lazer–McKenna result, since condition (1.7) fails. Nevertheless, the uniqueness is valid as we can derive from either \[3, Theorem 2.3\] or Theorem 2.3. But it is not transparent through (1.3) that the large solution fulfills $`lim_{d(x)0}\frac{u(x)}{\mathrm{log}_2\left(\frac{1}{d(x)}\right)}=1`$, as Corollary 2.7 proves.
###### Remark 2.4.
We point out that $`𝔏^{}NRV_1`$ with $`lim_u\mathrm{}𝔏(u)=\mathrm{}`$ if and only if
(2.3)
$$𝔏(u)=C\mathrm{exp}\left\{_B^u\frac{\mathrm{}(t)}{t}𝑑t\right\},uB>0$$
where $`C>0`$ is a constant and $`\mathrm{}`$ is a normalised slowly varying function satisfying $`lim_u\mathrm{}\mathrm{}(u)=0`$ and $`lim_u\mathrm{}_B^u\frac{\mathrm{}(t)}{t}𝑑t=\mathrm{}`$. Nontrivial examples of functions $`𝔏`$ are: $`\mathrm{exp}\{(\mathrm{log}u)^\gamma \}`$, where $`\gamma (0,1)`$, $`\mathrm{exp}\left\{\frac{\mathrm{log}u}{\mathrm{log}\mathrm{log}u}\right\}`$, and $`(\mathrm{log}_mu)^\alpha `$ with $`\alpha >0`$.
The hypothesis $`f𝔏RV_\rho `$ ($`\rho >0`$) is equivalent to the existence of $`gRV_\rho `$ so that $`f(u)=g(𝔏^{}(u))`$, for $`u`$ large (where $`𝔏^{}`$ denotes the inverse of $`𝔏`$). By Proposition 0.8 (v), $`𝔏^{}`$ is rapidly varying with index $`\mathrm{}`$ ($`𝔏^{}RV_{\mathrm{}}`$), i.e.,
$$\underset{u\mathrm{}}{lim}\frac{𝔏^{}(\lambda u)}{𝔏^{}(u)}=\{\begin{array}{cc}0\hfill & \text{if }\lambda (0,1),\hfill \\ 1\hfill & \text{if }\lambda =1,\hfill \\ \mathrm{}\hfill & \text{if }\lambda >1.\hfill \end{array}$$
Therefore, for $`g(u)=u^\rho `$, $`f(u)=[𝔏^{}(u)]^\rho `$ is rapidly varying with index $`\mathrm{}`$.
If $`gNRV_\rho `$, then $`𝑳_f`$ (which appears in (2.2)) can be taken as $`\frac{f(𝔏(u))}{u^\rho }`$. Moreover, $`\frac{f(u)}{u}`$ is increasing in a neighbourhood of infinity. For this, it is enough to see that $`lim_u\mathrm{}\frac{uf^{}(u)}{f(u)}>1`$. Indeed, using (2.3), we derive that
$$\underset{y\mathrm{}}{lim}\frac{f^{}(𝔏(y))𝔏(y)}{f(𝔏(y))}=\underset{y\mathrm{}}{lim}\frac{yg^{}(y)}{g(y)}\frac{𝔏(y)}{y𝔏^{}(y)}=\rho \underset{y\mathrm{}}{lim}\frac{𝔏(y)}{y𝔏^{}(y)}=\mathrm{}.$$
Proposition 2.2 will provide countless functions $`gNRV_\rho `$ and $`𝔏`$ as in (2.3). Hence, by taking $`f(u)=g(𝔏^{}(u))`$ ($`uB>0`$), the assumptions of Theorem 2.3 are fulfilled. It remains only to extend the definition of $`f`$ to the remaining part of $`(0,\mathrm{})`$ such that the smoothness of $`f`$ and $`(A_1)`$ hold.
Regarding the assumption $`(H)`$, $`𝔎NRV_\theta (0+)`$ if and only if there exists a normalised slowly varying function $`L_𝔎`$ such that
(2.4)
$$𝔎(t)=t^\theta L_𝔎(1/t),t(0,\nu )\text{with}\nu >0.$$
Therefore (2.4) is equivalent to saying that for some constants $`c,d>0`$ and $`\phi C(0,\nu )`$ with $`lim_{t0^+}\phi (t)=0`$ we have
$$𝔎(t)=ct^\theta \mathrm{exp}\left(_t^d\frac{\phi (y)}{y}𝑑y\right),\text{for }t(0,d).$$
Some examples of $`𝔎`$ as in $`(H)`$ are: $`t^\theta `$, $`(\mathrm{sin}t)^\theta `$, $`t^\theta /\mathrm{exp}\left[\frac{\mathrm{log}(1/t)}{\mathrm{log}\mathrm{log}(1/t)}\right]`$, $`t^\theta /\mathrm{exp}[(\mathrm{log}t)^\gamma ]`$ with $`\gamma (0,1)`$, $`t^\theta [\mathrm{log}(t+1)]^\alpha `$ or $`t^\theta [\mathrm{log}_m(1/t)]^\alpha `$ with $`\alpha >0`$ and $`m1`$ an integer.
###### Remark 2.5.
If in Theorem 2.3 we replace $`f𝔏RV_\rho `$ by the hypothesis $`f^{}RV_\rho `$ ($`\rho >0`$), then ($`P`$) still has a unique large solution $`u_a`$, $`a<\lambda _{\mathrm{},1}`$. However, the blow-up rate of $`u_a`$ near $`\mathrm{\Omega }`$ is as follows (see \[8, Theorem 1\])
(2.5)
$$\underset{d(x)0}{lim}\frac{u_a(x)}{\left(\frac{2\theta +\rho +2}{(2+\rho )(\theta +1)}\right)^{1/\rho }h(d(x))}=1,a<\lambda _{\mathrm{},1}$$
where $`h`$ is defined by
(2.6)
$$_{h(t)}^{\mathrm{}}\frac{ds}{\sqrt{2F(s)}}=_0^t𝔎(s)𝑑s,t(0,\nu ).$$
###### Remark 2.6.
The variation of $`f`$ at $`\mathrm{}`$ is not regular in Theorem 2.3 (i.e., $`fRV_\gamma `$, for any $`\gamma `$) in contrast to Remark 2.5 where $`fNRV_{\rho +1}`$. This fact will bring a significant change in the explosion speed of the large solution of ($`P`$).
By Lemma 3.4 we know that $`\mathrm{\Phi }NRV_{\frac{2(\theta +1)}{\rho }}(0+)`$. Since $`𝔏`$ varies slowly at infinity, we can invoke \[29, Proposition 0.8 (iv)\] to conclude that $`𝔏\mathrm{\Phi }RV_0(0+)`$.
We show that, in the setting of Remark 2.5, $`hRV_{\frac{2(\theta +1)}{\rho }}(0+)`$. It is easy to check that $`T(u)=(_0^{1/u}𝔎(s)𝑑s)^1RV_{\theta +1}`$. Set $`Y(u)=(_u^{\mathrm{}}\frac{ds}{\sqrt{2F(s)}})^1`$, for $`u>0`$. Clearly, $`Y`$ is increasing on $`(0,\mathrm{})`$, $`Y(\mathrm{})=\mathrm{}`$ and $`YRV_{\rho /2}`$. By (2.6), we find $`h(1/u)=Y^{}(T(u))`$ for $`u`$ sufficiently large, where $`Y^{}(u)`$ is the inverse of $`Y(u)`$. By Proposition 0.8(v) in , $`Y^{}RV_{2/\rho }`$ so that $`h(1/u)RV_{\frac{2(1+\theta )}{\rho }}`$.
As a consequence of Theorem 2.3 and (3.1), we obtain
###### Corollary 2.7.
Let ($`A_1`$) and $`(H)`$ hold. Assume that there exists $`\alpha ,\rho >0`$ and an integer $`m1`$ such that $`f((\mathrm{log}_mu)^\alpha )RV_\rho `$.
Then Eq. ($`P`$) has a unique large solution $`u_a`$, for any $`a<\lambda _{\mathrm{},1}`$. Moreover,
(2.7)
$$\underset{d(x)0}{lim}\frac{u_a(x)}{\left[\mathrm{log}_m\left(\frac{1}{d(x)}\right)\right]^\alpha }=\{\begin{array}{cc}\left(\frac{2(1+\theta )}{\rho }\right)^\alpha \hfill & \text{if }m=1,\hfill \\ 1\hfill & \text{if }m2.\hfill \end{array}$$
###### Remark 2.8.
For $`m=1`$ the influence of $`f`$ (resp., $`𝔎`$) into the blow-up rate (2.7) of the large solution can be seen through $`\alpha `$ and $`\rho `$ (resp., $`\theta `$). Nevertheless, if $`m2`$, then the order of iteration for logarithm changes accordingly in the asymptotic behaviour (2.7) that proves to be independent of the index of regular variation $`\rho `$ (for $`f((\mathrm{log}_mu)^\alpha )`$) and $`\theta `$ (for $`𝔎`$).
The assumption $`f((\mathrm{log}_mu)^\alpha )RV_\rho `$ holds if and only if there exists a slowly varying function $`L`$ such that
(2.8)
$$f(u)=\left[\mathrm{exp}_m(u^{\frac{1}{\alpha }})\right]^\rho L(\mathrm{exp}_m(u^{\frac{1}{\alpha }})),uB>0.$$
Such examples are given below:
* $`f(u)=u^\beta \mathrm{exp}\{\rho u^{\frac{1}{\alpha }}\}`$, $`f(u)=(\mathrm{log}u)^\beta \mathrm{exp}\{\rho u^{\frac{1}{\alpha }}\}`$, where $`\beta `$ is arbitrary;
* $`f(u)=\mathrm{exp}\left\{u^{\frac{1}{\alpha }}(\rho +\frac{\alpha }{\mathrm{log}u})\right\}`$, $`f(u)=\mathrm{exp}\left\{u^{\frac{1}{\alpha }}[\rho +u^{\frac{2}{3\alpha }}\mathrm{cos}(u^{\frac{1}{3\alpha }})]\right\}`$;
* $`f(u)=\mathrm{exp}\left\{u^{\frac{1}{\alpha }}(\rho +u^{\frac{\alpha _11}{\alpha }})\right\}`$ with $`\alpha _1(0,1)`$;
* $`f(u)=\mathrm{exp}\{u^{\frac{1}{\alpha }}+\rho \mathrm{exp}\{u^{\frac{1}{\alpha }}\}\}`$, $`f(u)=\mathrm{exp}\{(u^{\frac{1}{\alpha }}+\rho )\mathrm{exp}\{u^{\frac{1}{\alpha }}\}\}`$
($`m=1`$ in (i)–(iii) and $`m=2`$ in (iv)).
###### Example 2.9.
Among functions $`f`$ which fulfill the hypotheses of Corollary 2.7, we illustrate: $`f(u)=\mathrm{exp}\{u\}1`$, $`f(u)=\mathrm{sinh}(u)`$, $`f(u)=\mathrm{cosh}(u)1`$, $`f(u)=\mathrm{exp}\{u\}\mathrm{log}(u+1)`$, $`f(u)=u^\beta \mathrm{exp}\{\rho u^{\frac{1}{\alpha }}\}`$ with $`\beta 1`$, $`\alpha ,\rho >0`$, $`f(u)=u^\beta \mathrm{exp}(\mathrm{exp}\{u\})`$ with $`\beta 1`$ and $`f(u)=\mathrm{exp}(\mathrm{exp}\{u\})e`$.
Boundary blow-up phenomena for $`(P)`$ with $`a=0`$, $`b=1`$ and $`f(u)=u^p`$, $`1<p2`$, appear in the analytical theory of a Markov process called superdiffusion. In this case, the uniqueness of the large solution was studied in Dynkin by probabilistic techniques. It is remarkable that Dynkin’s papers realize, on one hand, a connection between superprocesses and singularity phenomena and, on the other hand, they contain a probabilistic representation of the minimal large solution. By means of a probabilistic representation, a uniqueness result in domains with non-smooth boundary was established by Le Gall in the case $`p=2`$. The existence of large solutions is usually deduced by comparison methods combined with Keller-Osserman a priori bounds, Calderon-Zygmund estimates, Agmon-Douglis-Nirenberg’s theory, or Alexandrov and Krylov-Safonov techniques.
Our interest falls here on the uniqueness of large solutions to $`(P)`$ when $`f`$ does not vary regularly at infinity (thus excluding the power case). Note that if $`f(u)=\mathrm{exp}(u)1`$ or $`f(u)=\mathrm{exp}(u)u1`$, then by Corollary 2.7 the equation $`\mathrm{\Delta }u=f(u)`$ in $`\mathrm{\Omega }`$ has a unique large solution which satisfies $`lim_{d(x)0}\frac{u(x)}{\mathrm{log}[d(x)]^2}=1`$. This asymptotic behaviour is exactly the same as for the unique large solution of $`\mathrm{\Delta }u=\mathrm{exp}(u)`$ in $`\mathrm{\Omega }`$, going back to the pioneering works of Bieberbach and Rademacher . For the two-term asymptotic expansion of the large solution of $`\mathrm{\Delta }u=\mathrm{exp}(u)`$ we refer to . We point out that our approach is completely different from the above papers for it relies exclusively on the regular variation theory (see for details) not only in the statement, but also in the proof of the main result.
## 3. Auxiliary results
For details about Propositions 3.1 and 3.3 we refer the reader to ( or ).
###### Proposition 3.1 (Elementary properties of slowly varying functions).
Assume that $`L`$ is a slowly varying function. Then the following hold
* $`\mathrm{log}L(u)/\mathrm{log}u0`$ as $`u\mathrm{}`$.
* For any $`m>0`$, $`u^mL(u)\mathrm{}`$, $`u^mL(u)0`$ as $`u\mathrm{}`$.
* $`(L(u))^m`$ varies slowly for every $`m`$.
* If $`L_1`$ varies slowly, so do $`L(u)L_1(u)`$ and $`L(u)+L_1(u)`$.
###### Remark 3.2.
If $`gRV_\rho `$ with $`\rho >0`$ ($`\rho <0`$), then $`lim_u\mathrm{}g(u)=\mathrm{}`$ ($`0`$). However, the behaviour at infinity for a slowly varying function cannot be predicted. We see that $`L(u)=\mathrm{exp}\{(\mathrm{log}u)^{1/3}\mathrm{cos}((\mathrm{log}u)^{1/3})\}`$ is a (normalised) slowly varying function (use $`lim_u\mathrm{}\frac{uL^{}(u)}{L(u)}=0`$) for which $`lim\; inf_u\mathrm{}L(u)=0`$ and $`lim\; sup_u\mathrm{}L(u)=\mathrm{}`$.
###### Proposition 3.3 (Karamata’s Theorem).
Let $`RRV_\rho `$ be locally bounded in $`[A,\mathrm{})`$. Then, for any $`j<(\rho +1)`$ (resp., $`j=(\rho +1)`$ if $`^{\mathrm{}}x^{(\rho +1)}R(x)𝑑x<\mathrm{}`$)
$$\underset{u\mathrm{}}{lim}\frac{u^{j+1}R(u)}{_u^{\mathrm{}}x^jR(x)𝑑x}=(j+\rho +1).$$
Under the assumptions of Theorem 2.3, we prove
###### Lemma 3.4.
The function $`\mathrm{\Phi }`$ given by (2.2) is well defined on some interval $`(0,\beta )`$. Furthermore, $`\mathrm{\Phi }NRV_{\frac{2(\theta +1)}{\rho }}(0+)`$ satisfies
(3.1)
$$\underset{t0^+}{lim}\frac{\mathrm{log}_m\mathrm{\Phi }(t)}{\mathrm{log}_m(\frac{1}{t})}=\{\begin{array}{cc}\frac{2(1+\theta )}{\rho }\hfill & \text{if }m=1,\hfill \\ 1\hfill & \text{if }m2.\hfill \end{array}$$
(3.2)
$$\underset{t0^+}{lim}\frac{\mathrm{\Phi }(t)\mathrm{\Phi }^{\prime \prime }(t)}{[\mathrm{\Phi }^{}(t)]^2}=1+\frac{\rho }{2(\theta +1)}\text{and}\underset{t0^+}{lim}\frac{𝔏(\mathrm{\Phi }(t))}{𝔏^{}(\mathrm{\Phi }(t))}\frac{\mathrm{\Phi }(t)}{[\mathrm{\Phi }^{}(t)]^2}=0.$$
###### Proof.
Let $`b>0`$ be such that $`𝑳_f`$ resp., $`𝔏^{}`$ is positive on $`[b,\mathrm{})`$. Since $`𝔏^{}RV_1`$ and $`𝑳_f`$ is slowly varying, Proposition 3.1 yields
$$\underset{u\mathrm{}}{lim}\frac{[𝔏^{}(u)]^{\frac{1}{2}}}{u^{\frac{\rho +1}{2}}[𝑳_f(u)]^{\frac{1}{2}}}u^{1+\tau }=0,\text{for any}\tau (0,\rho /2).$$
Therefore, there exists $`B>b`$ large so that
$$\zeta (x)=_x^{\mathrm{}}\frac{[𝔏^{}(y)]^{\frac{1}{2}}}{y^{\frac{\rho +1}{2}}[𝑳_f(y)]^{\frac{1}{2}}}𝑑y<\mathrm{},x>B.$$
It follows that $`\mathrm{\Phi }`$ is well defined on $`(0,\beta )`$, for some $`\beta >0`$. Moreover, $`\mathrm{\Phi }C^2(0,\beta )`$ and $`lim_{t0^+}\mathrm{\Phi }(t)=\mathrm{}`$. Using (2.2), we find
(3.3)
$$\frac{\mathrm{\Phi }^{}(t)[𝔏^{}(\mathrm{\Phi }(t))]^{\frac{1}{2}}}{[\mathrm{\Phi }(t)]^{\frac{\rho +1}{2}}[𝑳_f(\mathrm{\Phi }(t))]^{\frac{1}{2}}}=𝔎(t),t(0,\beta ).$$
In view of Proposition 3.3, we have
$$\underset{u\mathrm{}}{lim}\frac{[𝔏^{}(u)]^{\frac{1}{2}}}{u^{\frac{\rho 1}{2}}[𝑳_f(u)]^{\frac{1}{2}}\zeta (u)}=\frac{\rho }{2}$$
which, together with (2.2), produces
(3.4)
$$\underset{t0^+}{lim}\frac{[𝔏^{}(\mathrm{\Phi }(t))]^{\frac{1}{2}}[\mathrm{\Phi }(t)]^{\frac{\rho +1}{2}}}{[𝑳_f(\mathrm{\Phi }(t))]^{\frac{1}{2}}_0^t𝔎(s)𝑑s}=\frac{\rho }{2}.$$
By (3.3), (3.4) and L’Hospital’s rule, we find
(3.5)
$$\underset{t0^+}{lim}\frac{\mathrm{log}\mathrm{\Phi }(t)}{\mathrm{log}(_0^t𝔎(s)𝑑s)}=\underset{t0^+}{lim}\frac{\mathrm{\Phi }^{}(t)}{\mathrm{\Phi }(t)}\frac{_0^t𝔎(s)𝑑s}{𝔎(t)}=\frac{2}{\rho }.$$
We differentiate (3.3) to obtain
(3.6) $`\mathrm{\Phi }^{\prime \prime }(t)={\displaystyle \frac{𝔎(t)\mathrm{\Phi }^{}(t)[𝑳_f(\mathrm{\Phi }(t))]^{\frac{1}{2}}}{[𝔏^{}(\mathrm{\Phi }(t))]^{\frac{1}{2}}[\mathrm{\Phi }(t)]^{\frac{1\rho }{2}}}}`$ $`\{{\displaystyle \frac{\rho +1}{2}}+{\displaystyle \frac{𝔎^{}(t)\mathrm{\Phi }(t)}{𝔎(t)\mathrm{\Phi }^{}(t)}}+{\displaystyle \frac{\mathrm{\Phi }(t)𝑳_f^{}(\mathrm{\Phi }(t))}{2𝑳_f(\mathrm{\Phi }(t))}}`$
$`{\displaystyle \frac{\mathrm{\Phi }(t)𝔏^{\prime \prime }(\mathrm{\Phi }(t))}{2𝔏^{}(\mathrm{\Phi }(t))}}\}`$
for each $`t(0,\beta )`$. By $`𝔎NRV_\theta (0+)`$ we mean $`\stackrel{~}{𝔎}(u)=𝔎(1/u)NRV_\theta `$. Hence, $`lim_{t0^+}\frac{t𝔎^{}(t)}{𝔎(t)}=\theta `$ and $`lim_{t0^+}\frac{_0^t𝔎(s)𝑑s}{t𝔎(t)}=\frac{1}{\theta +1}`$. This, combined with (3.5), yields
(3.7)
$$\underset{t0^+}{lim}\frac{𝔎^{}(t)}{𝔎(t)}\frac{\mathrm{\Phi }(t)}{\mathrm{\Phi }^{}(t)}=\frac{\rho \theta }{2(\theta +1)}\text{and}\underset{t0^+}{lim}\frac{t\mathrm{\Phi }^{}(t)}{\mathrm{\Phi }(t)}=\frac{2(\theta +1)}{\rho }.$$
Thus, $`\mathrm{\Phi }NRV_{\frac{2(\theta +1)}{\rho }}(0+)`$. By (3.7) and L’Hospital’s rule, we obtain
(3.8)
$$\underset{t0^+}{lim}\frac{\mathrm{log}\mathrm{\Phi }(t)}{\mathrm{log}t}=\underset{t0^+}{lim}\frac{t\mathrm{\Phi }^{}(t)}{\mathrm{\Phi }(t)}=\frac{2}{\rho }(1+\theta ).$$
Proceeding by induction, we conclude (3.1). Since $`𝑳_f`$ is a normalised slowly varying function and $`𝔏^{}NRV_1`$, we have
(3.9)
$$\underset{t0^+}{lim}\frac{\mathrm{\Phi }(t)𝑳_f^{}(\mathrm{\Phi }(t))}{𝑳_f(\mathrm{\Phi }(t))}=0\text{and}\underset{t0^+}{lim}\frac{\mathrm{\Phi }(t)𝔏^{\prime \prime }(\mathrm{\Phi }(t))}{𝔏^{}(\mathrm{\Phi }(t))}=1.$$
By (3.6), (3.7) and (3.9), we infer that
$$\underset{t0^+}{lim}\frac{\mathrm{\Phi }^{\prime \prime }(t)[𝔏^{}(\mathrm{\Phi }(t))]^{\frac{1}{2}}}{𝔎(t)\mathrm{\Phi }^{}(t)[\mathrm{\Phi }(t)]^{\frac{\rho 1}{2}}[𝑳_f(\mathrm{\Phi })]^{\frac{1}{2}}}=\left(1+\frac{\rho }{2(\theta +1)}\right).$$
Replacing $`𝔎(t)`$ by its value in (3.3), we obtain the first assertion of (3.2). Moreover,
(3.10)
$$\underset{t0^+}{lim}\frac{\mathrm{log}(\mathrm{\Phi }^{}(t))}{\mathrm{log}\mathrm{\Phi }(t)}=\underset{t0^+}{lim}\frac{\mathrm{\Phi }^{\prime \prime }(t)\mathrm{\Phi }(t)}{[\mathrm{\Phi }^{}(t)]^2}=1+\frac{\rho }{2(1+\theta )}.$$
Since $`𝔏`$ varies slowly at $`\mathrm{}`$ and $`𝔏^{}RV_1`$, we use Proposition 3.1 (i) to obtain
(3.11)
$$\underset{t0^+}{lim}\frac{\mathrm{log}𝔏(\mathrm{\Phi }(t))}{\mathrm{log}\mathrm{\Phi }(t)}=0\text{and}\underset{t0^+}{lim}\frac{\mathrm{log}𝔏^{}(\mathrm{\Phi }(t))}{\mathrm{log}\mathrm{\Phi }(t)}=1.$$
We notice that
$$\mathrm{log}\left(\frac{𝔏(\mathrm{\Phi }(t))\mathrm{\Phi }(t)}{𝔏^{}(\mathrm{\Phi }(t))[\mathrm{\Phi }^{}(t)]^2}\right)=\mathrm{log}\mathrm{\Phi }(t)\left[1+\frac{\mathrm{log}𝔏(\mathrm{\Phi }(t))}{\mathrm{log}\mathrm{\Phi }(t)}\frac{2\mathrm{log}|\mathrm{\Phi }^{}(t)|}{\mathrm{log}\mathrm{\Phi }(t)}\frac{\mathrm{log}𝔏^{}(\mathrm{\Phi }(t))}{\mathrm{log}\mathrm{\Phi }(t)}\right]$$
which, together with (3.10) and (3.11), leads to
$$\underset{t0^+}{lim}\mathrm{log}\left(\frac{𝔏(\mathrm{\Phi }(t))\mathrm{\Phi }(t)}{𝔏^{}(\mathrm{\Phi }(t))[\mathrm{\Phi }^{}(t)]^2}\right)=\mathrm{}.$$
Thus the second claim of (3.2) is proved. ∎
## 4. Proof of Theorem 2.3
Let us first remark that the Keller–Osserman condition $`(A_2)`$ holds. Indeed, by using Proposition 3.1, we arrive at
$$\underset{z\mathrm{}}{lim}\frac{f(z)}{z^p}=\underset{u\mathrm{}}{lim}\frac{f(𝔏(u))}{u^\rho 𝑳_f(u)}\frac{u^\rho 𝑳_f(u)}{[𝔏(u)]^p}=\underset{u\mathrm{}}{lim}\frac{u^\rho 𝑳_f(u)}{[𝔏(u)]^p}=\mathrm{},p>1.$$
Thus, Eq. ($`P`$) has at least a large solution when $`a<\lambda _{\mathrm{},1}`$ and no large solution provided that $`a\lambda _{\mathrm{},1}`$ (see \[9, Theorem 1.1\]).
We now prove that, for $`a<\lambda _{\mathrm{},1}`$ fixed, every large solution of ($`P`$) exhibits the same asymptotic behaviour near $`\mathrm{\Omega }`$, namely (2.1). Set $`\vartheta ^\pm =\left(\frac{\rho }{2(1+\theta )(12ϵ_0)}\right)^{\frac{1}{\rho }}`$, where $`ϵ_0(0,1/2)`$ is arbitrary. Let $`\delta (0,\beta /2)`$ be small such that
* $`d(x)`$ is a $`C^2`$-function on the set $`\{x^N:d(x)<2\delta \}`$.
* $`𝔎`$ is nondecreasing on $`(0,2\delta )`$.
* $`(1ϵ_0)𝔎^2(d(x))<b(x)<(1+ϵ_0)𝔎^2(d(x))`$, for all $`x\mathrm{\Omega }`$ with $`d(x)<2\delta `$.
* $`𝔏(\vartheta ^\pm \mathrm{\Phi }(2\delta ))>0`$.
Let $`\sigma (0,\delta )`$ be arbitrary. We define $`u_\sigma ^\pm (x)=𝔏(\vartheta ^\pm \mathrm{\Phi }(d(x)\sigma ))`$, where $`d(x)(\sigma ,2\delta )`$ (resp., $`d(x)+\sigma <2\delta `$) for $`u_\sigma ^+(x)`$ (resp., $`u_\sigma ^{}(x)`$). It follows that
$`\mathrm{\Delta }u_\sigma ^\pm `$ $`=\mathrm{div}\left(\vartheta ^\pm 𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(d(x)\sigma ))\mathrm{\Phi }^{}(d(x)\sigma )d(x)\right)`$
$`=(\vartheta ^\pm )^2𝔏^{\prime \prime }(\vartheta ^\pm \mathrm{\Phi }(d(x)\sigma ))[\mathrm{\Phi }^{}(d(x)\sigma )]^2|d(x)|^2`$
$`+\vartheta ^\pm 𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(d(x)\sigma ))\mathrm{\Phi }^{\prime \prime }(d(x)\sigma )|d(x)|^2`$
$`+\vartheta ^\pm 𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(d(x)\sigma ))\mathrm{\Phi }^{}(d(x)\sigma )\mathrm{\Delta }d(x).`$
In view of (i)–(iii), when $`\sigma <d(x)<2\delta `$ we obtain (since $`|d(x)|=1`$)
(4.1) $`\mathrm{\Delta }u_\sigma ^++au_\sigma ^+b(x)f(u_\sigma ^+){\displaystyle \frac{\vartheta ^+𝔏^{}(\vartheta ^+\mathrm{\Phi }(d(x)\sigma ))[\mathrm{\Phi }^{}(d(x)\sigma )]^2}{\mathrm{\Phi }(d(x)\sigma )}}\times `$
$`\times \left\{{\displaystyle \frac{\mathrm{\Phi }(d(x)\sigma )}{\mathrm{\Phi }^{}(d(x)\sigma )}}\mathrm{\Delta }d+^+(d(x)\sigma )\right\}`$
respectively, when $`d(x)+\sigma <2\delta `$,
(4.2) $`\mathrm{\Delta }u_\sigma ^{}+au_\sigma ^{}b(x)f(u_\sigma ^{}){\displaystyle \frac{\vartheta ^{}𝔏^{}(\vartheta ^{}\mathrm{\Phi }(d(x)+\sigma ))[\mathrm{\Phi }^{}(d(x)+\sigma )]^2}{\mathrm{\Phi }(d(x)+\sigma )}}\times `$
$`\times \left\{{\displaystyle \frac{\mathrm{\Phi }(d(x)+\sigma )}{\mathrm{\Phi }^{}(d(x)+\sigma )}}\mathrm{\Delta }d+^{}(d(x)+\sigma )\right\}`$
Here $`^\pm `$ are real functions defined on $`(0,2\delta )`$ as follows
(4.3)
$$^\pm (t):=\frac{\mathrm{\Phi }^{\prime \prime }(t)\mathrm{\Phi }(t)}{[\mathrm{\Phi }^{}(t)]^2}+\frac{\vartheta ^\pm \mathrm{\Phi }(t)𝔏^{\prime \prime }(\vartheta ^\pm \mathrm{\Phi }(t))}{𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(t))}+\frac{a𝔏(\vartheta ^\pm \mathrm{\Phi }(t))}{\vartheta ^\pm 𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(t)))}\frac{\mathrm{\Phi }(t)}{[\mathrm{\Phi }^{}(t)]^2}𝒟^\pm (t)$$
where we denote
$$𝒟^\pm (t)=(1ϵ_0)\frac{𝔎^2(t)f(𝔏(\vartheta ^\pm \mathrm{\Phi }(t))}{\vartheta ^\pm 𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(t))}\frac{\mathrm{\Phi }(t)}{[\mathrm{\Phi }^{}(t)]^2}.$$
By virtue of (3.3), we may rewrite $`𝒟^\pm (t)`$ as
$$𝒟^\pm (t)=(1ϵ_0)(\vartheta ^\pm )^\rho \frac{f(𝔏(\vartheta ^\pm \mathrm{\Phi }(t)))}{(\vartheta ^\pm \mathrm{\Phi }(t))^\rho 𝑳_f(\vartheta ^\pm \mathrm{\Phi }(t))}\frac{𝑳_f(\vartheta ^\pm \mathrm{\Phi }(t))}{𝑳_f(\mathrm{\Phi }(t))}\frac{𝔏^{}(\mathrm{\Phi }(t))(\vartheta ^\pm )^1}{𝔏^{}(\vartheta ^\pm \mathrm{\Phi }(t))}.$$
It follows that $`lim_{t0^+}𝒟^\pm (t)=(1ϵ_0)(\vartheta ^\pm )^\rho `$. Note that $`lim_u\mathrm{}\frac{u𝔏^{\prime \prime }(u)}{𝔏^{}(u)}=1`$ (since $`𝔏^{}NRV_1`$). Moreover, by (4.3) and Lemma 3.4, we find
$$\underset{t0^+}{lim}^\pm (t)=\frac{\rho }{2(1+\theta )}(1ϵ_0)(\vartheta ^\pm )^\rho =\frac{\rho }{2(1+\theta )}\frac{\pm ϵ_0}{12ϵ_0}.$$
Hence, using (4.1) and (4.2), we can choose $`\delta >0`$ small enough so that
(4.4)
$$\begin{array}{cc}\hfill \mathrm{\Delta }u_\sigma ^++au_\sigma ^+b(x)f(u_\sigma ^+)& 0,x\text{with}\sigma <d(x)<2\delta \hfill \\ \hfill \mathrm{\Delta }u_\sigma ^{}+au_\sigma ^{}b(x)f(u_\sigma ^{})& 0,x\text{with}d(x)+\sigma <2\delta .\hfill \end{array}$$
For $`\eta >0`$, we define $`C_\eta =\{x\mathrm{\Omega }:d(x)>\eta \}`$ and $`D_\eta =\{x^N:\mathrm{dist}(x,\mathrm{\Omega })<\eta \}`$.
Let $`\eta >0`$ be small such that $`a<\lambda _1(\mathrm{\Delta },D_\eta \overline{\mathrm{\Omega }})`$, where $`\lambda _1(\mathrm{\Delta },D_\eta \overline{\mathrm{\Omega }})`$ denotes the first Dirichlet eigenvalue of $`(\mathrm{\Delta })`$ in the domain $`D_\eta \overline{\mathrm{\Omega }}`$. Let $`pC^{0,\mu }(\overline{D}_{2\eta })`$ satisfy $`0<p(x)b(x)`$ for $`x\mathrm{\Omega }C_{2\delta }`$, $`p=0`$ on $`\overline{D}_\eta \mathrm{\Omega }`$ and $`p>0`$ on $`\overline{D}_{2\eta }\overline{D}_\eta `$.
We denote by $`w`$ a positive large solution of $`\mathrm{\Delta }w+aw=p(x)f(w)`$ in $`D_{2\eta }\overline{C}_\delta `$ (see \[9, Theorem 1.1\]). Set $`U:=u_a+w`$ and $`V_\sigma :=u_\sigma ^++w`$, where $`u_a`$ is an arbitrary large solution of ($`P`$). It follows that
$$\mathrm{\Delta }U+aUb(x)f(U)0\text{in}\mathrm{\Omega }\overline{C}_\delta \text{and}\mathrm{\Delta }V_\sigma +aV_\sigma b(x)f(V_\sigma )0\text{in}C_\sigma \overline{C}_\delta .$$
Notice that $`U|_\mathrm{\Omega }=\mathrm{}>u_\sigma ^{}|_\mathrm{\Omega }`$, $`U|_{C_\delta }=\mathrm{}>u_\sigma ^{}|_{C_\delta }`$, resp., $`V_\sigma |_{C_\sigma }=\mathrm{}>u_a|_{C_\sigma }`$ and $`V_\sigma |_{C_\delta }=\mathrm{}>u_a|_{C_\delta }`$. Thus, by \[9, Lemma 2.1\], we deduce $`u_a+wu_\sigma ^{}`$ on $`\mathrm{\Omega }\overline{C}_\delta `$ and $`u_\sigma ^++wu_a`$ on $`C_\sigma \overline{C}_\delta `$. Letting $`\sigma 0^+`$, we find
$$𝔏(\vartheta ^{}\mathrm{\Phi }(d))wu_a𝔏(\vartheta ^+\mathrm{\Phi }(d))+w\text{on}\mathrm{\Omega }\overline{C}_\delta .$$
Since $`w`$ is uniformly bounded on $`\mathrm{\Omega }`$ and $`𝔏`$ varies slowly at $`\mathrm{}`$, we conclude (2.1).
Thus, $`lim_{d(x)0}\frac{u_1(x)}{u_2(x)}=1`$ for any two large solutions $`u_1`$, $`u_2`$ of ($`P`$). From now on, we can use the same line of reasoning as in the proof of \[8, Theorem 1\] to obtain $`u_1u_2`$ on $`\mathrm{\Omega }`$. ∎
## Acknowledgment
F. Cîrstea wishes to thank Prof. N. Dancer for bringing this problem into her attention. She also thanks Prof. N. S. Trudinger for fruitful discussions, hospitality and support received during a short visit to the ANU in January 2004.
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# NUHEP Report 989 UW Report MADPH-04-1382 May 8, 2005 New evidence for the saturation of the Froissart bound
## Acknowledgments
The work of FH is supported in part by the U.S. Department of Energy under Grant No. DE-FG02-95ER40896 and in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation.
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# Two early-stage inverse power-law dynamics in nonlinear complex systems far-from-equilibrium
## I Introduction
Slow relaxation phenomena in random composite materials made of components having widely different generalized susceptibilities (e.g., permeability, dielectric constant, electrical/ thermal conductivity, viscosity, elastic modulii etc.), continue to remain intriguing and hence a topic of intense research. The susceptibility, being a measure of the response of a system to an appropriate external perturbation or a driving force, is typically assumed to be linear under a vanishingly small external force. But, many natural systems/ phenomena do not show any measurable response until an appropriate driving force exceeds a measurable threshold. For example, a rigid body on a rough surface does not move until the driving force exceeds a finite frictional force. More intriguingly, it is almost a law of nature that in such driven, macroscopic systems, the response characteristic is nonlinear kardar with a concommitant critical behaviour at the threshold. Besides the case of frictional motion, nonlinear response above a threshold is also found in pinned charge-density-wave systems, disordered granular superconductors doniach , etc.
In studying the dynamical behaviour of a material, one usually measures its appropriate response property, say, $`\varphi (t)`$, as a function of time $`t`$, from a non-equilibrium to an equilibrium (for a closed system) or to a steady (for an open system) state. In general, this relaxation is classified into two groups: (i) a purely Debye type with an exponential relaxation function, $`\varphi (t)=\mathrm{exp}(t/\tau )`$, $`\tau `$ being a characteristic time-scale, called the relaxation time; or (ii) a non-Debye type where $`\varphi (t)`$ is either a linear superposition of exponential functions, or a sub-exponential function (as in some glassy systems), with multiple relaxation times. In more intriguing cases, this non-Debye relaxation may even be a power-law or a logarithmic function, without any time-scale, or one with $`\tau \mathrm{}`$. The appearence of such a scale-free, slow dynamics with one or more power-laws from the early stages of evolution, is what concerns us here because of the fundamental issue of the failure of Boltzmann’s relaxation time approximation zim in driven systems far-from-equilibrium.
## II Experiments and some related models
In a lucidly written review Scher, Shlesinger and Bendler scher focus on experimental observations (1970’s onwards) and the origin of two power-law kinetics. For a few examples, we cite some transient photocurrent measurements tie on $`a`$-Si:H, $`a`$-As<sub>2</sub>Se<sub>3</sub> etc. and a couple of hole transport hole data on PVK and Si-MOS devices, where two consecutive power-law decays, of the forms $`t^\alpha `$ and $`t^\beta `$ ($`\alpha ,\beta >0`$), covering one or more decades in time each (with a crossover in-between) were observed. Based on the continuous time random walk with a long-tailed power-law probability density function for the random waiting times (release time of trapped carriers by tunneling), Scher et al. scher formulated a theory regarding the above results. The latter long-tailed power-law function violates the Central Limit Theorem, since all of its moments including the first (mean waiting-time) diverge. The unifying feature of the above random walk is the scale-invariance of the shape of the relaxation current $`I(t)`$, if one normalizes the time by a transit time $`t_r`$, which is a sample dependent parameter. This stochastic theory by Scher et al. scher explains the results of many early experiments, following the relation $`\alpha +\beta =2`$.
But, there is a huge variety of relatively recent, more intriguing experiments on soft-condensed or complex systems where the couple of exponents $`\alpha `$ and $`\beta `$ do not seem to follow any simple algebraic relation. Weron and Jurlewicz note that an experiment on dielectric relaxation in a system of dipoles hill involves the crossover between different forms of self-similarity. They argue that wj a couple of power-law decays appear due to the coupling of micro-clusters of dipoles with a distribution of $`\tau `$’s (i.e., multiple time-scales). An experiment on fluorescence intermittency of single ZnS overcoated CdSe quantum dots kuno , the distribution of on and off times (blinking kinetics) is reported to follow a single power-law <sup>2</sup><sup>2</sup>2Our analysis suggests that two inverse power-law kinetics are present in this experiment kuno , as well.. In a biological system of photo-dissociated heme-proteins, the rebinding of the ligands of iron (i.e., the CO and the O<sub>2</sub> molecules) is observed to follow an inverse power-law dynamics pf . In a theoretical study of the same work, Tsallis et al. tsal claim two inverse power-law regimes and demonstrated its possible connection with nonextensive thermo-statistics (entropy). Naudts and Czachor naud analyzed the data of many experiments including those of kuno ; pf , and maintain that these two power-law decays result from some choice of parameters of their probability density function. All of the above theories lead to the result that $`\alpha <\beta `$, if two power-law relaxations exist. In a different biological system, the dynamics of Ca channels in living cells bez , the distribution of the survival-times of the channels has been studied. A stochastic dynamical model, with one dimensional geometry, was proposed specifically to explain only the later power-law bar dynamics <sup>3</sup><sup>3</sup>3Two inverse power-law kinetics are clearly present in the Fig.2, a typical result, of this 1-d model bar ..
Two power-law growths are also claimed to have been found, using atomic force microscopy (AFM), in the early dynamics of sputtered Ag particles on Si(0 0 1) substrate sangam . These authors relate the two different growth-dynamics to two competing structural rearrangements at two different length-scales. In the case of the growth of a large single DLA (diffusion limited aggregation) cluster, using upto 10<sup>8</sup> particles in a computer experiment, two power-laws seem to dictate the growth process as a function mandel of a time-like entity. A similar theme has been reflected in various methods (e.g., the sol-gel method) of synthesis of nanomaterials edcam . This involves a slow relaxation of restructuring many local clusters by crossing the local barriers due to a caging effect (typical of liquid-like, amorphous or dense granular materials) followed by (or, may be, added to) another slow relaxation of global restructuring of the local clusters themselves. Even in some relatively recent studies of cellular automata models of earthquakes naka , one power-law dynamics seems to result under certain choice of parameters.
## III A tunneling percolation model
To study the basic physics behind the transport properties within a minimal model, Sen and co-workers dcac ; pre1 constructed a lattice based bond percolation model. In this model, in addition to the randomly placed ohmic bonds (linear response) in the Random Resistor Network (RRN), tunneling bonds (t-bond) were introduced between two nearest neighbour ($`nn`$) ohmic bonds (o-bonds) separated by one insulator only. This tunneling is actually semi-quantum or semi-classical in the sense that no quantum mechanical phase-information of the charge carriers appears in the model. In the steady-state, a t-bond is insulating if the microscopic voltage difference across it ($`v`$) is less than a fixed voltage threshold ($`v_g`$, identical for all the t-bonds) and conducts linearly a current if $`vv_g`$. Indeed, the phenomenological parameters $`p`$ and $`v_g`$ ensure that both the disorder and the coulomb interaction (nonlinear response mehran as an outcome) is in-built in this model.
Historically dcac , we call this tunneling-enhanced percolation model as the Random Resistor cum Tunneling Network (RRTN). The appearence of the t-bonds in this perfectly correlated, i.e., deterministic fashion is the origin of a very low percolation threshold, in the RRTN model. For example, in the fig.1(a), the RRN with $`p`$=0.3, does not percolate. In this respect, one may note that the percolation threshold for the RRN (square lattice) is $`p_c`$ = 0.50. But, the maximal RRTN (where all the possible t-bonds, both along and across the field, are active, i.e., conducts a maximum possible current in the steady state) generated from that particular RRN, percolates as shown in the fig.1(b). Standard calculations with finite-size scaling analysis, establishes a new percolation threshold of $`p_{ct}`$ = 0.18 pre1 , for the maximal RRTN.
Further, because of the finite threshold of the t-bonds, the bulk dc conductance $`G(V)`$ (defined either as $`\frac{I}{V}`$ or as $`\frac{dI}{dV}`$) is a strongly nonlinear curve, as stated earlier. At a sufficiently low voltage (below $`v_g`$), no t-bond is active, and hence the effective bulk conductance is only due to the percolating ohmic backbone. So the $`G(V)`$ vs $`V`$ graph is parallel to the $`V`$-axis (falls on the $`V`$-axis if the ohmic backbone is non-percolating). This region is termed as the lower linear/ohmic regime. Thereafter, under further increment of the external voltage $`V`$, yet more microscopic t-bonds cross the threshold ($`v_g`$), i.e., become active, and give rise to extra parallel paths in the network. At an very high voltage foot2 (asymptotically infinite), all the possible t-bonds become active (i.e., creates a maximal RRTN) and hence the conductance saturates to its largest possible value (for that particular configuration). Obviously, this region belongs to the upper linear/ohmic regime. Other than the success of this discrete model in understanding various experiments on an ultra-low percolation threshold, and various aspects of nonlinear dc and ac responses, it has also been quite useful for understanding some very unusual aspects observed in low-temperature variable range hopping (VRH) conduction vrh , and some interesting aspects of breakdown phenomena break . Thus, even though, time enters in an implicit fashion in some of the above studies, an explicit characterization of the relaxation dynamics in the RRTN model in the perspective of various experiments and models (as described above) was considered necessary.
## IV Non-equilibrium dynamics in the RRTN model
In this paper, we study the current relaxation dynamics towards a steady-state in the RRTN model. Away from equilibrium, a t-bond with a microscopic $`v<v_g`$ behaves like a dielectric material between two metals (o-bonds); and the resulting charging effect gives rise to a displacement current ($`C\frac{dv(t)}{dt}`$, where $`C`$ is the capacitance). Thus, a t-bond gives rise to a displacement current if $`v(t)<v_g`$, and an ‘ohmic’ current if $`v(t)v_g`$ sust . For our calculations, we use the values of the microscopic conductance $`g_o=1.0`$ (o-bonds), $`g_t=10^2`$ (t-bonds), $`v_g=0.5`$, and $`C=10^5`$ for the t-bonds (in some arbitrary units).
In our numerical study, we apply an uniform electric field across RRTN’s of different system sizes ($`L`$) and ohmic bond concentrations ($`p`$). We study the evolution of the current in a RRTN starting from the switching on state until it approaches its asymptotic steady-state. To do this, we follow basically the current conservation (Kirchhoff’s laws) locally at each node of the lattice. The aim is to study the achievement of a global current conservation as an outcome of the local current conservation (hence, the dynamics). A discrete, scaled time unit has been chosen as one complete scan through each site of the lattice. This local conservation or the equation of continuity reads as,
$$i_{nn}(t)=0,t.$$
(1)
Here the sum has been taken over currents $`i_{nn}(t)`$ through various types of nearest neighbour ($`nn`$) microscopic bonds around any node/site of the lattice. For the case of a square lattice, one considers the four $`nn`$’s around a node inside the bulk (three and two $`nn`$’s respectively at any boundary or a corner). If eq. (1) were true simultaneously for each site of the lattice, then the global conservation (the steady state) for the entire network would automatically be achieved. As we need to start with an initial (arbitrary) microscopic voltage distribution, the eq. (1) would not hold for all the sites of the lattice. Some correction term would be required at each site and this requirement leads to the following time evolution algorithm which we call as the lattice Kirchhoff’s dynamics:
$$v(j,k,t+1)=v(j,k,t)+\frac{i_{nn}(t)}{g_{nn}},$$
(2)
where $`g_{nn}`$ are the various microscopic conductances of the $`nn`$ bonds around the node $`(j,k)`$. Then we numerically solve a set of coupled difference equations on the lattice. The move towards a macroscopic steady-state implies that the difference of currents through the first and the last layers tends to zero as a function of time. In practice, the system is considered to have reached its steady state when this difference decreases to a pre-assigned smallness.
## V Results and discussions
As shown in the figs. 2, we observe a non-Debye type current relaxation with two consecutive initial power-laws (and a crossover in-between), each of them spanning at least a decade. This asymptotic steady state current (whether insulating or conducting) for any randomly chosen RRTN, is found to be very robust against any initiallly chosen voltage distribution on the lattice. To analyse the current evolution of an insulating or a conducting sample on the same footing, we subtract the corresponding steady current, $`I(t\mathrm{})`$, from the evolving current $`I(t)`$ at the time $`t`$. Our observation on the transient current response, indicates clearly the existence of a couple of initial power-laws, whose exponents differ significantly for systems with different configurations, $`p`$, $`L`$ and external voltage $`V`$. So, we choose to work with one sample at a time and analyze its results within our numerical accuracy.
For example, in the fig. 2(a), we show the dynamics for a sample with $`L=60,p=0.55,V=2.0`$ and in the fig 2(b) another with $`L=80,p=0.35,V=20.0`$. The first figure (i.e., fig. 2(a)) represents the class of relaxation, where the second exponent ($`0.72`$) is larger than the first ($`0.44`$) (the only class reported in our previous work at a particular $`p=p_c`$ sust ). There are quite a few theoretical works scher ; wj ; tsal in this regard. Some other experiments bez ; bar ; sangam find a second exponent smaller than the first. Since our earlier report sust , we have been able to reproduce this other class of relaxation (e.g., fig. 2(b) with the exponents $`0.78`$, $`0.57`$), as well, within the context of our RRTN. In special cases, we do find only one power-law relaxation, which may be considered to be the borderline between the above two, or as the merging of the two power-law exponents. Further, we do not find any particular relation between the exponents. Next, the existence or the lack of any asymptotic exponential kinetics is not explicitly stated in most of them. Indeed, in some of the theoretical studies (e.g. in scher ; wj ; mandel ), the second power-law persists upto asymptotically infinite times. This trend can not describe the possiblity for these systems to reach an appropriate steady state. In contrast, beyond the power-law relaxations (whether one or two), our model acquires a relaxation time $`\tau `$ and the system follows a fast exponential dynamics as a precursor of a final steady state which is effectively diffusive (or ohmic). Obviously, the power-law relaxations at times $`t<\tau `$ imply a strong deviation from the Boltzmann’s relaxation time approximation (i.e., strongly non-Debye type).
As far as the origin of the two power-law dynamics are concerned, we have already outlined the main content of some of them scher ; wj ; tsal ; bar ; naud ; naka in the section on experiments. In most of them they occur due to local structural rearrangements preceding the final global structural rearrangements. In our case, the basic structure in a particular sample is created once for all, and it is the fields across the bonds which keep changing in such a fashion that the local conservation (Kirchhoff’s laws at each node) dominates the first power-law regime and the global current conservation is dominated in the latter power-law regime (of course, there are the required structural rearrangements of the active t-bonds). So, it is interesting that these two very different mechanisms give rise to a qualitatively identical outcome. Further, since the power-law dynamics occurs even for $`p`$’s away from $`p_c`$ or $`p_{ct}`$, it clearly demonstrates that they are not organized by any type of criticality. Finally, as discussed above, while most of the other theoretical works, are destined to get only one class of two-power-law relaxation behaviour (namely, $`\alpha <\beta `$), the RRTN dynamics has the ability to capture both of the classes for different sets of parameters.
###### Acknowledgements.
Sincere thanks are due for warm supports to the organizers of many International Conferences where earlier results on relaxation in the RRTN were presented. Foremost among them are Profs. P. Sheng and P.M. Hui for one named ETOPIM5 (AKS in 1999) was held at Hong Kong. Two more on Nonlinear Dynamics and Earthquake (SB in SMR1322 in 2001 and SMR1519 in 2003) were held at the AS-ICTP, Trieste, Italy. The last one named ICAMMP2002 (AKS and SB in 2003) were held at the SUST, Sylhet, Bangladesh. AKS also acknowledges fruitful discussions with Prof. H. Levine in 2003 on their Ref.bar and the support (a Visiting Professor in 2004) plus discussions with Prof. A. Hansen of the NTNU, Trondheim, Norway.
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# Singlet-triplet decoherence due to nuclear spins in a double quantum dot
## I Introduction
Decoherence due to the coupling of a qubit to its environment is widely regarded as the major obstacle to quantum computing and quantum information processing in solid-state systems. Electron spins confined in semiconductor quantum dotsLoss and DiVincenzo (1998) couple to their environments primarily through the spin-orbit interaction and hyperfine interaction with nuclear spins in the surrounding lattice.Burkard et al. (1999); Cerletti et al. (2005) To reach the next step in coherent electron spin state manipulation, the strongest decoherence effects in this system must be understood and reduced, if possible.
The effects of spin-orbit interaction are reduced in confined quantum dots at low temperatures.Khaetskii and Nazarov (2000) Indeed, recent experiments give longitudinal relaxation times $`T_1`$ for quantum-dot-confined electrons that reach $`T_120\mathrm{ms}`$Kroutvar et al. (2004) in self-assembled dots and $`T_10.85\mathrm{ms}`$ in gated dotsElzerman et al. (2004), in agreement with theory.Golovach et al. (2004) These times suggest that the spin-orbit interaction is a relatively weak source of decoherence in these structures since theory predicts that the transverse spin decay time $`T_2`$ due to spin-orbit interaction alone (neglecting other sources of decoherence) would be given by $`T_2=2T_1`$.Golovach et al. (2004) Other strategies for reducing the effects of spin-orbit interaction may include using hole (instead of electron) spin, where a recent study has found that $`T_2=2T_1`$ also applies, and the hole spin relaxation time can be made even longer than that for the electron spin.Bulaev and Loss (2005)
Unlike the spin-orbit interaction, the hyperfine interaction of a single electron spin with a random nuclear spin environment can lead to pure dephasing, giving a transverse spin decay time on the order of $`5\mathrm{ns}`$,Khaetskii et al. (2002); Merkulov et al. (2002); Coish and Loss (2004) six orders of magnitude shorter than the measured longitudinal decay times $`T_1`$. To minimize errors during qubit gating operations in these proposed devices, this decay must be fully understood. The hyperfine interaction in a single quantum dot is described by a Hamiltonian $`H=𝐡𝐒`$, where $`𝐒`$ is the electron spin operator and $`𝐡`$ is a collective quantum nuclear spin operator, which we will refer to as the “Overhauser operator”. A common assumption in the literature is to replace the Overhauser operator by a classical effective magnetic field $`𝐡𝐁_N`$.Schulten and Wolynes (1978); Merkulov et al. (2002); Khaetskii et al. (2002); Erlingsson et al. (2003); Erlingsson and Nazarov (2004); Yuzbashyan et al. (2004); Bracker et al. (2005); Braun et al. (2005); Dutt et al. (2005); Taylor et al. (2005); Petta et al. (2005); Koppens et al. (2005); Greilich et al. (2005) Since a classical magnetic field only induces precession (not decoherence), the classical-field picture necessitates an ensemble of nuclear spin configurations to induce decay of the electron spin expectation value.Khaetskii et al. (2002); Merkulov et al. (2002) For experiments performed on a large bulk sample of electron spins, or experiments performed over timescales that are longer than the typical timescale for variation of $`𝐁_N`$, the source of the ensemble averaging is clear. However, one conclusion of this model is that single-electron-spin experiments performed over a timescale shorter than the nuclear spin correlation time should show no decay. This conclusion is contradicted by numerical Schliemann et al. (2002); Shenvi et al. (2005) and analytical Zurek et al. (2003); Coish and Loss (2004) results, which show that the quantum nature of the Overhauser operator can lead to rapid decay of a single electron spin, even for a fully static nuclear spin system. This rapid decay is, however, reversible with a standard Hahn spin-echo sequence in an applied magnetic field and the timescale of the decay can be increased by squeezing the nuclear spin state.Coish and Loss (2004)
Another potential solution to the hyperfine decoherence problem is to polarize the nuclear spins. Polarizing the nuclear spin system in zero applied magnetic field reduces the longitudinal spin-flip probability by the factor $`1/p^2N`$, where $`p`$ is the nuclear spin polarization and $`N`$ is the number of nuclear spins within the quantum dot.Burkard et al. (1999); Coish and Loss (2004) The effect on the transverse components of electron spin is different. Unless the nuclear spin state is squeezed or a spin-echo sequence is performed, the transverse components of electron spin will decay to zero in a time $`t_c5\mathrm{ns}`$ in a typical GaAs quantum dot. Polarizing the nuclear spin system increases $`t_c`$ by reducing the phase-space available for fluctuations in the Overhauser operator, resulting in $`t_c5\mathrm{ns}/\sqrt{1p^2}`$.Coish and Loss (2004) Recent experiments show that the nuclear spin system can be polarized by as much as 60%.Bracker et al. (2005) However, to achieve an order-of-magnitude increase in $`t_c`$, the polarization degree would have to be on the order of 99%,Cerletti et al. (2005) for which more ambitious polarization schemes have been proposed.Imamoglu et al. (2003)
If electron spins in quantum dots are to be used as quantum information processors, the two-electron states of double quantum dots must also be coherent during rapid two-qubit switching times.<sup>1</sup><sup>1</sup>endnote: 1For exchange gates with spin-1/2 qubitsLoss and DiVincenzo (1998), the relevant requirement is that the qubit switching time $`t_S`$ should be much smaller than the singlet-triplet decoherence time.Burkard et al. (1999) Measurements of singlet-triplet relaxation times $`t_{ST}`$ in vertical double dots $`(t_{ST}200\mu \mathrm{s})`$,Fujisawa et al. (2002) gated lateral double dots $`(t_{ST}70\mu \mathrm{s})`$,Petta et al. (2004a) and single dots $`(t_{ST}2.58\mathrm{ms})`$Hanson et al. (2005) suggest that these states may be very long-lived. Recent experiments have now probed the decoherence time of such states, which is believed to be limited by the hyperfine interaction with surrounding nuclear spins.Petta et al. (2005) The dramatic effect of the hyperfine interaction on two-electron states in a double quantum dot has previously been illustrated in experiments that show slow time-dependent current oscillations in transport current through a double dot in the spin blockade regime.Ono and Tarucha (2004)
It may be possible to circumvent some of the complications associated with single-spin decoherence by considering an encoded qubit, composed of the two-dimensional subspace of states with total $`z`$-projection of spin equal to zero for two electrons in a double quantum dot.Taylor et al. (2005) One potential advantage of such a setup is that it may be possible to reduce the strength of hyperfine coupling to the encoded state space for a symmetric double-dot (see Appendix A). A potential disadvantage of this scheme is that coupling to the orbital (charge) degree of freedom can then lead to additional decoherence, but we find that orbital dephasing can be made negligible under appropriate conditions (see Sec. IV). To achieve control of the singlet-triplet subspace, however, the decoherence process for the two-electron system should be understood in detail.
In this paper we give a fully quantum mechanical solution for the spin dynamics of a two-electron system coupled to a nuclear-spin environment via the hyperfine interaction in a double quantum dot. Although we focus our attention here on quantum dots, decoherence due to a spin bath is also an important problem for, e.g., proposals to use molecular magnets for quantum information processing.Cerletti et al. (2005); Meier et al. (2003a, b); Troiani et al. (2005) In fact, the problem of a pair of electrons interacting with a bath of nuclear spins via the contact hyperfine interaction has been addressed long ago to describe spin-dependent reaction rates in radicals.Schulten and Wolynes (1978); Werner et al. (1977) A semiclassical theory has been developed,Schulten and Wolynes (1978) in which electron spins in radicals experience a randomly oriented effective classical magnetic field due to the contact hyperfine interaction between electron and nuclear spins. In this semiclassical theory, random hopping events of the electrons were envisioned to induce a randomly fluctuating local magnetic field at the site of the electron spin, resulting in decay of a singlet-triplet correlator. Here, we solve a different problem. Ensemble averaging over nuclear spin configurations is natural for a large sample of $`10^{23}`$ radicals. In contrast, we consider the coherent dynamics of two-electron spin states within a single double quantum dot. More importantly, the Heisenberg exchange interaction, which was found to be negligible in Ref. Schulten and Wolynes, 1978, can be any value (large or small) in our system of interest. We find that a nonzero exchange interaction can lead to a drastic change in the form and timescale of decoherence. Moreover, this paper is of direct relevance to very recent experiments Johnson et al. (2005); Petta et al. (2005); Koppens et al. (2005) related to such double-dot systems.
The rest of this paper is organized as follows. In Sec. II we solve the problem for electron spin dynamics in the subspace of total spin $`z`$-component $`S^z=0`$ with an exact solution for the projected effective Hamiltonian. In Sec. III we show that a perturbative solution is possible for electron spin dynamics in the subspace of singlet and $`S^z=+1`$ triplet states. Sec. IV contains a discussion of the contributions to singlet-triplet decoherence from orbital dephasing. In Sec. V we review our most important results. Technical details are given in Appendixes A to C.
## II Dynamics in the $`S^z=0`$ subspace
We consider two electrons confined to a double quantum dot, of the type considered, for example, in Refs. Johnson et al., 2005; Petta et al., 2005; Koppens et al., 2005. Each electron spin experiences a Zeeman splitting $`ϵ_z=g\mu _BB`$ due to an applied magnetic field $`𝐁=(0,0,B)`$, $`B>0`$, defining the spin quantization axis $`z`$, which can be along or perpendicular to the quantum dot axis. In addition, each electron interacts with an independent quantum nuclear field $`𝐡_l,l=1,2`$, due to the contact hyperfine interaction with surrounding nuclear spins. The nuclear field experienced by an electron in orbital state $`l`$ is $`𝐡_l=_kA_k^l𝐈_k`$, where $`𝐈_k`$ is the nuclear spin operator for a nucleus of total spin $`I`$ at lattice site $`k`$, and the hyperfine coupling constants are given by $`A_k^l=vA\left|\psi _0^l(𝐫_k)\right|^2`$, with $`v`$ the volume of a unit cell containing one nuclear spin, $`A`$ characterizes the hyperfine coupling strength, and $`\psi _0^l(𝐫_k)`$ is the single-particle envelope wavefunction for orbital state $`l`$, evaluated at site $`k`$. This problem simplifies considerably in a moderately large magnetic field ($`B\mathrm{max}\{\delta 𝐡_{\mathrm{rms}}/\mathrm{g}\mu _\mathrm{B},𝐡_{\mathrm{rms}}/\mathrm{g}\mu _\mathrm{B}\}`$, where $`𝒪_{\mathrm{rms}}=\psi _I\left|𝒪^2\right|\psi _I^{1/2}`$ is the root-mean-square expectation value of the operator $`𝒪`$ with respect to the nuclear spin state $`|\psi _I`$, $`\delta 𝐡=\frac{1}{2}\left(𝐡_1𝐡_2\right)`$, and $`𝐡=\frac{1}{2}\left(𝐡_1+𝐡_2\right)`$). In a typical unpolarized GaAs quantum dot, this condition is $`BIA/\sqrt{N}g\mu _B10\mathrm{mT}`$ (see Appendix A). For this estimate, we have used $`IA/g\mu _\mathrm{B}5\mathrm{T}`$, based on a sum over all three nuclear spin isotopes (all three hyperfine coupling constants) present in GaAsPaget et al. (1977) and $`N10^5`$ nuclei within each quantum dot. In this section, we also require $`BJ/g\mu _B`$, where $`J`$ is the Heisenberg exchange coupling between the two electron spins. For definiteness we take $`J>0`$, but all results are valid for either sign of $`J`$, with $`J`$ replaced by its absolute value. In the above limits, the electron Zeeman energy dominates all other energy scales and the relevant spin Hamiltonian becomes block-diagonal, with blocks labeled by the total spin projection along the magnetic field $`S^z`$ (see Appendix B). In the subspace of $`S^z=0`$ we write the projected two-electron spin Hamiltonian in the subspace of singlet and $`S^z=0`$ triplet states $`(|S,|T_0)`$ to zeroth order in the inverse Zeeman splitting $`1/ϵ_z`$ as $`H_0=\frac{J}{2}𝐒𝐒+\delta h^z\delta S^z`$, where $`𝐒=𝐒_1+𝐒_2`$ is the total spin operator in the double dot and $`\delta 𝐒=𝐒_1𝐒_2`$ is the spin difference operator. In terms of the vector of Pauli matrices $`𝝉=(\tau ^x,\tau ^y,\tau ^z)`$:$`|S|\tau ^z=1,|T_0|\tau ^z=+1`$ $`H_0`$ can be rewritten as:
$$H_0=\frac{J}{2}\left(1+\tau ^z\right)+\delta h^z\tau ^x.$$
(1)
Diagonalizing this two-dimensional Hamiltonian gives eigenvalues and eigenvectors
$`E_n^\pm `$ $`=`$ $`{\displaystyle \frac{J}{2}}\pm {\displaystyle \frac{1}{2}}\sqrt{J^2+4\left(\delta h_n^z\right)^2},`$ (2)
$`|E_n^\pm `$ $`=`$ $`{\displaystyle \frac{\delta h_n^z|S+E_n^\pm |T_0}{\sqrt{\left(E_n^\pm \right)^2+\left(\delta h_n^z\right)^2}}}|n,`$ (3)
where $`|n`$ is an eigenstate of the operator $`\delta h^z`$ with eigenvalue $`\delta h_n^z`$. Since the eigenstates $`|E_n^\pm `$ are simultaneous eigenstates of the operator $`\delta h^z`$, we note that there will be no dynamics induced in the nuclear system under the Hamiltonian $`H_0`$. In other words, the nuclear system remains static under the influence of $`H_0`$ alone, and there is consequently no back action on the electron spin due to nuclear dynamics.
We fix the electron system in the singlet state $`|S`$ at time $`t=0`$:
$$|\psi (t=0)=|S|\psi _I;|\psi _I=\underset{n}{}a_n|n,$$
(4)
where $`a_n`$ is an arbitrary set of (normalized) coefficients $`(_n\left|a_n\right|^2=1)`$. The initial nuclear spin state $`|\psi _I`$ is, in general, not an eigenstate $`|n`$. The probability to find the electron spins in the state $`|T_0`$ at $`t>0`$ is then given by the correlation function (setting $`\mathrm{}=1`$):
$$C_{T_0}(t)=\underset{n}{}\rho _I(n)\left|n|T_0\left|e^{iH_0t}\right|S|n\right|^2,$$
(5)
where $`\rho _I(n)=|a_n|^2`$ gives the diagonal matrix elements of the nuclear-spin density operator, which describes a pure (not mixed) state of the nuclear system: $`\rho _I=|\psi _I\psi _I|=_n\rho _I(n)|nn\left|+_{nn^{}}a_n^{}a_n^{}\right|n^{}n|`$. $`C_{T_0}(t)`$ is the sum of a time-independent piece $`\overline{C_n}`$ and an interference term $`C_{T_0}^{\mathrm{int}}(t)`$:
$`C_{T_0}(t)`$ $`=`$ $`\overline{C_n}+C_{T_0}^{\mathrm{int}}(t),`$ (6)
$`C_n`$ $`=`$ $`{\displaystyle \frac{2\left(\delta h_n^z\right)^2}{J^2+4\left(\delta h_n^z\right)^2}},`$ (7)
$`C_{T_0}^{\mathrm{int}}(t)`$ $`=`$ $`\overline{C_n\mathrm{cos}\left(\left[E_n^+E_n^{}\right]t\right)}.`$ (8)
Here, the overbar is defined by $`\overline{f(n)}=_n\rho _I(n)f(n)`$. Note that $`C_n`$ depends only on the exchange and Overhauser field inhomogeneity $`\delta h_n^z`$ through the ratio $`\delta h_n^z/J`$.
For a large number of nuclear spins $`N1`$ in a superposition of $`\delta h^z`$-eigenstates $`|n`$, we assume that $`\rho _I(n)`$ describes a continuous Gaussian distribution of $`\delta h_n^z`$ values, with mean $`\overline{\delta h_n^z}=0`$ (for the case $`\overline{\delta h_n^z}0`$, see Sec. II.1) and variance $`\sigma _0^2=\overline{\left(\delta h_n^z\overline{\delta h_n^z}\right)^2}=\overline{\left(\delta h_n^z\right)^2}`$ (i.e. $`\sigma _0=\delta h^z_{\mathrm{rms}}`$). The approach to a Gaussian distribution in the limit of large $`N`$ for a sufficiently randomized nuclear system is guaranteed by the central limit theorem.Coish and Loss (2004) The assumption of a continuous distribution of $`\delta h_n^z`$ precludes any possibility of recurrence in the correlator we calculate.<sup>2</sup><sup>2</sup>endnote: 2We recall that a superposition $`f(t)`$ of oscillating functions with different periods leads to quasiperiodic behavior, i.e., after the so-called Poincaré recurrence time $`t_p`$, the function $`f(t)`$ will return back arbitrarily close to its initial value (see, e.g., Ref. Fick and Sauermann, 1990). A lower-bound for the Poincaré recurrence time in this system is given by the inverse mean level spacing for the fully-polarized problemKhaetskii et al. (2002): $`t_\mathrm{p}N^2/A`$. In a GaAs double quantum dot containing $`N10^5`$ nuclear spins, this estimate gives $`t_{\mathrm{rec}}0.1\mathrm{s}`$. Moreover, by performing the continuum limit, we restrict ourselves to the free-induction signal (without spin-echo). In fact, we remark that *all* decay in the correlator given by (8) can be recovered with a suitable $`\pi `$-pulse, defined by the unitary operation $`U_\pi |E_n^\pm =|E_n^{}`$. This statement follows directly from the sequence
$$e^{iJt}|E_n^\pm =U_\pi e^{iH_0t}U_\pi e^{iH_0t}|E_n^\pm .$$
(9)
Thus, under the above sequence of echoes and free induction, all eigenstates are recovered up to a common phase factor. Only higher-order corrections to the effective Hamiltonian $`H_0`$ may induce completely irreversible decay. This irreversible decay could be due, for example, to the variation in hyperfine coupling constants, leading to decay on a timescale $`tN/A`$, as in the case of a single electron spin in Refs. Khaetskii et al., 2002; Coish and Loss, 2004. Another source of decay is orbital dephasing (see Sec. IV).
We perform the continuum limit for the average of an arbitrary function $`f(n)`$ according to the prescription
$`{\displaystyle \underset{n}{}}\rho _I(n)f(n)`$ $``$ $`{\displaystyle 𝑑xP_{\sigma ;\overline{x}}(x)f(n(x))},`$ (10)
$`P_{\sigma ;\overline{x}}(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\sigma }}\mathrm{exp}\left({\displaystyle \frac{\left(x\overline{x}\right)^2}{2\sigma ^2}}\right),`$ (11)
with $`\overline{x}=0`$, $`\sigma ^2=\overline{x^2}`$, and here we take $`x=\delta h_n^z,\sigma =\sigma _0`$. Using
$$C_n=C(\delta h_n^z)=C(x)=\frac{2x^2}{J^2+4x^2},$$
(12)
we evaluate $`C_{T_0}^{\mathrm{int}}(t)=\mathrm{Re}\left[\stackrel{~}{C}_{T_0}^{\mathrm{int}}(t)\right]`$, where the complex interference term is given by the integral
$`\stackrel{~}{C}_{T_0}^{\mathrm{int}}(t)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xC(x)P_{\sigma _0;0}(x)e^{it\sqrt{J^2+4x^2}}.`$ (13)
In general, the interference term given by Eq. (13) will decay to zero after the singlet-triplet decoherence time. We note that the interference term decays *even* for a *purely static* nuclear spin configuration with no ensemble averaging performed over initial conditions, as is the case for an isolated electron spin.Schliemann et al. (2002); Zurek et al. (2003); Coish and Loss (2004) The total $`z`$-component of the nuclear spins will be essentially static in any experiment performed over a timescale less than the nuclear spin diffusion time (the diffusion time is several seconds for nuclei surrounding donors in GaAsPaget (1982)). We stress that the relevant timescale in the present case is the spin diffusion time, and not the dipolar correlation time, since nonsecular corrections to the dipole-dipole interaction are strongly suppressed by the nuclear Zeeman energy in an applied magnetic field of a few GaussSlichter (1980) (as assumed here). Without preparation of the initial nuclear state or implementation of a spin-echo technique, this decoherence process therefore cannot be eliminated with fast measurement, and in general cannot be modeled by a classical nuclear field moving due to slow internal dynamics; a classical nuclear field that does not move cannot induce decay.
At times longer than the singlet-triplet decoherence time the interference term vanishes, leaving $`C_{T_0}(\mathrm{})=\overline{C_n}`$, which depends only on the ratio $`\delta h_n^z/J`$, and could therefore be used to trace-out the slow adiabatic dynamics $`\delta h_n^z(t)`$ of the nuclear spins, or to measure the exchange coupling $`J`$ when the size of the hyperfine field fluctuations is known. We evaluate $`C_{T_0}(\mathrm{})`$ from
$$C_{T_0}(\mathrm{})=\overline{C_n}=_{\mathrm{}}^{\mathrm{}}𝑑xC(x)P_{\sigma _0;0}(x).$$
(14)
In two limiting cases, we find the saturation value is given by (see Appendix C)
$$C_{T_0}(\mathrm{})\{\begin{array}{c}\frac{1}{2}\sqrt{\frac{\pi }{2}}\frac{J}{4\sigma _0},\sigma _0J,\\ 2\left(\frac{\sigma _0}{J}\right)^2,\sigma _0J.\end{array}$$
(15)
We recover the semiclassical high-magnetic-field limitSchulten and Wolynes (1978) $`(C_{T_0}(\mathrm{})=1/2)`$ *only* when the exchange $`J`$ is much smaller than $`\sigma _0`$. Furthermore, due to the average over $`\delta h_n^z`$ eigenstates, the approach to the semiclassical value of $`\frac{1}{2}`$ is a slowly-varying (linear) function of the ratio $`J/\sigma _0`$, in spite of the fact that $`C_n\left(J/\delta h_n^z\right)^2`$ as $`J0`$. In Figure 1 we plot the correlator saturation value $`C_{T_0}(\mathrm{})`$ as a function of the ratio $`\delta h^z_{\mathrm{rms}}/J`$ for a nuclear spin system described by a fixed eigenstate of $`\delta h^z`$ (i.e. $`\rho _I=|nn|`$), and for a nuclear spin system that describes a Gaussian distribution of $`\delta h^z`$ eigenstates with variance $`\sigma _0^2=\overline{\left(\delta h_n^z\right)^2}=\delta h^z_{\mathrm{rms}}^2`$. We also show the asymptotic expression for $`\sigma _0J`$, as given in Eq. (15).
Now we turn to the interference term $`C_{T_0}^{\mathrm{int}}(t)`$ given by Eq. (13), which can be evaluated explicitly in several interesting limits. First, in the limiting case of vanishing exchange $`(J=0)`$, we have $`C(x)=\frac{1}{2}`$ from (12). Direct integration of Eq. (13) then gives
$$C_{T_0}^{\mathrm{int}}(t)=\frac{1}{2}\mathrm{exp}\left(\frac{t^2}{2t_0^2}\right),t_0=\frac{1}{2\sigma _0},J=0.$$
(16)
For zero exchange interaction, the correlator decays purely as a Gaussian, with decoherence time $`t_0=\frac{1}{2\sigma _0}\frac{\sqrt{N}}{IA}`$ for a typical asymmetric double quantum dot (see Appendix A). However, for arbitrary nonzero exchange interaction $`J0`$, we find the asymptotic form of the correlator at long times is given by (see Appendix C):
$`C_{T_0}^{\mathrm{int}}\left(t\right)`$ $``$ $`{\displaystyle \frac{\mathrm{cos}\left(Jt+\frac{3\pi }{4}\right)}{4\sigma _0\sqrt{J}t^{3/2}}},`$ (18)
$`t\mathrm{max}({\displaystyle \frac{1}{J}},{\displaystyle \frac{1}{2\sigma _0}},{\displaystyle \frac{J}{4\sigma _0^2}}).`$
Thus, for arbitrarily small exchange interaction $`J`$, the asymptotic decay law of the correlator is modified from the Gaussian behavior of Eq. (16) to a (much slower) power law ($`1/t^{3/2}`$). We also note that the long-time correlator has a universal phase shift of $`\frac{3\pi }{4}`$, which is independent of any microscopic parameters. Our calculation therefore provides an example of interesting non-Markovian decay in an experimentally accessible system. Furthermore, the slow-down of the asymptotic decay suggests that the exchange interaction can be used to modify the *form* of decay, in addition to the decoherence time, through a narrowing of the distribution of eigenstates (see the discussion following Eq. (20) below). We have evaluated the full correlator $`C_{T_0}(t)`$ by numerical integration of Eq. (13) and plotted the results in Figure 2 along with the analytical asymptotic forms from (18).
We now investigate the relevant singlet-triplet correlator $`C_{T_0}(t)`$ in the limit of large exchange $`J`$. In this case, we have $`x\sigma _0J`$ for the typical $`x`$ contributing to the integral in Eq. (13). Thus, we can expand the prefactor $`C(x)`$ and frequency term in the integrand:
$`C(x)`$ $``$ $`2{\displaystyle \frac{x^2}{J^2}},`$ (19)
$`\sqrt{J^2+4x^2}`$ $``$ $`J+2{\displaystyle \frac{x^2}{J}}.`$ (20)
From Eq. (20) it is evident that the range of frequencies that contribute to the correlator is suppressed by $`\sigma _0/J`$ (increasing the exchange narrows the distribution of eigenenergies that can contribute to decay). This narrowing of the linewidth will increase the decoherence time. Moreover, the leading-order $`x^2`$-dependence in (20) collaborates with the Gaussian distribution of $`\delta h^z`$ eigenstates to induce a power-law decay. With the approximations in Eqs. (19) and (20), we find an expression for the correlator that is valid for all times in the limit of large exchange $`J`$ by direct evaluation of the integral in Eq. (13):
$`C_{T_0}^{\mathrm{int}}(t)`$ $`=`$ $`2\left({\displaystyle \frac{\sigma _0}{J}}\right)^2{\displaystyle \frac{\mathrm{cos}\left(Jt+\frac{3}{2}\mathrm{arctan}\left(\frac{t}{t_0^{}}\right)\right)}{\left(1+\left(\frac{t}{t_0^{}}\right)^2\right)^{3/4}}},`$ (22)
$`t_0^{}={\displaystyle \frac{J}{4\sigma _0^2}},J\sigma _0.`$
There is a new timescale ($`t_0^{}=J/4\sigma _0^2`$) that appears for large $`J`$ due to dynamical narrowing; increasing the exchange $`J`$ results in rapid precession of the pseudospin $`𝝉`$ about the $`z`$-axis, which makes transverse fluctuations along $`\tau ^x`$ due to $`\delta h^z`$ progressively unimportant. Explicitly, we have $`t_0^{}JN/4A^2\sqrt{N}/A`$ for $`J\sigma _0A/\sqrt{N}`$.
Eq. (22) provides a potentially useful means of extracting the relevant microscopic parameters from an experiment. $`J`$ and $`\sigma _0`$ can be determined independent of each other exclusively from a measurement of the oscillation frequency and phase shift of $`C_{T_0}^{\mathrm{int}}(t)`$. In particular, any loss of oscillation amplitude (visibility) due to systematic error in the experiment can be ignored for the purposes of finding $`\sigma _0`$ and $`J`$. The loss in visibility can then be quantified by comparison with the amplitude expected from Eq. (22). We illustrate the two types of decay that occur for large and small $`J`$ in Figure 3.
### II.1 Inhomogeneous polarization, $`\overline{\delta h_n^z}0`$
It is possible that a nonequilibrium inhomogeneous average polarization could be generated in the nuclear spin system, in which case $`\overline{\delta h_n^z}0`$. Pumping of nuclear spin polarization occurs naturally, for example, at donor impurities in GaAs during electron spin resonance (ESR), resulting in a shift of the ESR resonance condition.Seck et al. (1997) It is therefore important to investigate the effects of a nonzero average Overhauser field inhomogeneity on the decay law and timescale of the singlet-triplet correlator. In this subsection we generalize our previous results for the case $`\overline{\delta h_n^z}0`$.
We set the mean Overhauser field inhomogeneity to $`\overline{\delta h_n^z}=x_0`$, in which case the complex singlet-triplet interference term is given by
$`\stackrel{~}{C}_{T_0}^{\mathrm{int}}(t)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xC(x)P_{\sigma _0;x_0}(x)e^{it\sqrt{J^2+4x^2}}.`$ (23)
When the mean value of the Overhauser field inhomogeneity $`x_0`$ is much larger than the fluctuations $`\sigma _0`$ ($`x_0\sigma _0`$), we approximate $`C(x)C(x_0)`$ and expand the frequency term $`\sqrt{J^2+4x^2}=\omega _0+\frac{4x_0}{\omega _0}(xx_0)+\frac{2J^2}{\omega _0^3}(xx_0)^2+\mathrm{}`$, where $`\omega _0=\sqrt{J^2+4x_0^2}`$. We retain only linear order in $`xx_0`$ for the frequency term, which is strictly valid for times $`t(J^2+4x_0^2)^{3/2}/2J^2\sigma _0^2`$. This time estimate is found by replacing $`\left(xx_0\right)^2\sigma _0^2`$ in the quadratic term and demanding that the quadratic term multiplied by time be much less than one. In this limit, the correlator and range of validity are then
$`C_{T_0}^{\mathrm{int}}(t)`$ $`=`$ $`{\displaystyle \frac{2x_0^2}{\omega _0^2}}e^{\frac{1}{2}\left(\frac{t}{t_0^{\prime \prime }}\right)^2}\mathrm{cos}\left(\omega _0t\right),`$ (26)
$`t_0^{\prime \prime }={\displaystyle \frac{\omega _0}{4x_0\sigma _0}},\omega _0=\sqrt{J^2+4x_0^2},`$
$`x_0\sigma _0,t{\displaystyle \frac{\left(J^2+4x_0^2\right)^{3/2}}{2J^2\sigma _0^2}}.`$
This expression is valid for any value of the exchange $`J`$, up to the timescale indicated.
In contrast with the previous result for $`x_0=0`$, from Eq. (26) we find that the long-time saturation value of the correlator deviates from the semiclassical result ($`C_{T_0}(\mathrm{})=C_{T_0}^{\mathrm{int}}(0)=1/2`$) by an amount that is quadratic in the exchange $`J`$ for $`Jx_0`$:
$$C_{T_0}(\mathrm{})=C_{T_0}^{\mathrm{int}}(0)\{\begin{array}{c}\frac{1}{2}\frac{1}{8}\left(\frac{J}{x_0}\right)^2,Jx_0,\\ 2\left(\frac{x_0}{J}\right)^2,Jx_0.\end{array},x_0\sigma _0.$$
(27)
In the limit of large exchange, $`J\mathrm{max}(\sigma _0,x_0)`$, we can once again apply the approximations given in Eqs. (19) and (20). Using these approximations in Eq. (23) and integrating then gives
$`\stackrel{~}{C}_{T_0}^{\mathrm{int}}(t)`$ $`=`$ $`2\left({\displaystyle \frac{\sigma _0}{J}}\right)^2\xi ^3(t)\left(1+\left({\displaystyle \frac{x_0}{\sigma _0}}\right)^2\xi ^2(t)\right)\mathrm{exp}\left\{iJt{\displaystyle \frac{x_0^2}{2\sigma _0^2}}\left(1\xi ^2(t)\right)\right\},`$ (29)
$`\xi (t)=\left(1i{\displaystyle \frac{t}{t_0^{}}}\right)^{1/2},t_0^{}={\displaystyle \frac{J}{4\sigma _0^2}},J\mathrm{max}(x_0,\sigma _0),t{\displaystyle \frac{J^3}{2\mathrm{max}(x_0^4,\sigma _0^4)}}.`$
We have found the limit on the time range of validity in Eq. (29) using the same estimate that was used for Eqs. (26-26). At short times, $`tt_0^{}=J/4\sigma _0^2`$, we expand $`\xi ^2(t)1+i\frac{t}{t_0^{}}\left(\frac{t}{t_0^{}}\right)^2`$ and find that this function decays initially as a Gaussian with timescale $`t_0^{\prime \prime }J/4x_0\sigma _0`$:
$`C_{T_0}^{\mathrm{int}}(t)`$ $``$ $`2{\displaystyle \frac{\sigma _0^2+x_0^2}{J^2}}e^{\frac{1}{2}\left(\frac{t}{t_0^{\prime \prime }}\right)^2}\mathrm{cos}\left(\omega _0^{}t\right),`$ (32)
$`t_0^{\prime \prime }{\displaystyle \frac{J}{4x_0\sigma _0}},\omega _0^{}=J+{\displaystyle \frac{2x_0^2}{J}},`$
$`tt_0^{}={\displaystyle \frac{J}{4\sigma _0^2}},J\mathrm{max}(x_0,\sigma _0).`$
This agrees with the result in Eq. (26) when $`Jx_0\sigma _0`$.
For sufficiently large exchange $`J`$, the expression given by Eq. (29) is valid for times longer than the previous expression, given by Eq. (26). We perform an asymptotic expansion of Eq. (29) for long times using $`\xi (tt_0^{})e^{i\pi /4}\sqrt{t_0^{}/t}`$. This gives
$`C_{T_0}^{\mathrm{int}}(t)`$ $``$ $`{\displaystyle \frac{e^{x_0^2/2\sigma _0^2}\mathrm{cos}(Jt+\frac{3\pi }{4})}{4\sigma _0\sqrt{J}t^{3/2}}},`$ (34)
$`tt_0^{}={\displaystyle \frac{J}{4\sigma _0^2}},J\mathrm{max}(x_0,\sigma _0).`$
As in the case of $`x_0=0`$, the long-time asymptotics of Eq. (29) once again give a power law $`1/t^{3/2}`$, although the amplitude of the long-time decay is exponentially suppressed in the ratio $`x_0^2/\sigma _0^2`$. When $`x_0=0`$, Eq. (34) recovers the previous result, given in Eq. (18).
### II.2 Reducing decoherence
The results of this section suggest a general strategy for increasing the amplitude of coherent oscillations between the singlet $`|S`$ and triplet $`|T_0`$ states, and for weakening the form of decay. To avoid a rapid Gaussian decay with a timescale $`t_0^{\prime \prime }=J/4x_0\sigma _0`$, the mean Overhauser field inhomogeneity should be made smaller than the fluctuations ($`\overline{\delta h_n^z}=x_0\sigma _0`$) and the exchange $`J`$ should be made larger than $`x_0`$ and $`\sigma _0`$ ($`J\mathrm{max}(x_0,\sigma _0)`$). Explicitly, the ideal condition for slow and weak (power-law) decay can be written as
$$J\sigma _0x_0.$$
(35)
The condition in Eq. (35) can be achieved equally well by increasing the exchange coupling $`J`$ for fixed hyperfine fluctuations $`\sigma _0`$ or by reducing the fluctuations $`\sigma _0`$ through state squeezing or by making the double-dot confining potential more symmetric (see Appendix A).
## III Dynamics in the subspace of $`|S`$ and $`|T_+`$
We now consider the case when the Zeeman energy of the $`S^z=1`$ triplet state approximately compensates the exchange ($`\left|\mathrm{\Delta }\right|J`$, where $`\mathrm{\Delta }=ϵ_z+J`$). In addition, we assume the exchange is much larger than the nuclear field energy scales $`J\mathrm{max}\{\delta 𝐡_{\mathrm{rms}},𝐡_{\mathrm{rms}}\}`$. Under these conditions, we consider the dynamics in a subspace formed by the singlet $`|S|\tau ^z=1`$ and the $`S^z=1`$ triplet state $`|T_+|\tau ^z=+1`$, governed by the Hamiltonian (to zeroth order in $`1/J`$, see Appendix B):
$$H_+=\frac{1}{2}(\mathrm{\Delta }+h^z)(1+\tau ^z)\frac{1}{\sqrt{2}}(\delta h^{}\tau ^++\mathrm{H}.\mathrm{c}.).$$
(36)
Here, $`\delta h^\pm =\delta h^x\pm i\delta h^y`$ and $`\tau ^\pm =\frac{1}{2}\left(\tau ^x\pm i\tau ^y\right)`$. The $`|T_+`$ probability at time $`t>0`$ is
$$C_{T_+}(t)=\underset{n,n^{}}{}\rho _I(n)\left|n^{}|T_+\left|e^{iH_+t}\right|S|n\right|^2.$$
(37)
This case is essentially different from the previous one, since the eigenstates of $`H_+`$ are no longer simply product states of electron and nuclear spin, implying a back-action of the electron on the nuclear system. Nevertheless, when $`h^z+\mathrm{\Delta }_{\mathrm{rms}}\delta 𝐡^\pm _{\mathrm{rms}}`$, we can evaluate the correlator in standard time-dependent perturbation theory to leading order in the term $`V=\frac{1}{\sqrt{2}}\left(\tau ^+\delta h^{}+\tau ^{}\delta h^+\right)`$. Neglecting corrections of order $`h_n^z/\mathrm{\Delta }1`$, this gives
$$C_{T_+}^{(2)}(t)\overline{\frac{\alpha _n^2}{\mathrm{\Delta }^2}\left(1\mathrm{cos}\left(\left[\left[h^z\right]_n+\mathrm{\Delta }\right]t\right)\right)},$$
(38)
where $`\alpha _n=_n^{}\left|n^{}\left|\delta h^{}\right|n\right|^2`$, and $`|n`$ is now an eigenstate of the operator $`h^z`$ with eigenvalue $`\left[h^z\right]_n`$. To estimate the size of $`\alpha _n`$, we assume identical completely decoupled dots and nuclear polarization $`p1`$, which gives $`\alpha _n^2\frac{1}{2}I(I+1)_kA_k^2`$, where $`A_k`$ is the hyperfine coupling constant to the nuclear spin at lattice site $`k`$ (with total nuclear spin $`I`$) and the sum $`_k`$ runs over all lattice sites in one of the dots. We estimate the typical size of $`\alpha _n`$ with the replacements $`A_k\frac{A}{N},_kN`$, which gives $`\alpha _n\alpha /\sqrt{2}=\sqrt{\frac{I(I+1)}{2N}}A`$, where $`N`$ characterizes the number of nuclear spins within the dot envelope wavefunction. If we assume the nuclear spin state is described by a continuous Gaussian distribution of $`h^z`$ eigenstates with mean $`\overline{h_n^z}=0`$ and variance $`\sigma _+^2`$, we find
$$C_{T_+}^{(2)}(t)\frac{1}{2}\left(\frac{\alpha }{\mathrm{\Delta }}\right)^2\left(1e^{t^2/2t_+^2}\mathrm{cos}\left(\mathrm{\Delta }t\right)\right),t_+=\frac{1}{2\sigma _+}.$$
(39)
Thus, if we ignore any possibility for recurrence, the distribution of $`h^z`$ eigenstates will lead to Gaussian decay of the two-electron spin state, as is the case for a single electron.Schliemann et al. (2002); Coish and Loss (2004) However, as in the case of a single electron, this decay can be reduced or eliminated altogether by narrowing the distribution of $`h^z`$ eigenstates $`|n`$ through measurement (squeezing the nuclear spin state).Coish and Loss (2004) We show these two cases (with and without squeezing of the nuclear state) in Figure 4.
## IV Singlet-triplet decoherence due to orbital dephasing
To this point we have neglected dephasing of the singlet $`|S`$ and triplet $`|T_j`$ $`(j=0,+)`$ states due to coupling in the orbital sector. The effective Hamiltonian description ignores the different character of the orbital states for singlet and triplet, and so it is tempting to assume that orbital dephasing is unimportant where the effective Hamiltonian is valid. However, the singlet and triplet do have different orbital states which can, in general, couple differently to the environment through the charge degree of freedom, and therefore acquire different phases. Examples of such environmental influences are charge fluctuators or measurement devices, such as quantum point contacts used for charge readout.Engel et al. (2004); Elzerman et al. (2004) Here we briefly step away from the effective Hamiltonians derived in Appendix B to give a physical picture of the effects of orbital dephasing in terms of the true double-dot wavefunctions. We then return to the effective Hamiltonian picture in order to give a more general estimate of the effects of orbital dephasing on singlet-triplet decoherence for a two-electron double dot.
We consider a double quantum dot containing a fixed (quantized) number of electrons $`N`$. Within the far-field approximation, the double-dot charge distribution couples to the environment first through a monopole, and then a dipole term. Since the charge on the double dot is quantized, the monopole term gives an equal contribution for both the singlet and triplet wavefunctions. The leading interaction that can distinguish singlet from triplet is the electric dipole term:
$$V_{\mathrm{orb}}(t)𝐩_N𝐄(t).$$
(40)
Here, $`𝐩_N`$ is the electric dipole moment operator for the charge distribution in a double dot containing $`N`$ electrons and $`𝐄(t)`$ is a fluctuating electric field due to the surrounding environment, which we model by a Gaussian random process. For a double quantum dot with well-localized single-particle eigenstates we denote the charge states by $`|(n,m)`$, indicating that the double-dot has $`n`$ electrons in dot $`1`$ and $`m`$ electrons in dot $`2`$, where $`n+m=N`$. If the double dot contains only a single electron $`(N=1)`$, the environment can distinguish the two localized states through the difference in the dipole moment operator, which has the size $`\left|\mathrm{\Delta }𝐩_1\right|=\left|(1,0)\left|𝐩_1\right|(1,0)(0,1)\left|𝐩_1\right|(0,1)\right|2\left|e\right|a`$, where $`e`$ is the electron charge and $`2a`$ is the inter-dot spacing. When $`N=2`$, for highly-localized states, only the states with double-occupancy ($`|(0,2)`$ and $`|(2,0)`$) contribute to the dipole moment. If the typical hyperfine energy scale is much smaller than the detuning from resonance $`\delta `$ of the $`|(1,1)`$ and $`|(0,2)`$ states ($`\mathrm{max}(\delta 𝐡_{\mathrm{rms}},𝐡_{\mathrm{rms}})\delta `$), only the $`|(1,1)`$ singlet state (not the triplets) will mix with the doubly-occupied states, so the singlet and triplet states will be energetically distinguishable through $`\left|\mathrm{\Delta }𝐩_2\right|=\left|S\left|𝐩_2\right|S\right|2\left|e\right|a\left|P_{(0,2)}P_{(2,0)}\right|2\left|e\right|aD`$, where $`P_{(0,2)}`$ $`\left(P_{(2,0)}\right)`$ is the probability to find the singlet $`|S`$ in the $`|(0,2)`$ $`(|(2,0))`$ state and $`D=P_{(0,2)}+P_{(2,0)}`$ is the double occupancy. In this discussion, we assume that the exchange is much larger than the hyperfine energy scales, $`J\mathrm{max}(𝐡_{\mathrm{rms}},\delta 𝐡_{\mathrm{rms}})`$, so that the singlet and triplet states are good approximates for the true two-electron eigenstates.
For weak coupling to the environment, and assuming the environment correlation time is much less than the orbital dephasing time $`t_\varphi ^{(N)}`$, we can apply standard techniques to determine the dephasing time for a two-level system described by the Bloch equations.Blum (1981) We find that the fluctuations in $`𝐄(t)`$ lead to exponential dephasing with the rate $`1/t_\varphi ^{(N)}=\frac{1}{4}\left|\mathrm{\Delta }𝐩_N\right|^2_{\mathrm{}}^{\mathrm{}}𝑑tE(t)E(0)`$, where the scalar $`E(t)`$ is the component of $`𝐄(t)`$ along $`\mathrm{\Delta }𝐩_N`$ and we assume $`lim_t\mathrm{}\frac{1}{t}_0^t𝑑t^{}E(t^{})=0`$. Assuming equivalent environments for the single-particle and two-particle cases, the ratio of the single-particle to two-particle dephasing times is then
$$\frac{t_\varphi ^{(1)}}{t_\varphi ^{(2)}}=\left|\frac{\mathrm{\Delta }𝐩_2}{\mathrm{\Delta }𝐩_1}\right|^2D^2.$$
(41)
The single-electron orbital dephasing rate has been measured to be $`t_\varphi ^{(1)}1\mathrm{ns}`$Hayashi et al. (2003) and $`t_\varphi ^{(1)}400\mathrm{ps}`$Petta et al. (2004b) in different gated double quantum dots. If the hyperfine interaction (which becomes important on the timescale $`t5\mathrm{ns}`$) is to provide the major source of decoherence in these two-electron structures, we therefore require $`t_\varphi ^{(2)}t_\varphi ^{(1)}`$. This condition can be achieved by ensuring a small double occupancy $`D1`$ of the singlet state. When the inter-dot tunnel coupling $`t_{12}`$ is much less than the detuning from resonance $`\delta `$ ($`t_{12}\delta U+U^{}`$, with on-site and nearest-neighbor charging energies $`U`$ and $`U^{}`$, respectively – see Appendix B) we find the double-occupancy of $`|S`$ in perturbation theory is
$$D2\left(\frac{t_{12}}{\delta }\right)^21.$$
(42)
Even in this regime, orbital dephasing may become the limiting timescale for singlet-triplet decoherence after the removal of hyperfine-induced decoherence by spin echo. A detailed analysis of the double-occupancy and its relation to the concurrence (an entanglement measure) for a symmetric double dot can be found in Refs. Golovach and Loss, 2003, 2004.
With this physical picture in mind, we can generalize the above results to the case when the electrons experience fluctuations due to any time-dependent classical fields. In particular, if the separation in single-particle energy eigenstates for $`N=1`$ is $`ϵ+\delta ϵ(t)`$, where $`\delta ϵ(t)`$ fluctuates randomly with amplitude $`\delta ϵ`$, and similarly, if for $`N=2`$ the singlet and triplet levels are separated by an exchange $`J+\delta J(t)`$, where $`\delta J(t)`$ has amplitude $`\delta J`$, we find
$$\frac{t_\varphi ^{(1)}}{t_\varphi ^{(2)}}=\left|\frac{\delta J}{\delta ϵ}\right|^2.$$
(43)
From this expression we conclude that the optimal operating point of the double dot is where the slope of $`J`$ vs. $`ϵ`$ vanishes, i.e., $`\delta J/\delta ϵ=0`$. At this optimal point, $`t_\varphi ^{(2)}\mathrm{}`$, within the approximations we have made. Eq. (43) is valid for weak coupling to the environment (i.e. $`\delta JJ`$ and $`\delta ϵϵ`$), and when the environment correlation time is small compared to the dephasing times. If, for example, we take $`J2t_{12}^2/\delta `$ for $`U+U^{}\delta t_{12}`$ from Eq. (58) and if $`\delta ϵ`$ corresponds to fluctuations in the single-particle charging energy difference ($`ϵ(V_{g1}V_{g2})\delta `$ from Eq. (56)), we find $`t_\varphi ^{(1)}/t_\varphi ^{(2)}4t_{12}^4/\delta ^4`$, in agreement with Eqs. (41) and (42). In particular, the hyperfine-dominated singlet-triplet decoherence becomes visible when $`t_\varphi ^{(2)}t_0^{},t_0^{\prime \prime }t_0,t_+`$. This regime is achievable by choosing $`\delta t_{12}`$, but still $`J2t_{12}^2/\delta \sigma _0`$, since $`t_\varphi ^{(2)}`$ is a much stronger function of $`\delta `$ than $`t_0^{},t_0^{\prime \prime }`$. That is, the two-particle dephasing time scales like $`t_\varphi ^{(2)}\delta ^4`$, but the typical hyperfine-induced decay times scale like $`t_0^{},t_0^{\prime \prime }J1/\delta `$. On the other hand, when $`t_{12}\delta `$, we have $`\left|\delta J/\delta ϵ\right|O(1)`$, which gives $`t_\varphi ^{(2)}t_\varphi ^{(1)}`$, and thus a very short singlet-triplet decoherence time ($`1\mathrm{ns}`$), which is dominated by orbital dephasing.
## V Conclusions
We have shown that a fully quantum mechanical solution is possible for the dynamics of a two-electron system interacting with an environment of nuclear spins under an applied magnetic field. Our solution shows that the singlet-triplet correlators $`C_{T_0}(t)`$ and $`C_{T_+}(t)`$ will decay due to the quantum distribution of the nuclear spin system, even for a nuclear system that is static. We have found that the asymptotic behavior of $`C_{T_0}(t)`$ undergoes a transition from Gaussian to power-law ($`1/t^{3/2}`$) when the Heisenberg exchange coupling $`J`$ becomes nonzero, and acquires a universal phase shift of $`3\pi /4`$. The oscillation frequency and phase shift as a function of time can be used to determine the exchange and Overhauser field fluctuations. We have also investigated the effects of an inhomogeneous polarization on $`C_{T_0}(t)`$, and have suggested a general strategy for reducing decoherence in this system. Finally, we have discussed orbital dephasing and its effect on singlet-triplet decoherence.
###### Acknowledgements.
We thank G. Burkard, J. C. Egues, J. A. Folk, V. N. Golovach, D. Klauser, F. H. L. Koppens, J. Lehmann, C. M. Marcus, and J. R. Petta for useful discussions. We acknowledge financial support from the Swiss NSF, the NCCR nanoscience, EU RTN Spintronics, EU RTN QuEMolNa, EU NoE MAGMANet, DARPA, ARO, ONR, and NSERC of Canada.
## Appendix A Estimating the Overhauser field
In this appendix we estimate the size of the Overhauser field inhomogeneity for a typical double quantum dot, and show that this quantity depends, in a sensitive way, on the form of the orbital wavefunctions.
As in the main text, we take the average Overhauser field and the Overhauser field inhomogeneity to be $`𝐡=\frac{1}{2}\left(𝐡_1+𝐡_2\right)`$ and $`\delta 𝐡=\frac{1}{2}\left(𝐡_1𝐡_2\right)`$ respectively, where $`𝐡_l=Av_k\left|\psi _0^l(𝐫_k)\right|^2𝐈_k`$, and $`\psi _0^l(𝐫)`$ is orbital eigenstate $`l`$ in the double quantum dot. In the presence of tunneling, the eigenstates of a symmetric double quantum dot will be well-describedBurkard et al. (1999); Golovach and Loss (2004) by the symmetric and antisymmetric linear combination of dot-localized states $`\varphi _l(𝐫),l=1,2`$: $`\psi _0^{1,2}(\mathrm{r})=\frac{1}{\sqrt{2}}\left(\varphi _1(\mathrm{r})\pm \varphi _2(\mathrm{r})\right)`$. In this case, we find
$$\delta 𝐡_{\mathrm{rms}}=Av\underset{k}{}\mathrm{Re}\left[\varphi _1^{}(𝐫_k)\varphi _2(𝐫_k)\right]𝐈_k_{\mathrm{rms}}.$$
(44)
We take $`\frac{1}{N}_k𝐈_k_{\mathrm{rms}}\sqrt{I(I+1)/N}`$ to be the r.m.s. value for a system of $`N`$ nuclear spins with uniform polarization $`p1`$. Changing the sum to an integral according to $`v_kd^3r`$ then gives
$$\delta 𝐡_{\mathrm{rms}}\gamma \sqrt{\frac{I(I+1)}{N}}A=\gamma \alpha ,$$
(45)
where $`\gamma =d^3r\mathrm{Re}\left[\varphi _1^{}(𝐫)\varphi _2(𝐫)\right]`$ is the overlap of the localized orbital dot states and we have introduced the energy scale $`\alpha =\sqrt{I(I+1)}A/\sqrt{N}`$. The result in Eq. (45) suggests that the Overhauser field inhomogeneity can be drastically reduced in a symmetric double quantum dot simply by separating the two dots, reducing the wavefunction overlap. If, however, the double dot is sufficiently asymmetric, the correct orbital eigenstates will be well-described by localized states $`\psi _0^l(\mathrm{r})=\varphi _l(\mathrm{r}),l=1,2`$, (with overlap $`\gamma 1`$), in which case we find
$$\delta 𝐡_{\mathrm{rms}}\sqrt{\frac{I(I+1)}{N}}A=\alpha .$$
(46)
Thus, great care should be taken in determining $`\delta 𝐡_{\mathrm{rms}}`$ based on microscopic parameters. In particular, for a symmetric double quantum dot, the overlap $`\gamma `$ must also be known to determine $`\delta 𝐡_{\mathrm{rms}}`$ based on $`N`$.
In contrast, for the total Overhauser operator $`𝐡`$, in both of the above cases ($`\psi _0^{1,2}(\mathrm{r})=\frac{1}{\sqrt{2}}\left(\varphi _1(\mathrm{r})\pm \varphi _2(\mathrm{r})\right)`$ or $`\psi _0^l(\mathrm{r})=\varphi _l(\mathrm{r}),l=1,2`$), we find
$$𝐡_{\mathrm{rms}}\sqrt{\frac{I(I+1)}{N}}A=\alpha .$$
(47)
## Appendix B Effective Hamiltonians for two-electron states in a double quantum dot
In this appendix we derive effective Hamiltonians for a two-electron system interacting with nuclear spins in a double quantum dot via the contact hyperfine interaction.
We begin from the two-electron Hamiltonian in second-quantized form,
$$H=H_{SP}+H_C+H_T+H_Z+H_{\mathrm{hf}},$$
(48)
where $`H_{SP}`$ describes the single-particle charging energy, $`H_C`$ models the Coulomb interaction between electrons in the double dot, $`H_T`$ describes tunneling between dot orbital states, $`H_Z`$ gives the electron Zeeman energy (we neglect the nuclear Zeeman energy, which is smaller by the ratio of nuclear to Bohr magneton: $`\mu _N/\mu _B10^3`$) and $`H_{\mathrm{hf}}`$ describes the Fermi contact hyperfine interaction between electrons on the double dot and nuclei in the surrounding lattice. Explicitly, these terms are given by
$`H_{SP}`$ $`=`$ $`{\displaystyle \underset{l\sigma }{}}V_{gl}n_{l\sigma };n_{l\sigma }=d_{l\sigma }^{}d_{l\sigma },`$ (49)
$`H_C`$ $`=`$ $`U{\displaystyle \underset{l}{}}n_ln_l+U^{}(n_1+n_1)(n_2+n_2),`$ (50)
$`H_T`$ $`=`$ $`t_{12}{\displaystyle \underset{\sigma }{}}\left(d_{1\sigma }^{}d_{2\sigma }+d_{2\sigma }^{}d_{1\sigma }\right),`$ (51)
$`H_Z`$ $`=`$ $`{\displaystyle \frac{ϵ_z}{2}}{\displaystyle \underset{l}{}}(n_ln_l),`$ (52)
$`H_{\mathrm{hf}}`$ $`=`$ $`{\displaystyle \underset{l}{}}𝐒_l𝐡_l;𝐒_l={\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma \sigma ^{}}{}}d_{l\sigma }^{}𝝈_{\sigma \sigma ^{}}d_{l\sigma ^{}}.`$ (53)
Here, $`d_{l\sigma }^{}`$ creates an electron with spin $`\sigma `$ in orbital state $`l`$ $`(l=1,2)`$, $`V_{gl}`$ is the single-particle charging energy for orbital state $`l`$, $`U`$ is the two-particle charging energy for two electrons in the same orbital state, and $`U^{}`$ is the two-particle charging energy when there is one electron in each orbital. When the orbital eigenstates are localized states in quantum dot $`l=1,2`$, $`V_{gl}`$ is supplied by the back-gate voltage on dot $`l`$ and $`U`$ $`(U^{})`$ is the on-site (nearest-neighbor) charging energy. $`t_{12}`$ is the hopping matrix element between the two orbital states, $`ϵ_z`$ is the electron Zeeman splitting, $`𝐡_l`$ is the nuclear field (Overhauser operator) for an electron in orbital $`l`$, and $`𝝈_{\sigma \sigma ^{}}`$ gives the matrix elements of the vector of Pauli matrices $`𝝈=(\sigma _x,\sigma _y,\sigma _z)`$. In the subspace of two electrons occupying two orbital states, the spectrum of $`H_{SP}+H_C`$ consists of four degenerate “delocalized” states with one electron in each orbital, all with unperturbed energy $`E_{(1,1)}`$ (a singlet $`|S(1,1)`$ and three triplets: $`|T_j(1,1);j=\pm ,0`$), and two non-degenerate “localized” singlet states $`|S(2,0)`$ and $`|S(0,2)`$, with two electrons in orbital $`l=1`$ or $`l=2`$, having energy $`E_{(2,0)}`$ and $`E_{(0,2)}`$, respectively.
To derive an effective Hamiltonian $`H_{\mathrm{eff}}`$ from a given Hamiltonian $`H`$, which has a set of nearly degenerate levels $`\left\{|i\right\}`$, we use the standard procedureStoneham (1975),
$$H_{\mathrm{eff}}=PHP+PHQ\frac{1}{EQHQ}QHP,$$
(54)
where $`P=_i|ii|`$ is a projection operator onto the relevant subspace and $`Q=1P`$ is its complement.
We choose the arbitrary zero of energy such that $`E_{(1,1)}=V_{g1}+V_{g2}+U^{}=0`$ and introduce the detuning parameters
$`\delta _1`$ $`=`$ $`E_{(1,1)}E_{(2,0)}=2V_{g1}U=\delta UU^{},`$ (55)
$`\delta _2`$ $`=`$ $`E_{(1,1)}E_{(0,2)}=2V_{g2}U=\delta .`$ (56)
We then project onto the four-dimensional subspace formed by the delocalized singlet $`|S(1,1)`$ and three delocalized triplet states $`|T_j(1,1),j=\pm ,0`$. That is, we choose $`Q=|S(0,2)S(0,2)|+|S(2,0)S(2,0)|`$, $`P=1Q`$. When $`\delta _1,\delta _2t_{12}`$, we have *$`EE_{(1,1)}=0`$* in the denominator of Eq. (54). This gives an effective spin Hamiltonian in the subspace of one electron in each orbital state:
$`H_{\mathrm{eff}}`$ $`=`$ $`ϵ_z{\displaystyle \underset{l}{}}S_l^z+{\displaystyle \underset{l}{}}𝐡_l𝐒_lJ\left({\displaystyle \frac{1}{4}}𝐒_1𝐒_2\right),`$ (58)
$`J2t_{12}^2\left({\displaystyle \frac{1}{\delta }}{\displaystyle \frac{1}{\delta +U+U^{}}}\right).`$
This Hamiltonian is more conveniently rewritten in terms of the sum and difference vectors of the electron spin and Overhauser operators $`𝐒=𝐒_1+𝐒_2,\delta 𝐒=𝐒_1𝐒_2`$ and $`𝐡=\frac{1}{2}\left(𝐡_1+𝐡_2\right),\delta 𝐡=\frac{1}{2}\left(𝐡_1𝐡_2\right)`$:
$`H_{\mathrm{eff}}`$ $`=`$ $`ϵ_zS^z+𝐡𝐒+\delta 𝐡\delta 𝐒+{\displaystyle \frac{J}{2}}𝐒𝐒J.`$ (59)
Neglecting the constant term, in the basis of singlet and three triplet states, $`\left\{|S(1,1)=|S,|T_j(1,1)=|T_j,j=\pm ,0\right\}`$, the Hamiltonian matrix for $`H_{\mathrm{eff}}`$ takes the form
$$\left(\begin{array}{cccc}0& \delta h^+/\sqrt{2}& \delta h^z& \delta h^{}/\sqrt{2}\\ \delta h^{}/\sqrt{2}& J+ϵ_z+h^z& h^{}/\sqrt{2}& 0\\ \delta h^z& h^+/\sqrt{2}& J& h^{}/\sqrt{2}\\ \delta h^+/\sqrt{2}& 0& h^+/\sqrt{2}& Jϵ_zh^z\end{array}\right),$$
(60)
where $`\delta h^\pm =\delta h^x\pm i\delta h^y`$ and $`h^\pm =h^x\pm ih^y`$. We are interested in this Hamiltonian in two limiting cases, where it becomes block-diagonal in a two-dimensional subspace.
### B.1 Effective Hamiltonian in the $`|S|T_0`$ subspace
Projecting $`H`$ onto the two-dimensional subspace spanned by $`|T_0|\tau ^z=+1`$ and $`|S|\tau ^z=1`$, we find
$$H_0=N_0+\frac{1}{2}𝐯_0𝝉,$$
(61)
where $`𝝉=(\tau ^x,\tau ^y,\tau ^z)`$ is a vector of Pauli matrices. The leading and first subleading corrections to $`H_0`$ in powers of $`1/ϵ_z`$ are ($`H_0=H_0^{(0)}+H_0^{(1)}+\mathrm{}`$, $`H_0^{(i)}=N_0^{(i)}+𝐯_0^{(i)}`$):
$`N_0^{(0)}`$ $`=`$ $`{\displaystyle \frac{J}{2}},`$ (62)
$`v_0^{z(0)}`$ $`=`$ $`J,`$ (63)
$`v_0^{+(0)}`$ $`=`$ $`2\delta h^z,`$ (64)
$`N_0^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{4ϵ_z}}\left([h^{},h^+]+[\delta h^{},\delta h^+]\right),`$ (65)
$`v_0^{z(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2ϵ_z}}\left([h^{},h^+][\delta h^{},\delta h^+]\right),`$ (66)
$`v_0^{+(1)}`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_z}}\left(\delta h^+h^{}+\delta h^{}h^+\right).`$ (67)
Here, $`𝐍_X=(N_X^x,N_X^y,N_X^z)`$, $`𝐯_X=(v_X^x,v_X^y,v_X^z)`$, $`N_X^\pm =N_X^x\pm iN_X^y`$, and $`v_X^\pm =v_X^x\pm iv_X^y`$. For a typical unpolarized system, we estimate the size of all subleading corrections from their r.m.s. expectation values, taken with respect to an unpolarized nuclear state. This gives
$$H_0^{(1)}_{\mathrm{rms}}=O\left(\frac{\alpha ^2}{ϵ_z}\right),$$
(68)
where $`\alpha `$ is given by $`\alpha =\sqrt{I(I+1)}A/\sqrt{N}`$ (for a GaAs quantum dot containing $`N10^5`$ nuclear spins, $`1/\alpha 5\mathrm{ns}`$). We therefore expect dynamics calculated under $`H_0^{(0)}`$ to be valid up to timescales on the order of $`ϵ_z/\alpha ^21/\alpha `$, when $`ϵ_z\alpha `$.
### B.2 Effective Hamiltonian in the $`|S|T_+`$ subspace
When the Zeeman energy of the $`|T_+`$ triplet state approximately compensates the exchange, $`\mathrm{max}(𝐡_{\mathrm{rms}},\delta 𝐡_{\mathrm{rms}},\left|\mathrm{\Delta }\right|)J`$ (where $`\mathrm{\Delta }=ϵ_z+J`$), we find an effective Hamiltonian in the subspace $`|T_+|\tau ^z=+1,|S|\tau ^z=1`$:
$$H_+=N_++\frac{1}{2}𝐯_+𝝉,$$
(69)
where the leading and subleading corrections in powers of $`1/J`$ are
$`N_+^{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Delta }+h^z\right),`$ (70)
$`v_+^{z(0)}`$ $`=`$ $`\mathrm{\Delta }+h^z,`$ (71)
$`v_+^{+(0)}`$ $`=`$ $`\sqrt{2}\delta h^+,`$ (72)
$`N_+^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2J}}\left(\left(\delta h^z\right)^2+{\displaystyle \frac{1}{4}}\delta h^{}\delta h^++{\displaystyle \frac{1}{2}}h^{}h^+\right),`$ (73)
$`v_+^{z(1)}`$ $`=`$ $`{\displaystyle \frac{1}{J}}\left(\left(\delta h^z\right)^2+{\displaystyle \frac{1}{4}}\delta h^{}\delta h^+{\displaystyle \frac{1}{2}}h^{}h^+\right),`$ (74)
$`v_+^{+(1)}`$ $`=`$ $`\sqrt{2}{\displaystyle \frac{\delta h^zh^+}{J}}.`$ (75)
Once again, we estimate the influence of the subleading corrections from their r.m.s. value with respect to a nuclear spin state of polarization $`p1`$, giving
$$H_+^{(1)}_{\mathrm{rms}}=O\left(\frac{\alpha ^2}{J}\right).$$
(76)
We therefore expect the dynamics under $`H_+^{(0)}`$ to be valid up to time scales on the order of $`tJ/\alpha ^21/\alpha `$ for $`J\alpha `$.
## Appendix C Asymptotics
### C.1 $`C_{T_0}(\mathrm{})`$ for $`J2\sigma _0,J2\sigma _0`$
In the limit of $`J0`$, we perform an asymptotic expansion of the integral in Eq. (14) by separating the prefactor into a constant piece and an unnormalized Lorentzian of width $`J/2`$:
$$C(x)=\frac{1}{2}\left(1\frac{(J/2)^2}{\left(J/2\right)^2+x^2}\right).$$
(77)
The Gaussian average over the constant term gives $`1/2`$ and when $`J/2\sigma _0`$, the typical $`x`$ contributing to the Lorentzian part of Eq. (14) is $`xJ/2\sigma _0`$, so we approximate $`\mathrm{exp}(\frac{1}{2}x^2/\sigma _0^2)1`$ in the integrand of this term. Integrating the Lorentzian then gives the result in Eq. (15) for $`J2\sigma _0`$. In the opposite limit of $`J2\sigma _0`$, the Lorentzian is slowly-varying with respect to the Gaussian, and the prefactor can be expanded within the integrand $`C(x)2x^2/J^2`$. Performing the remaining Gaussian integral gives the result in Eq. (15) for $`J2\sigma _0`$.
### C.2 $`C_{T_0}^{\mathrm{int}}(t)`$ for $`t\mathrm{}`$
To evaluate the integral in Eq. (13) at long times when $`J0`$, we make the change of variables $`u=\sqrt{\lambda ^2+\left(x/\sigma _0\right)^2}\lambda `$, $`\lambda =J/2\sigma _0`$, $`\stackrel{~}{t}=2\sigma _0t`$, which gives
$`\stackrel{~}{C}_{T_0}^{\mathrm{int}}(\stackrel{~}{t}/2\sigma _0)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑u{\displaystyle \frac{\sqrt{u(u+2\lambda )}}{u+\lambda }}\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}\left(u^2+2u\lambda \right)+i(u+\lambda )\stackrel{~}{t}\right\},`$ (79)
$`\lambda =J/2\sigma _0,\stackrel{~}{t}=2\sigma _0t.`$
At long times, the major contributions to this integral come from a region near the lower limit, where $`u1/\stackrel{~}{t}`$. For $`\stackrel{~}{t}\mathrm{max}(1/\lambda ,1)`$ (i.e. $`t\mathrm{max}(1/J,1/2\sigma _0)`$), we approximate the integrand by its form for $`u\mathrm{max}(\lambda ,1)`$, retaining the exponential term as a cutoff. This gives
$$\stackrel{~}{C}_{T_0}^{\mathrm{int}}(\stackrel{~}{t}/2\sigma _0)\frac{e^{i\lambda \stackrel{~}{t}}}{\sqrt{\pi \lambda }}_0^{\mathrm{}}𝑑u\sqrt{u}e^{(\lambda i\stackrel{~}{t})u}=\frac{e^{i\lambda \stackrel{~}{t}}}{2\sqrt{\lambda }\left(\lambda i\stackrel{~}{t}\right)^{3/2}}.$$
(80)
When $`\stackrel{~}{t}\lambda `$ (i.e. $`tJ/4\sigma _0^2`$), we expand the denominator of the above expression, which gives the result in Eq. (18).
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# X-ray emission from NGC 1808: more than a complex starburstPartially based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA Member States and NASA.
## 1 Introduction
NGC 1808 ($`\alpha `$(J2000)=05:07:42.34; $`\delta `$(J2000)=–37:30:47.0 with an uncertainty of 1$`\stackrel{}{.}`$25, obtained from 2MASS<sup>1</sup><sup>1</sup>1 http://www.ipac.caltech.edu/2mass/), classified as an SABb Seyfert 2 galaxy, is located at a distance of 10.9 Mpc (H<sub>0</sub> = 75 km s <sup>-1</sup> Mpc<sup>-1</sup>, 1″= 53 pc). Images in different wavebands suggest a high star formation (SF) intensity in the central region with a diameter of $`20`$″ (i.e., $`1`$ kpc). There are several optical hot spots associated with H II-regions (Sérsic & Pastoriza 1965). The luminous and compact knots detected in radio and infra-red (IR) images (Saikia et al. 1990, Kotilainen et al. 1996) that do not coincide with the optical hot spots are probably supernova remnants (SNRs) or complexes of unresolved SNRs. Dust filaments were found perpendicular to the central disk plane and can be explained as outflowing material driven by supernovae (SNe; Heckman et al. 1990). A recent interaction of NGC 1808 with its companion NGC 1792 could explain both the intense star-formation activity and its peculiar morphology (Dahlem et al. 1990, Koribalski et al. 1993).
The nature of the nucleus is still unclear. It is classified as Seyfert 2 based on the optical emission lines detected by Véron-Cetty & Véron (1985). However, there is also evidence from optical observations of a hidden starburst (SB), based on polarization measurements by Scarrott et al. (1993). The nucleus is claimed to have a strong non-stellar component, because only 10% of the IR radiation observed by ISO could be interpreted as emission related to star formation (Siebenmorgen et al. 2001). On the other hand, soft X-ray (0.1–2.4 keV) data obtained with ROSAT (Dahlem et al. 1994 and Junkes et al. 1995) favor an interpretation in terms of SF-related emission, but do not discard other hypotheses. Awaki & Koyama (1993) interpret their Ginga X-ray data (1.5–37 keV) as an obscured active galactic nucleus (AGN). However, Awaki et al. (1996), using ASCA (2–10 keV) observations, point out that the hard X-ray spectrum could also be the result of starburst activity, while the long-term variability from the Ginga and ASCA observations again suggests a Seyfert nature of the nucleus.
With the advent of the XMM-Newton and Chandra satellites, the coexistence of starburst and AGN activity has been detected in X-rays in the nuclei of several nearby galaxies. NGC 4303 presents evidence of the composite SB/AGN nature of the nuclear region: Chandra and UV-HST observations show that a superstellar cluster coexists with a hard X-ray source, possibly an AGN, in the 3 central pc of the galaxy (Jiménez-Bailón et al. 2003, Colina et al. 2002). At larger scales, Persic et al. (2004a) claim that starburst and AGN coexist in NGC 4666 by detecting evidence of diffuse thermal emission, high-luminousity X-ray sources, and the presence of an AGN in the X-ray band. The detailed analysis of our XMM-Newton data shown here, combined with other observations, sheds new light on the nature of the nuclear activity in NGC 1808 and the associated outflow of gas.
The observations and data reduction are described in Sect. 2, our results are presented in Sect. 3, followed by our interpretation of the inner region and circumnuclear emission in Sect. 4.1 and of the more extended X-ray emission outside the nuclear area in Sect. 4.2. The conclusions of this work are given in Sect. 5.
## 2 Observations and data reduction
### 2.1 XMM-Newton data
Information on the XMM-Newton instrumentation is provided by Jansen et al. (2001; XMM-Newton mission), Strüder et al. (2001; EPIC-pn), Turner et al. (2001; EPIC-MOS), den Herder et al. (2001; RGS), and Mason et al. (2001; OM). Our XMM-Newton observation (Obs-Id 0110980801) was performed on April 6, 2002. The scheduled exposure time was 40 ks, while effective exposure times per instrument are listed below. The EPIC-pn exposure was in extended full frame mode and with the thin filter. Both RGS cameras were in spectroscopic mode and the two MOS cameras in full frame, with the thin filter. OM was used in full frame low resolution mode with the U, UVW1 (219 nm) and UVW2 (212 nm) broad band filters, plus the two (optical and UV) grisms.
The data were processed with the Science Analysis Subsystem, SAS, v.5.4.1 (Gabriel et al. 2004) and the most up-to-date calibration available in February 2003. The standard tasks to process data, ep/mchain and rgsproc, were run with default parameters. For the EPIC instruments, checks against pile-up were performed with the epatplot task. No sign of pile-up was found for any of EPIC-pn, MOS1, or MOS2 data. Additionally, background flaring time intervals, for which the signal-to-noise does not improve by including these events, were filtered following the method described in Piconcelli et al. (2004b). Only single and double events were considered in the analysis. For RGS the selection was for periods with less than 0.12 c/s in the background region of the CCD number 9. The effective exposure times after each filter was applied were 31.4 ks, 39.0 ks, 38.6 ks, and 37.4 ks for EPIC-pn, the two MOS cameras, and RGS, respectively. The OM broad band data were processed with the SAS omichain task, which performed all necessary corrections. Unfortunately, about half of the UVW1 filter image was lost due to ground-station problems, and the UV grism data was also lost due to telemetry drops.
### 2.2 Chandra data
The Chandra ACIS (AXAF CCD Imaging Camera) archival image, sequence number 700451, was extracted from the Chandra archive. The ACIS observation took place on December 19, 2002 with an exposure time of 43 ks. We used the science level 2 files, generated by the pipeline standard processing on December 19, 2002, as stored in the archive. They were processed with Ciao 3.0.2. According to the information available on the Chandra pages, further corrections should be applied (i.e. reduction of tap-ringing distortions). However, for the purpose of the present work, this was found to be unnecessary. No additional corrections were therefore applied.
### 2.3 Ground-based H$`\alpha `$ imagery
The H$`\alpha `$ data used here were obtained with the CTIO 4-m telescope on November 18, 1993. A $`2048\times 2048`$ chip with a $`0\stackrel{}{.}431`$ pixel scale was used, leading to a total field-of-view of $`14\stackrel{}{.}7\times 14\stackrel{}{.}7`$.
Three exposures of 20 min. each were taken with a redshifted H$`\alpha `$ filter ($`\lambda _0`$ = 657.8 nm; $`\mathrm{\Delta }\lambda `$ = 0.64 nm), with a total integration time of 1 hr. The continuum subtraction was performed using an $`R`$-band image with a total integration time of 8 min. The data are extremely sensitive to low surface brightness emission, but as they are uncalibrated, will be used only for qualitative studies.
## 3 Results
### 3.1 X-ray imaging
XMM-Newton EPIC X-ray images of NGC 1808 were produced by grouping the photons from the cleaned event lists into the sub-bands defined by the XMM-Newton Survey Science Centre (SSC; Watson et al. 2001) and using the merge task in SAS v.5.4.1, which takes exposure time corrections into account. These images are displayed in Fig. 1. The maximum emission is located at $`\alpha `$<sup>EPIC</sup>(J2000)=05:07:42.40; $`\delta `$<sup>EPIC</sup>(J2000)=–37:30:46.2 with an uncertainty of 1$`\stackrel{}{.}`$5–3″. The position of the 2MASS nucleus is indicated by a cross. An adaptive filter was used to smooth the images in areas of low signal-to-noise (S/N) ratio, while not smearing out the emission from point sources.
NGC 1808 shows extended emission in the soft X-ray images up to 4.5 keV (Fig. 1). Fig. 2 shows the comparison of the radial profiles of NGC 1808 and MCG–06–30–15 in two energy bands: 0.2–4.5 keV and 4.5–10 keV. The QSO MCG–06–30–15 can be considered as a point-like source in all the XMM-Newton energy bands (Kirsch et al. 2005). The comparison shows that the emission of NGC 1808 in the lowest energy band is extended while its radial profile in the 4.5–10 keV band does not differ from the QSO profile, taking the errors into account. Therefore, NGC 1808 can be considered as a point-like source above 4.5 keV. This points out that the origin of the X-ray emission is not due to the compact central source alone. The elongated X-ray emission of NGC 1808 follows the same orientation as the H$`\alpha `$ emission (see Sect. 3.4). The extended morphology is detected at energies below 4.5 keV, with a maximum in the 0.5–2.0 keV band (see Fig. 1). At these low energies, thermal emission processes tracing high-mass SF play an important role. Thus, our data suggest that the stellar activity in the nuclear region is significantly contributing to the observed soft X-ray emission. With increasing energy, the central emission maximum becomes more prominent and the non-nuclear sources fade away.
In order to resolve more details, we have compared our 0.2–12 keV EPIC image with the Chandra ACIS data. Both were smoothed with an adaptive filter and are displayed in Fig. 3. The Chandra ACIS image shows how the nucleus is surrounded by a population of discrete sources. The image (left panel of Fig. 3) shows the detailed structure in the central part of NGC 1808: at least four point-like sources are detected in the image, apart from the diffuse emission. Two of them are located at $``$4″ (212 pc) NW and SE from the maximum of the X-ray emission in the Chandra image and have also been detected in a Chandra HRC observation. Zezas et al. (2001) reported a luminosity for each of them of around $`10^{39}`$ erg s<sup>-1</sup>.
Although the nucleus is observed as a point-like source in the HRC image (see Zezas el al. 2001), the ACIS image reveals the presence of a double-peak structure (see Fig. 3). The hard (E$`>`$2 keV) and soft (E$`<`$1.5 keV) X-ray images of the two peaks presented in Fig. 4 reveal that the eastern source, S1, is significantly softer than the western one, S2. The locations of the two peaks are $`\alpha `$$`{}_{1}{}^{}{}_{}{}^{\mathrm{Chandra}}`$(J2000)=05:07:42.35; $`\delta `$$`{}_{1}{}^{}{}_{}{}^{\mathrm{Chandra}}`$(J2000)= –37:30:45.8 and $`\alpha `$$`{}_{2}{}^{}{}_{}{}^{\mathrm{Chandra}}`$(J2000)=05:07:42.19; $`\delta `$$`{}_{2}{}^{}{}_{}{}^{\mathrm{Chandra}}`$(J2000)= –37:30:45.8 with an uncertainty of 0$`\stackrel{}{.}`$6. The nuclear location from 2MASS is indicated by an arrow in the Chandra image, closer to the eastern source found by Chandra, S1, but compatible within the errors with both of them. The remaining two point-like sources are situated on the ring of radio emission as imaged by Saikia et al. (1990), close to the detected SNR and H II-regions.
### 3.2 X-ray spectral analysis
We performed an extensive X-ray spectral analysis of NGC 1808 using the EPIC and RGS instruments on board XMM-Newton, as well as, with Chandra’s ACIS camera.
The XMM-Newton EPIC data were used to perform a spectral analysis of NGC 1808 in the 0.35–9 keV energy band with moderate resolution (80 eV at 1 keV). In order to search for differences in the X-ray spectrum of the nucleus and its surroundings, we defined three circular regions and one annulus. The circles have radii of 16″ (850 pc), 35″ (1.9 kpc), and 2′ (6.4 kpc), while the annulus has external and internal radii of 20″ and 50″ (1.0–2.7 kpc), respectively (excluding one extra-nuclear source). Background regions were selected to be on the same CCD as the source but far enough to prevent source contamination. Figure 5 shows the spectra of each region. We simultaneously analysed the EPIC-pn and RGS spectra, and finally we also studied the Chandra spectra of the two nuclear point-like sources S1 and S2 (see Fig. 4), both located within the 16″ inner region extracted in the XMM-Newton analysis.
In order to apply the modified $`\chi ^2`$ minimization technique in the spectral analysis, all EPIC spectra were grouped such that each spectral bin contains at least 50 counts. The Chandra spectra were grouped with at least 20 counts in each bin. The spectra were analysed using XSPEC v.11.3.0 (Arnaud 1996).
The spectra of the 35″ and 2′ regions are very similar, both in shape and intensity (see Fig. ). Only the spectrum of the 16″ region, although in good agreement with the two previous ones above 2 keV, is clearly weaker in the soft band.
#### 3.2.1 EPIC spectrum of the inner region
The EPIC spectrum of the inner circular region of NGC 1808 is displayed in Fig. 6. The extraction region has a radius of 16″ (850 pc), centred on the maximum peak of the X-ray emission. MOS1 and MOS2 spectra and response matrices were combined by merging these to maximize the signal-to-noise level. Subsequently, the EPIC-pn and combined MOS spectra were fitted simultaneously. Figure 6 shows the 0.35–9 keV spectrum, the best fit model with a statistical quality of $`\chi _\nu ^2=0.99`$ for 174 degrees of freedom (dof), and the residuals.
Simple models as a power law, a Raymond-Smith (Raymond and Smith, 1976 ), mekal (Mewe, Lemen, & van den Oord 1986), and bremsstrahlung components provide unacceptable fits. The best fit model includes, apart from the fixed Galactic foreground absorption with a column density of $`N_H=3.23\times 10^{20}`$ cm<sup>-2</sup> (Dickey & Lockman, 1990), a power law with index $`\mathrm{\Gamma }=0.79_{0.06}^{+0.07}`$ and a mekal thermal component with a temperature of kT=0.580$`{}_{0.016}{}^{}{}_{}{}^{+0.019}`$ keV absorbed by a Hydrogen column of N<sub>H</sub>=2.0$`{}_{0.4}{}^{+0.3}\times 10^{21}`$cm<sup>-2</sup>. The value of $`\chi ^2`$ is unacceptable, but an intrinsic absorption to the power law component does not improve the fit; the inferred upper limit for the Hydrogen column density is $`8\times 10^{20}`$cm<sup>-2</sup>. The fit only improved, in terms of $`\chi ^2`$, when some abundances were allowed to differ from the solar values.The abundances of Ne, Mg, Si, and Fe, left free in the fits, are 1.6$`{}_{0.5}{}^{}{}_{}{}^{+1.0}`$, 2.2$`{}_{0.5}{}^{}{}_{}{}^{+1.2}`$, 2.1$`{}_{0.5}{}^{}{}_{}{}^{+1.0}`$, and 0.63$`{}_{0.11}{}^{}{}_{}{}^{+0.2}`$ times the solar value ($`Z_{}`$), respectively. Fits with other thermal models such as Raymond-Smith or bremsstrahlung emission lead to higher values of $`\chi ^2`$. No Fe K$`\alpha `$ line was significantly detected. An upper limit of 80 eV was calculated for the equivalent width of a narrow iron line. Table 1 shows the goodness of the fit together with the values and uncertainties for each free parameter.
#### 3.2.2 The Chandra spectra of sources S1 and S2
The ACIS image in Fig. 3 shows a double-peak nuclear structure surrounded by a few point-like sources inside the inner, r=16″, region defined for EPIC and analysed in the previous sub-section. We analysed the standard extracted spectra obtained with psextract of the two circular, r=1$`\stackrel{}{.}`$2, regions S1 and S2 (see Fig. 4). The background was extracted from a clean circular region with a radius of 10″ located $``$1′ North-East from the sources in the same CCD. For the two regions, models with a single component were rejected due to the high $`\chi ^2`$ value obtained.
The best fit model obtained for the observed spectrum of the hard source, S2, includes an absorbed mekal component and an absorbed power law. The values of the parameters and the goodness of the fit are presented in Table 2 and the observed data, the best fit model, and the residuals are shown in Fig. 7b. It is worth noting the high values of the absorbing column densities for both components, N$`{}_{\mathrm{H}}{}^{}=(23)\times 10^{22}`$ cm<sup>-2</sup>. The index of the power law found is $`\mathrm{\Gamma }=1.2\pm 0.3`$, and the temperature of the thermal component is kT=$`0.087_{0.006}^{+0.005}`$ keV. The thermal component in S2 is significantly colder than the one detected in the EPIC spectrum of the inner region, with kT=0.58$`\pm 0.02`$ keV. The fit does not improve by adding an extra thermal component that could account for the thermal emission observed in the inner region of XMM-Newton, kT=0.58 keV. Neither does it improve by varying the abundances from the solar value.
Figure. 4 shows that above 2 keV, source S1 is dimmer than source S2, while below 1.5 keV it is the opposite. Assuming that the low hard emission of S1 could be due to contamination of the harder source, S2, we tested a model which includes an absorbed power law component with photon index and N<sub>H</sub> fixed to the values obtained for the best fit model of source S2. In order to model the soft range, we considered an absorbed mekal thermal model with free temperature. The resulting fit is unacceptable, $`\chi _\nu ^2=1.9`$. The addition of an extra thermal component with kT=0.087 keV equal to the one detected in source S2 does not significantly improve the fit. The best fit model is obtained with an unabsorbed power law with $`\mathrm{\Gamma }=0.99_{0.14}^{+0.18}`$ and a moderately absorbed, N$`{}_{H}{}^{}=5.2_{0.7}^{+1.0}\times 10^{21}`$ cm<sup>-2</sup>, mekal component with a temperature of kT=$`0.58_{0.09}^{+0.07}`$ keV and the abundances of Ne, Mg, Si, and Fe fixed to the values obtained in the fit of the EPIC spectrum of the inner region. If an absorption component above the Galactic value is applied to the power law, the value of N<sub>H</sub> for that additional component is compatible with zero. The observed data, the best fit model, and the residuals are shown in Fig. 7a.
#### 3.2.3 Simultaneous analysis of EPIC-pn and RGS spectra
In order to study the nuclear emission of NGC 1808 in greater detail, we took advantage of the high spectral resolution of the RGS data (0.35–2.5 keV).
Figure 8 compares the combined RGS1 and RGS2 spectrum of NGC 1808 with that of M 82, a proto-typical starburst galaxy. Both spectra were generated with the SAS task rgsfluxer applying the standard extraction techniques. The RGS spectrum of NGC 1808 shows no continuum emission above the noise level, but it does exhibit a number of emission lines. Figure 8 shows that these have wavelengths and relative intensity ratios very similar to the strongest lines identified in M 82 (Read and Stevens, 2002).
We determined the wavelength of the emission lines observed with the RGS by fitting Gaussian profiles to them. The redshift-corrected locations of the fitted Gaussians are less than 0.02 Å from the theoretical value of the lines that we have identified. The fitted line widths are within the range of 0.02–0.04 Å, compatible with the instrumental resolution of the RGS. The detected lines are thus identified as: the Ly$`\alpha `$ emission line from O VIII , transitions of He-like Ne IX ions, and Fe L emission lines from Fe XVII and Fe XVIII . The weaker lines visible in the M 82 spectrum were not detected in NGC 1808, most probably due to the lower S/N. The similarity in the line ratios suggests that there is a common origin of the soft X-ray emission in both galaxies, i.e. thermal emission from a hot, extended gas component, as shown for M 82 (Read and Stevens, 2002). This result confirms the detection of the starburst in NGC 1808 in our data, which dominates the soft X-ray emission.
The RGS spectra of NGC 1808 and M 82 in Fig. 8 show only a qualitative comparison. In particular, in the reduction process of the data, calibration effects affecting the line profile were not considered, and therefore the wings of the lines are not properly determined. The most important effect of this issue is that the flux of the lines could be underestimated. The ratio of fluxes of the detected lines can be used to derive physical properties of the gas where the lines originate. In order to perform an accurate study of the properties of the emission lines, we analysed the RGS spectrum, taking the response of the detector into account and therefore avoiding the uncertainties explained before. Using XSPEC v.11.3, we simultaneously analysed the RGS and the 16″ region EPIC-pn spectra. Figure 9 shows that the soft X-ray emission of NGC 1808 can be explained solely by the contribution of emission lines. Consequently, we tried to model the pn-RGS spectra with an absorbed (N$`{}_{\mathrm{H}}{}^{}=1_1^{+17}\times 10^{19}`$ cm<sup>-2</sup>) power law ($`\mathrm{\Gamma }=1.1_{0.5}^{+0.6}`$ and A$`{}_{\mathrm{pl}}{}^{}=9.4_{0.7}^{+0.9}\times 10^5`$ photons/keV/cm<sup>2</sup>/s at 1 keV) which accounts for the hard band emission and a number of lines accounting for the soft band emission. The emission lines were modeled with Gaussian profiles with widths fixed to their instrumental value, both in pn and RGS. The wavelengths were also fixed to the theoretical values of the most likely identifications in order to reduce the degrees of freedom. The best fit model consists of 13 emission lines and a power law with a value for $`\chi ^2=307`$ for 209 dof. The energy, intensity, and identification of each line are shown in Table 3 and the fit in Fig. 9.
The line ratios provide a diagnostic of the conditions of the gas. In particular, the ratios between lines of iron in different degrees of ionization degrees allow for the best temperature diagnostics. Table 3 gives all the ratios relative to the strongest isolated line, i.e. FeXVII (15.01 Å). We therefore use the observed FeXVII line and compare it with FeXVIII and others with higher ionization states to infer the plasma temperature. We use the web version of the ATOMDB v.1.3, http://cxc.harvard.edu/atomdb and conclude that a collisionally ionized plasma with kT$``$0.4 keV ($`5\times 10^6`$ K) has a line spectrum compatible with the observed spectrum of NGC 1808. This is not surprising, given our previous results using EPIC data alone.
#### 3.2.4 EPIC spectrum of the non-nuclear extended region
In order to investigate the non-nuclear extended emission of NGC 1808, the EPIC-pn spectrum of an annular region was extracted and analysed. The inner and the outer radii, 16$`\stackrel{}{.}`$5 (875 pc) and 35″ (1.9 kpc), were chosen to minimize the contribution of the so-called inner region, i.e. r=16″, and to exclude the CCD edges.
For the spectral fit, we have assumed that the emission in this region is contaminated by the emission from the inner region. Therefore, the fitting model includes a power law with an index fixed to the value obtained in the fit of the inner region spectrum, i.e. $`\mathrm{\Gamma }=0.79`$, and an absorbed mekal component with N<sub>H</sub> fixed to 2$`\times 10^{21}`$cm<sup>-2</sup> and kT<sub>I</sub> fixed to 0.58 keV. The normalisations of both components were left free. An excess with respect to the model is observed at low energies. This feature can be well fitted with an extra mekal component with a temperature kT$`{}_{\mathrm{II}}{}^{}0.1`$ keV. The abundances of Ne, Mg, Si, and Fe have been left free in the fits and their values are compatible within the errors with the abundances obtained for the inner region. The values of all parameters and the goodness of the fit can be found in Table 1. Figure 10 shows the data, the best fit model, and the residuals.
Comparison of the inner and annular regions shows that the annular spectrum includes a component not detected in the inner region, the low temperature gas, which can therefore be associated with a non-nuclear, extended emission component.
The results of the analysis confirm our expectation that the spectrum of the annular region is contaminated by some emission from the inner region. This is especially true for the power law, which is fainter by a factor of $``$6.5 compared to that in the inner region, while the hotter thermal gas contribution may be partially due to either extended emission or discrete sources in the annulus, as its normalisation is only a factor of 3.4 weaker than in the inner region. Indeed, the encircled energy function for EPIC-pn predicts that about 20% of the total energy emitted by the unresolved source will be imaged between r=16$`\stackrel{}{.}`$5 and r=35″ from the centre of the point spread function.
#### 3.2.5 Integrated EPIC spectrum of the total galaxy
A circular region with $`r=62\stackrel{}{.}5`$ (3.3 kpc) was analysed. The extracted EPIC spectrum contains the emission of the whole galaxy and the region coincides with the extraction radius used in the ROSAT analysis (Dahlem et al. 1994; Junkes et al. 1995).
Analogously to the annular region, the components of the best fit model ( $`\chi _\nu ^2=1.00`$ for 222 dof) are a power law and two absorbed thermal components. The values of the various parameters and the goodness of the fit can be found in Table 1. Figure 11 shows the data, the best fit model, and the residuals. In the lower panel of Fig. 11, we also show the best fit model and its various additive components in $`Ef(E)`$. The values of the parameters suggest that the emission has the same origin as for the annular region; i.e. the power law and higher temperature component are associated with the nuclear emission, and the soft thermal emission is associated with the non-nuclear extended emission.
### 3.3 X-ray fluxes and luminosities
Various absorbed X-ray fluxes and unabsorbed luminosities, measured in the three regions of the XMM-Newton analysis and the two regions of the Chandra analysis, are collated in Table 4.
For all cases, the thermal components dominate in the soft X-ray band, while the power law dominates in the 2–10 keV band. The innermost part of the galaxy (r=16″=850 pc) is responsible for emission in the hard band, while in the soft energy band it contributes only half of the total luminosity of the galaxy ($`r=62\stackrel{}{.}5`$=3.3 kpc).
In order to study the long term variability in NGC 1808, we compared the fluxes observed with XMM-Newton in April 2002 with the ones observed with Chandra eight months later. To allow comparison, the spectra of the nuclear, the annular, and the inner regions were extracted from the Chandra ACIS observation. These regions are equivalent to the ones analysed with XMM-Newton. The spectra of the inner regions were fitted using the best fit model obtained from the XMM-Newton analysis, i.e. a power law plus an absorbed mekal component. The index of the power law, hydrogen column density, and the temperature and abundances of the mekal component of the best fit model ($`\chi ^2=170`$ for 137 dof) are compatible with the values obtained with XMM-Newton. Therefore, no spectral variability was detected between both observations. The spectrum of the annular region was also fitted using the best fit model obtained with XMM-Newton. The values of the relevant parameters of the best fit model ($`\chi ^2=45`$ for 36 dof) are again compatible with the ones derived from the XMM-Newton observation.
In the inner region, the comparison of the XMM-Newton and Chandra values reveals no sign of variability in the soft energy range Chandra flux is only 5% lower than the XMM-Newton flux, well within the statistical errors and calibration uncertainties of both instruments. In the hard band, a decrease of 16$`\%`$ in flux is observed at a 2.5$`\sigma `$ level (see Table 4).
The Chandra analysis of the two discrete sources detected within the r=16″ region reveals that source S2 is responsible for $``$50% of the hard emission from the inner Chandra region, while the combined emission of sources S1 and S2 only explains $``$15% of the observed soft emission in this inner region. It is worth noting that the upper limit of the intrinsic, i.e. absorption-corrected soft X-ray luminosity of S2 obtained with Chandra, is more than three orders of magnitude higher than the luminosity of the whole galaxy. However, considering the uncertainties of this measurement, it is compatible with lower values. Moreover, taking the absorbed fluxes into account, no discrepancy is found, indicating that the intrinsic absorption components are responsible for the high luminosity measured. Although the statistical test shows that the model is acceptable, alternative scenarios cannot be ruled out.
In the annular region, the Chandra flux measurement in the soft band is significantly lower than those of XMM-Newton. However, Chandra is less efficient than XMM-Newton to detect weak extended emission and probably the discrepancy is due to the extended emission not determined by Chandra.
We also measured the fluxes and luminosities with Chandra within an r=62$`\stackrel{}{.}`$5 aperture to compare the XMM-Newton data with Chandra. The Chandra spectrum does not require the complexity needed to describe the XMM-Newton spectrum: only a combination of a power law and one mekal component is sufficient. The lowest temperature mekal, kT$`0.1`$ keV, found with XMM-Newton is not detected in the Chandra spectrum due to the small effective area of the instrument at such low energies. The values of the parameters of the two components included in the best-fit model ($`\chi ^2`$=220 for 200 dof) are compatible with the XMM-Newton results.
As expected, taking the results for the inner region into account, a mild variability in the hard band was observed between the XMM-Newton and Chandra observations, while no significant variation of the flux was measured in the soft energy range. This variation is dominated by the power law component of the inner region.
In addition to the newly derived luminosity values, Table 4 also lists the luminosity measured with ROSAT (Junkes et al. 1995) in this region. The values indicate an important decrease in the luminosity between the ROSAT, XMM-Newton, and Chandra observations. Comparing our results ($`L_{0.510keV}3\times 10^{40}`$ ergs<sup>-1</sup>) for the r=62$`\stackrel{}{.}`$5 region with the ASCA observations of February 1994 obtained by Awaki et al. 1996 ($`L_{0.510keV}=2\times 10^{40}`$ ergs<sup>-1</sup>), a small increase in the luminosity is observed. The comparison with the Ginga results (October 1990) indicates a decrease of the luminosity from $`L_{210keV}=5\times 10^{40}`$ ergs<sup>-1</sup> for Ginga to $`L_{210keV}=1.7\times 10^{40}`$ ergs<sup>-1</sup> for XMM-Newton. Nevertheless, existing uncertainties in the relative calibration of various X-ray detectors could explain the discrepancy of XMM-Newton and Chandra with the ASCA and Ginga values, but the luminosity measured by ROSAT is 5 times that observed with XMM-Newton and Chandra, too high to be explained by calibration uncertainties.
Short-term variability during the XMM-Newton observation was not detected.
### 3.4 OM optical/UV data and H$`\alpha `$ imagery
OM UV images of NGC 1808 with the U (350 nm), UVW1 (291 nm) and UVW2 (212 nm) filters were taken in parallel with the X-ray observations. Due to technical problems, part of the UVW1 image was lost. All images are displayed in Fig. 12. The U and UVW2 images show a double-peak structure in the centre of the galaxy. The locations of the peaks in UVW2 image are $`\alpha _1^{\mathrm{OM}}`$(J2000)=05:07:41.97; $`\delta `$$`{}_{1}{}^{}{}_{}{}^{\mathrm{OM}}`$(J2000)=–37:30:43.3 and $`\alpha `$$`{}_{2}{}^{}{}_{}{}^{\mathrm{OM}}`$(J2000)=05:07:42.62; $`\delta `$$`{}_{2}{}^{}{}_{}{}^{\mathrm{OM}}`$(J2000)=–37:30:48.1 with uncertainties of $``$2″. Figure 12.c shows the central region of the UVW2 image with the Chandra contours overlaid. The locations of the two point-like sources detected by Chandra (at $``$4″ from the nucleus and aligned in the NW and SE direction) coincide with the positions of the UV peaks observed in the OM image. However, there is no UV emission associated with any of the two bright X-ray sources at the very nucleus of the galaxy.
In Fig. 13 we have overlaid the OM 212 nm image of NGC 1808 on the EPIC 0.5–2.0 keV frame from Fig. 1 and on an extremely sensitive H$`\alpha `$ image. An RGB composition of the images in optical (H$`\alpha `$), UV (OM UVW2) and X-rays (EPIC) is presented in Fig. 14.
The correspondence of UV and H$`\alpha `$ emission maxima indicates that the OM data trace the massive stars ionizing the gas in H II regions in NGC 1808. The H$`\alpha `$ image also shows several knots located in an outer spiral arm. SF is traced in the circumnuclear starburst but also further out in the galaxy disk. Similarly, the soft X-ray emission observed in the same regions is expected to come from hot gas in or near these SF regions. Optical/UV emission also tracks the bar.
## 4 Discussion
ASCA and Ginga observations in the past missed the starburst’s soft X-ray emission, favouring an interpretation in terms of an AGN (Awaki et al. 1996, Awaki & Koyama 1993), as suggested first by Véron-Cetty and Véron (1985). Similarly, the IR emission observed by ISO is dominated by this AGN (Siebenmorgen et al. 2001).
ROSAT, on the other hand, detected the soft emission from the starburst, being less sensitive to the hard emission of an AGN (e.g. Junkes et al. 1995), thus suggesting the predominance of the starburst.
Based on the data presented above, the most likely interpretation is that NGC 1808 hosts both, a circumnuclear starburst and an unresolved nuclear X-ray source. Both in the imaging (Figs. 1,3, and 4) and spectral (Figs. 6, 7, and 9) domains, XMM-Newton and Chandra detect the presence of the starburst and the hard unresolved nuclear source. In the following we discuss the physical origin of the emission of both types of activity: the starburst and the unresolved X-ray source.
### 4.1 X-ray emission of the inner region
#### 4.1.1 Hard component
Spectral characteristics
The slope of the power law fitted to the 16″ region, $`\mathrm{\Gamma }=0.79_{0.06}^{+0.07}`$, is flatter in comparison to results obtained for radio-quiet AGN and LINERs. Works by Georgantopoulos et al. (2002) and Terashima et al. (2002) find values for the power law spectral indices ranging from 1.7 to 2.3 for these objects. However, there are also objects with lower spectral indices (e.g. PKS 2251$`+`$113, with $`\mathrm{\Gamma }=0.95\pm 0.24`$, Reeves & Turner, 2001). Analysis of the Chandra data reveals that this hard emission is dominated by one of the two detected nuclear sources, S2 in Fig. 4, which is harder. In addition, source S2 is absorbed by a column density of N$`{}_{H}{}^{}=3.1_{0.7}^{+0.8}\times 10^{22}`$ cm<sup>-2</sup>, compatible with values observed in Compton-thin Seyfert 2 galaxies. This fact suggests that in case one of the nuclear peaks observed with Chandra is associated with an AGN, the most likely possibility is that it is source S2. Interestingly, the best-fit model of the XMM-Newton 16″ region spectrum does not require absorption in excess of the Galactic value. However, this XMM-Newton region includes several emitting sources that could blur the presence of an absorbing component in the best fit model.
Luminosity
The total luminosity of the central region with a radius of 16″= 850 pc of NGC 1808 is L$`{}_{210\mathrm{keV}}{}^{}=(1.61\pm 0.06)\times 10^{40}`$ erg s<sup>-1</sup>. The luminosity in the same energy band associated with the harder peak with a radius of 1$`\stackrel{}{.}`$2=65 pc observed with Chandra, i.e. source S2, is L$`{}_{210\mathrm{keV}}{}^{}=0.92_{0.7}^{+0.08}\times 10^{40}`$ erg s<sup>-1</sup>, which is $`60\%`$ of the emission of the XMM-Newton inner region. These values are in the range of Low Luminosity AGN (LLAGN) and LINERS (Terashima et al. 2002, Georgantopoulos et al. 2002). Moreover, in both cases the luminosity associated with the power law component is in good agreement with the typical values of low luminosity Seyfert galaxies and LINERs. Terashima et al. (2002) found in a sample of 53 objects of these types observed by ASCA that the luminosity varies within the range $`L_{210keV}=5\times 10^{3940}`$ ergs<sup>-1</sup>. If only LINERS are considered, the mean luminosity obtained is $`L_{210keV}=7\times 10^{39}`$ ergs<sup>-1</sup> (Ho et al. 2001), which is lower than the one derived for NGC 1808. In the soft band, a sample of LINERS observed with ROSAT (Komossa et al. 1999) shows luminosities in the range of $`L_{0.242keV}=10^{3841}`$ ergs<sup>-1</sup>, in agreement with luminosity measured for the NGC 1808.
Variability
We studied the variability of NGC 1808. While no short-term variation of flux was detected during the XMM-Newton observation, medium- and long-term variability has been suggested in both soft and hard bands.
NGC 1808 was observed by XMM-Newton and Chandra with an interval of eight months. The measurements indicate the soft X-ray flux remained constant, taking the uncertainties into account. A decrease of 16% (at a 2.5$`\sigma `$ level) in the measured flux is observed in the hard X-ray band. In particular, the power law component of the innermost region spectrum of NGC 1808 is responsible for this variability.
The long- and medium-term flux variations measured over the years by various satellites are most probably caused by variability of the nuclear source. In the soft energy band, a significant decrease of a factor of five has been observed when the ROSAT and XMM-Newton (or Chandra) luminosities are compared. It should be noted that the ROSAT spectral fit is particularly uncertain because of the high absorbing column density, which leaves a very narrow effective bandpass for the X-ray emission, from about 0.5 to 2.4 keV.
In summary, mild medium-term variability of the hard band luminosity was measured between the XMM-Newton and Chandra observations. This variation can be attributed to an AGN or a ULX. Furthermore, the uncertainties introduced by comparing different satellites does not allow us to firmly conclude that other changes in flux are intrinsically due to the nucleus of NGC 1808.
Ultraluminous X-ray sources
The study of the hard X-rays emanating from the inner region indicates that the nuclear X-ray emission itself could be due to a LLAGN or a LINER. The observed hard band X-ray luminosity, L$`_{210keV}`$=$`(1.61\pm 0.06)\times 10^{40}`$ erg s<sup>-1</sup>, is very difficult to explain by standard stellar processes. Nevertheless, non-nuclear point-like sources with X-ray luminosities on the order of $`10^{3841}`$ erg s<sup>-1</sup> have been detected in nearby galaxies. An alternative origin to an AGN nucleus of NGC 1808 could be one of the so-called ultraluminous X-ray sources (ULX).
The two central sources S1 and S2 were studied in detail in Sect. 3.2.2. The nuclear location from 2MASS is compatible with either of them being the nucleus. The luminosity of source S2 of L$`{}_{210keV}{}^{}=9.2_7^{+0.8}\times 10^{39}`$ erg s<sup>-1</sup> is compatible with either a LLAGN or a ULX. Strickland et al. (2001) and Roberts et al. (2002) show that the analysis of Chandra ULX spectra favours simple power laws with indices in the range of 1.8–2.9, higher than the value obtained in our analysis of 0.79$`{}_{0.06}{}^{}{}_{}{}^{+0.07}`$. Interestingly, one of the compact radio sources reported by Saika et al. (1990) coincides with the X-ray point-like source S1.
Although more data analysis is necessary to reach any firm conclusion, this finding leaves open the possibility that the nucleus of NGC 1808 could be a ULX.
#### 4.1.2 Soft component
Emission distribution
The soft X-ray emission of the inner region of NGC 1808 is extended (see e.g. Fig. 1). The elongated emission is detected at energies up to 4.5 keV. Moreover, the Chandra images clearly show the simultaneous presence of a double emission peak and diffuse, or at least unresolved, X-ray emission around the two peaks. The extended emission is likely to have been created by star formation processes. Our XMM-Newton data permit a clear spectral distinction between the emission associated with the starburst and the additional nuclear sources in the r=16″ (850 pc) region, see Fig. 6. The starburst, associated with the mekal component in the spectral fit, dominates the soft energy band (below 1 keV), see Tables 1 and 4. The starburst also shows up in the form of radio continuum and IR knots (Saikia et al. 1990, Kotilainen et al. 1996), bright H$`\alpha `$ emission, and bright circumnuclear UV continuum from massive stars, as displayed in Fig. 13. In particular, the position of the 2MASS nucleus coincides perfectly with one of the radio compact sources detected by Saikia et al (1990), explained as SNR or complexes of unresolved SNR. The nucleus is prominent in X-rays and in the optical regime (see Figs. 1 and 12.a).
However, while the OM UVW2 image shows that the UV emission close to the position of the nucleus is very weak (see Fig. 12.c). These facts suggest that the nucleus is highly absorbed. The spectral analysis of source S2 shows a high value of the equivalent hydrogen column, N$`{}_{H}{}^{}10^{22}`$ cm<sup>-2</sup>, which can explain this absence of UV and soft X-ray emission associated with the hard X-ray peak.
Plasma temperature
The best fit to the EPIC data contains a thermal plasma component with a temperature of kT = $`0.580_{0.016}^{+0.019}`$ keV. This corresponds roughly to the “medium” component identified in the spectra of other starburst galaxies (Dahlem et al. 1998; see also Weaver et al. 2000 and Dahlem et al. 2000). It is consistent with $`0.5\pm 0.2`$ keV (Junkes et al. 1995), which was measured from ROSAT data, based only on about 600 photons. It is interesting to note that the surprisingly high Hydrogen column density, on the order of $`8\times 10^{21}`$ cm<sup>-2</sup>, in the best-fitting model to the ROSAT data is also corroborated by our results.
X-ray luminosity and SFR
The luminosity of the starburst component, which we identify with the thermal spectral components, is dominated by the emission in the soft energy band. The observed X-ray luminosity, $`L_{0.24keV}1.3\times 10^{40}`$ erg s<sup>-1</sup>, which represents $``$70% of the emission associated with the starburst component of the whole galaxy, leads to a far-infrared-to-X-ray luminosity ratio of log($`L_\mathrm{X}/L_{\mathrm{FIR}})`$ = –3.8 <sup>2</sup><sup>2</sup>2$`L_{\mathrm{FIR}}=2\pi D^2FIR=9.3\times 10^{43}`$ erg s<sup>-1</sup>(Dahlem et al. 1992), where $`FIR=1.26(2.58f_{60}+f_{100}`$), $`f_{60}`$ being the IRAS 60 $`\mu `$m flux density and $`f_{100}`$ the 100 $`\mu `$m flux density.. This value is consistent with those obtained by Heckman et al. (1990) for a sample of six starburst galaxies, which are in the range between –3.7 and –4.3, and also with the results in Mas-Hesse et al. (1995), who obtained a mean value of –3.33 for a sample of starforming galaxies.
Ranalli, Comastri & Setti (2003) deduced that the X-ray luminosity is a tracer of SFR, based on a sample of nearby galaxies. As was pointed out by Persic et al. (2004b), the luminosity in the 2–10 keV band is not a precise SFR indicator, due to the contamination by low-mass X-ray binaries (LMXB). These objects are bright in the X-ray band over long timescales, t$`{}_{X}{}^{}10^7`$ yr, and therefore do not trace the instantaneous SFR. For this reason, we use the soft X-ray luminosity as a tracer of the SRF, with the SFR-L$`_{0.52keV}`$ relation in Ranalli, Comastri & Setti (2003):
$$SFR_X(M_{}yr^1)=2.2\times 10^{40}L_{0.52keV}$$
(1)
The L$`_{0.52keV}`$ associated with the starburst in the r=16″(850 pc) region, as derived from the mekal component contribution, is on the order of L$`{}_{0.52keV}{}^{SB}1\times 10^{40}`$ erg s<sup>-1</sup> and the SFR<sub>X</sub>, therefore, is $`2.5`$ $`M_{}`$ yr<sup>-1</sup>.
Metallicities
We detect several key emission lines, which for the first time provide us with a good measure of the abundances of some elements (Sect. 3.2.1). The RGS spectrum shows a very weak continuum emission with several emission lines. The values obtained indicate slightly super-solar metallicities for Ne, Mg, and Si and solar or sub-solar abundance for Fe, which is consistent with the results by Dahlem et al. (1998) and Weaver at al (2000).
In this context the absence of a thermal Fe K$`\alpha `$ line at $``$6.5 keV energy is noteworthy. From our data in the annular region, we derive a 3-$`\sigma `$ upper limit on the equivalent width of 170 eV and on the flux of $`4\times 10^7`$ photons cm<sup>-2</sup>s<sup>-1</sup>. This result is compatible with any possible iron line emission entirely due to the mekal component. We have used ATOMDB v.1.3 to derive the expected flux of the emission lines for a gas with temperature and normalisation as obtained from the single mekal component fitted to the annular spectrum, i.e. kT=$`0.53_{0.03}^{+0.05}`$ keV and normalisation A$`{}_{\mathrm{mekal}}{}^{}=2.2\times 10^4`$. The total flux of the most prominent iron lines (FeXXI-6.505 keV and FeXXII-6.504 & 6.57 keV ) is $`2\times 10^4`$ photons cm<sup>-2</sup>s<sup>-1</sup> (with an emissivity of $`5\times 10^{22}`$ photons cm<sup>3</sup>s<sup>-1</sup>), consistent with the upper limit measured.
### 4.2 X-ray emission outside the central region
In Sect. 3.2, we show the results of the spectral analysis of the 875 pc–1.9 kpc annular region. The best fit model indicates that, apart from the contribution of the nuclear emission, the annular region includes an extra thermal component which can be explained by a mekal model with kT=$`0.10_{0.01}^{+0.02}`$ keV. The analysis of the total galaxy X-ray emission also indicates the presence of a soft thermal model with a similar temperature, kT=$`0.11_{0.02}^{+0.03}`$ keV.
Although weak, the contribution in the soft band of this component (L$`{}_{0.12.4keV}{}^{}=7.7\times 10^{39}`$ erg s<sup>-1</sup>) to the soft X-ray luminosity of the whole galaxy is significant, representing around 30% of the total soft X-ray emission of NGC 1808. In the hard band, the contribution of this low temperature thermal component is negligible. Using Eq. 1 (see Sect. 4.1.2 for further details), the observed X-ray luminosity associated with both mekal components of the annular region, L$`{}_{0.52keV}{}^{}4\times 10^{39}`$, can be translated into a SFR of $``$$`M_{}`$ yr<sup>-1</sup>.
Far-infrared emission is also a good tracer of the star formation; using the relationship by Kennicutt et al. (1998), we derive for NGC 1808:
$$SFR_{FIR}(M_{}yr^1)=\frac{1}{e}4.5\times 10^{44}L_{FIR}$$
(2)
where $`e`$ is the fraction of UV/optical flux emitted by the stars and absorbed by the dust and remitted in the IR, which is close to unity even for low values of reddening (Mas-Hesse & Kunth 1991). Assuming that the whole FIR emission is due to the starburst, we obtain SFR$`{}_{FIR}{}^{}`$$`M_{}`$ yr<sup>-1</sup>. This result is consistent with the one obtained through the X-ray emission, combining the values derived from the nuclear region (SFR$`{}_{X}{}^{}2.5`$ $`M_{}`$ yr<sup>-1</sup>, see Sect. 4.1.2) and from the annular region (SFR$`{}_{X}{}^{}1`$ $`M_{}`$ yr<sup>-1</sup>). It suggests that the contribution of the AGN to the FIR emission is small.
The second, soft mekal component detected in both the annular and the r=3.3 kpc region is most likely thermal emission associated with the halo. Such a thermal component (kT=0.2–0.3 keV) is found in all starburst galaxies in the sample by Dahlem et al. (1998) and Weaver et al. (2000) and also by us in other starburst galaxies observed with XMM-Newton (e.g. Dahlem et al. 2003 and Ehle & Dahlem 2004). The presence of both, the prominent vertical dust filaments in NGC 1808 and the gaseous outflow inferred by the kinematics of H I gas (Koribalski et al. 1993), suggested that this component would most likely be present in NGC 1808.
If the diffuse X-ray emission in the soft energy band is assumed to be due to thermal cooling of a hot gas in collisional ionization equilibrium<sup>3</sup><sup>3</sup>3See, however, the discussion in the validity of this assupntion in Breitschwerdt & Schmutzler (1999), it is possible to estimate the gas density $`n_\mathrm{e}`$, mass $`m_{\mathrm{gas}}`$, and cooling time $`\tau `$ of that plasma. To this end we make use of a formula given by Nulsen et al. (1984), i.e. $`L_\mathrm{x}(\mathrm{soft})=0.812\mathrm{\Lambda }(T)n_\mathrm{e}^2V\eta `$. The unknown filling factor $`\eta `$ allows for some clumpiness of the gas filling the emitting volume V, and was found in simulations (Avillez & Breitschwerdt, 2004 and references therein) to be in the range of 0.2–0.4, only modestly increasing with higher SN rate.
As explained in Sect. 3.2.2, the spectrum of the non-nuclear extended X-ray emission is complex and a best fit was obtained with a model consisting of a foreground absorbed power law plus (at least) two internally absorbed mekal components, corroborating the fact that the hot gas is a multi-temperature mixture. Nevertheless for the purpose of estimating hot gas parameters, we fitted this emission component with a single ’characteristic’ temperature that was found to be $`0.53_{0.03}^{+0.05}`$ keV. For such a gas temperature of the starburst component ($`6.2\times 10^6`$ K) Raymond et al. (1976) give a cooling coefficient $`\mathrm{\Lambda }(T)`$ of $`5\times 10^{23}`$ erg cm<sup>3</sup> s<sup>-1</sup>. For the unabsorbed luminosity of the hot gas component we make use of the total (Thermal I + II) $`L_{0.12.4keV}=7.7\times 10^{39}`$ erg s<sup>-1</sup> (see Table 4). Assuming for the emitting volume a torus geometry with an inner radius of 875 pc and outer radius of 1.9 kpc, the calculated gas parameters are presented in Table 5.
The mass of the detected hot gas is only $`0.5\%`$ of the total gas mass<sup>4</sup><sup>4</sup>4The total gas mass calculated using H<sub>I</sub> observations is $`>1.8\times 10^9`$$`M_{}`$(Dahlem et al. 2001).. The cooling time, $`\tau `$, is comparable to typical duty cycles of starbursts, which are estimated to be on the order of a few times $`10^7`$ yr (Rieke et al. 1988).
No significant X-ray emission was found, at the sensitivity of our current data, in the intergalactic space around NGC 1808. Similarly, no H I and H$`\alpha `$ emission was found by us (Dahlem et al. 2001).
## 5 Summary and conclusions
The work presented in this paper shows evidence of the presence of starburst activity and a hard unresolved source within the inner r=850 pc region in NGC 1808. Analysis of the optical-to-X-ray emission proves the co-existence of thermal diffuse plasma, non nuclear point-like sources and, of a LLAGN or a ULX.
The XMM-Newton EPIC images show extended X-ray emission for energies below 4.5 keV. The elongated shape observed in the X-ray band follows the same orientation as the H$`\alpha `$ and optical-UV emission, suggesting that SF contributes significantly to the soft X-ray emission. At least four non-nuclear point-like sources were detected in the Chandra ACIS image of the inner 250 pc of NGC 1808. The high spatial resolution Chandra image allows the detection of a double-peak nuclear structure. The locations of both sources are compatible with the position of the nucleus as derived from 2MASS. The spectral characteristics of these sources show that one is dominated by hot gas emission, indicating starburst activity, while the other is a highly absorbed hard X-ray source.
The spectral analysis of the XMM-Newton data completes the scenario outlined by X-ray imaging. Comparing EPIC-pn spectra of several regions in the central part of NGC 1808, it is inferred that the bulk of the emission originates from the nucleus, although in the soft X-ray band the contribution of the circumnuclear regions is not negligible.
The X-ray EPIC spectrum of the inner nuclear region (r=16″, 850 pc) is explained by a power law, $`\mathrm{\Gamma }=0.79_{0.06}^{+0.07}`$, accounting for the hard X-ray emission and associated with an unresolved point-like source. The absorption measured in the Chandra spectral analysis of source S2 is compatible with typical values of Compton-thin Seyfert 2 galaxies. However, this absorbing material does not show up in the XMM-Newton spectrum, probably because it includes contributions from different regions, making it difficult to separate absorption and soft emission components.
In order to explain the soft X-ray emission of the inner region, it is necessary to include emission due to a diffuse plasma with kT=0.580$`{}_{0.016}{}^{}{}_{}{}^{+0.019}`$ keV. The data obtained with RGS exhibit emission lines similar in wavelength and relative intensity ratios to the ones found for the prototypical starburst galaxy M 82. The temperature, the luminosity, and the various ratios of emission lines are consistent with being originated in a collisionally ionized plasma associated with SF regions, which dominates the total emission spectrum including the continuum in the soft X-ray regime.
The luminosities measured with XMM-Newton for the r=850 pc inner region (L$`{}_{210keV}{}^{}=(1.61\pm 0.06)\times 10^{40}`$ erg s<sup>-1</sup>) and with Chandra for source S2 (L$`{}_{210keV}{}^{}=0.92_{0.7}^{+0.08}\times 10^{40}`$ erg s<sup>-1</sup>) are both in the range of LLAGNs but are also compatible with the values of ULX sources. A decrease in the 2-10 keV flux on the order of 15$`\%`$ in an interval of eight months was measured when comparing XMM-Newton and Chandra observations, while no change in flux was detected in the soft energy band. Both the luminosities and the variability detected are compatible with the unresolved source being a LLAGN or a ULX.
Halo emission of NGC 1808 could also be detected: a softer thermal component, kT$``$0.1 keV, with L$`{}_{0.12.4keV}{}^{}7\times 10^{39}`$ erg s<sup>-1</sup> ($``$30% of the total soft X-ray emission) is necessary to explain the spectrum at distances larger than 875 pc from the centre.
The multiwavelength analysis of the nuclear region of NGC 1808 performed in this work allows us to detect and for the first time to unambiguously disentangle the contributions of an unresolved nuclear X-ray source and the starburst regions within the r=850 pc region. Further analysis of better quality data is required to determine the exact nature of the unresolved nuclear source.
###### Acknowledgements.
We specially thank N. Stuhrmann for reducing the H$`\alpha `$ data. EJB, MSLl and JMMH have been supported by Spanish MEC under grant AYA2001-3939-C03-02
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# Déformations des cônes de vecteurs primitifs
## Introduction
Soit $`G`$ un groupe réductif connexe complexe, et $`V`$ un $`G`$-module rationnel de dimension finie. On dit qu’un sous-schéma fermé $`G`$-stable $`XV`$ est à multiplicités finies si son algèbre affine est somme directe de modules simples avec des multiplicités finies.
Alexeev et Brion ont montré récemment dans \[AlBr\] que les sous-schémas $`X`$ ayant des multiplicités finies fixées sont paramétrés par un schéma quasi-projectif: le schéma de Hilbert invariant. Leur travail est basé sur celui de Haiman et Sturmfels (\[HaSt\]), qui correspond au cas particulier où le groupe réductif $`G`$ est un tore.
On se propose ici de déterminer le schéma de Hilbert invariant dans le cas “le plus simple”. L’espace ambiant est le $`G`$-module simple de plus grand poids $`\lambda `$:
$$V=V(\lambda ).$$
On cherche à paramétrer les plus petits sous-schémas fermés $`G`$-stables de $`V`$ de dimension positive: on va donc choisir les multiplicités les plus petites possibles. Pour cela, on remarque que l’algèbre affine d’un tel schéma $`X`$ contient le dual de $`V(\lambda )`$; on note $`\lambda ^{}`$ son plus grand poids. Notons $`f`$ un vecteur de plus grand poids de $`V(\lambda ^{})`$. Les puissances de $`f`$ sont non nulles dans l’algèbre affine de $`X`$, donc celle-ci contient tous les modules simples $`V(d\lambda ^{})`$, où $`d`$ est un entier positif.
Ainsi, on va prendre pour multiplicité $`1`$ pour les modules simples $`V(d\lambda ^{})`$, et $`0`$ pour les autres modules simples. Un point particulier du schéma de Hilbert invariant correspondant $`H_\lambda `$ est alors donné par le cône des vecteurs primitifs de $`V(\lambda )`$ (réunion de l’orbite des vecteurs de plus grand poids, et de l’origine). Ces cônes ne sont autres que les cônes affines sur les variétés de drapeaux plongées par un système linéaire complet.
On montrera que pour la plupart des poids $`\lambda `$, le schéma $`H_\lambda `$ est en fait réduit à ce point. Dans les autres cas, $`H_\lambda `$ est la droite affine. On obtient ainsi une classification des $`G`$-modules simples dont le cône des vecteurs primitifs admet une déformation invariante.
Cette classification rappelle celle (obtenue par Akhiezer dans \[Ak1\]) des variétés projectives lisses dont les orbites sous l’action d’un groupe algébrique affine connexe sont un diviseur ample et son complémentaire. On décrira comment relier ces deux classifications et on les reliera aussi à celle des algèbres de Jordan simples. Cette dernière se trouve être plus “petite”, mais on verra que lorsqu’un cône de vecteurs primitifs admet une déformation invariante dans $`V(\lambda )`$, cette déformation provient d’une algèbre de Jordan simple.
On montrera enfin que pour la plupart des poids $`\lambda `$, le cône des vecteurs primitifs de $`V(\lambda )`$ est en fait rigide: les seuls cônes de vecteurs primitifs admettant des déformations infinitésimales non triviales sont, outre ceux qui admettent une déformation invariante, le cône affine $`C_m`$ au dessus de la courbe rationnelle normale de degré $`m`$ dans $`^m`$ (les déformations de ce cône ont été étudiées dans \[Pi\]), et le cône affine $`C_{mn}`$ au dessus de $`^1\times ^n`$ dans le plongement de bidegré $`(m\mathrm{,1})`$ avec $`m2`$. On verra que la déformation verselle de $`C_{mn}`$ est lisse, contrairement à celle de $`C_m`$.
Remerciements. Je suis très vivement reconnaissant envers mon directeur de thèse Michel Brion pour son aide tout au long de ce travail.
## 1 Une classe de schémas de Hilbert invariants
### 1.1 Notations et résultat principal
On considère des schémas et des groupes algébriques sur $``$. Les références utilisées sont \[Ha\] pour la théorie des schémas et \[PoVi\] pour celle des groupes algébriques de transformations.
Soit $`G`$ un groupe réductif connexe. On en choisit un sous-groupe de Borel $`B`$, et un tore maximal $`T`$ inclus dans $`B`$. On considère le radical unipotent $`U`$ de $`B`$ : on a : $`B=TU`$. Les algèbres de Lie respectives de $`G`$, $`B`$, $`T`$ et $`U`$ sont notées : $`𝔤`$, $`𝔟`$, $`𝔱`$, et $`𝔲`$. Le système de racines de $`G`$ relativement à $`T`$ est noté $`R`$. Le choix de $`B`$ nous en fournit une base $`S`$, et on a : $`R=R_+R_{}`$$`R_+`$ est l’ensemble des racines positives, et $`R_{}`$ celui des racines négatives.
On note $`\mathrm{\Lambda }`$ le groupe des caractères de $`T`$. On a un ordre partiel sur $`\mathrm{\Lambda }`$ : $`\mu \lambda `$ si et seulement si $`\lambda \mu `$ est une somme de racines positives. On note $`\mathrm{\Lambda }^+`$ l’ensemble des éléments de $`\mathrm{\Lambda }`$ qui sont des poids dominants (relativement à la base $`S`$ du système de racines $`R`$). On sait que $`\mathrm{\Lambda }^+`$ est en bijection avec l’ensemble des classes d’isomorphisme de $`G`$-modules rationnels simples. Si $`\lambda `$ est un élément de $`\mathrm{\Lambda }^+`$, on notera $`V(\lambda )`$ un $`G`$-module simple correspondant, c’est-à-dire de plus grand poids $`\lambda `$, et $`v_\lambda `$ un vecteur de $`V(\lambda )`$ de poids $`\lambda `$. Si $`\lambda `$ est un poids qui n’est pas dominant, on pose $`V(\lambda )=0`$.
Les $`G`$-modules simples peuvent être construits de la façon suivante : soit $`\lambda `$ un poids dominant, et $`P`$ un sous-groupe parabolique de $`G`$ contenant $`B`$ tel que $`\lambda `$ se prolonge en un caractère de $`P`$. Notons $`\pi :GG/P`$ la surjection canonique, et $`_\lambda `$ le faisceau inversible sur $`G/P`$ qui associe à un ouvert $`\mathrm{\Omega }G/P`$ :
$$_\lambda (\mathrm{\Omega }):=\{f𝒪_G(\pi ^1(\mathrm{\Omega }))\text{ }|\text{ }gG\text{}pP\text{}f(gp)=\lambda (p)f(g)\}.$$
Le faisceau $`_\lambda `$ est alors $`G`$-linéarisé (via l’action de $`G`$ sur ses fonctions régulières par translation à gauche), et l’espace des sections globales de $`_\lambda `$ est un $`G`$-module simple :
$$\mathrm{\Gamma }(G/P,_\lambda )V(\lambda )^{}.$$
Si $`V`$ est un $`T`$-module rationnel (éventuellement de dimension infinie), on note $`V=_{\lambda \mathrm{\Lambda }}V_\lambda `$ sa décomposition en sous-espaces propres. Par exemple, l’algèbre de Lie de $`G`$ admet la décomposition :
$$𝔤=𝔱\underset{\alpha R}{}𝔤_\alpha ,$$
où chaque $`𝔤_\alpha `$ est de dimension 1. On choisit pour tout $`\alpha R`$ un générateur $`e_\alpha `$ de $`𝔤_\alpha `$.
Si $`V`$ est un $`G`$-module rationnel, on note $`V_{(\lambda )}`$ sa composante isotypique de type $`\lambda `$, c’est-à-dire le sous-module de $`V`$ somme des sous-modules isomorphes à $`V(\lambda )`$. On a alors la décomposition $`V=_{\lambda \mathrm{\Lambda }^+}V_{(\lambda )}`$.
Soit $`V`$ un $`G`$-module rationnel de dimension finie, et $`h:\mathrm{\Lambda }^+`$ une fonction. On appelle famille de sous-schémas fermés $`G`$-stables de $`V`$ un sous-schéma fermé $`G`$-stable de $`𝔛S\times V`$, où $`S`$ est un schéma avec action triviale de $`G`$. On note $`\pi :𝔛S`$ le morphisme induit par la projection $`S\times VS`$, et $`\pi _{}𝒪_𝔛`$ l’image directe par $`\pi `$ du faisceau structural de $`𝔛`$ . La famille $`𝔛`$ est dite de fonction de Hilbert $`h`$ si on a un isomorphisme de $`𝒪_S`$\- $`G`$-modules
$$\pi _{}𝒪_𝔛\underset{\lambda \mathrm{\Lambda }^+}{}_\lambda V(\lambda ),$$
où chaque $`_\lambda `$ est un $`𝒪_S`$-module localement libre de rang $`h(\lambda )`$. (Le morphisme $`\pi `$ est alors plat.)
Le foncteur contravariant : $`(\text{Schémas})^{}(\text{Ensembles})`$ qui associe à tout schéma $`S`$ l’ensemble des familles $`𝔛S\times V`$ de fonction de Hilbert $`h`$ est représenté par un schéma quasi-projectif noté $`Hilb_h^G(V)`$. (On renvoie à \[AlBr\]§1.2 pour plus de détails.)
On fixe désormais un poids dominant $`\lambda `$. On note $`\lambda ^{}`$ le plus grand poids du $`G`$-module $`V(\lambda )^{}`$ dual de $`V(\lambda )`$. Soit $`h_\lambda :\mathrm{\Lambda }^+`$ la fonction valant 1 sur $`\lambda ^{}`$ et $`0`$ ailleurs. On note dans la suite
$$H_\lambda :=Hilb_{h_\lambda }^G(V(\lambda ))$$
le schéma de Hilbert invariant associé à ce choix.
Si $`E`$ est un espace vectoriel de dimension finie, on note $`(E)`$ l’espace de ses droites. On a une action régulière de $`G`$ sur l’espace $`(V(\lambda ))`$. Notons $`[v_\lambda ](V(\lambda ))`$ la droite engendrée par $`v_\lambda `$ et
$$P_\lambda :=G_{[v_\lambda ]}$$
son stabilisateur dans $`G`$ : c’est le plus grand sous-groupe parabolique de $`G`$ qui contient $`B`$ et tel que $`\lambda `$ se prolonge en un caractère de $`P_\lambda `$. L’orbite de $`[v_\lambda ]`$ est la seule orbite fermée de $`(V(\lambda ))`$ (donc l’unique orbite de plus petite dimension). L’espace homogène projectif $`G/P_\lambda `$ se plonge ainsi dans $`(V(\lambda ))`$, et le faisceau inversible très ample associé à ce plongement est en fait $`_\lambda `$. Le cône affine au dessus de $`G/P_\lambda `$ dans $`V(\lambda )`$ est le cône
$$C_\lambda :=G.v_\lambda \{0\}=\overline{G.v_\lambda }$$
des vecteurs primitifs de $`V(\lambda )`$. C’est une variété normale (cf. \[Kr\], III.3.5). La variété $`G/P_\lambda (V(\lambda ))`$ est donc projectivement normale, et l’algèbre affine graduée du cône $`C_\lambda `$ est
$$\underset{d}{}\mathrm{\Gamma }(G/P_\lambda ,_{d\lambda })=\underset{d}{}V(d\lambda ^{}).$$
On peut donc voir $`C_\lambda `$ comme un point fermé de $`H_\lambda `$.
L’objet de cette partie est de montrer le théorème suivant, énoncé avec les notations de \[Bo\] :
###### Théorème 1.1.
Le schéma de Hilbert invariant $`H_\lambda `$ est un point réduit, sauf dans les cas suivants où $`H_\lambda `$ est la droite affine :
(H1) $`G`$ est simple de type $`A_1`$, et $`\lambda =2\omega _1\text{ ou }4\omega _1`$.
(H2) $`G`$ est simple de type $`A_n`$, $`n2`$ et $`\lambda =\omega _1+\omega _n`$.
(H3) $`G`$ est simple de type $`B_3`$ et $`\lambda =\omega _3\text{ ou }2\omega _3`$.
(H4) $`G`$ est simple de type $`B_n`$, $`n2`$ et $`\lambda =\omega _1\text{ ou }2\omega _1`$.
(H5) $`G`$ est simple de type $`C_n`$, $`n3`$ et $`\lambda =\omega _2`$.
(H6) $`G`$ est simple de type $`D_n`$, $`n3`$ et $`\lambda =\omega _1\text{ ou }2\omega _1`$.
(H7) $`G`$ est simple de type $`F_4`$ et $`\lambda =\omega _4`$.
(H8) $`G`$ est simple de type $`G_2`$ et $`\lambda =\omega _1\text{ ou }2\omega _1`$.
(H9) $`G`$ est semi-simple de type $`A_1\times A_1`$ et $`\lambda =(\omega _1,\omega _1)\text{ ou }(2\omega _1\mathrm{,2}\omega _1)`$.
et dans les cas $`(G,V(\lambda ))`$ obtenus à partir d’un cas $`(G_0,V_0)`$ parmi les précédents par factorisation : $`GG_0GL(V_0)`$.
###### Remarque 1.2.
* Le cas (H6) avec $`n=3`$ revient à un groupe $`G`$ simple de type $`A_3`$, avec $`\lambda =\omega _2\text{ ou }2\omega _2`$.
* On pourrait voir le cas (H9) comme étant le cas (H6) avec $`n=2`$.
### 1.2 Action du groupe multiplicatif sur le schéma de Hilbert invariant
On a une opération naturelle du groupe multiplicatif $`𝔾_m`$ sur le schéma de Hilbert invariant $`H_\lambda `$ : elle provient de l’action de $`𝔾_m`$ sur $`V(\lambda )`$ par homothéties (qui commute avec l’action de $`G`$).
Dans cette partie, on montre que cette action admet pour unique point fixe le cône $`C_\lambda `$ des vecteurs primitifs. On montre aussi que $`C_\lambda `$ est dans l’adhérence de toutes les orbites de $`𝔾_m`$, et on en déduit que le schéma de Hilbert invariant est affine. Ces propriétés peuvent être déduites de ce qui est fait dans \[AlBr\], §2.1 à 2.3 ; on a préféré donner ici des preuves directes.
###### Proposition 1.3.
(a) Le cône $`C_\lambda `$ est l’unique point fermé de $`H_\lambda `$ fixé par $`𝔾_m`$.
(b) Soit $`X`$ un point fermé de $`H_\lambda `$. Le morphisme : $`𝔾_mH_\lambda `$, $`tt.X`$ se prolonge en un morphisme $`𝔸^1H_\lambda `$, $`0C_\lambda `$.
Preuve. (a) On note $`S^eV(\lambda )^{}`$ la puissance symétrique d’ordre $`e`$ de $`V(\lambda )^{}`$. On identifie l’algèbre des fonctions régulières sur $`V(\lambda )`$ à l’algèbre symétrique de $`V(\lambda )^{}`$ :
$$\mathrm{Sym}V(\lambda )^{}:=\underset{e}{}S^eV(\lambda )^{}.$$
Les points fermés fixés par $`𝔾_m`$ correspondent aux idéaux homogènes
$$I=\underset{e}{}I_e\mathrm{Sym}V(\lambda )^{}=\underset{e}{}S^eV(\lambda )^{}$$
qui sont stables par $`G`$ et de fonction de Hilbert $`h`$.
On sait que $`S^eV(\lambda )^{}`$ contient un unique sous-$`G`$-module isomorphe à $`V(e\lambda )^{}`$, et que ses autres composantes isotypiques non nulles sont de type inférieur à $`e\lambda ^{}`$ :
$$S^eV(\lambda )^{}V(e\lambda )^{}\underset{\mu <e\lambda ^{}}{}[S^eV(\lambda )^{}]_{(\mu )}.$$
(1)
On va montrer par récurrence sur $`e`$, qu’un tel idéal $`I`$ vérifie :
$$I_e=\underset{\mu <e\lambda ^{}}{}[S^eV(\lambda )^{}]_{(\mu )}.$$
(2)
En effet, comme $`I`$ ne contient pas les constantes, on a $`I_0=0`$. Puis, si $`(2)`$ est satisfait pour tout $`d<e`$, il faut que :
$$\underset{\mu <e\lambda ^{}}{}[S^eV(\lambda )^{}]_{(\mu )}I_e$$
pour que la fonction de Hilbert de $`I`$ soit $`h_\lambda `$. Cette dernière inclusion est en fait une égalité, car sinon on aurait $`I_d=S^dV(\lambda )^{}`$ pour tout $`de`$.
Il n’y a donc pas d’autre point fermé fixé par $`𝔾_m`$ que $`C_\lambda `$.
(b) Pour mieux comprendre l’action de $`𝔾_m`$ sur $`H_\lambda `$, on reprend la construction du schéma de Hilbert invariant (voir \[HaSt\], §1,2,3 et \[AlBr\], §1.2).
Considérons l’action naturelle de $`𝔾_m`$ sur l’algèbre symétrique de $`V(\lambda )^{}`$, où $`𝔾_m`$ opère sur la composante $`S^eV(\lambda )^{}`$ avec le poids $`e`$, de sorte que $`\mathrm{Sym}V(\lambda )^{}`$ est une $`G\times 𝔾_m`$-algèbre rationnelle.
La sous-algèbre $`[\mathrm{Sym}V(\lambda )^{}]^U`$ des invariants par $`U`$ est alors une $`T\times 𝔾_m`$\- algèbre rationnelle de type fini, selon \[Gros\],Thm 9.4.
On en choisit un système fini de générateurs $`f_1,\mathrm{}f_n`$ formé de $`T\times 𝔾_m`$-vecteurs propres, et on note $`S=[x_1,\mathrm{},x_n]`$ l’algèbre de polynômes correspondante. L’algèbre $`S`$ est naturellement une $`T\times 𝔾_m`$\- algèbre rationnelle, et on a un morphisme surjectif de $`T\times 𝔾_m`$\- algèbres rationnelles :
$$\pi :S(\mathrm{Sym}V(\lambda )^{})^U.$$
L’action de $`T`$ sur $`S`$ fournit une graduation de $`S`$ par le groupe abélien $`\mathrm{\Lambda }`$.
On peut alors identifier $`H_\lambda `$ à un sous-schéma localement fermé d’un produit de Grassmanniennes, donc d’un produit d’espaces projectifs. Plus précisément, on sait (\[HaSt\]) qu’il existe une partie finie $`D`$ de $`\mathrm{\Lambda }`$, et pour tout $`\mu D`$, un sous-espace vectoriel de dimension finie $`N_\mu `$ de $`S_\mu `$ que l’on peut choisir stable par $`𝔾_m`$, tels que l’on ait un plongement :
$$H_\lambda \underset{\mu D}{}(\stackrel{r_\mu }{}N_\mu )$$
(3)
$`r_\mu :=dimN_\mu h_\lambda (\mu )`$. Décrivons l’image d’un point fermé par ce plongement : si $`I`$ est l’idéal d’un sous-schéma fermé de $`V(\lambda )`$ correspondant à un point fermé de $`H_\lambda `$, on lui associe pour tout $`\mu D`$ :
$$J_\mu :=\pi ^1(I^U)N_\mu .$$
Les $`N_\mu `$ sont des modules rationnels pour l’action de $`𝔾_m`$, donc les $`(^{r_\mu }N_\mu )`$ sont munis d’une action régulière de $`𝔾_m`$, pour laquelle le plongement $`(3)`$ est équivariant.
On peut maintenant vérifier le point (b) de la proposition :
Soit $`\mu D`$. Si $`\mu \lambda ^{}`$, alors $`r_\mu =dimN_\mu `$, et $`(^{r_\mu }N_\mu )`$ est réduit à un point. Sinon, écrivons $`\mu =e\lambda ^{}`$. On a $`(^{r_\mu }N_\mu )(N_\mu ^{})`$. Notons $`K:=N_\mu \mathrm{ker}\pi `$ et $`L`$ un supplémentaire $`𝔾_m`$-stable de $`K`$ dans $`N_\mu `$ :
$$N_\mu =KL.$$
Selon la décomposition (1) (considérée à tous les ordres), le plus grand poids de l’action de $`𝔾_m`$ sur $`L`$ est $`e`$ :
$$L=L_e\underset{c>e}{}L_c.$$
L’espace vectoriel $`J_\mu `$ est un hyperplan de $`KL`$ qui contient $`K`$ (par définition de $`J_\mu `$ et $`K`$) mais qui ne contient pas $`L_e`$ selon le lemme qui suit.
Montrons alors que le morphisme $`𝔾_m(N_\mu ^{})`$, $`tt.J_\mu `$ se prolonge en un morphisme $`f:𝔸^1(N_\mu ^{})`$ en posant $`f(0):=K_{c>e}L_c`$. Choisissons une base de $`N_\mu ^{}`$ compatible avec la décomposition $`N_\mu =KL_e_{c>e}L_c`$. Notons $`d`$ la dimension de $`L`$. Les coordonnées homogènes de $`J_\mu `$ dans $`(N_\mu ^{})`$ sont $`[\underset{\mathrm{dim}(K)\text{ fois}}{\underset{}{0:\mathrm{}:0}}:x_1:x_2:\mathrm{}:x_d]`$ et on a $`x_10`$. Celles de $`t.J_\mu `$ sont donc $`[0:\mathrm{}:0:t^ex_1:t^{c_2}x_2:\mathrm{}:t^{c_d}x_d]`$, où les $`c_j`$ sont des entiers strictement supérieurs à $`e`$. D’où l’assertion.
Comme $`f(0)`$ ne dépend pas de l’idéal $`I`$ considéré, il s’agit de $`\pi ^1(I_0^U)N_\mu `$$`I_0`$ est l’idéal du cône $`C_\lambda `$, d’où (b). $`\mathrm{}`$
###### Lemme 1.4.
Soit $`X`$ un sous-schéma fermé de $`V(\lambda )`$ de fonction de Hilbert $`h_\lambda `$, et $`I\mathrm{Sym}V(\lambda )^{}`$ son idéal.
Alors pour tout $`e`$, le sous-$`G`$-module de $`S^eV(\lambda )^{}`$ isomorphe à $`V(e\lambda )^{}`$ n’est pas inclus dans $`I`$.
Preuve. Notons $`f`$ un vecteur de plus grand poids de $`V(\lambda )^{}`$.
Le $`G`$-module $`[S^eV(\lambda )^{}]_{(e\lambda ^{})}`$ est simple, et $`f^e`$ en est un vecteur de plus grand poids. Supposons par l’absurde : $`f^eI`$. Alors $`f`$ appartient à l’idéal du sous-schéma réduit $`X_{red}`$ associé à $`X`$. Le sous-espace vectoriel de $`V(\lambda )`$ engendré par $`X_{red}`$ est un sous-$`G`$-module de $`V(\lambda )`$ inclus dans l’hyperplan défini par $`f`$ : il est donc réduit à $`\{0\}`$, et l’espace vectoriel $`[X]`$ est de dimension finie : une contradiction. $`\mathrm{}`$
###### Corollaire 1.5.
Le schéma $`H_\lambda `$ est affine. Son algèbre affine $`A`$ est graduée par l’action du groupe multiplicatif sur $`H_\lambda `$, en degrés négatifs : $`A=_dA_d`$, et l’anneau $`A_0`$ est local.
Preuve. On rappelle que $`H_\lambda `$ s’identifie à un sous-schéma localement fermé $`𝔾_m`$-stable de $`(M)`$, où $`M`$ est un $`𝔾_m`$-module rationnel de dimension finie. Le sous-schéma réduit $`\overline{H_\lambda }H_\lambda `$ est donc aussi $`𝔾_m`$-stable, et son idéal homogène aussi. Il existe donc un élément homogène $`f\mathrm{Sym}(M^{})`$, $`𝔾_m`$-vecteur propre, définissant un ouvert $`U`$ contenant $`C_\lambda `$. Cet ouvert $`U`$ est $`𝔾_m`$-stable, et contient donc $`H_\lambda `$, selon le point (b) de la proposition précédente. Ainsi, $`H_\lambda `$ est fermé dans $`U`$, et $`U`$ est un ouvert affine : $`H_\lambda `$ est affine.
Montrons maintenant que $`A_e=0`$ pour tout $`e0`$. Par l’absurde, soit $`fA\{0\}`$ de degré $`e>0`$. Soit $`X`$ un point de $`H_\lambda `$ tel que $`f(X)0`$. La fonction $`𝔾_m`$, $`tf(tX)=t^ef(X)`$ se prolonge en une fonction régulière sur $`𝔸^1`$ : une contradiction.
Enfin, montrons que $`A_0=A^{𝔾_m}`$ n’a qu’un seul idéal maximal. On sait (\[Kr\], II.3.2) que le morphisme $`\mathrm{Spec}(A)\mathrm{Spec}(A^{𝔾_m})`$ est surjectif. Donc si $`A^{𝔾_m}`$ avait deux points fermés distincts, $`H_\lambda `$ aurait deux fermés $`𝔾_m`$-stables disjoints : une contradiction avec le point (b) de la proposition précédente. $`\mathrm{}`$
### 1.3 Calcul de l’espace tangent au point fixe
L’objet de ce paragraphe est de montrer le résultat suivant :
###### Proposition 1.6.
L’espace tangent $`T_{C_\lambda }H_\lambda `$ est nul, sauf dans les cas (H1) à (H9) du théorème 1.1 où il est de dimension 1.
Notons $`G_{v_\lambda }`$ le stabilisateur de $`v_\lambda `$ dans $`G`$. L’espace tangent en $`v_\lambda `$ à $`C_\lambda =G.v_\lambda \{0\}`$ est $`𝔤.v_\lambda V(\lambda )`$. Il est stabilisé par l’action de $`G_{v_\lambda }`$. On note enfin $`[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}`$ l’espace des invariants du quotient par $`G_{v_\lambda }`$.
Le point de départ de la démonstration est l’isomorphisme canonique :
$$T_{C_\lambda }H_\lambda [V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}.$$
Cet isomorphisme découle de la proposition 1.5 (iii) de \[AlBr\]. En effet on peut supposer que l’espace vectoriel $`V(\lambda )`$ n’est pas une droite : la variété $`C_\lambda `$ est alors de dimension supérieure ou égale à $`2`$, et normale. La codimension de $`C_\lambda \backslash G.v_\lambda =\{0\}`$ est donc supérieure ou égale à $`2`$, et la proposition s’applique.
###### Lemme 1.7.
On a $`G_{v_\lambda }=T_{v_\lambda }.G_{v_\lambda }^{}`$, en notant $`T_{v_\lambda }`$ le stabilisateur de $`v_\lambda `$ dans $`T`$ (on a $`T_{v_\lambda }=\mathrm{ker}(\lambda )`$ et $`G_{v_\lambda }^{}`$ la composante neutre de $`G_{v_\lambda }`$.
Preuve. Considèrons la décomposition de Lévi de $`P_\lambda =G_{[v_\lambda ]}`$ relative à $`T`$ :
$$P_\lambda =L_\lambda .U_\lambda .$$
Comme $`T`$ est un tore maximal du groupe réductif $`L_\lambda `$, on a $`L_\lambda =T.[L_\lambda ,L_\lambda ]`$. D’où $`P_\lambda =T.[L_\lambda ,L_\lambda ].U_\lambda `$ et $`G_{v_\lambda }=T_{v_\lambda }.[L_\lambda ,L_\lambda ].U_\lambda `$ (car $`[L_\lambda ,L_\lambda ]`$ et $`U_\lambda `$ stabilisent $`v_\lambda `$), d’où le résultat, car $`[L_\lambda ,L_\lambda ]`$ et $`U_\lambda `$ sont connexes. $`\mathrm{}`$
###### Proposition 1.8.
On a une action du tore $`T`$ sur l’espace $`[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}`$. Sa décomposition en sous-espaces propres est :
$$[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}=[V(\lambda )/𝔤.v_\lambda ]_0^U[V(\lambda )/𝔤.v_\lambda ]_\lambda ^U$$
Preuve. On observe que les poids de $`V(\lambda )`$ qui sont des multiples de $`\lambda `$ sont $`\lambda `$ et éventuellement $`0`$ et $`\lambda `$.
Comme $`P_\lambda `$ stabilise $`𝔤.v_\lambda `$, il agit sur $`V(\lambda )/𝔤.v_\lambda `$ et sur $`[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}`$ et $`[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }^{}}`$ (puisque $`G_{v_\lambda }`$ et $`G_{v_\lambda }^{}`$ sont des sous-groupes distingués de $`P_\lambda `$).
On a donc une action du tore $`T`$ sur $`[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}`$, et ses poids sont des poids de $`V(\lambda )`$ qui sont multiples de $`\lambda `$, car la restriction de l’action à $`T_{v_\lambda }`$ est triviale :
$$[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}=[V(\lambda )/𝔤.v_\lambda ]_0^{G_{v_\lambda }}[V(\lambda )/𝔤.v_\lambda ]_\lambda ^{G_{v_\lambda }}$$
(le poids $`\lambda `$ n’apparaît pas car $`V(\lambda )_\lambda `$ est inclus dans $`𝔤.v_\lambda `$). D’où, selon le lemme précédent
$$[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}=[V(\lambda )/𝔤.v_\lambda ]_0^{G_{v_\lambda }^{}}[V(\lambda )/𝔤.v_\lambda ]_\lambda ^{G_{v_\lambda }^{}}$$
car $`T_{v_\lambda }`$ agit trivialement sur le membre de droite.
Ainsi, on a $`[V(\lambda )/𝔤.v_\lambda ]^{G_{v_\lambda }}=[V(\lambda )/𝔤.v_\lambda ]_0^{𝔤_{v_\lambda }}[V(\lambda )/𝔤.v_\lambda ]_\lambda ^{𝔤_{v_\lambda }}`$
On en déduit alors la proposition. En effet, on a :
$$𝔤_{v_\lambda }=𝔲𝔱_{v_\lambda }\underset{\begin{array}{c}\alpha R_+,\lambda ,\alpha ^{}=0\end{array}}{}𝔤_\alpha $$
Tout vecteur de poids $`\lambda `$ est invariant par les algèbres de Lie $`𝔱_{v_\lambda }`$ et
$`{\displaystyle \underset{\begin{array}{c}\alpha R_+,\lambda ,\alpha ^{}=0\end{array}}{}}𝔤_\alpha `$, et tout vecteur de poids $`0`$ aussi s’il est invariant par $`𝔲`$. $`\mathrm{}`$
On obtient maintenant une condition nécessaire pour que l’espace tangent soit non nul :
###### Proposition 1.9.
Si $`[V(\lambda )/𝔤.v_\lambda ]_0^U0`$, alors $`\lambda `$ s’écrit $`\lambda =\alpha +\beta `$$`\alpha S`$ et $`\beta R_+`$.
Si $`[V(\lambda )/𝔤.v_\lambda ]_\lambda ^U0`$, alors $`\lambda `$ s’écrit $`\lambda =\frac{\alpha +\beta }{2}`$$`\alpha S`$ et $`\beta R_+`$.
Preuve. Soit $`vV(\lambda )_0`$ dont la classe dans $`V(\lambda )/𝔤.v_\lambda `$ est un $`U`$-invariant non nul. On exprime cela à l’aide de l’algèbre de Lie $`𝔲`$ de $`U`$ :
$$𝔲v𝔤v_\lambda \text{ et }v𝔤v_\lambda .$$
Comme $`vV(\lambda )^U`$, il existe une racine simple $`\alpha `$ telle que $`e_\alpha v0`$. Donc $`e_\alpha v`$ est un $`T`$-vecteur propre de $`𝔤.v_\lambda `$. Si $`e_\alpha v`$ était proportionnel à $`v_\lambda `$, alors $`vV(\lambda )_{\lambda \alpha }=𝔤_\alpha .v_\lambda `$ : une contradiction. Donc il existe une racine positive $`\beta `$, telle que $`e_\alpha v`$ est proportionnel à $`e_\beta v_\lambda `$. En considérant les poids, on obtient : $`\alpha =\beta +\lambda `$.
On vérifie de même la seconde implication. $`\mathrm{}`$
###### Proposition 1.10.
Supposons l’espace tangent à $`H_\lambda `$ en $`C_\lambda `$ non nul, et $`G`$ semi-simple. Alors l’image de $`G`$ dans $`\mathrm{GL}(V(\lambda ))`$ est simple ou de type $`A_1\times A_1`$.
Preuve.
L’algèbre de Lie de $`G`$ est un produit d’algèbres de Lie simples : $`𝔤=𝔤_1\times \mathrm{}\times 𝔤_r`$, avec $`r1`$. Celle de $`U`$ s’écrit : $`𝔲=𝔲_1\times \mathrm{}\times 𝔲_r`$, avec $`𝔲_j𝔤_j`$.
La donnée du poids dominant $`\lambda `$ de $`𝔤`$ revient à celle d’un poids dominant $`\lambda _i`$ de $`𝔤_i`$ pour tout $`i`$, et
$$V(\lambda )=V(\lambda _1)_{}\mathrm{}_{}V(\lambda _r).$$
On peut supposer que tous les $`\lambda _i`$ sont non nuls.
D’apres la proposition précédente, $`\lambda `$ est somme ou demi-somme de deux racines, on peut donc se limiter au cas où $`r=1\text{ ou }2`$. Le cas $`r=1`$ correspond au cas où $`G`$ est simple ; supposons donc que $`r=2`$, et que l’espace tangent en $`C_\lambda `$ est non nul.
Selon la proposition 1.8, il existe un vecteur non nul $`v`$ appartenant à $`V(\lambda )_0\text{ ou }V(\lambda )_\lambda `$ qui est $`U`$-invariant modulo $`𝔤.v_\lambda `$.
Le vecteur $`v`$ n’est invariant ni par $`𝔲_1`$, ni par $`𝔲_2`$. Donc il existe une racine simple $`\alpha _1`$ de $`𝔤_1`$ (resp. $`\alpha _2`$ de $`𝔤_2`$) telle que : $`e_{\alpha _1}.v0`$ (resp. $`e_{\alpha _2}.v0`$). Comme $`e_{\alpha _1}.v`$ (resp. $`e_{\alpha _2}.v`$) est dans $`𝔤.v_\lambda `$, il existe une racine $`\beta _1`$ telle que $`e_{\alpha _1}.v`$ est proportionnel à $`e_{\beta _1}.v_\lambda `$ (resp. une racine $`\beta _2`$ telle que $`e_{\alpha _2}.v`$ est proportionnel à $`e_{\beta _2}.v_\lambda `$).
Supposons que $`v`$ appartient à $`V(\lambda )_0`$. En considérant les poids, on a
$$\alpha _1+\beta _1=\alpha _2+\beta _2=\lambda ,$$
donc $`\beta _1`$ est en fait une racine de $`𝔤_2`$, et $`\beta _2`$ une racine de $`𝔤_1`$. On a donc, avec des notations évidentes
$$(\alpha _1\mathrm{,0})+(0,\beta _1)=(\lambda _1,\lambda _2)\text{ et }(0,\alpha _2)+(\beta _2\mathrm{,0})=(\lambda _1,\lambda _2).$$
Donc
$$\lambda _1=\alpha _1\text{ et }\lambda _2=\alpha _2.$$
Lorsque $`v`$ appartient à $`V(\lambda )_\lambda `$ on en déduit de même
$$\lambda _1=\alpha _1/2\text{ et }\lambda _2=\alpha _2/2.$$
Dans les deux cas, les algèbres de Lie simples $`𝔤_1`$ et $`𝔤_2`$ admettent une racine simple qui est un poids dominant : elles sont donc de type $`A_1`$. $`\mathrm{}`$
#### 1.3.1 Cas restant à étudier
Selon la proposition 1.9, pour que l’espace tangent en $`C_\lambda `$ soit non nul, il faut que $`\lambda `$ soit somme ou demi-somme d’une racine simple et d’une racine positive. On dresse ci-dessous la liste des cas où l’espace tangent peut être non nul, obtenue en calculant toutes les sommes et demi-sommes d’une racine simple et d’une racine positive, puis en ne gardant que celles qui sont des poids dominants.
Les notations sont celles de \[Bo\].
On donne $`\lambda `$ sous la forme d’une somme de poids fondamentaux, et sous la forme de somme ou demi-somme d’une racine simple et d’une racine positive de toutes les façons possibles.
1. $`G`$ est simple de type $`A_1`$,
$`\lambda =2\omega _1=\frac{1}{2}(\alpha _1+\alpha _1)`$
2. $`G`$ est simple de type $`A_1`$,
$`\lambda =4\omega _1=\alpha _1+\alpha _1`$
3. $`G`$ est simple de type $`A_2`$,
$`\lambda =3\omega _1=\alpha _1+(\alpha _1+\alpha _2)`$
4. $`G`$ est simple de type $`A_2`$,
$`\lambda =3\omega _2=\alpha _2+(\alpha _1+\alpha _2)`$
5. $`G`$ est simple de type $`A_3`$,
$`\lambda =\omega _2=\frac{1}{2}[\alpha _2+(\alpha _1+\alpha _2+\alpha _3)]`$
6. $`G`$ est simple de type $`A_3`$,
$`\lambda =2\omega _2=\alpha _2+(\alpha _1+\alpha _2+\alpha _3)`$
7. $`G`$ est simple de type $`A_n`$, $`n2`$,
$`\lambda =\omega _1+\omega _n=\alpha _1+(\alpha _2+\mathrm{}+\alpha _n)=\alpha _n+(\alpha _1+\mathrm{}+\alpha _{n1})`$
8. $`G`$ est simple de type $`B_2`$,
$`\lambda =\omega _2=\frac{1}{2}[\alpha _2+(\alpha _1+\alpha _2)]`$
9. $`G`$ est simple de type $`B_2`$,
$`\lambda =2\omega _2=\alpha _2+(\alpha _1+\alpha _2)`$
10. $`G`$ est simple de type $`B_3`$,
$`\lambda =\omega _3=\frac{1}{2}[\alpha _3+(\alpha _1+2\alpha _2+2\alpha _3)]`$
11. $`G`$ est simple de type $`B_3`$,
$`\lambda =2\omega _3=\alpha _3+(\alpha _1+2\alpha _2+2\alpha _3)`$
12. $`G`$ est simple de type $`B_n`$, $`n3`$,
$`\lambda =\omega _2=\alpha _2+(\alpha _1+\alpha _2+2(\alpha _3+\mathrm{}+\alpha _n))`$
13. $`G`$ est simple de type $`B_n`$, $`n2`$,
$`\lambda =\omega _1=\alpha _1+(\alpha _2+\mathrm{}+\alpha _n)=\alpha _n+(\alpha _1+\mathrm{}+\alpha _{n1})=\frac{1}{2}[\alpha _1+(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _n))]`$
14. $`G`$ est simple de type $`B_n`$, $`n2`$,
$`\lambda =2\omega _1=\alpha _1+(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _n))`$
15. $`G`$ est simple de type $`C_n`$, $`n3`$,
$`\lambda =\omega _1=\frac{1}{2}[\alpha _1+(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _{n1})+\alpha _n)]`$
16. $`G`$ est simple de type $`C_n`$, $`n3`$,
$`\lambda =2\omega _1=\alpha _1+(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _{n1})+\alpha _n)`$
17. $`G`$ est simple de type $`C_n`$, $`n3`$,
$`\lambda =\omega _2=\alpha _1+(2(\alpha _2+\mathrm{}+\alpha _{n1})+\alpha _n)=\alpha _2+(\alpha _1+\alpha _2+2(\alpha _3+\mathrm{}+\alpha _{n1})+\alpha _n)`$
18. $`G`$ est simple de type $`D_4`$,
$`\lambda =\omega _3=\frac{1}{2}[\alpha _3+(\alpha _1+2\alpha _2+\alpha _3+\alpha _4)]`$
19. $`G`$ est simple de type $`D_4`$,
$`\lambda =\omega _4=\frac{1}{2}[\alpha _4+(\alpha _1+2\alpha _2+\alpha _3+\alpha _4)]`$
20. $`G`$ est simple de type $`D_4`$,
$`\lambda =2\omega _3=\alpha _3+(\alpha _1+2\alpha _2+\alpha _3+\alpha _4)`$
21. $`G`$ est simple de type $`D_4`$,
$`\lambda =2\omega _4=\alpha _4+(\alpha _1+2\alpha _2+\alpha _3+\alpha _4)`$
22. $`G`$ est simple de type $`D_n`$, $`n4`$,
$`\lambda =\omega _1=\frac{1}{2}[\alpha _1+(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _{n2})+\alpha _{n1}+\alpha _n)]`$
23. $`G`$ est simple de type $`D_n`$, $`n4`$,
$`\lambda =2\omega _1=\alpha _1+(\alpha _1+2(\alpha _2+\mathrm{}+\alpha _{n2})+\alpha _{n1}+\alpha _n)`$
24. $`G`$ est simple de type $`D_n`$, $`n4`$,
$`\lambda =\omega _2=\alpha _2+(\alpha _1+\alpha _2+2(\alpha _3+\mathrm{}+\alpha _{n2})+\alpha _{n1}+\alpha _n)`$
25. $`G`$ est simple de type $`E_6`$,
$`\lambda =\omega _2=\alpha _2+(\alpha _1+\alpha _2+2\alpha _3+3\alpha _4+2\alpha _5+\alpha _6)`$
26. $`G`$ est simple de type $`E_7`$,
$`\lambda =\omega _1=\alpha _1+(\alpha _1+2\alpha _2+3\alpha _3+4\alpha _4+3\alpha _5+2\alpha _6+\alpha _7)`$
27. $`G`$ est simple de type $`E_8`$,
$`\lambda =\omega _8=\alpha _8+(2\alpha _1+3\alpha _2+4\alpha _3+6\alpha _4+5\alpha _5+4\alpha _6+3\alpha _7+\alpha _8)`$
28. $`G`$ est simple de type $`F_4`$,
$`\lambda =\omega _1=\alpha _1+(\alpha _1+3\alpha _2+4\alpha _3+2\alpha _4)`$
29. $`G`$ est simple de type $`F_4`$,
$`\lambda =\omega _4=\alpha _3+(\alpha _1+2\alpha _2+2\alpha _3+2\alpha _4)=\alpha _4+(\alpha _1+2\alpha _2+3\alpha _3+\alpha _4)`$
30. $`G`$ est simple de type $`G_2`$,
$`\lambda =\omega _1=\alpha _1+(\alpha _1+\alpha _2)=\frac{1}{2}[\alpha _1+(3\alpha _1+2\alpha _2)]`$
31. $`G`$ est simple de type $`G_2`$,
$`\lambda =2\omega _1=\alpha _1+(3\alpha _1+2\alpha _2)`$
32. $`G`$ est simple de type $`G_2`$,
$`\lambda =\omega _2=\alpha _2+(3\alpha _1+\alpha _2)`$
33. $`G`$ est semi-simple de type $`A_1\times A_1`$,
$`\lambda =(\omega _1,\omega _1)=\frac{1}{2}[(\alpha _1\mathrm{,0})+(0,\alpha _1)]`$
34. $`G`$ est semi-simple de type $`A_1\times A_1`$,
$`\lambda =(2\omega _1\mathrm{,2}\omega _1)=(\alpha _1\mathrm{,0})+(0,\alpha _1)`$
Tous les cas se traitent de la même manière : selon la proposition 1.8, il s’agit de calculer
$$dim[V(\lambda )/𝔤.v_\lambda ]_0^U\text{ et }dim[V(\lambda )/𝔤.v_\lambda ]_\lambda ^U.$$
Des connaissances élémentaires sur les modules irréductibles (pour lesquelles on renvoie à \[Serr\]) permettent de mener à bien le calcul. A titre d’exemples, on traite quelques cas représentatifs dans les sections 1.3.2 à 1.3.4.
#### 1.3.2 Cas de la représentation adjointe d’un groupe simple
Il s’agit des cas (C1), (C7), (C9), (C12), (C16), (C25), (C26), (C27), (C28) et (C32).
On a $`V(\lambda )𝔤`$ et $`\lambda `$ est la plus grande racine.
1) Le cas (C1) où $`G`$ est de type $`A_1`$ se traite à part : $`𝔤`$ s’identifie à l’algèbre de Lie des matrices $`2\times 2`$ de trace nulle. On pose :
$$x:=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)h:=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)y:=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)$$
La matrice $`x`$ est un vecteur de plus grand poids : l’espace tangent est isomorphe à
$$[𝔤/𝔤.x]_\lambda ^U=[𝔤/(xh)]_\lambda ^Uy.$$
Il est donc de dimension $`1`$.
2) Dans les autres cas, $`\lambda `$ n’est pas demi-somme d’une racine simple et d’une racine positive. L’espace tangent est donc isomorphe à $`[𝔤/𝔤.v_\lambda ]_0^U`$.
Posons
$$E:=\{\gamma R\text{ }|\text{ }\delta R,\gamma +\delta =\lambda \}.$$
On a alors, en notant $`h_\lambda `$ l’élément de $`[𝔤_\lambda ,𝔤_\lambda ]`$ tel que $`\lambda (h_\lambda )=2`$ :
$$𝔤v_\lambda =[𝔤,𝔤_\lambda ]=𝔤_\lambda h_\lambda \underset{\gamma E}{}𝔤_\gamma .$$
Le sous-espace de poids nul de $`𝔤/𝔤.v_\lambda `$ est isomorphe à $`𝔱/h_\lambda `$, et l’espace tangent est isomorphe à
$$\{t𝔱\text{ }|\text{ }\alpha S,[𝔤_\alpha ,𝔱]𝔤v_\lambda \}/h_\lambda ,$$
donc à
$$\{t𝔱\text{ }|\text{ }\alpha SE,\alpha (t)=0\}/h_\lambda .$$
Il est donc de dimension $`dim𝔱\mathrm{card}(SE)1`$. Or $`dim𝔱=\mathrm{card}S`$. La dimension de l’espace tangent est donc $`\mathrm{card}(SE)1`$.
Comme $`\mathrm{card}(SE)`$ est le nombre de façons d’écrire $`\lambda `$ comme somme d’une racine simple et d’une racine positive, on conclut à l’aide de la liste 1.3.4 que l’espace tangent est de dimension $`1`$ dans le cas où $`G`$ est de type $`A_n`$, et $`0`$ dans les autres cas.
#### 1.3.3 Cas (C23)
Ici, $`𝔤`$ s’identifie à l’algèbre de Lie $`𝔰𝔬(2n)`$ des matrices de taille $`2n\times 2n`$ antisymétriques par rapport à la seconde diagonale. On a une action naturelle de $`𝔰𝔬(2n)`$ sur $`^{2n}`$, dont la base canonique est notée $`(e_1,\mathrm{},e_n,e_n,\mathrm{},e_1)`$. Le module simple $`V(2\omega _1)`$ peut être vu comme un quotient du carré symétrique de $`^{2n}`$ :
$$V(2\omega _1)=S^2^{2n}/e_1e_1+\mathrm{}+e_ne_n.$$
Il s’agit de calculer la dimension de $`[V(2\omega _1)/𝔤v_{2\omega _1}]_0^U`$.
Notons $`\pi :V(2\omega _1)V(2\omega _1)/𝔤v_{2\omega _1}`$ la surjection canonique. On remarque que $`\pi `$ induit des isomorphismes :
$$\pi (V(2\omega _1)_0)V(2\omega _1)_0,\text{ }\pi (V(2\omega _1)_{\alpha _2})V(2\omega _1)_{\alpha _2},\text{ … , }\pi (V(2\omega _1)_{\alpha _n})V(2\omega _1)_{\alpha _n}.$$
Par contre, $`\pi (V(2\omega _1)_{\alpha _1})=0`$. On a donc, en notant $`A`$ la sous-algèbre de $`𝔤`$ engendrée par une partie $`A`$ de $`𝔤`$ :
$$[\pi (V(2\omega _1))]_0^UV(2\omega _1)_0^{e_{\alpha _2},\mathrm{},e_{\alpha _n}}=V(2\omega _1)_0^{e_{\alpha _2},\mathrm{},e_{\alpha _{n1}}}$$
car $`e_{\alpha _{n1}}\text{ et }e_{\alpha _n}`$ stabilisent les mêmes éléments dans $`V(2\omega _1)_0`$.
Or on sait que $`V(2\omega _1)_0`$ est de dimension $`n1`$, et $`V(2\omega _1)_{\alpha _1},\mathrm{},V(2\omega _1)_{\alpha _{n1}}`$ sont de dimension $`1`$.
Comme $`V(2\omega _1)_0^{e_{\alpha _1},\mathrm{},e_{\alpha _{n1}}}=V(2\omega _1)_0^U=\{0\}`$, on en déduit que l’espace vectoriel $`V(2\omega _1)_0^{e_{\alpha _2},\mathrm{},e_{\alpha _{n1}}}`$ est de dimension $`1`$.
Donc l’espace tangent est de dimension $`1`$.
#### 1.3.4 Cas (C31)
Les poids de $`V(2\omega _1)`$ sont :
* $`0`$ avec multiplicité $`3`$
* $`\alpha _1`$ et ses conjugués sous l’action du groupe de Weyl, chacun avec multiplicité $`2`$
* $`\alpha _2`$ et ses conjugués, chacun avec multiplicité $`1`$
* $`2\omega _1=4\alpha _1+2\alpha _2`$ et ses conjugués, chacun avec multiplicité $`1`$.
Comme $`V(2\omega _1)_0^U=\{0\}`$, l’application linéaire
$$\begin{array}{ccc}\hfill \varphi :V(2\omega _1)_0& \hfill & V(2\omega _1)_{\alpha _1}V(2\omega _1)_{\alpha _2}\\ \hfill v& \hfill & (e_{\alpha _1}v,e_{\alpha _2}v)\end{array}$$
est injective, donc bijective.
Notons $`\pi :V(2\omega _1)V(2\omega _1)/𝔤.v_{2\omega _1}`$ la surjection canonique. On remarque que $`\pi `$ induit des isomorphismes
$$\pi (V(2\omega _1)_0)V(2\omega _1)_0\text{ et }\pi (V(2\omega _1)_{\alpha _2})V(2\omega _1)_{\alpha _2}.$$
Par contre, $`dim\pi (V(2\omega _1)_{\alpha _1})=1`$.
On en déduit que l’application quotient
$$\overline{\varphi }:\pi (V(2\omega _1)_0)\pi (V(2\omega _1)_{\alpha _1})\pi (V(2\omega _1)_{\alpha _2})$$
a pour noyau $`\pi (V(2\omega _1)_0)^U`$, de dimension $`1`$.
Donc l’espace tangent est de dimension $`1`$.
### 1.4 Conclusion
Dans ce paragraphe, on ramène la démonstration du théorème 1.1, à des résultats obtenus au §2.4 (de façon indépendante). On commence par énoncer sans démonstration une conséquence immédiate du lemme de Nakayama :
###### Lemme 1.11.
Soit $`A=_dA_d`$ un anneau gradué par $``$. On suppose l’anneau $`A_0`$ local d’idéal maximal $`_0`$. Notons $``$ l’idéal maximal $`_0_{d=1}^{\mathrm{}}A_d`$ de $`A`$. Soit $`M=_dM_d`$ un $`A`$-module gradué de type fini. Si $`.M=M`$, alors $`M=0`$.
Le schéma $`H_\lambda `$ est presque déterminé par le corollaire suivant du corollaire 1.5 et de la proposition 1.6 :
###### Corollaire 1.12.
Le schéma $`H_\lambda `$ est soit une droite affine, soit un point épaissi $`\mathrm{Spec}[t]/(t^N)`$, pour un entier $`N`$.
Preuve. On garde les notations du corollaire 1.5. L’idéal $``$ de $`A`$ correspondant au point $`C_\lambda `$ de $`H_\lambda `$ est l’unique idéal maximal de $`A`$ fixé par $`𝔾_m`$, donc on a
$$=_0\underset{d=1}{\overset{\mathrm{}}{}}A_d$$
$`_0`$ est l’idéal maximal de $`A_0`$. Selon la proposition 1.6, l’espace vectoriel $`/^2`$ est de dimension inférieure ou égale à 1. Il existe donc un élément homogène $`f`$ de $``$ tel que $`=Af+^2`$. Selon le lemme précédent, $`=Af`$. Or $`A`$ s’écrit, comme espace vectoriel $`A=`$, donc on a $`A=ff^2\mathrm{}=[f]`$. Ainsi, l’algèbre $`A`$ est monogène, et comme son spectre est connexe (proposition 1.3), on en déduit le corollaire. $`\mathrm{}`$
Le théorème 1.1 est donc démontré dans les cas où l’espace tangent est nul. Dans les cas où il est de dimension 1, il reste à exclure les points épaissis. On va voir (§2.4) qu’à chaque fois que l’espace tangent à $`H_\lambda `$ en $`C_\lambda `$ est de dimension 1, il existe d’autres points fermés que $`C_\lambda `$ dans $`H_\lambda `$, et celui-ci est donc une droite affine.
## 2 Algèbres de Jordan simples et familles universelles
On va maintenant relier la classification (H) du théorème 1.1 à deux classifications déjà connues :
* celle notée (J) des algèbres de Jordan simples (théorème 2.1)
* et une classification notée (A) de variétés projectives à deux orbites (théorème 2.2).
Pour cela, on va associer de façon naturelle
* aux objets de (J) des objets de (A) dans §2.3 (en considérant le cône des éléments de rang $`1`$ des algèbres de Jordan simples)
* et aux objets de (A) des objets de (H) dans §2.4
On constatera alors que lors de ces deux opérations, tous les cas sont atteints. Ainsi, pour chacun des cas du théorème 1.1, on obtiendra à l’aide d’une algèbre de Jordan simple une déformation non triviale du cône $`C_\lambda `$ dans $`V(\lambda )`$, qui sera en fait la famille universelle au dessus de $`H_\lambda `$ (§2.5).
### 2.1 Classification des algèbres de Jordan simples
Pour plus de détails concernant les algèbres de Jordan, on renvoie à \[Jac\] ou \[FaKo\]. On appellera algèbre de Jordan (complexe) une $``$-algèbre $`(A,)`$ commutative unitaire de dimension finie (non nécessairemant associative) telle que
$$a,bA,a^2(ab)=a(a^2b).$$
On peut montrer que $`A`$ est associative relativement aux puissances (c’est-à-dire : toutes ses sous-algèbres monogènes sont associatives).
On peut alors définir le polynôme minimal d’un élément $`a`$ de $`A`$ : c’est le générateur unitaire de l’idéal $`\{P[X]\text{ tels que }P(a)=0\}`$ de $`[X]`$. Un élément de $`A`$ est dit régulier si le degré de son polynôme minimal est maximal. Les éléments réguliers forment un ouvert dense de $`A`$.
Il existe alors (\[FaKo\], prop II.2.1) des fonctions polynômiales $`p_1,\mathrm{},p_r`$ sur $`A`$ telles que si $`a`$ est un élément régulier de $`A`$, son polynôme minimal est $`X^r+p_1(a)X^{r1}+\mathrm{}+p_r(a)`$. Chaque $`p_i`$ est alors homogène de degré $`i`$. On définit la trace et le déterminant sur $`A`$ par :
$$\mathrm{tr}:=p_1\text{ et }\mathrm{det}:=(1)^rp_r.$$
Dans la suite, on s’intéresse aux algèbres de Jordan simples, que l’on définit maintenant. Si $`a`$ est un élément de $`A`$, on note $`L(a):AA`$, $`bab`$ la multiplication par $`a`$. On note $`\mathrm{Tr}`$ la trace d’un endomorphisme $`AA`$. Une algèbre de Jordan $`A`$ est semi-simple si la forme bilinéaire
$$\begin{array}{ccc}A\times A& \hfill & \hfill \\ (a,b)& \hfill & \mathrm{Tr}L(ab)\hfill \end{array}$$
est non dégénérée. Elle est dite simple si de plus elle n’admet pas d’idéaux non triviaux. (Dans ce cas, les formes linéaires $`a\mathrm{Tr}L(a)`$ et $`atr(a)`$ sont en fait proportionnelles.)
Les matrices hermitiennes sur les complexifiées
$$R=_{}\text{}_{}\text{}_{}\text{}_{}𝕆$$
des algèbres de Hurwitz donnent des exemples d’algèbres de Jordan simples : on note $`H_n(R)`$ l’espace vectoriel des matrices de taille $`n\times n`$ qui sont égales à la transposée de leur conjuguée. On munit $`H_n(R)`$ d’une structure d’algèbre en posant :
$$M_1M_2=\frac{1}{2}(M_1M_2+M_2M_1).$$
La classification des algèbres de Jordan simples est connue (\[Jac\], théorème 8 p 203) :
###### Théorème 2.1 (P.Jordan, J.von Neumann, E.Wigner).
Toute algèbre de Jordan simple est isomorphe à l’une des suivantes :
(J1) $`W`$, où $`W`$ est un espace vectoriel de dimension finie muni d’une forme bilinéaire non dégénérée $`.,.`$. La loi de l’algèbre est donnée par la formule : $`(t_1,w_1)(t_2,w_2)=(t_1t_2+w_1,w_2,t_1w_2+t_2w_1).`$
(J2) $`H_n(_{})`$ , $`n3`$
(J3) $`H_n(_{})`$ , $`n3`$
(J4) $`H_n(_{})`$ , $`n3`$
(J5) $`H_3(_{}𝕆)`$.
On définit enfin le groupe de structure d’une algèbre de Jordan simple $`A`$ : c’est le groupe des automorphismes (d’espace vectoriel) de $`A`$ qui conservent à un scalaire près le déterminant :
$$\mathrm{Str}(A):=\{gGL(A)\text{ }|\text{ }u^{},aA,\mathrm{det}(ga)=u\mathrm{det}(a)\}.$$
Il contient le groupe $`\mathrm{Aut}(A)`$ des automorphismes d’algèbre de $`A`$. Les deux groupes $`\mathrm{Aut}(A)`$ et $`\mathrm{Str}(A)`$ sont des groupes algébriques réductifs (mais non connexes). En fait, $`\mathrm{Aut}(A)^{}`$ est semi-simple, et $`\mathrm{Str}(A)^{}`$ est de centre les homothéties.
### 2.2 Classification de variétés à deux orbites
Les espaces homogènes sous l’action d’un groupe réductif admettant une complétion équivariante par un diviseur homogène ont été classifiés par D. Akhiezer (voir \[Ak1\] ou \[HuSn\] ; la classification est retrouvée dans \[Bri\] par des méthodes algébriques). Dans le cas où le diviseur est ample, on a :
###### Théorème 2.2 (D.Akhiezer).
Soit $`Z`$ une variété projective lisse. Soit $`D`$ un diviseur ample de $`Z`$, et $`\mathrm{\Omega }`$ son complémentaire. On suppose qu’il existe une action régulière d’un groupe algébrique affine connexe $`\mathrm{\Gamma }`$ sur $`Z`$ sous laquelle $`\mathrm{\Omega }`$ et $`D`$ sont les orbites de $`Z`$.
Soit $`G`$ l’image de $`\mathrm{\Gamma }`$ dans $`\mathrm{Aut}(Z)`$, et $`H`$ le stabilisateur d’un point de $`\mathrm{\Omega }`$. Alors, à revêtement fini de $`G`$ près, on est dans un des cas suivants :
(A1) $`G=SL(n+1),n1`$ , $`H=GL(n)`$ et $`Z=^n\times (^n)^{}`$.
(A2) $`G=SO(n),n3`$ , $`H=SO(n1)`$ et $`Z=Q(n1)`$.
(A3) $`G=SO(n),n3`$ , $`H=O(n1)`$ et $`Z=^{n1}`$.
(A4) $`G=Sp(2n),n2`$ , $`H=Sp(2)\times Sp(2n2)`$ et $`Z`$ est la grassmanienne des 2-plans de $`^{2n}`$.
(A5) $`G=F_4`$ , $`H=Spin(9)`$ et $`Z=E_6/P`$ en notant $`P`$ le sous-groupe parabolique maximal de $`E_6`$ dont les racines simples sont $`\alpha _2,\mathrm{},\alpha _n`$ (notations de \[Bo\]).
(A6) $`G=G_2`$ , $`H=SL(3)`$ et $`Z=Q(6)`$.
(A7) $`G=G_2`$ , $`H=N_G(SL(3))`$ et $`Z=^6`$.
(A8) $`G=Spin(7)`$ , $`H=G_2`$ et $`Z=Q(7)`$.
(A9) $`G=SO(7)`$ , $`H=G_2`$ et $`Z=^7`$.
On a noté $`Q(n)`$ une quadrique projective lisse de dimension $`n`$.
Les actions de $`SL(n+1)`$, $`SO(n)`$ et $`Sp(2n)`$ sont les actions naturelles. Les actions de $`G_2`$, $`Spin(7)`$ et $`SO(7)`$ sont déduites des plongements $`G_2SO(7)`$ via la représentation de dimension 7 de $`G_2`$, et $`Spin(7)SO(8)`$ via la représentation spinorielle de $`Spin(7)`$.
### 2.3 Un lien entre les classifications (J) et (A)
Dans ce paragraphe, on rappelle une définition du cône des éléménts de rang $`1`$ d’une algèbre de Jordan simple. La variété projective formée des droites de ce cône nous donne alors une variété à 2 orbites de la classification (A) ; l’orbite ouverte de cette variété correspond aux éléments de trace non nulle.
On obtient au passage une bijection entre les classes d’isomorphisme des algèbres de Jordan simples complexes et les classes d’isomorphisme des espaces symétriques de rang 1 complexes (en effet, ceux-ci sont les quotients $`G/H`$, où $`G`$ et $`H`$ sont les groupes donnés dans les cas (A1) à (A5)).
La correspondance analogue dans le cas réel (entre les algèbres de Jordan simples réelles et les espaces symétriques de rang 1 réels compacts) est établie dans \[Hi\] ; dans ce cas, tout élément de rang 1 est de trace non nulle, c’est pourquoi l’espace symétrique des droites d’éléments de rang 1 (et de trace non nulle) est compact.
Soit $`A`$ une algèbre de Jordan simple.
Notons $`\mathrm{\Gamma }`$ la composante neutre du groupe des automorphismes de $`A`$, et $`\mathrm{\Gamma }^{}`$ celle du groupe de structure de $`A`$ :
$$\mathrm{\Gamma }:=\text{Aut}(A)^{}\mathrm{\Gamma }^{}:=\text{Str}(A)^{}.$$
Notons $`.1`$ la droite engendrée par l’élément unité de $`A`$, et $`V`$ le sous-espace vectoriel de $`A`$ formé des éléments de trace nulle. Alors $`A`$ est somme directe des deux sous-espaces vectoriels :
$$A=.1V,$$
(4)
et cette décomposition est stable par $`\mathrm{\Gamma }`$. On vérifie (grâce à la classification) que $`V`$ est un $`\mathrm{\Gamma }`$-module simple.
Notons $`D`$ la $`\mathrm{\Gamma }`$-orbite fermée dans $`(V)`$ et $`\stackrel{~}{D}`$ le cône des vecteurs primitifs de $`V`$, c’est-à-dire le cône affine sur $`D`$.
On vérifie également que comme $`\mathrm{\Gamma }^{}`$-module rationnel, $`A`$ est simple ; notons $`Z`$ la $`\mathrm{\Gamma }^{}`$-orbite fermée dans $`(A)`$ et $`\stackrel{~}{Z}`$ son cône des vecteurs primitifs de $`A`$. Les éléments (non nuls) de $`\stackrel{~}{Z}`$ sont appelés les éléments de rang $`1`$ de l’algèbre de Jordan $`A`$. Les éléments (non nuls) de $`\stackrel{~}{D}`$ sont les éléments de rang $`1`$ et de trace nulle.
On remarque que $`D`$ est un diviseur ample de $`Z`$ :
$$D=Z(V)Z(.1V)$$
et on vérifie enfin que $`\mathrm{\Gamma }`$ agit transitivement sur $`D`$ et sur $`ZD`$.
Ainsi, à toute algèbre de Jordan simple on fait correspondre un élément de la classification (A) en prenant comme groupe $`G`$ le groupe $`\mathrm{\Gamma }`$ ; on peut aussi prendre comme groupe $`G`$ un sous-groupe fermé de $`\mathrm{\Gamma }`$ pourvu qu’il agisse transitivement sur $`D`$ et sur $`ZD`$.
Précisément :
* à partir de (J1), on obtient le cas (A2) quand $`G`$ est le groupe $`\mathrm{\Gamma }=SO(W)`$, mais aussi, en considérant des sous-groupes stricts de $`\mathrm{\Gamma }`$, le cas (A6) quand $`W`$ est de dimension $`7`$ et $`G=G_2`$ et le cas (A8) quand $`W`$ est de dimension $`8`$ et $`G=\mathrm{Spin}(7)`$.
* à partir de (J2), on obtient le cas (A3) quand $`G`$ est le groupe $`\mathrm{\Gamma }=SO(n)`$, mais aussi, en considérant des sous-groupes stricts de $`\mathrm{\Gamma }`$, les cas (A7) quand $`n=7`$ et (A9) quand $`n=8`$.
* à partir de (J3), on obtient le cas (A1) avec $`G=\mathrm{\Gamma }=\mathrm{PGL}(n)`$.
* à partir de (J4), on obtient le cas (A4) avec $`G=\mathrm{\Gamma }=\mathrm{Sp}(2n)`$.
* à partir de (J5), on obtient le cas (A5) avec $`G=\mathrm{\Gamma }=F_4`$.
On voit donc que tous les cas du théorème 2.2 peuvent être obtenus à partir des algèbres de Jordan simples.
### 2.4 Un lien entre les classifications (A) et (H)
Supposons que le groupe réductif connexe $`G`$ agit sur une variété projective lisse $`Z`$, et que ses orbites sont un diviseur ample $`D`$ et son complémentaire $`\mathrm{\Omega }`$ (de sorte que l’on est dans la situation du théorème 2.2).
Alors $`D`$ est en fait très ample, et si l’on plonge $`Z`$ dans $`(\mathrm{\Gamma }(Z,𝒪_Z(D))^{})`$ en associant à tout $`zZ`$ l’hyperplan des sections globales qui s’annulent en $`z`$, le cône affine au-dessus de $`Z`$ dans $`\mathrm{\Gamma }(Z,𝒪_Z(D))^{}`$ est normal (car selon le §2.3, c’est le cône des vecteurs primitifs d’un $`G`$-module simple). Ce cône affine est donc le spectre de l’algèbre graduée
$$R:=\underset{d}{}\mathrm{\Gamma }(Z,𝒪_Z(dD)),$$
où l’on note $`𝒪_Z(dD)`$ le faisceau inversible sur $`Z`$ associé au diviseur $`dD`$, et $`\mathrm{\Gamma }(Z,𝒪_Z(dD))=:R_d`$ l’espace de ses sections globales. On note $`\sigma _DR_1`$ la section canonique de $`𝒪_Z(D)`$. L’algèbre $`R`$ est naturellement munie d’une structure de $`G`$-algèbre rationnelle ; on munit $`\stackrel{~}{Z}`$ de l’action de $`G`$ correspondante (qui induit celle de $`G`$ sur $`Z=\mathrm{Proj}R`$).
###### Proposition 2.3.
Il existe un poids dominant $`\lambda `$ tel que $`R_1`$ se décompose comme $`G`$-module sous la forme
$$R_1=\sigma _DV(\lambda )^{}.$$
Notons $`f`$ l’immersion fermée correspondant au morphisme surjectif d’algèbres $`\mathrm{Sym}(\sigma _DV(\lambda )^{})R`$ et $`\pi `$ le morphisme donné par la fonction régulière $`\sigma _D`$ : on a un diagramme commutatif de morphismes équivariants
$`𝔸^1`$ est muni de l’action triviale de $`G`$. De plus $`\pi :\stackrel{~}{Z}𝔸^1`$ est une famille de fonction de Hilbert $`h_\lambda `$ et la fibre de $`\pi `$ en $`0𝔸^1`$ est le cône des vecteurs primitifs de $`V(\lambda )`$. Les autres fibres $`\pi ^1(t)`$, $`t0`$ sont isomorphes à l’orbite ouverte $`\mathrm{\Omega }`$.
Preuve.
Comme $`D`$ est complet et homogène sous l’action de $`G`$, il est isomorphe à un quotient $`G/P`$, où $`P`$ est un sous-groupe parabolique de $`G`$ contenant $`B`$. On a donc un plongement $`i:G/PZ`$. L’image réciproque de $`𝒪_Z(D)`$ par $`i`$ est un faisceau inversible ample sur $`G/P`$ : elle est donc isomorphe au faisceau $`_\lambda `$ pour un certain poids dominant $`\lambda `$.
On a une suite exacte de $`𝒪_Z`$-modules :
$$0𝒪_Z(D)\stackrel{\sigma _D}{}𝒪_Zi_{}𝒪_{G/P}0.$$
On la tensorise par $`𝒪_Z(dD)`$ :
$$0𝒪_Z((d1)D)\stackrel{\sigma _D}{}𝒪_Z(dD)i_{}_{d\lambda }0.$$
On a donc une suite exacte de $`G`$-modules de sections globales :
$$0R_{d1}\stackrel{\sigma _D}{}R_d\stackrel{f_d}{}V(d\lambda ^{}).$$
(5)
Lorsque $`d=1`$, la suite (5) est
$$0\stackrel{\sigma _D}{}R_1\stackrel{f_1}{}V(\lambda ^{}).$$
Comme $`Z`$ se plonge dans $`(R_1^{})`$, l’espace vectoriel $`R_1`$ n’est pas de dimension $`1`$. De plus $`V(\lambda )^{}`$ est un $`G`$-module simple, donc le morphisme $`f_1`$ est surjectif, et on en déduit le premier point de la proposition.
Montrons que le morphisme $`f_d`$ est surjectif pour tout entier $`d`$. Comme $`R_1`$ contient un $`B`$-vecteur propre de poids $`\lambda ^{}`$, $`R_d`$ contient un $`B`$-vecteur propre de poids $`d\lambda ^{}`$, et il contient donc un $`G`$-module simple isomorphe à $`V(d\lambda ^{})`$. Or en considérant la suite exacte (5) pour tout $`d^{}<d`$, on remarque que $`R_{d1}`$ est un sous-$`G`$-module de $`_{d=0}^{d1}V(d\lambda ^{})`$, d’où l’assertion.
(On peut retrouver la surjectivité des $`f_d`$ par un argument cohomologique. En effet, selon le §2.3, la variété $`Z`$ est une variété de drapeaux pour l’action d’un groupe $`G^{}`$ réductif connexe, que l’on peut supposer simplement connexe quitte à le remplacer par un revêtement fini. L’algèbre des fonctions régulières sur $`G^{}`$ est alors factorielle, et le groupe de Picard de $`G^{}`$ est nul. Selon \[KnKrVu\], prop 3.2 (i), tout faisceau inversible sur $`Z`$ est donc linéarisable. On peut donc appliquer le théorème de Borel-Weil-Bott (\[Ak2\] p 113 ) au faisceau inversible $`𝒪_Z((d1)D)`$. Comme celui-ci est ample, on obtient $`H^1(Z,𝒪_Z((d1)D))=0`$, d’où le résultat.)
La fibre de $`\pi `$ au dessus de $`0`$ est le sous-cône de $`\stackrel{~}{Z}`$ d’algèbre affine graduée $`R/\sigma _DR`$. D’après ce qui précède, on a un isomorphisme de $`G`$-modules
$$R/\sigma _DR\underset{d}{}V(d\lambda ^{}).$$
La fibre au dessus de $`0`$ est donc le cône des vecteurs primitifs de $`V(\lambda )`$, selon la proposition 1.3(a).
La fibre au dessus de $`t0`$ est la section de $`Z`$ par l’hyperplan affine $`\{\sigma _D=t\}`$, donc est isomorphe à l’ouvert $`\mathrm{\Omega }`$. Enfin le morphisme $`\pi `$ est plat car $`R`$ est un $`[\sigma _D]`$-module sans torsion : le morphisme $`\pi `$ est donc bien une famille de fonction de Hilbert $`h_\lambda `$. $`\mathrm{}`$
###### Remarque 2.4.
La déformation $`\pi `$ ainsi obtenue est toujours non triviale, car les fibres $`\{\pi ^1(t),t0\}`$ sont homogènes pour l’action de $`G`$, contrairement à $`\pi ^1(0)`$.
On associe ainsi à chaque objet de (A) un schéma de Hilbert invariant de la classification (H), et on constate que l’on obtient ainsi toute cette classification :
* Le cas (H2) provient du cas (A1), avec $`n2`$.
* Le cas (H1) (resp. (H4), (H6), (H9)) provient des cas (A2) et (A3), avec $`n=3`$ (resp. $`n\text{ impair supérieur à }5`$, $`n\text{ pair supérieur à }6`$, $`n=4`$).
* Le cas (H3) provient des cas (A8) et (A9).
* Le cas (H5) provient du cas (A4), avec $`n3`$.
* Le cas (H7) provient du cas (A5).
* Le cas (H8) provient des cas (A6) et (A7).
En particulier, on en déduit dans chacun des cas l’existence d’un point de $`H_\lambda `$ distinct de $`C_\lambda `$, comme annoncé au §1.4.
### 2.5 Construction des familles universelles
Plaçons-nous dans l’un des cas du théorème 1.1 : le schéma de Hilbert invariant $`H_\lambda `$ est isomorphe à la droite affine. Il résulte des §2.3 et 2.4 qu’on obtient une déformation de $`C_\lambda `$ à partir d’une (unique) algèbre de Jordan simple $`A`$, de la façon suivante. On note $`\stackrel{~}{Z}`$ le cône des éléments de rang $`1`$ de $`A`$, et $`i`$ l’inclusion $`\stackrel{~}{Z}A`$. Le morphisme donné par la restriction de la trace de $`A`$ à $`\stackrel{~}{Z}`$ est noté $`\pi _\lambda `$. Le diagramme analogue à celui de la proposition 2.3 est alors :
###### Proposition 2.5.
Si $`A`$ n’est pas de type (J1), la famille $`\pi _\lambda `$ est la famille universelle au dessus de $`H_\lambda 𝔸^1`$.
Sinon on a $`A=W`$ ; la famille universelle est alors, avec des notations évidentes
Preuve. Les familles de la proposition sont les images inverses par un (unique) morphisme $`f:𝔸^1H_\lambda `$ de la famille universelle (car ce sont bien des familles de fonction de Hilbert $`h_\lambda `$).
On va montrer que $`f`$ est injectif ; comme $`H_\lambda `$ est isomorphe à $`𝔸^1`$, on en conclura que $`f`$ est un isomorphisme, et la proposition sera démontrée.
L’injectivité de $`f`$ signifie que les fibres des familles de sous-schémas considérées sont deux à deux distinctes.
Cela est clair dans le cas où $`A`$ est de type (J1) : la fibre de la famille au dessus de $`t𝔸^1`$ est la sous-variété $`\{wW|t=w,w\}`$.
Dans le cas où $`A`$ n’est pas de type (J1), on vérifie que son élément unité n’est pas somme de deux éléments de rang $`1`$.
Le module simple $`V(\lambda )`$ est l’espace $`V`$ de la décomposition $`(2)`$. La fibre de $`\pi _\lambda `$ au dessus de $`t𝔸^1`$ est
$$\{at.1|a\stackrel{~}{Z}\text{ et }\mathrm{tr}(a)=t\}V.$$
Supposons que les fibres de $`\pi _\lambda `$ au dessus de $`t`$ et $`t^{}`$ soient égales : on peut alors écrire $`at.1=a^{}t^{}.1`$, donc $`(tt^{}).1=a^{}a`$, donc $`t=t^{}`$.
Le morphisme $`f`$ est donc bien injectif. $`\mathrm{}`$
Le second cas de la proposition correspond aux cas (H1),(H4),(H6),(H9)$`G=\mathrm{SO}(W)`$ et $`V(\lambda )=W`$, ainsi qu’au cas (H8)$`G=G_2`$ et $`V(\lambda )=W`$ est de dimension $`7`$, et au cas (H3)$`G=\mathrm{Spin}(7)`$ et $`V(\lambda )=W`$ est de dimension $`8`$.
## 3 Rigidité des cônes de vecteurs primitifs
On sait (\[Ha\], ex 9.8 p 267) que les déformations infinitésimales de $`C_\lambda `$ sont classifiées par un $``$-espace vectoriel noté $`T^1(C_\lambda )`$. Comme $`C_\lambda `$ est une $`G`$-variété, l’espace vectoriel $`T^1(C_\lambda )`$ est un $`G`$-module rationnel (\[Ri\]) ; on le note dans la suite $`T_\lambda ^1`$.
On voit facilement que l’espace des éléments $`G`$-invariants de $`T_\lambda ^1`$ est en fait l’espace tangent en $`C_\lambda `$ au schéma de Hilbert invariant (proposition 3.5) ; il est donc déterminé par le théorème 1.1. Dans cette partie, on détermine complètement le $`G`$-module $`T_\lambda ^1`$ :
###### Théorème 3.1.
L’espace $`T_\lambda ^1`$ des déformations infinitésimales de $`C_\lambda `$ est nul, sauf dans les cas suivants :
1. Si l’on est dans les cas (H2) à (H9) du théorème 1.1, alors $`T_\lambda ^1=V(0)`$.
2. Si $`G=\mathrm{SL}(2)`$ et $`m2`$ est un entier, alors $`T_m^1=V(m2)V(m4)`$ (on indexe les poids de $`\mathrm{SL}(2)`$ par les entiers).
3. Si le groupe s’écrit $`G=\mathrm{SL}(2)\times H`$ et le module simple $`V(\lambda )=V_{\mathrm{SL}(2)}(m)W,`$$`V_{\mathrm{SL}(2)}(m)`$ est le $`\mathrm{SL}(2)`$-module simple de plus grand poids $`m`$, et $`(H,W)=(\mathrm{SL}(V),V)`$, $`(\mathrm{SL}(V),V^{})`$ ou $`(\mathrm{Sp}(V),V)`$ (dans ce dernier cas, $`V`$ est un espace vectoriel de dimension paire supérieure ou égale à $`2`$), alors $`T_\lambda ^1=V_{\mathrm{SL}(2)}(m2)W`$.
et dans les cas obtenus à partir d’un cas parmi les précédents par factorisation.
###### Remarques 3.2.
1. On retrouve ainsi des faits déjà connus :
* Le cône affine sur le plongement de Segre de la variété $`^m\times ^n`$ (dans $`^{(m+1)(n+1)1}`$) est rigide quand $`m+n3`$ (voir par exemple \[KlLa\] thm 2.2.8).
* Pinkham a déterminé dans \[Pi\] l’espace $`T_m^1`$ quand $`C_m`$ est le cône affine sur la courbe rationnelle normale de degré $`m`$ dans $`^m`$ ; en particulier, il a montré que $`\mathrm{dim}(T_m^1)=2m4`$ (pour $`m4`$). Il a aussi montré que si $`m5`$, la déformation verselle est irréductible, de dimension $`m1`$, et lisse hors de l’origine. Si $`m=4`$, elle a deux composantes de dimensions $`3`$ et $`1`$ qui se rencontrent transversalement à l’origine.
Le cône de vecteurs primitifs $`C_4`$ est exceptionnel, car c’est le seul dont l’espace des déformations infinitésimales admet à la fois une partie $`G`$-invariante et une partie non invariante : $`T_4^1=V(0)V(2)`$. La direction $`G`$-invariante de l’espace $`T_4^1`$ correspond à la composante de dimension $`1`$ de la déformation verselle.
Ainsi, on constate que dans tous les cas où le schéma de Hilbert invariant $`H_\lambda `$ n’est pas réduit à un point, il donne une composante irréductible de la déformation verselle de $`C_\lambda `$.
* Svanes a montré dans \[Sv1\] et \[Sv2\] que les cônes affines sur les plus petits plongements des variétés de drapeaux de $`\mathrm{SL}(n)`$ (qui correspondent au cas où $`G`$ est simple de type $`A_n`$ et $`\lambda `$ est une somme de poids fondamentaux) sont rigides, à l’exception des cas (H2) et, pour $`n=3`$, (H6) du théorème 1.1.
2. Les couples $`(H,W)`$ du troisième cas du théorème peuvent être décrits géométriquement : ce sont ceux où le groupe $`H`$ agit transitivement sur les droites du module $`W`$ (cela résulte par exemple de \[Ak2\] thm2 p75).
Enfin on détermine les déformations verselles des cônes $`C_\lambda `$. Si on est dans le cas (R1) du théorème 3.1, la déformation verselle de $`C_\lambda `$ est donnée par le schéma de Hilbert invariant $`H_\lambda `$. Dans le cas (R2), elle a été déterminée dans \[Pi\] par des équations explicites.
Plaçons-nous maintenant dans le cas (R3) du théorème 3.1. On va décrire la déformation verselle de $`C_\lambda `$ à l’aide d’équations analogues à celles de \[Pi\]. La dimension de l’espace vectoriel $`W`$ est $`n+1`$, pour un entier non nul $`n`$. Le cône $`C_\lambda `$ est le cône affine dans $`𝔸^{(m+1)(n+1)}`$ au dessus de l’image de $`^1\times ^n`$ par le plongement
$$\begin{array}{ccc}^1\times ^n& \hfill & ^{(m+1)(n+1)1}\hfill \\ ([a:b],[c_0,\mathrm{},c_n])& \hfill & [a^jb^{mj}c_i]_{_{i,j}}\hfill \end{array}$$
On note ici ce cône $`C_{mn}`$. Son idéal homogène est engendré par les mineurs $`2\times 2`$ de la matrice
$$\left(\begin{array}{ccccccccccccc}x_{00}& x_{01}& \mathrm{}& x_{0m1}& x_{10}& x_{11}& \mathrm{}& x_{1m1}& \mathrm{}& x_{n0}& x_{n1}& \mathrm{}& x_{nm1}\\ x_{01}& x_{02}& \mathrm{}& x_{0m}& x_{11}& x_{12}& \mathrm{}& x_{1m}& \mathrm{}& x_{n1}& x_{n2}& \mathrm{}& x_{nm}\end{array}\right).$$
###### Proposition 3.3.
La déformation verselle de $`C_{mn}`$ est la déformation $`𝔙`$ au dessus du spectre de l’anneau $`[[𝐭]]`$ des séries formelles en les $`t_{i,j}`$ (où $`0in`$ et $`1jm1`$) définie par les mineurs $`2\times 2`$ de la matrice
$$\left(\begin{array}{ccccccccc}x_{00}& \mathrm{}& x_{0m2}& x_{0m1}& \mathrm{}& x_{n0}& \mathrm{}& x_{nm2}& x_{nm1}\\ x_{01}t_{01}& \mathrm{}& x_{0m1}t_{0m1}& x_{0m}& \mathrm{}& x_{n1}t_{n1}& \mathrm{}& x_{nm1}t_{nm1}& x_{nm}\end{array}\right)$$
(6)
On démontre le théorème 3.1 dans les parties 3.1 à 3.3, et la proposition 3.3 dans la partie 3.4. On rappelle d’abord quelques faits connus.
Notons $`𝒯_{C_\lambda }`$ et $`𝒯_{V(\lambda )}`$ les faisceaux tangents respectifs de $`C_\lambda `$ et $`V(\lambda )`$, et $`𝒩_{C_\lambda }`$ le faisceau normal de $`C_\lambda `$ dans $`V(\lambda )`$. On a $`𝒯_{V(\lambda )}=𝒪_{V(\lambda )}V(\lambda )`$. Comme $`C_\lambda `$ est normal, $`𝒯_{C_\lambda }`$ et $`𝒩_{C_\lambda }`$ sont des faisceaux réflexifs.
Les déformations infinitésimales de $`C_\lambda `$ se plongent en fait toutes dans $`V(\lambda )`$, et l’on a une suite exacte (\[Ha\], ex 9.8 p 267)
$$0H^0(C_\lambda ,𝒯_{C_\lambda })H^0(C_\lambda ,𝒯_{V(\lambda )}|_{C_\lambda })H^0(C_\lambda ,𝒩_{C_\lambda })T_\lambda ^10,$$
c’est-à-dire
$$0H^0(C_\lambda ,𝒯_{C_\lambda })H^0(C_\lambda ,𝒪_{C_\lambda })_{}V(\lambda )H^0(C_\lambda ,𝒩_{C_\lambda })T_\lambda ^10.$$
(7)
On peut supposer $`C_\lambda `$ de dimension supérieure ou égale à $`2`$ ; on note $`E_\lambda :=C_\lambda \{0\}`$ le cône épointé. La suite exacte ci-dessus s’identifie alors à la suivante (avec des notations analogues)
$$0H^0(E_\lambda ,𝒯_{E_\lambda })H^0(E_\lambda ,𝒪_{E_\lambda })_{}V(\lambda )H^0(E_\lambda ,𝒩_{E_\lambda })T_\lambda ^10.$$
(8)
Or comme $`E_\lambda `$ est lisse, on a la suite exacte courte
$$0𝒯_{E_\lambda }𝒪_{E_\lambda }_{}V(\lambda )𝒩_{E_\lambda }0.$$
On en déduit la proposition suivante, due à Schlessinger (\[Sc\]) :
###### Proposition 3.4.
On a une suite exacte :
$$0T_\lambda ^1H^1(E_\lambda ,𝒯_{E_\lambda })H^1(E_\lambda ,𝒪_{E_\lambda })_{}V(\lambda ).$$
### 3.1 Préliminaires
On commence par déterminer la partie invariante de l’espace $`T_\lambda ^1`$ :
###### Proposition 3.5.
On a un isomorphisme canonique
$$(T_\lambda ^1)^GT_{C_\lambda }H_\lambda .$$
Ainsi, l’espace $`(T_\lambda ^1)^G`$ est nul, sauf dans les cas (H1) à (H9) du théorème 1.1, où il est de dimension 1.
Preuve.
En prenant les $`G`$-invariants de (7), on obtient la suite exacte de $`G`$-modules de \[AlBr\], prop 1.13 :
$$0H^0(C_\lambda ,𝒯_{C_\lambda })^G(H^0(C_\lambda ,𝒪_{C_\lambda })V(\lambda ))^GT_{C_\lambda }H_\lambda (T_\lambda ^1)^G0,$$
qui s’écrit dans notre cas (\[AlBr\], prop 1.15 (iii)) :
$$0[𝔤.v_\lambda ]^{G_{v_\lambda }}V(\lambda )^{G_{v_\lambda }}T_{C_\lambda }H_\lambda (T_\lambda ^1)^G0.$$
Comme ses deux premiers termes sont de dimension $`1`$ :
$$[𝔤.v_\lambda ]^{G_{v_\lambda }}=V(\lambda )^{G_{v_\lambda }}=v_\lambda ,$$
on en déduit le résultat. $`\mathrm{}`$
Ainsi, on a montré que la partie $`G`$-invariante était bien celle annoncée dans le théorème 3.1. On va maintenant déterminer les autres composantes isotypiques de $`T_\lambda ^1`$.
On note $`X_\lambda `$ la variété de drapeaux $`G/P_\lambda `$, et $`\pi :E_\lambda X_\lambda `$ la surjection naturelle. On remarque que $`E_\lambda `$ est l’espace total du faisceau $`_\lambda `$ privé de la section nulle, donc $`\pi `$ est un morphisme affine lisse.
###### Proposition 3.6.
On a une suite exacte de faisceaux $`G`$-linéarisés sur $`X_\lambda `$ :
$$0\underset{d}{}_{d\lambda }\pi _{}𝒯_{E_\lambda }\underset{d}{}_{d\lambda }𝒯_{X_\lambda }0,$$
donc une suite exacte de $`G`$-modules :
$$\underset{d}{}H^1(X_\lambda ,_{d\lambda })H^1(E_\lambda ,𝒯_{E_\lambda })\underset{d}{}H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })\underset{d}{}H^2(X_\lambda ,_{d\lambda }).$$
(9)
Preuve.
Comme le morphisme $`\pi `$ est lisse, on a la suite exacte courte
$$0𝒯_\pi 𝒯_{E_\lambda }\pi ^{}𝒯_{X_\lambda }0.$$
On remarque que le faisceau $`𝒯_\pi `$ tangent à $`\pi `$ est isomorphe à $`𝒪_{E_\lambda }`$. Puis, comme $`\pi `$ est affine, on a
$$0\pi _{}𝒯_\pi \pi _{}𝒯_{E_\lambda }\pi _{}\pi ^{}𝒯_{X_\lambda }0,$$
d’où la suite exacte de faisceaux annoncée, car $`\pi _{}𝒯_\pi _d_{d\lambda }`$, et, selon la formule de projection, $`\pi _{}\pi ^{}𝒯_{X_\lambda }_d_{d\lambda }𝒯_{X_\lambda }`$.
La suite exacte de $`G`$-modules donnée en découle aussi, car, comme $`\pi `$ est affine,
$`H^1(E_\lambda ,𝒯_{E_\lambda })H^1(X_\lambda ,\pi _{}𝒯_{E_\lambda })`$. $`\mathrm{}`$
La suite exacte (9) nous permettra de démontrer le théorème 3.1 dans certains cas, à l’aide du théorème de Borel-Weil-Bott (\[Ak2\], thm p113).
Notons $`Q_\lambda `$ l’unique sous-groupe parabolique de $`G`$ conjugué à $`P_\lambda `$ et contenant le sous-groupe de Borel $`B^{}`$ opposé à $`B`$ ; on voit naturellement $`Q_\lambda `$ comme un point de $`X_\lambda =G/P_\lambda `$. Regardons, pour appliquer le théorème de Borel-Weil-Bott, quels sont les poids de l’action de $`T`$ sur les fibres $`_{d\lambda }|_{\{Q_\lambda \}}`$ et $`𝒯_{X_\lambda }|_{\{Q_\lambda \}}𝔤/𝔮_\lambda `$ : le tore $`T`$ agit sur le premier espace avec le poids $`d\lambda ^{}`$ ; les poids du second sont les racines positives de $`G`$ qui ne sont pas des racines de $`Q_\lambda `$.
Si $`W`$ est un $`Q_\lambda `$-module rationnel, on note $`𝒱(W)`$ le $`𝒪_{X_\lambda }`$-module $`G`$-linéarisé dont la fibre en $`Q_\lambda `$ est $`W`$. On rappelle (voir \[Gros\] ou \[Jan\]) que l’espace des sections globales de $`𝒱(W)`$ est le $`G`$-module induit par le $`Q_\lambda `$-module $`W`$ :
$$H^0(X_\lambda ,𝒱(W))=\mathrm{Ind}_{Q_\lambda }^G(W)$$
et les groupes de cohomologie de $`𝒱(W)`$ donnent les foncteurs dérivés à droite du foncteur $`\mathrm{Ind}_{Q_\lambda }^G`$ :
$$H^j(X_\lambda ,𝒱(W))=R^j\mathrm{Ind}_{Q_\lambda }^G(W).$$
Dans toute la suite, on note simplement $`\mathrm{Ind}`$ le foncteur $`\mathrm{Ind}_{Q_\lambda }^G`$.
Si $`\mu `$ est un caractère de $`Q_\lambda `$, on définit le $`Q_\lambda `$-module
$$W[\mu ]:=_\mu _{}W,$$
$`_\mu `$ est la droite où $`Q_\lambda `$ opère avec le poids $`\mu `$.
On note enfin $``$ l’action tordue du groupe de Weyl sur $`\mathrm{\Lambda }`$ : si $`w`$ est un élément du groupe de Weyl, et $`\mu `$ un poids, on pose $`w\mu :=w(\mu +\rho )\rho `$, où $`\rho `$ est la demi-somme des racines positives.
###### Proposition 3.7.
Soit $`d`$. Si l’espace $`H^1(X_\lambda ,_{d\lambda })`$ est non nul, alors $`V(\lambda )`$ est en fait un $`\mathrm{SL}(2)`$-module, et l’action de $`G`$ se factorise sous la forme $`G\mathrm{SL}(2)\mathrm{GL}(V(\lambda )).`$
Preuve. Si $`d0`$, on sait que tous les groupes de cohomologie de $`_{d\lambda }`$ sont nuls sauf en degré $`0`$.
Supposons $`d<0`$. Selon le théorème de Borel-Weil-Bott, $`H^1(X_\lambda ,_{d\lambda })`$ est nul ou irréductible, et on a $`H^1(X_\lambda ,_{d\lambda })V(\mu )`$ si et seulement si il existe une racine simple $`\alpha `$ et un poids dominant $`\mu `$ tels que
$$s_\alpha \mu =d\lambda ^{},$$
en notant $`s_\alpha `$ la réflexion simple associée à $`\alpha `$. On en déduit
$$\mu d\lambda ^{}=(1+\alpha ^{},\mu )\alpha .$$
La racine simple $`\alpha `$ est donc un poids dominant : c’est une racine simple de $`G`$ isolée dans son diagramme de Dynkin. Notons $`\omega _\alpha =\alpha /2`$ le poids fondamental associé à $`\alpha `$. Le poids $`d\lambda ^{}`$ est proportionnel à $`\omega _\alpha `$, et $`\lambda =\lambda ^{}`$ aussi, d’où le résultat. $`\mathrm{}`$
###### Proposition 3.8.
Lorsque $`d0`$, on a $`H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })=0`$.
Preuve. C’est une conséquence du fait suivant (cf. \[Bro\], thm2.2) : notons $`T^{}X_\lambda `$ le fibré cotangent de $`X_\lambda =G/P_\lambda `$, et $`p:T^{}X_\lambda X_\lambda `$ la projection canonique. On a, pour tout $`i1`$,
$$H^i(T^{}X_\lambda ,p^{}_{d\lambda })=0.$$
Selon la formule de projection,
$$f_{}f^{}_{d\lambda }=_{d\lambda }f_{}𝒪_{T^{}X_\lambda }=_{d\lambda }\mathrm{Sym}(𝒯_{X_\lambda }),$$
$`\mathrm{Sym}(𝒯_{X_\lambda })=_k\mathrm{S}^k𝒯_{X_\lambda }`$ est l’algèbre symétrique du $`𝒪_{X_\lambda }`$-module $`𝒯_{X_\lambda }`$. Ainsi,
$$\underset{k}{}H^i(X_\lambda ,_{d\lambda }\mathrm{S}^k𝒯_{X_\lambda })=0,$$
et la proposition en découle en prenant $`i=k=1`$. $`\mathrm{}`$
### 3.2 Cas où le groupe $`G`$ est simple
Dans ce paragraphe, on établit le théorème 3.1 dans le cas où $`G`$ est un groupe simple.
Le §3.2.1 concerne le cas où $`G=\mathrm{SL}(2)`$. L’espace des déformations infinitésimales a alors été déterminé dans \[Pi\] ; on donne cependant une preuve simple de ce résultat, qui a l’avantage de fournir la structure de $`\mathrm{SL}(2)`$-module de $`T_\lambda ^1`$.
Le cas des autres groupes simples est traité dans le §3.2.2.
#### 3.2.1 Cas où $`G=\mathrm{SL}(2)`$
On note $`H`$ le sous-groupe des matrices unipotentes triangulaires inférieures $`\left(\begin{array}{cc}1& 0\\ x& 1\end{array}\right)`$. On identifie le groupe des poids de $`\mathrm{SL}(2)`$ à $``$, de sorte que les poids dominants sont les éléments de $``$. Le poids dominant $`\lambda `$ est donc un entier, que l’on note ici $`m`$. On suppose $`m1`$ ; on a donc $`Q_\lambda =B^{}`$.
Pour déterminer $`T_\lambda ^1`$, on utilise la suite exacte (8). Comme $`\pi `$ est affine, on a la suite exacte
$$0H^0(X_\lambda ,\pi _{}𝒯_{E_\lambda })\stackrel{𝑓}{}H^0(X_\lambda ,\pi _{}𝒪_{E_\lambda })V(\lambda )\stackrel{𝑔}{}H^0(X_\lambda ,\pi _{}𝒩_{E_\lambda })T_\lambda ^10.$$
On remarque que les fibres respectives en $`Q_\lambda `$ des faisceaux $`G`$-linéarisés $`\pi _{}𝒯_{E_\lambda }`$, $`\pi _{}𝒪_{E_\lambda }V(\lambda )`$ et $`\pi _{}𝒩_{E_\lambda }`$ sont les $`Q_\lambda `$-modules gradués
$$\underset{d}{}𝔤v_\lambda [md]\text{ , }\underset{d}{}V(m)[md]\text{ et }\underset{d}{}V(m)/𝔤v_\lambda [md],$$
et que les morphismes $`f`$ et $`g`$ sont les images par le foncteur $`\mathrm{Ind}`$ des morphismes de modules gradués qui forment en chaque degré $`d`$ une suite exacte courte :
$$0𝔤v_\lambda [md]V(m)[md]V(m)/𝔤v_\lambda [md]0.$$
Considérons la suite exacte longue associée à cette dernière
$$0\mathrm{Ind}(𝔤v_\lambda [md])\stackrel{f_d}{}\mathrm{Ind}(V(m)[md])\stackrel{g_d}{}\mathrm{Ind}(V(m)/𝔤v_\lambda [md])R^1\mathrm{Ind}(𝔤v_\lambda [md])\mathrm{}$$
(10)
On a
$$T_\lambda ^1\underset{d}{}\mathrm{coker}g_d.$$
Pour connaître $`\mathrm{coker}g_d`$, on remarque qu’on a des isomorphismes de $`Q_\lambda `$-modules :
$$𝔤v_\lambda V(1)[m+1]\text{ et }V(m)/𝔤v_\lambda V(m2)[2],$$
donc
$$𝔤v_\lambda [md]V(1)[m(d1)+1]$$
et
$$V(m)/𝔤v_\lambda [md]V(m2)[md+2].$$
La suite exacte longue (10) devient alors
$$0V(1)V(m(d1)+1)\stackrel{f_d}{}V(m)V(md)\stackrel{g_d}{}V(m2)V(md+2)$$
$$R^1\mathrm{Ind}(V(1)[(m(d1)+1])\mathrm{}$$
On peut alors conclure :
Supposons $`m3`$ :
* si $`d<0`$, on a $`V(m2)V(md+2)=0`$, donc $`\mathrm{coker}g_d=0`$.
* si $`d1`$, montrons que $`R^1\mathrm{Ind}(V(1)[m(d1)+1])=0`$ : considérons la suite exacte de $`Q_\lambda `$-modules suivante
$$0[m(d1)]V(1)[m(d1)+1][m(d1)+2]0.$$
La suite exacte longue associée (relativement au foncteur $`\mathrm{Ind}`$) s’écrit
$$\mathrm{}0R^1\mathrm{Ind}(V(1)[m(d1)+1])0\mathrm{}$$
Ainsi le morphisme $`g_d`$ est surjectif : $`\mathrm{coker}g_d=0`$.
* si enfin $`d=0`$, on remarque que $`\mathrm{coker}g_0=V(m2)V(m4)`$.
Si $`m=2`$, le conoyau de $`g_d`$ est nul si $`d1`$, et on a $`\mathrm{coker}g_1=V(0)`$. Donc on a toujours $`T_\lambda ^1=V(m2)V(m4)`$.
#### 3.2.2 Autres groupes simples
On suppose maintenant que $`G`$ est un groupe simple de type autre que $`A_1`$, et il s’agit de montrer que les seules déformations infinitésimales de $`C_\lambda `$ sont celles qui proviennent du théorème 1.1.
Selon la proposition 3.5, il ne reste qu’à montrer que le groupe $`G`$ agit trivialement sur l’espace $`T_\lambda ^1`$ (ie tous les éléments de l’espace sont $`G`$-invariants).
En vertu de la proposition 3.7, la suite exacte (9) donne ici
$$H^1(E_\lambda ,𝒯_{E_\lambda })\underset{d}{}H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda }).$$
Il ne reste donc plus qu’à montrer que pour tout $`d`$ et pour tout poids dominant $`\mu `$ non nul, la composante isotypique $`H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })_{(\mu )}`$ est nulle.
Pour cela on va utiliser le lemme et la proposition suivants. Les notations sont celles de \[Bo\].
###### Lemme 3.9.
Soit $`R`$ un système de racines irréductible muni d’une base $`S`$. Soit $`\alpha `$ un élément de $`S`$. Alors il existe une racine positive longue $`\gamma `$ telle que $`\gamma ^{},\alpha =1`$, sauf dans les cas suivants :
* si $`R`$ est de type $`A_1`$.
* si $`R`$ est de type $`B_2`$ et $`\alpha =\alpha _1`$ est la racine simple longue.
* si $`R`$ est de type $`C_n`$, avec $`n3`$ et $`\alpha =\alpha _n`$ est la racine simple longue.
Lorsqu’une telle racine $`\gamma `$ existe, elle n’est pas unique, sauf dans les cas suivants :
* si $`R`$ est de type $`A_2`$.
* si $`R`$ est de type $`B_2`$ et $`\alpha =\alpha _2`$ est la racine simple courte : seule $`\gamma =\alpha _1`$ convient.
* si $`R`$ est de type $`C_n`$, avec $`n3`$ et $`\alpha =\alpha _i`$, $`i=1\mathrm{}n1`$ est une racine simple courte : seule $`\gamma =2\alpha _{i+1}+\mathrm{}+2\alpha _{n1}+\alpha _n`$ convient .
* si $`R`$ est de type $`G_2`$ : pour $`\alpha =\alpha _1`$, seule $`\gamma =\alpha _2`$ convient ; pour $`\alpha =\alpha _2`$, seule $`\gamma =3\alpha _1+\alpha _2`$ convient.
Preuve. Traitons d’abord le cas où toutes les racines de $`R`$ sont de même longueur. Si $`R`$ est de type $`A_1`$, il n’y a rien à prouver ; on suppose donc que $`R`$ est de type $`A_n`$ $`(n2)`$, $`D_n`$ $`(n4)`$, $`E_6`$, $`E_7`$ ou $`E_8`$. L’existence de $`\gamma `$ est claire : toute racine simple reliée à $`\alpha `$ dans le diagramme de Dynkin de $`R`$ convient. Montrons que si $`R`$ n’est pas de type $`A_2`$, $`\gamma `$ n’est pas unique. Cela est clair si $`\alpha `$ est reliée à plusieurs racines simples dans le diagramme de Dynkin. Sinon, notons $`\beta `$ la seule racine simple reliée à $`\alpha `$. Le sous-système de racines $`R^{}`$ de $`R`$ de base $`S^{}:=S\{\alpha \}`$ est un système de racines irréductible, de rang supérieur ou égal à $`2`$. On sait qu’il admet plusieurs racines $`\gamma `$ telles que $`\beta `$ a pour coefficient $`1`$ dans l’écriture de $`\gamma `$ dans la base $`S^{}`$ (voir par exemple \[Ak2\], prop 1 p 126). Toutes ces racines $`\gamma `$ conviennent clairement.
Supposons maintenant que $`R`$ est de type $`B_n`$, avec $`n3`$. Si $`\alpha `$ est une racine simple longue, on est ramené au premier cas, car les racines longues de $`R`$ forment un système de racines de type $`A_3`$ si $`n=3`$, et $`D_n`$ sinon. Sinon, $`\alpha =\alpha _n`$, et la racine $`\gamma =\alpha _i+\mathrm{}+\alpha _{n1}`$ convient pour tout $`i=1\mathrm{}n1`$.
Supposons que $`R`$ est de type $`F_4`$. Si $`\alpha `$ est une racine simple longue, on est ramené au premier cas, car les racines longues de $`R`$ forment un système de racines de type $`D_4`$. Si $`\alpha =\alpha _3=ϵ_4,`$ alors $`\gamma =ϵ_iϵ_4`$ convient pour tout $`i=\mathrm{1,2,3}.`$ Si $`\alpha =\alpha _4=(ϵ_1ϵ_2ϵ_3ϵ_4)/2,`$ alors $`\gamma =ϵ_i+ϵ_j`$ convient pour tout $`2i<j4`$.
Enfin, on vérifie aisément les assertions du lemme concernant $`B_2`$, $`C_n`$ ($`n3`$) et $`G_2`$. $`\mathrm{}`$
###### Proposition 3.10.
Soit $`R`$ un système de racines irréductible muni d’une base. On suppose que $`R`$ n’est pas de type $`A_1`$.
Soient $`\alpha `$ une racine simple, $`\beta `$ une racine positive, et $`N2`$ un entier, tels que $`N\alpha +\beta `$ est un poids dominant. Alors $`N=2`$, et on est dans l’un des cas suivants :
* si $`R`$ est de type $`A_2`$, on a $`2\alpha _1+\alpha _2=3\omega _1`$ et $`2\alpha _2+\alpha _1=3\omega _2`$.
* si $`R`$ est de type $`B_2`$, on a $`2\alpha _2+\alpha _1=2\omega _2`$.
* si $`R`$ est de type $`C_n`$, on a $`2\alpha _1+(2\alpha _2+2\alpha _3+\mathrm{}+2\alpha _{n1}+\alpha _n)=2\omega _1`$.
* si $`R`$ est de type $`G_2`$, on a $`2\alpha _1+\alpha _2=3\omega _1`$.
Preuve. Traitons d’abord le cas où il existe plusieurs racines positives longues $`\gamma `$ telles que $`\gamma ^{},\alpha =1`$. Soit $`\gamma `$ une telle racine, que l’on suppose distincte de $`\beta `$. On a $`\gamma ^{},N\alpha +\beta =N+\gamma ^{},\beta 0.`$ D’où $`N\gamma ^{},\beta 1`$, car $`\gamma `$ est une racine longue distincte de $`\beta `$.
Pour conclure, on étudie un à un les cas du lemme précédent où il n’existe pas de racine $`\gamma `$, ainsi que ceux où il existe une unique racine $`\gamma `$ (selon le premier point de la démonstration, on peut alors supposer $`\beta =\gamma `$). $`\mathrm{}`$
###### Proposition 3.11.
Le groupe $`G`$ agit trivialement sur $`H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })`$.
Preuve. Selon la proposition 3.8, on peut supposer $`d<0`$.
Afin d’appliquer le théorème de Borel-Weil-Bott, on considère une suite de Jordan-Hölder du $`Q_\lambda `$-module $`𝔤/𝔮_\lambda `$, c’est-à-dire une suite décroissante
$$𝔤/𝔮_\lambda =W_0W_1\mathrm{}W_r=0$$
de $`Q_\lambda `$-modules telle que les quotients $`W_i/W_{i+1}`$ sont des modules simples.
Soit $`\mu `$ un poids dominant non nul. On veut montrer que la composante isotypique
$$H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })_{(\mu )}=R^1\mathrm{Ind}(𝔤/𝔮_\lambda [d\lambda ^{}])_{(\mu )}$$
est nulle.
Pour tout $`i`$, on a une suite exacte
$$0W_{i+1}[d\lambda ^{}]W_i[d\lambda ^{}](W_i/W_{i+1})[d\lambda ^{}]0,$$
donc une suite exacte sur les composantes isotypiques
$$R^1\mathrm{Ind}(W_{i+1}[d\lambda ^{}])_{(\mu )}R^1\mathrm{Ind}(W_i[d\lambda ^{}])_{(\mu )}R^1\mathrm{Ind}((W_i/W_{i+1})[d\lambda ^{}])_{(\mu )}.$$
Il suffit donc de montrer que pour tout $`i`$, on a
$$R^1\mathrm{Ind}((W_i/W_{i+1})[d\lambda ^{}])_{(\mu )}=0$$
et la proposition sera prouvée.
Supposons le contraire : selon le théorème de Borel-Weil-Bott, il existe une racine simple $`\alpha `$ telle que
$$s_\alpha \mu =d\lambda ^{}+\beta ,$$
$`\beta `$ est le plus grand poids de $`W_i/W_{i+1}`$ (c’est donc une racine de $`G`$ qui n’est pas une racine de $`Q_\lambda `$). On en déduit
$$\mu d\lambda ^{}=N\alpha +\beta ,$$
en posant
$$N:=1+\alpha ^{},\mu 1$$
Comme $`d<0`$, le poids $`N\alpha +\beta `$ est dominant : la racine simple $`\alpha `$ et la racine positive $`\beta `$ sont donc données dans la liste du §1.3.1 (si $`N=1`$) ou de la proposition 3.10 (si $`N2`$). Ici, on remarque de plus que :
* Le poids $`N\alpha +\beta `$ est la somme de deux poids dominants $`\mu `$ et $`d\lambda ^{}`$ non nuls.
* On a $`\alpha ^{},\mu =0`$ si et seulement si $`N=1`$.
* On a $`\beta ^{},\lambda ^{}0`$ (car $`\beta `$ n’est pas une racine de $`Q_\lambda `$).
En consultant les deux listes, on constate immédiatement que cela est impossible (si $`G`$ n’est pas de type $`A_1`$). $`\mathrm{}`$
### 3.3 Cas où le groupe $`G`$ n’est pas simple
On suppose dans cette partie que le groupe $`G`$ est de la forme $`G^1\times G^2`$. Le sous-groupe de Borel $`B`$ de $`G`$ s’écrit $`B=B^1\times B^2`$ ; de même pour le tore maximal $`T=T^1\times T^2`$. La représentation $`V(\lambda )`$ de $`G`$ s’écrit $`V(\lambda )=V(\lambda _1)V(\lambda _2)`$, où l’on suppose les poids dominants respectifs $`\lambda _1`$ et $`\lambda _2`$ de $`G^1`$ et $`G^2`$ tous les deux non nuls.
On note $`P_{\lambda _i}`$ le stabilisateur dans $`G^i`$ de la droite des vecteurs de plus grand poids de $`V(\lambda _i)`$. On a $`P_\lambda =P_{\lambda _1}\times P_{\lambda _2}`$.
Notre variété de drapeaux est donc un produit $`X_\lambda =X_{\lambda _1}\times X_{\lambda _2}`$, où l’on note $`X_{\lambda _i}`$ la variété de drapeaux $`G^i/P_{\lambda _i}`$. On note $`p_i:X_\lambda X_{\lambda _i}`$ les projections canoniques.
On a
$$_\lambda =p_1^{}_{\lambda _1}p_2^{}_{\lambda _2},$$
et
$$𝒯_{X_\lambda }=p_1^{}𝒯_{X_{\lambda _1}}p_2^{}𝒯_{X_{\lambda _2}}.$$
On remarque que selon la proposition 3.7
$$H^1(E_\lambda ,𝒪_{E_\lambda })=\underset{d}{}H^1(X_\lambda ,_{d\lambda })=0.$$
La proposition 3.4 donne donc
$$T_\lambda ^1H^1(E_\lambda ,𝒯_{E_\lambda }).$$
Cet isomorphisme est aussi conséquence de \[Sern\] prop II.5.8 (ii). En effet, $`\mathrm{dim}(C_\lambda )=\mathrm{dim}(X_{\lambda _1})+\mathrm{dim}(X_{\lambda _2})+13`$. Comme $`C_\lambda `$ est Cohen-Macaulay \[Ra\], sa profondeur en $`0`$ est supérieure ou égale à $`3`$.
#### 3.3.1 Cas où $`G=\mathrm{SL}(2)\times \mathrm{SL}(2)`$
On note encore $`H`$ le groupe des matrices unipotentes triangulaires inférieures de taille $`2\times 2`$.
On écrit le poids dominant de $`G`$ sous la forme $`\lambda =(m,n)`$, où l’on peut supposer les entiers $`m`$ et $`n`$ tels que $`mn1`$.
Pour calculer $`H^1(E_\lambda ,𝒯_{E_\lambda })H^1(X_\lambda ,\pi _{}𝒯_{E_\lambda })`$, on va utiliser une résolution du faisceau $`G`$-linéarisé $`\pi _{}𝒯_{E_\lambda }`$. Sa fibre en $`Q_\lambda `$ est le $`Q_\lambda `$-module
$$\underset{d}{}𝔤/𝔤_{v_\lambda }[md,nd].$$
D’où
$$H^1(E_\lambda ,𝒯_{E_\lambda })=\underset{d}{}R^1\mathrm{Ind}(𝔤/𝔤_{v_\lambda }[md,nd]).$$
Soit $`d`$. On remarque que
$$𝔤_{v_\lambda }=(𝔥\times 𝔥)𝔱_{v_\lambda },$$
où le stabilisateur $`𝔱_{v_\lambda }`$ de $`v_\lambda `$ dans $`𝔱`$ est une droite $`Q_\lambda `$-invariante.
On a donc une suite exacte de $`Q_\lambda `$-modules
$$0[md,nd]𝔤/(𝔥\times 𝔥)[md,nd]𝔤/𝔤_{v_\lambda }[md,nd]0.$$
D’où une suite exacte longue
$$\mathrm{}R^1\mathrm{Ind}([md,nd])R^1\mathrm{Ind}(𝔤/(𝔥\times 𝔥)[md,nd])R^1\mathrm{Ind}(𝔤/𝔤_{v_\lambda }[md,nd])$$
$$R^2\mathrm{Ind}([md,nd])\stackrel{h_d}{}R^2\mathrm{Ind}(𝔤/(𝔥\times 𝔥)[md,nd])\mathrm{}$$
###### Proposition 3.12.
1. L’espace $`R^1\mathrm{Ind}([md,nd])`$ est nul pour tout $`d`$.
2. L’espace $`R^1\mathrm{Ind}(𝔤/(𝔥\times 𝔥)[md,nd])`$ est nul sauf si $`d=1`$ et $`n=1`$. Dans ce cas il vaut $`V(m\mathrm{2,1})`$.
3. Le noyau de $`h_d`$ est nul, sauf si $`d=1`$ et $`(m,n)=(\mathrm{2,2})`$ et si $`d=2`$ et $`(m,n)=(\mathrm{1,1})`$. Dans ces deux cas (qui correspondent au cas (H9) du théorème 1.1) il vaut $`V(\mathrm{0,0})`$.
Preuve.
(1) Cela découle immédiatement du théorème de Borel-Weil-Bott (ou simplement de la cohomologie des faisceaux inversibles sur $`^1\times ^1`$).
(2) On remarque qu’on a l’isomorphisme de $`Q_\lambda `$-modules
$$𝔤/(𝔥\times 𝔥)[md,nd]V(\mathrm{0,1})[md,nd+1]V(\mathrm{1,0})[md+1,nd],$$
avec par exemple
$$R^1\mathrm{Ind}(V(\mathrm{0,1})[md,nd+1])=V(\mathrm{0,1})R^1\mathrm{Ind}([md,nd+1]).$$
Selon le théorème de Borel-Weil-Bott, pour que l’espace $`R^1\mathrm{Ind}([md,nd+1])`$ soit non nul, il faut que les entiers $`md`$ et $`nd+1`$ soient l’un positif, l’autre strictement négatif. Comme $`md`$ et $`nd`$ sont de même signe, on a nécessairement $`md<0`$ et $`nd+10`$ (donc $`d=1`$ et $`n=1`$). Dans ce cas,
$$R^1\mathrm{Ind}(V(\mathrm{0,1})[md,nd+1])=V(\mathrm{0,1})V(m\mathrm{2,0})=V(m\mathrm{2,1}).$$
Il faut donc que l’on ait $`m2`$.
De la même façon, on obtient que $`R^1\mathrm{Ind}(V(\mathrm{1,0})[md+1,nd])`$ est toujours nul, car on a supposé $`mn`$.
(3) On va en fait déterminer le conoyau de l’application transposée$`{}_{}{}^{t}h_{d}^{}`$. Selon le théorème de dualité de Serre (\[Ha\] III.7), on a un isomorphisme fonctoriel
$$R^2\mathrm{Ind}(W)^{}\mathrm{Ind}(W^{}[2,2])$$
pour tout $`Q_\lambda `$-module $`W`$ (car la fibre en $`Q_\lambda `$ du faisceau anticanonique de $`X_\lambda `$ est $`[2,2]`$).
Supposons que le conoyau de l’application
$$\mathrm{Ind}([𝔤/(𝔥\times 𝔥)]^{}[md2,nd2])\stackrel{{}_{}{}^{t}h_{d}^{}}{}\mathrm{Ind}([md2,nd2]).$$
est non nul.
Pour que l’espace d’arrivée de $`{}_{}{}^{t}h_{d}^{}`$ soit non nul, il faut que $`M:=md2`$ et $`N:=nd2`$ soient positifs. Cet espace est alors le module simple $`V(M,N)`$, et il faut que l’application $`{}_{}{}^{t}h_{d}^{}`$ soit nulle.
Comme dans la démonstration de (2), on remarque que l’espace de départ de $`{}_{}{}^{t}h_{d}^{}`$ est une somme directe :
$$\mathrm{Ind}([𝔤/(𝔥\times 𝔥)]^{}[md2,nd2])V(\mathrm{0,1})[M,N1]V(\mathrm{1,0})[M1,N].$$
Il faut donc que les deux composantes de $`{}_{}{}^{t}h_{d}^{}`$ soient nulles. La première composante est l’image par $`\mathrm{Ind}`$ du morphisme $`j`$ de la suite exacte courte suivante
$$0\stackrel{}{}[M,N2]\stackrel{}{}V(\mathrm{0,1})[M,N1]\stackrel{𝑗}{}[M,N]\stackrel{}{}0.$$
Le conoyau de la première composante se plonge donc dans l’espace $`R^1\mathrm{Ind}([M,N2])`$. Selon le théorème de Borel-Weil-Bott, pour que ce dernier espace soit non nul, il faut $`N22`$, donc $`N=0`$. On montre de même, à l’aide de la seconde composante de $`{}_{}{}^{t}h_{d}^{}`$, que $`M=0`$.
Enfin, dans le cas où $`(M,N)=(\mathrm{0,0})`$, le conoyau de $`0\stackrel{{}_{}{}^{t}h_{d}^{}}{}V(0)`$ est bien $`V(0)`$. $`\mathrm{}`$
#### 3.3.2 Autres cas
Lorsque les actions de $`G^1`$ et $`G^2`$ se factorisent par $`\mathrm{SL}(2)`$ :
$$G^i\mathrm{SL}(2)\mathrm{GL}(V(\lambda _i)),$$
on est dans la situation du §3.3.1 ; on suppose donc que l’action de $`G^2`$ ne se factorise pas par $`\mathrm{SL}(2)`$.
###### Proposition 3.13.
1. L’espace $`H^2(X_\lambda ,_{d\lambda })`$ est nul pour tout entier $`d`$.
2. On a donc un isomorphisme $`T_\lambda ^1_dH^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })`$.
Preuve.
(1) On rappelle que si $`d0`$, tous les groupes de cohomologie de $`_{d\lambda }`$ sont nuls sauf en degré $`0`$, et si $`d<0`$, le groupe de cohomologie de degré $`0`$ est nul. On peut donc supposer $`d<0`$, et on a, selon la formule de Künneth (\[Da\] p32)
$$H^2(X_\lambda ,_{d\lambda })=H^1(X_{\lambda _1},_{d\lambda _1})H^1(X_{\lambda _2},_{d\lambda _2}).$$
Or selon la proposition 3.7 et l’hypothèse faite sur $`\lambda _2`$, on a $`H^1(X_{\lambda _2},_{d\lambda _2})=0`$, d’où le résultat.
(2) Selon le point (1) et la proposition 3.7, la suite exacte (9) s’écrit
$$0H^1(E_\lambda ,𝒯_{E_\lambda })\underset{d}{}H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })0,$$
d’où le résultat. $`\mathrm{}`$
La proposition suivante achève donc la démonstration du théorème 3.1 :
###### Proposition 3.14.
L’espace $`H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })`$ est nul, sauf dans les cas suivants (à factorisation près) :
* Si $`G=\mathrm{SL}(2)\times \mathrm{SL}(n)`$ et $`d=1`$ et $`\lambda =(m,\omega _1)`$, alors il vaut $`V_{\mathrm{SL}(2)}(m2)V(\omega _1)`$.
* Si $`G=\mathrm{SL}(2)\times \mathrm{SL}(n)`$ et $`d=1`$ et $`\lambda =(m,\omega _n)`$, alors il vaut $`V_{\mathrm{SL}(2)}(m2)V(\omega _n)`$.
* Si $`G=\mathrm{SL}(2)\times \mathrm{Sp}(2n)`$ et $`d=1`$ et $`\lambda =(m,\omega _1)`$, alors il vaut $`V_{\mathrm{SL}(2)}(m2)V(\omega _1)`$.
Preuve. Selon la proposition 3.8, on peut supposer $`d<0`$. On a
$$_{d\lambda }𝒯_{X_\lambda }=(p_1^{}(_{d\lambda _1}𝒯_{X_{\lambda _1}})p_2^{}(_{d\lambda _2}))(p_1^{}(_{d\lambda _1})p_2^{}(_{d\lambda _2}𝒯_{X_{\lambda _2}})).$$
Donc, selon la formule de Künneth,
$$H^1(X_\lambda ,_{d\lambda }𝒯_{X_\lambda })=H^1(X_{\lambda _1},_{d\lambda _1})H^0(X_{\lambda _2},_{d\lambda _2}𝒯_{X_{\lambda _2}}).$$
En effet, comme $`d<0`$, on a
$$H^0(X_{\lambda _1},_{d\lambda _1})=H^0(X_{\lambda _2},_{d\lambda _2})=0$$
et selon la proposition 3.7 (grâce à l’hypothèse faite sur $`\lambda _2`$), on a
$$H^1(X_{\lambda _2},_{d\lambda _2})=0.$$
Si l’action de $`G^1`$ sur $`V(\lambda _1)`$ ne se factorise pas par $`\mathrm{SL}(2)`$, on aura également $`H^1(X_{\lambda _1},_{d\lambda _1})=0`$. On peut donc supposer que $`G^1=\mathrm{SL}(2)`$, de sorte que $`H^1(X_{\lambda _1},_{d\lambda _1})`$ est non nul, et vaut $`V(d\lambda _12)`$ (on considère désormais $`\lambda _1`$ comme un entier).
Il reste à calculer l’espace des sections globales $`H^0(X_{\lambda _2},_{d\lambda _2}𝒯_{X_{\lambda _2}})`$.
Remarquons tout d’abord que sa partie $`G`$-invariante est nulle : en effet, selon l’isomorphisme (où $``$ est la droite munie de l’action triviale de $`G`$)
$$\mathrm{Hom}^G(,\mathrm{Ind}(𝔤^2/𝔮_{\lambda _2}[d\lambda _2^{}]))\mathrm{Hom}^{Q_{\lambda _2}}(,𝔤^2/𝔮_{\lambda _2}[d\lambda _2^{}]),$$
la partie $`G`$-invariante est isomorphe à l’espace des $`Q_{\lambda _2}`$-invariants suivant :
$$(𝔤^2/𝔮_{\lambda _2}[d\lambda _2^{}])^{Q_{\lambda _2}}.$$
Supposons par l’absurde ce dernier espace non nul. Ses éléments sont en particulier de poids nul pour le tore $`T`$ : ils admettent donc un représentant dans $`𝔤^2`$ dont le poids est une racine positive $`\beta `$ de $`𝔤^2`$ telle que
$$\beta +d\lambda _2^{}=0.$$
On en déduit d’une part que l’on peut supposer que le groupe $`G^2`$ est simple (car son action sur $`V(\lambda _2)`$ se factorise par celle d’un groupe simple), et d’autre part que la racine positive $`\beta `$ est une racine dominante. Or on sait qu’alors le sous-$`𝔮_{\lambda _2}`$-module de $`𝔤^2`$ engendré par $`𝔤_\beta ^2`$ contient $`𝔤_\alpha ^2`$ pour toute racine simple $`\alpha `$ de $`𝔤^2`$.
Ainsi, toute racine simple de $`𝔤^2`$ est une racine de $`𝔮_{\lambda _2}`$, et $`𝔮_{\lambda _2}=𝔤^2`$ : une contradiction.
Déterminons maintenant les autres composantes isotypiques de l’espace des sections globales : soit $`\mu `$ un poids dominant de $`G^2`$ non nul. Supposons que $`H^0(X_{\lambda _2},_{d\lambda _2}𝒯_{X_{\lambda _2}})_{(\mu )}`$ est non nulle.
Comme dans la démonstration de la proposition 3.11, appliquons le théorème de Borel-Weil-Bott à l’aide d’une suite de Jordan-Hölder du $`Q_{\lambda _2}`$-module $`𝔤^2/𝔮_{\lambda _2}`$.
On obtient qu’il existe une racine positive $`\beta `$ de $`G^2`$ telle que
$$\beta +d\lambda _2^{}=\mu .$$
Comme $`d<0`$, pour que $`\beta +d\lambda _2^{}`$ soit dominant, il faut que l’action de $`G^2`$ sur $`V(\lambda _2)`$ se factorise par un groupe simple ; on suppose donc que $`G^2`$ est un groupe simple.
On remarque que $`\beta =\mu d\lambda _2^{}`$ est un poids dominant, non fondamental (car $`\mu `$ et $`d\lambda _2^{}`$ sont tous les deux non nuls).
Or les seules racines dominantes des systèmes de racines irréductibles qui ne sont pas des poids fondamentaux sont
* la plus grande racine $`\omega _1+\omega _n`$ des systèmes de racines de type $`A_n`$ ($`n1`$).
* la plus grande racine $`2\omega _1`$ des systèmes de racines de type $`C_n`$ ($`n2`$).
On a donc $`d=1`$, et l’on est dans l’une des deux situations suivantes (à revêtement fini de $`G^2`$ près) :
* On a $`G^2=\mathrm{SL}(n)`$ ($`n2`$), et $`\mu =\lambda _2=\omega _1`$ ou $`\mu =\lambda _2=\omega _n`$.
* On a $`G^2=\mathrm{Sp}(2n)`$ ($`n2`$), et $`\mu =\lambda _2=\omega _1`$.
Il ne reste plus qu’à vérifier que dans ces deux cas, l’espace $`T_\lambda ^1`$ est celui annoncé (la seule autre possibilité est qu’il soit nul).
En utilisant les notations analogues à celles de la démonstration de la proposition 3.11 (en remplaçant $`G`$ par $`G^2`$), on a pour tout $`i`$ une suite exacte
$$0\mathrm{Ind}(W_{i+1}[d\lambda _2^{}])_{(\mu )}\mathrm{Ind}(W_i[d\lambda _2^{}])_{(\mu )}\mathrm{Ind}(W_i/W_{i+1}[d\lambda _2^{}])_{(\mu )}R^1\mathrm{Ind}(W_{i+1}[d\lambda _2^{}])_{(\mu )}.$$
Selon cette même démonstration, on a
$$R^1\mathrm{Ind}(W_{i+1}[d\lambda _2^{}])_{(\mu )}=0$$
pour tout $`i`$. On en conclut facilement que
$$H^0(X_{\lambda _2},_{d\lambda _2}𝒯_{X_{\lambda _2}})=\mathrm{Ind}(W_0[d\lambda _2^{}])=V(\lambda _2)$$
dans les deux situations. $`\mathrm{}`$
### 3.4 Démonstration de la proposition 3.3
On va d’abord démontrer que la déformation $`𝔙`$ de la proposition 3.3 est plate (proposition 3.16). Elle est donc déduite de la déformation verselle de $`C_{mn}`$ par un changement de base. On montrera ensuite (proposition 3.17) que la différentielle du changement de base est un isomorphisme (de l’espace tangent à $`\mathrm{Spec}[[𝐭]]`$ vers l’espace des déformations infinitésimales de $`C_{mn}`$), ce qui démontrera la proposition 3.3.
Soit $`N`$ un entier supérieur ou égal à $`2`$. On munit l’ensemble des monômes de l’anneau de polynomes $`[x_1,\mathrm{},x_N,z]`$ de l’ordre lexicographique : si $`m_1=x_1^{\alpha _1}\mathrm{}x_N^{\alpha _N}z^{\alpha _{N+1}}`$ et $`m_2=x_1^{\beta _1}\mathrm{}x_N^{\beta _N}z^{\beta _{N+1}}`$ sont deux monômes distincts, alors $`m_1<m_2`$ si et seulement si $`\alpha _i<\beta _i`$ pour le plus petit indice $`i`$ tel que $`\alpha _i\beta _i`$.
On renvoie à \[Ei\] p 332 pour la définition d’une base de Gröbner. Le lemme suivant est une application immédiate du critère de Buchberger (\[Ei\] p 338):
###### Lemme 3.15.
Soient $`s_1,\mathrm{},s_N`$ des nombres complexes. Notons $`g_{ij}`$ le mineur $`2\times 2`$ de la matrice
$$\left(\begin{array}{cccc}x_0& x_1& \mathrm{}& x_{N1}\\ x_1s_1z& x_2s_2z& \mathrm{}& x_Ns_Nz\end{array}\right)$$
obtenu en prenant les colonnes $`i`$ et $`j`$. La famille $`(g_{ij})_{i<j}`$ est une base de Gröbner de l’idéal qu’elle engendre.
On reprend maintenant les notations de la proposition 3.3.
###### Proposition 3.16.
Notons $`𝔛`$ la famille de sous schémas fermés
définie par les mineurs $`2\times 2`$ de (6). La famille $`𝔛`$ est plate au dessus de $`𝔸^{(m1)(n+1)}`$.
Par conséquent la famille $`𝔙`$ est plate.
Preuve. On va montrer que l’adhérence schématique $`\overline{𝔛}`$ de $`𝔛`$ dans $`𝔸^{(m1)(n+1)}\times ^{(m+1)(n+1)}`$ est une famille plate sur $`𝔸^{(m1)(n+1)}`$, ce qui donnera le résultat.
Pour cela, il suffit (\[Ha\] thm 9.9 p 261) de montrer que le polynôme de Hilbert de la fibre $`\overline{𝔛}_s`$ de $`\overline{𝔛}`$ au dessus de $`s=(s_{ij})_{i,j}`$ (vue comme un sous-schéma fermé de $`^{(m+1)(n+1)}`$) est indépendant de $`s`$.
L’idéal homogène $`I_s`$ de $`\overline{𝔛}_s`$ est engendré par les mineurs $`2\times 2`$ de la matrice
$$\left(\begin{array}{ccccccccc}x_{00}& \mathrm{}& x_{0m2}& x_{0m1}& \mathrm{}& x_{n0}& \mathrm{}& x_{nm2}& x_{nm1}\\ x_{01}s_{01}z& \mathrm{}& x_{0m1}s_{0m1}z& x_{0m}& \mathrm{}& x_{n1}s_{n1}z& \mathrm{}& x_{nm1}s_{nm1}z& x_{nm}\end{array}\right)$$
$`z`$ est une indéterminée supplémentaire. Selon le lemme précédent, ceux-ci forment une base de Gröbner pour l’ordre lexicographique tel que $`x_{ij}<x_{kl}`$ si $`i<k`$ ou si $`i=k`$ et $`j<l`$, et tel que $`z`$ soit la plus petite indéterminée.
On remarque que les termes initiaux de ces mineurs $`2\times 2`$ ne dépendent pas de $`s`$. Selon \[Ei\] Thm 15.3 p 329, le $``$-espace vectoriel quotient $`[𝐱,z]/I_s`$ admet donc pour base les classes modulo $`I_s`$ d’une famille de monômes indépendante de $`s`$. La dimension des composantes homogènes de $`[𝐱,z]/I_s`$ est donc indépendante de $`s`$, et le polynôme de Hilbert de $`\overline{𝔛}_s`$ aussi, ce qui démontre la proposition. $`\mathrm{}`$
Remarquons que l’espace tangent à $`\mathrm{Spec}[[𝐭]]`$ est de même dimension que $`T^1(C_{mn})`$. Pour montrer que la différentielle du morphisme correspondant à $`𝔙`$ (de $`\mathrm{Spec}[[𝐭]]`$ vers la déformation verselle de $`C_{mn}`$) est un isomorphisme entre les espaces tangents, il suffit donc de montrer qu’elle est injective, c’est-à-dire que tout vecteur tangent non nul à $`\mathrm{Spec}[[𝐭]]`$ correspond à une déformation infinitésimale de $`C_{mn}`$ non triviale.
Notons $`ϵ`$ la classe de $`y`$ dans l’algèbre $`[y]/y^2`$.
###### Proposition 3.17.
Soit un vecteur tangent à $`\mathrm{Spec}[[𝐭]]`$, c’est-à-dire un morphisme de $`\mathrm{Spec}[ϵ]`$ vers $`\mathrm{Spec}[[𝐭]]`$. On note $`s_{ij}`$ les nombres complexes tels que le morphisme correspondant
$`[[𝐭]][ϵ]`$ envoie $`t_{ij}`$ vers $`ϵs_{ij}`$.
La déformation induite par $`𝔙`$ sur $`\mathrm{Spec}[ϵ]`$ est le sous-schéma fermé de $`\mathrm{Spec}[𝐱][ϵ]`$ défini par les mineurs $`2\times 2`$ de la matrice
$$\left(\begin{array}{ccccccccc}x_{00}& \mathrm{}& x_{0m2}& x_{0m1}& \mathrm{}& x_{n0}& \mathrm{}& x_{nm2}& x_{nm1}\\ x_{01}ϵs_{01}& \mathrm{}& x_{0m1}ϵs_{0m1}& x_{0m}& \mathrm{}& x_{n1}ϵs_{n1}& \mathrm{}& x_{nm1}ϵs_{nm1}& x_{nm}\end{array}\right).$$
Elle est triviale si et seulement si les $`s_{ij}`$ sont tous nuls.
Preuve. On va appliquer \[Ha\] ex 9.8 p 267. Notons $`J`$ l’idéal du sous-schéma fermé $`C_{mn}`$ de $`𝔸^{(m+1)(n+1)}`$, et $`A:=\mathrm{Spec}[𝐱]/J`$ son algèbre affine. La déformation infinitésimale de la proposition correspond au morphisme de $`A`$-modules $`\varphi :J/J^2A`$ tel que
$$\varphi \left(\begin{array}{cccc}& x_{ij}& x_{kl}& \\ & x_{ij+1}& x_{kl+1}& \end{array}\right)=\begin{array}{cc}x_{ij}& x_{kl}\\ s_{ij+1}& s_{kl+1}\end{array}.$$
Elle est triviale si et seulement si le morphisme $`\varphi `$ est induit par un champ de vecteurs
$$\underset{i,j}{}h_{ij}\frac{}{x_{ij}},$$
c’est-à-dire s’il existe des des éléments $`h_{ij}`$ de $`A`$ tels que
$$\varphi \left(\begin{array}{cccc}& x_{ij}& x_{kl}& \\ & x_{ij+1}& x_{kl+1}& \end{array}\right)=\begin{array}{cc}x_{ij}& x_{kl}\\ h_{ij+1}& h_{kl+1}\end{array}+\begin{array}{cc}h_{ij}& h_{kl}\\ x_{ij+1}& x_{kl+1}\end{array}.$$
En prenant alors $`i`$ et $`k`$ distincts (ce qui est possible car $`n`$ est non nul), on obtient que la composante homogène de degré $`1`$ de $`h_{ij}`$ est nulle dès que $`j<m`$, puis que les $`s_{ij}`$ sont tous nuls, d’où le résultat. $`\mathrm{}`$
## Appendice
On note $`\mathrm{Hilb}((V(\lambda )))`$ le schéma de Hilbert (construit dans \[Grot\]) des sous-schémas fermés de $`(V(\lambda ))`$. Le sous-schéma $`\mathrm{Hilb}^G((V(\lambda )))`$ des points fixes de $`G`$ dans $`\mathrm{Hilb}((V(\lambda )))`$ paramètre les sous-schémas fermés de $`(V(\lambda ))`$ qui sont stables par $`G`$.
On répond dans cet appendice à la question naturelle suivante : quelles déformations locales du cône des vecteurs primitifs $`C(\lambda )`$ peut-on obtenir à l’aide de $`\mathrm{Hilb}((V(\lambda )))`$?
Comme $`C(\lambda )`$ est le cône affine dans $`V(\lambda )`$ au dessus de la variété de drapeaux
$$X_\lambda :=G/P_\lambda (V(\lambda )),$$
on peut être tenté de déformer $`X_\lambda `$ dans $`(V(\lambda ))`$ à l’aide du schéma de Hilbert pour en déduire naturellement une déformation de $`C(\lambda )`$.
La proposition suivante montre que l’on n’obtient ainsi que des déformations triviales, c’est-à-dire provenant de l’action du groupe $`\mathrm{GL}(V(\lambda ))`$ des automorphismes d’espace vectoriel de $`V(\lambda )`$.
On note $`z`$ le point de $`\mathrm{Hilb}((V(\lambda )))`$ correspondant à $`X_\lambda `$. Le groupe $`\mathrm{GL}(V(\lambda ))`$ agit naturellement sur $`\mathrm{Hilb}((V(\lambda )))`$.
###### Proposition 3.18.
L’orbite $`\mathrm{GL}(V(\lambda ))z`$ est ouverte dans $`\mathrm{Hilb}((V(\lambda )))`$.
Preuve.
Il suffit de montrer que l’espace tangent à l’orbite est égal à l’espace tangent au schéma de Hilbert :
$$T_z(\mathrm{GL}(V(\lambda ))z)=T_z\mathrm{Hilb}((V(\lambda ))),$$
c’est-à-dire que l’application
$$𝔤𝔩(V(\lambda ))\stackrel{\mathit{\varphi }}{}T_z\mathrm{Hilb}((V(\lambda )))$$
obtenue en différentiant l’application naturelle
$$\begin{array}{ccc}\mathrm{GL}(V(\lambda ))& \hfill & \mathrm{Hilb}((V(\lambda ))),\\ u& \hfill & uz\end{array}$$
est surjective.
Notons $`𝒩_{X_\lambda }`$ le faisceau normal à $`X_\lambda `$ dans $`(V(\lambda ))`$.
Il est donné par la suite exacte courte de $`𝒪_{X_\lambda }`$-modules :
$$0𝒯_{X_\lambda }𝒯_{(V(\lambda ))}|_{X_\lambda }𝒩_{X_\lambda }0.$$
(11)
L’espace tangent au schéma de Hilbert en $`z`$ est canoniquement isomorphe à l’espace des sections globales de $`𝒩_{X_\lambda }`$ :
$$T_z\mathrm{Hilb}((V(\lambda )))\mathrm{H}^0(X_\lambda ,𝒩_{X_\lambda }).$$
On sait que l’espace
$$\mathrm{H}^1(X_\lambda ,𝒯_{X_\lambda })$$
est nul (cela résulte par exemple de la proposition 3.8).
En utilisant la suite exacte courte (11), on en déduit que l’application canonique
$$\mathrm{H}^0(X_\lambda ,𝒯_{(V(\lambda ))}|_{X_\lambda })\stackrel{\varphi _1}{}\mathrm{H}^0(X_\lambda ,𝒩_{X_\lambda })$$
est surjective.
On utilise ensuite la suite exacte de $`𝒪_{(V(\lambda ))}`$-modules (\[Ha\] Example II.8.20.1) :
$$0𝒪_{(V(\lambda ))}𝒪_{(V(\lambda ))}(1)_{}V(\lambda )𝒯_{(V(\lambda ))}0.$$
Comme les termes de cette suite sont des faisceaux localement libres, on obtient encore une suite exacte si on la restreint à $`X_\lambda `$ :
$$0𝒪_{X_\lambda }_\lambda _{}V(\lambda )𝒯_{(V(\lambda ))}|_{X_\lambda }0.$$
L’application associée
$$V(\lambda )^{}_{}V(\lambda )=\mathrm{H}^0(X_\lambda ,_\lambda )_{}V(\lambda )\stackrel{\varphi _2}{}\mathrm{H}^0(X_\lambda ,𝒯_{(V(\lambda ))}|_{X_\lambda })$$
est surjective, car l’espace $`\mathrm{H}^1(X_\lambda ,𝒪_{X_\lambda })`$ est nul (cela découle par exemple du théorème de Borel-Weil-Bott).
On remarque enfin que l’application $`\varphi `$ s’identifie à la composée $`\varphi _1\varphi _2`$ : elle est donc surjective, d’où la proposition. $`\mathrm{}`$
On en déduit immédiatement le corollaire suivant, qui montre que l’on n’obtient aucune déformation $`G`$-invariante à l’aide de $`\mathrm{Hilb}((V(\lambda )))`$.
###### Corollaire 3.19.
Le point $`z`$ est un point isolé réduit de $`\mathrm{Hilb}^G((V(\lambda )))`$.
Preuve. Il suffit de montrer que l’espace tangent à $`\mathrm{Hilb}^G((V(\lambda )))`$ en $`z`$ est nul. On reprend les notations de la démonstration de la proposition précédente. On a vu que l’application
$$𝔤𝔩(V(\lambda ))\stackrel{\mathit{\varphi }}{}T_z\mathrm{Hilb}((V(\lambda )))$$
est surjective. Comme le groupe $`G`$ est réductif, on en déduit une surjection sur les espaces des $`G`$-invariants :
$$𝔤𝔩(V(\lambda ))^G\stackrel{\mathit{\varphi }}{}T_z\mathrm{Hilb}^G((V(\lambda ))).$$
Or l’application $`\varphi `$ est nulle sur l’espace $`𝔤𝔩(V(\lambda ))^G`$ (qui est le centre de $`𝔤𝔩(V(\lambda ))`$), d’où le corollaire. $`\mathrm{}`$
Pinkham utilise dans \[Pi\] §4-5 de manière plus concluante le schéma de Hilbert $`\mathrm{Hilb}(\overline{V(\lambda )})`$ des sous-schémas fermés de l’espace projectif $`\overline{V(\lambda )}`$ obtenu en complétant $`V(\lambda )`$. Il étudie ainsi plus généralement les déformations des cônes affines sur les variétés projectives lisses, et montre sous certaines hypothèses (qui sont vérifiées dans notre situation) qu’on les obtient toutes en déformant à l’aide de $`\mathrm{Hilb}(\overline{V(\lambda )})`$ le complété du cône.
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# A study of the Gribov copies in linear covariant gauges in Euclidean Yang-Mills theories
## 1 Introduction
Gribov ambiguities and their relevance for the nonperturbative aspects of Euclidean Yang-Mills theories have witnessed growing interest in recent years. These ambiguities, affecting the Faddeev-Popov quantization formula, deeply modify the infrared behavior of Yang-Mills theories, as one learns from the large amount of results obtained in the Landau gauge as well as in the Coulomb gauge .
As pointed out in , the existence of these ambiguities is due to the lack of a globally well defined gauge-fixing procedure. Among the class of covariant gauges, the Gribov ambiguities have been much investigated in the Landau gauge, where the gauge field is transverse, $`_\mu A_\mu ^a=0`$. This property plays an important role here. It ensures that the Faddeev-Popov operator, $`^{ab}(A)=_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)`$, is hermitian, $`=^{}`$. Its eigenvalues are thus real. Concerning now the quantization of Yang-Mills theories in the Landau gauge, it turns out that, as a consequence of the existence of Gribov copies, the domain of integration in the Feynman path integral has to be restricted to the so called Gribov region $`\mathrm{\Omega }`$ , which is the set of all transverse fields for which the Faddeev-Popov operator is positive definite, i.e. $`\mathrm{\Omega }=\{A_\mu ^a,_\mu A_\mu ^a=0,^{ab}(A)=_\mu (_\mu \delta ^{ab}gf^{abc}A_\mu ^c)>0\}`$. The boundary $`\mathrm{\Omega },`$where the first vanishing eigenvalue of the Faddeev-Popov operator appears, is called the first Gribov horizon. For the partition function of Yang-Mills theories in the Landau gauge, one has
$`𝒵`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}DA\delta (A)\left(det\left(^2\delta ^{ab}+gf^{abc}A_\mu ^c_\mu \right)\right)e^{\frac{1}{4}{\scriptscriptstyle d^4xF_{\mu \nu }^aF_{\mu \nu }^a}}`$ (1.1)
$`=`$ $`{\displaystyle _\mathrm{\Omega }}DAD\overline{c}Dc\delta (A)e^{\frac{1}{4}{\scriptscriptstyle d^4xF_{\mu \nu }^aF_{\mu \nu }^a}+{\scriptscriptstyle d^4x\overline{c}^a_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)c^b}}.`$
The restriction of the domain of integration to the region $`\mathrm{\Omega }`$ has important consequences on the infrared behavior of the gluon and ghost propagators . More precisely, in the Landau gauge, the gluon propagator $`A_\mu ^a(k)A_\nu ^b(k)`$ is found to be suppressed in the infrared, while the ghost propagator $`\overline{c}^a(k)c^b(k)`$ is enhanced, according to
$$A_\mu ^a(k)A_\nu ^b(k)=\delta ^{ab}\left(\delta _{\mu \nu }\frac{k_\mu k_\nu }{k^2}\right)\frac{k^2}{k^4+\gamma ^4},$$
(1.2)
and
$`𝒢_{gh}(k)`$ $`=`$ $`{\displaystyle \frac{1}{N^21}}{\displaystyle \underset{ab}{}}\delta ^{ab}\overline{c}^a(k)c^b(k),`$
$`𝒢_{gh}(k)_{k0}`$ $``$ $`{\displaystyle \frac{1}{k^4}}.`$ (1.3)
The Gribov parameter $`\gamma `$ in eq.$`\left(\text{1.2}\right)`$ has the dimension of a mass and is determined by a gap equation which, to the first order approximation, reads
$$\frac{3Ng^2}{4}\frac{d^4k}{\left(2\pi \right)^4}\frac{1}{k^4+\gamma ^4}=1.$$
(1.4)
From equation $`\left(\text{1.3}\right)`$, one sees that the infrared behavior of the ghost propagator is more singular than the perturbative behavior $`1/k^2`$. This infrared enhancement is known as the Gribov-Zwanziger horizon condition , generally stated as $`lim_{k0}\left(k^2𝒢_{gh}(k)\right)^1=0`$. Remarkably, this behavior of the gluon and ghost propagators in the Landau gauge has received many confirmations from lattice simulations as well as from recent studies of the Schwinger-Dyson equations .
Several results have been established about the Gribov region $`\mathrm{\Omega }`$. It has been proven that $`\mathrm{\Omega }`$ is convex and bounded in every direction . Moreover, every gauge orbit passes inside $`\mathrm{\Omega }`$ . The latter result is deeply related to the possibility of introducing an auxiliary functional, $`(A)=d^4xA_\mu ^aA_\mu ^a`$, whose minimization with respect to the gauge transformations provides a characterization of the region $`\mathrm{\Omega }`$ . It can be checked in fact that the Gribov region $`\mathrm{\Omega }`$ can be identified as the set of all relative minima of the functional $`(A)`$. It should also be mentioned here that the region $`\mathrm{\Omega }`$ is not free from Gribov copies, i.e. Gribov copies still exist inside $`\mathrm{\Omega }`$ . To avoid the presence of these additional copies, a further restriction to a smaller region $`\mathrm{\Lambda }`$, known as the fundamental modular region, should be implemented . The region $`\mathrm{\Lambda }`$ is contained in the Gribov region $`\mathrm{\Omega }`$, being defined as the set of all absolute minima of the auxiliary functional $`(A)`$. However, it is difficult to give an explicit description of $`\mathrm{\Lambda }`$. Recently, it has been argued in that the additional copies existing inside $`\mathrm{\Omega }`$ have no influence on the expectation values, so that averages calculated over $`\mathrm{\Lambda }`$ or $`\mathrm{\Omega }`$ are expected to give the same result.
Besides the Landau gauge, the effects of the Gribov copies on the infrared behavior of Yang-Mills theories have been studied to a great extent in the noncovariant Coulomb gauge , $`_iA_i^a=0`$, $`i=1,2,3`$. In particular, the Coulomb gauge allows for a direct study of the confining properties of the static potential $`V(r)`$ between an external quark pair at spatial separation $`r`$. It turns out that $`V(r)`$ can be accessed by means of the $`44`$-component, $`A_4^a(\stackrel{}{x},t)A_4^b(0)`$, of the gluon propagator. Moreover, in analogy with the Landau gauge, the spatial components of the gluon propagator, $`A_i^a(k)A_j^b(k)`$, are found to be suppressed in the infrared, a feature confirmed by lattice simulations .
Concerning now other covariant gauges, although the presence of the Gribov copies is certainly to be expected , as explicitly confirmed by a recent study of the zero modes of the Faddeev-Popov operator in the maximal Abelian gauge , very little is known about their consequences on the infrared behavior of the gluon and ghost propagators and, more generally, on the whole Yang-Mills theory. The need for such an investigation is motivated by the increasing belief that the Gribov copies might have a crucial role for the infrared region of Yang-Mills theories as well as for color confinement. It would be desirable thus to improve our understanding on how the Gribov copies manifest themselves in different gauges and how they modify the infrared behavior of the theory.
This work aims at studying this issue in the class of covariant linear gauges, $`_\mu A_\mu ^a=\alpha b^a`$, where $`\alpha `$ is the gauge parameter and $`b^a`$ is the Lagrange multiplier. The task is far from being trivial. Many features of the Landau gauge are lost for a generic nonvanishing $`\alpha `$. The Faddeev-Popov operator, $`^{ab}(A)=_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)`$, is now not hermitian. Moreover, a suitable minimizing functional in these gauges is not yet at our disposal. As a consequence, the identification of the region to which the domain of integration in the path-integral should be restricted becomes difficult to be handled. All this makes a complete treatment of the Gribov copies in linear covariant gauges far beyond our present capabilities. Nevertheless, we shall be able to establish some preliminary results which enable us to obtain a characterization of the infrared behavior of the gluon propagator, at least for small values of the gauge parameter $`\alpha `$. Considering in fact small values of the gauge parameter $`\alpha `$, will allow us to stay close to the Landau gauge fixing. The present covariant gauge can be considered thus as a kind of deformation of the Landau gauge, whose results have to be recovered in the limit $`\alpha 0`$, thereby providing a useful consistency check.
The output of our findings can be summarized as follows. As in the case of the Landau gauge, the transverse component of the gluon propagator turns out to be suppressed in the infrared. Moreover, its longitudinal part remains unchanged.
Concerning now the behavior of the ghost fields in linear covariant gauges, it turns out that, instead of the ghost propagator, the Green function which is enhanced in the infrared is given by the quantity $`𝒢_{tr}(k)`$, defined as
$`𝒢_{tr}(k)`$ $`=`$ $`{\displaystyle \frac{1}{N^21}}{\displaystyle \underset{ab}{}}\delta ^{ab}𝒢_{tr}^{ab}(k),`$
$`𝒢_{tr}^{ab}(k)`$ $`=`$ $`k\left|\left(^1(A^T)\right)^{ab}\right|k,`$ (1.5)
where
$$A_\mu ^{aT}=\left(\delta _{\mu \nu }\frac{_\mu _\nu }{^2}\right)A_\nu ^a,$$
(1.6)
is the transverse component of the gauge field and $`\left(^1(A^T)\right)^{ab}`$ stands for the inverse of the Faddeev-Popov operator $`^{ab}(A^T)`$,
$$^{ab}(A^T)=_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{Tc}\right).$$
(1.7)
The Green function $`𝒢_{tr}(k)`$ is found to obey the Gribov-Zwanziger horizon condition, i.e. $`lim_{k0}\left(k^2𝒢_{tr}(k)\right)^1=0`$. It should be remarked that $`𝒢_{tr}(k)`$ does not coincide with the ghost propagator for a generic value of the gauge parameter $`\alpha `$. However, $`𝒢_{tr}(k)`$ reduces to the ghost two-point function for vanishing $`\alpha `$, so that our results turn out to coincide with those of the Landau gauge in the limit $`\alpha 0`$.
The infrared behavior of the gluon propagator and of $`𝒢_{tr}(k)`$ will be investigated also in the presence of the dimension two gluon condensate $`A_\mu ^aA_\mu ^a`$. A detailed study of this condensate in linear covariant gauges can be found in . In the presence of $`A_\mu ^aA_\mu ^a`$, the infrared suppression of the transverse component of the gluon propagator is enforced. Furthermore, its longitudinal component turns out to be suppressed as well. Concerning the Green function $`𝒢_{tr}(k)`$, its infrared enhancement is not modified by the condensate $`A_\mu ^aA_\mu ^a`$.
The work is organized as follows. In Sect.2 a few properties of the Gribov copies in the linear covariant gauges, and for small values of the gauge parameter $`\alpha `$, are presented. In Sect.3 the infrared behavior of the gluon propagator, of the Green function $`𝒢_{tr}(k)`$ and of the ghost propagator is worked out. Sect.4 accounts for the inclusion of the gluon condensate $`A_\mu ^aA_\mu ^a`$. Sect.5 is devoted to a comparison of our results with those obtained from the analysis of the Schwinger-Dyson equations and from lattice simulations. In Sect.6 we present our conclusion.
## 2 A few properties of the Gribov copies in the linear covariant gauges
In this section we shall discuss a few properties of the Gribov copies in the linear covariant gauges. Let us begin by considering the expression of the Euclidean Yang-Mills action quantized in the linear covariant $`\alpha `$gauges, namely
$$S_{YM}+S_{gf}=\frac{1}{4}d^4xF_{\mu \nu }^aF_{\mu \nu }^ad^4x\left(b^aA^a+\frac{\alpha }{2}b^ab^a+\overline{c}^a_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)c^b\right),$$
(2.8)
where $`\overline{c}^a`$, $`c^a`$ stand for the Faddeev-Popov ghosts. The Lagrange multiplier $`b^a`$ has been introduced to implement the gauge condition
$$_\mu A_\mu ^a=\alpha b^a,$$
(2.9)
from which it follows
$$d^4x\left(b^aA^a+\frac{\alpha }{2}b^ab^a\right)\frac{1}{2\alpha }d^4x\left(_\mu A_\mu ^a\right)^2.$$
(2.10)
From the relation $`\left(\text{2.9}\right)`$ we see that the field $`A_\mu ^a`$ is not transverse, due to the presence of the gauge parameter $`\alpha `$. Of course, in the limit $`\alpha 0`$, we recover the Landau gauge, $`_\mu A_\mu ^a=0`$. In what follows, it will be useful to decompose the gauge field $`A_\mu ^a`$ into transverse and longitudinal components, according to
$$A_\mu ^a=A_\mu ^{aT}+A_\mu ^{aL},$$
(2.11)
with
$`A_\mu ^{aT}`$ $`=`$ $`\left(\delta _{\mu \nu }{\displaystyle \frac{_\mu _\nu }{^2}}\right)A_\nu ^a,`$
$`_\mu A_\mu ^{aT}`$ $`=`$ $`0,`$ (2.12)
and
$`A_\mu ^{aL}`$ $`=`$ $`{\displaystyle \frac{_\mu _\nu }{^2}}A_\nu ^a=\alpha {\displaystyle \frac{_\mu }{^2}}b^a,`$
$`_\mu A_\mu ^{aL}`$ $`=`$ $`\alpha b^a.`$ (2.13)
As already underlined, we shall restrict ourselves to the case in which $`\alpha `$ is small, i.e. $`\alpha 1`$, so that the longitudinal component $`A_\mu ^{aL}`$ in eq.$`\left(\text{2.11}\right)`$ can be considered as a small perturbation with respect to the transverse part $`A_\mu ^{aT}`$.
Let us proceed by proving the following statement.
* Statement
If the transverse component $`A_\mu ^{aT}`$ of the gauge field $`A_\mu ^a=\left(A_\mu ^{aT}+A_\mu ^{aL}\right)`$ belongs to the Gribov region $`\mathrm{\Omega }`$, i.e. $`A_\mu ^{aT}\mathrm{\Omega }`$, $`\mathrm{\Omega }=\{A_\mu ^{aT},_\mu A_\mu ^{aT}=0,_\mu (_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT})>0\}`$, then the Faddeev-Popov operator $`^{ab}(A)=_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)`$ has no zero modes.
Proof
The proof of this statement is done by assuming the converse. Suppose indeed that the operator $`^{ab}(A)`$ has a zero mode, i.e. it exists a $`\stackrel{~}{\phi }^a(x,\alpha )`$ such that
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}gf^{abc}A_\mu ^{cL}\right)\stackrel{~}{\phi }^b(x,\alpha )=0,$$
(2.14)
which, from eq.$`\left(\text{2.13}\right)`$, becomes
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}+\alpha gf^{abc}\left(\frac{_\mu }{^2}b^c\right)\right)\stackrel{~}{\phi }^b(x,\alpha )=0.$$
(2.15)
Since $`\alpha `$ is small we perform a Taylor expansion of $`\stackrel{~}{\phi }^a(x,\alpha )`$, namely
$$\stackrel{~}{\phi }^a(x,\alpha )=\underset{n=0}{\overset{\mathrm{}}{}}\alpha ^n\stackrel{~}{\phi }_n^a(x),$$
(2.16)
with
$$\stackrel{~}{\phi }_0^a(x)=\stackrel{~}{\phi }^a(x,\alpha )|_{\alpha =0}.$$
(2.17)
Equation $`\left(\text{2.15}\right)`$ splits thus in a tower of equations, given by
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)\stackrel{~}{\phi }_0^b(x)=0,$$
(2.18)
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)\stackrel{~}{\phi }_1^b(x)_\mu \left(gf^{abc}\left(\frac{_\mu }{^2}b^c\right)\stackrel{~}{\phi }_0^b(x)\right)=0,$$
(2.19)
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)\stackrel{~}{\phi }_2^b(x)_\mu \left(gf^{abc}\left(\frac{_\mu }{^2}b^c\right)\stackrel{~}{\phi }_1^b(x)\right)=0,$$
(2.20)
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)\stackrel{~}{\phi }_3^b(x)_\mu \left(gf^{abc}\left(\frac{_\mu }{^2}b^c\right)\stackrel{~}{\phi }_2^b(x)\right)=0,$$
(2.21)
and so on.
Moreover, since $`A_\mu ^{aT}`$ belongs to the Gribov region $`\mathrm{\Omega }`$, $`A_\mu ^{aT}\mathrm{\Omega }`$, it follows that $`\stackrel{~}{\phi }_0^b(x)`$ in the first equation $`\left(\text{2.18}\right)`$ necessarily vanishes, i.e. $`\stackrel{~}{\phi }_0^b(x)=0`$, since the operator $`_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)`$ has no zero modes in $`\mathrm{\Omega }`$. Furthermore, setting $`\stackrel{~}{\phi }_0^b(x)=0`$ in the second equation $`\left(\text{2.19}\right)`$, we get
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)\stackrel{~}{\phi }_1^b(x)=0,$$
(2.22)
which, in turn, implies the vanishing of $`\stackrel{~}{\phi }_1^b(x)`$, i.e. $`\stackrel{~}{\phi }_1^b(x)=0`$. As a consequence, eq.$`\left(\text{2.20}\right)`$ reads
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT}\right)\stackrel{~}{\phi }_2^b(x)=0,$$
(2.23)
from which $`\stackrel{~}{\phi }_2^b(x)=0`$. It is apparent thus that the argument can be easily iterated to the whole tower of equations, implying that $`\stackrel{~}{\phi }^b(x,\alpha )`$ vanishes, $`\stackrel{~}{\phi }^b(x,\alpha )=0`$. This concludes the proof of the statement.
In particular, from this result it follows that if $`A_\mu ^{aT}`$ belongs to the Gribov region, $`A_\mu ^{aT}\mathrm{\Omega }`$, the field $`A_\mu ^a=\left(A_\mu ^{aT}+A_\mu ^{aL}\right)`$ does not possess Gribov copies which are closely related, i.e. obtained from $`A_\mu ^a`$ through an infinitesimal gauge transformation
$$\stackrel{~}{A}_\mu ^a=A_\mu ^a+\left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)\omega ^b,$$
(2.24)
where $`\omega ^a(x)`$ denotes the parameter of the gauge transformation. Indeed, from the condition
$$_\mu \stackrel{~}{A}_\mu ^a=_\mu A_\mu ^a,$$
(2.25)
we should have
$$_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)\omega ^b=0,$$
(2.26)
which, however, has no solution for $`\omega ^a(x)`$.
### 2.1 Restriction of the domain of integration in the path-integral
The results obtained in the previous sections suggest to restrict the domain of integration in the path integral to the region $`\stackrel{~}{\mathrm{\Omega }}`$ defined as
$$\stackrel{~}{\mathrm{\Omega }}\{A_\mu ^a;A_\mu ^a=A_\mu ^{aT}+A_\mu ^{aL},A_\mu ^{aT}\mathrm{\Omega }\},$$
(2.27)
i.e. $`\stackrel{~}{\mathrm{\Omega }}`$ is the set of all connections whose transverse part $`A_\mu ^{aT}`$ belongs to the Gribov region $`\mathrm{\Omega }=\{A_\mu ^{aT},_\mu A_\mu ^{aT}=0,_\mu (_\mu \delta ^{ab}gf^{abc}A_\mu ^{cT})>0\}`$. This choice is motivated by the following considerations:
* Since the gauge parameter $`\alpha `$ is small, $`\alpha 1`$, the region $`\stackrel{~}{\mathrm{\Omega }}`$ can be regarded as a deformation of the Gribov region $`\mathrm{\Omega }`$. It is apparent in fact that in the limit $`\alpha 0`$ the region $`\mathrm{\Omega }`$ is recovered, namely
$$\underset{\alpha 0}{lim}\stackrel{~}{\mathrm{\Omega }}=\mathrm{\Omega }.$$
(2.28)
* As discussed before, the Faddeev-Popov operator, $`^{ab}(A)=_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^c\right)`$, has no zero modes within $`\stackrel{~}{\mathrm{\Omega }}`$. Therefore, the restriction to $`\stackrel{~}{\mathrm{\Omega }}`$ allows us to get rid of the Gribov copies which can be obtained through infinitesimal gauge transformations.
* The effective implementation in the path-integral of the restriction of the domain of integration to the region $`\stackrel{~}{\mathrm{\Omega }}`$ can be done by repeating the no-pole prescription of Gribov’s original work . Indeed, since the transverse component $`A_\mu ^{aT}`$ of any field belonging to $`\stackrel{~}{\mathrm{\Omega }}`$ lies in the Gribov region $`\mathrm{\Omega }`$, to implement the restriction to $`\stackrel{~}{\mathrm{\Omega }}`$ it will be sufficient to require that the Green function $`𝒢_{tr}(k)`$ of eq.$`\left(\text{1.5}\right)`$ has no poles for any $`A_\mu ^{aT}\mathrm{\Omega }`$, except for a singularity at $`k^2=0`$, corresponding to the Gribov horizon $`\mathrm{\Omega }`$.
Thus, for the partition function of Yang-Mills theories in linear covariant gauges, we write
$`𝒵`$ $`=`$ $`{\displaystyle DADb\delta (A+\alpha b)det\left(^{ab}(A)\right)e^{\left(\frac{1}{4}{\scriptscriptstyle d^4xF_{\mu \nu }^aF_{\mu \nu }^a}{\scriptscriptstyle d^4x\left(b^aA^a+{\scriptscriptstyle \frac{\alpha }{2}}b^ab^a\right)}\right)}𝒱(\stackrel{~}{\mathrm{\Omega }})}`$ (2.29)
$`=`$ $`{\displaystyle DAdet\left(_\mu \left(\delta ^{ab}_\mu gf^{abc}A_\mu ^c\right)\right)e^{\left(\frac{1}{4}{\scriptscriptstyle d^4xF_{\mu \nu }^aF_{\mu \nu }^a}+\frac{1}{2\alpha }{\scriptscriptstyle d^4x\left(A^a\right)^2}\right)}𝒱(\stackrel{~}{\mathrm{\Omega }})},`$
where the factor $`𝒱(\stackrel{~}{\mathrm{\Omega }})`$ implements the restriction to the region $`\stackrel{~}{\mathrm{\Omega }}`$. An explicit characterization of $`𝒱(\stackrel{~}{\mathrm{\Omega }})`$ and of its consequences on the infrared behavior of the gluon propagator and of $`𝒢_{tr}(k)`$ will be discussed in the next section.
Finally, it is worth to spend a few words on the aspects which remain uncovered by the present investigation. Even if the restriction to the region $`\stackrel{~}{\mathrm{\Omega }}`$ allows us to get rid of the Gribov copies which are closely related, i.e. attainable by infinitesimal gauge transformations, we still lack a treatment of the copies which cannot be attained by infinitesimal transformations. This task is beyond our present possibilities, as the knowledge of a suitable auxiliary functional associated to the linear covariant gauges would be required. Nevertheless, since we are limiting ourselves to small values of $`\alpha `$, the restriction to the region $`\stackrel{~}{\mathrm{\Omega }}`$ looks quite natural.
### 2.2 Characterization of the factor $`𝒱(\stackrel{~}{\mathrm{\Omega }})`$
As already remarked, the factor $`𝒱(\stackrel{~}{\mathrm{\Omega }})`$ can be accommodated for by requiring that the Green function $`𝒢_{tr}(k)`$ of eq.$`\left(\text{1.5}\right)`$ has no poles for a given nonvanishing value of the momentum $`k`$. This condition can be understood by recalling that the region $`\mathrm{\Omega }`$ is defined as the set of all transverse gauge connections $`\left\{A_\mu ^{Ta}\right\}`$, $`_\mu A_\mu ^{Ta}=0,`$ for which the operator $`^{ab}(A^T)`$ is positive definite, i.e. $`^{ab}(A^T)=_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{Tc}\right)>0`$. This implies that the inverse operator $`\left[_\mu \left(_\mu \delta ^{ab}gf^{abc}A_\mu ^{Tc}\right)\right]^1`$, and thus $`𝒢_{tr}(k)`$, can become large only when approaching the horizon $`\mathrm{\Omega }`$, which corresponds in fact to $`k=0`$ . The quantity $`𝒢_{tr}(k)`$ can be evaluated order by order in perturbation theory. Repeating the same calculation of , we find that, up to the second order
$$𝒢_{tr}(k)\frac{1}{k^2}\frac{1}{1\rho (k,A^T)},$$
(2.30)
with $`\rho (k,A^T)`$ given by
$`\rho (k,A^T)`$ $`=`$ $`{\displaystyle \frac{N}{N^21}}{\displaystyle \frac{g^2}{V}}{\displaystyle \frac{1}{k^2}}{\displaystyle \underset{q}{}}{\displaystyle \frac{k_\mu (kq)_\nu }{\left(kq\right)^2}}A_\mu ^{Ta}(q)A_\nu ^{Ta}(q)`$ (2.31)
$`=`$ $`{\displaystyle \frac{N}{N^21}}{\displaystyle \frac{g^2}{V}}{\displaystyle \frac{k_\mu k_\nu }{k^2}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{\left(kq\right)^2}}A_\mu ^{Ta}(q)A_\nu ^{Ta}(q),`$
$`V`$ being the Euclidean space-time volume. According to , the no-pole condition for $`𝒢_{tr}(k)`$ reads
$`\rho (0,A^T)`$ $`<`$ $`1,`$
$`\rho (0,A^T)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{N}{N^21}}{\displaystyle \frac{g^2}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{q^2}}\left(A_\lambda ^{Ta}(q)A_\lambda ^{Ta}(q)\right).`$ (2.32)
Therefore, for the factor $`𝒱(\stackrel{~}{\mathrm{\Omega }})`$ in eq.$`\left(\text{2.29}\right)`$ we have
$$𝒱(\stackrel{~}{\mathrm{\Omega }})=\theta (1\rho (0,A^T)),$$
(2.33)
where $`\theta (x)`$ stands for the step function$`\theta (x)=1`$ if $`x>0`$, $`\theta (x)=0`$ if $`x<0`$..
## 3 The gluon propagator
In order to study the gluon propagator, it is sufficient to retain only the quadratic terms in expression $`\left(\text{2.29}\right)`$ which contribute to the two-point correlation function $`A_\mu ^a(k)A_\nu ^b(k)`$. Making use of the integral representation for the step function
$$\theta (1\rho (0,A^T))=_{i\mathrm{}+\epsilon }^{i\mathrm{}+\epsilon }\frac{d\eta }{2\pi i\eta }e^{\eta (1\rho (0,A^T))},$$
(3.34)
for the partition function $`\left(\text{2.29}\right)`$ one get
$$𝒵_{\mathrm{quadr}}=𝒩DA\frac{d\eta }{2\pi i}e^{\eta \mathrm{log}\eta }e^{\left(S_{\mathrm{quadr}}+\eta \rho (0,A^T)\right)},$$
(3.35)
where $`𝒩`$ is a constant factor and $`S_{\mathrm{quadr}}`$ stands for the quadratic part of the quantized Yang-Mills action
$$S_{\mathrm{quadr}}=\frac{1}{2}\underset{q}{}\left(q^2A_\mu ^a(q)A_\mu ^a(q)\left(1\frac{1}{\alpha }\right)q_\mu q_\nu A_\mu ^a(q)A_\nu ^a(q)\right).$$
(3.36)
From
$`A_\mu ^{Ta}(q)A_\mu ^{Ta}(q)`$ $`=`$ $`\left(\delta _{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right)\left(\delta _{\mu \sigma }{\displaystyle \frac{q_\mu q_\sigma }{q^2}}\right)A_\nu ^a(q)A_\sigma ^a(q)`$ (3.37)
$`=`$ $`A_\mu ^a(q)A_\mu ^a(q){\displaystyle \frac{q_\mu q_\nu }{q^2}}A_\mu ^a(q)A_\nu ^a(q),`$
it follows that
$$𝒵_{\mathrm{quadr}}=𝒩DA\frac{d\eta }{2\pi i}e^{\eta \mathrm{log}\eta }e^{\frac{1}{2}_qA_\mu ^a(q)𝒬_{\mu \nu }^{ab}A_\nu ^b(q)},$$
(3.38)
with
$$𝒬_{\mu \nu }^{ab}=\left(\left(q^2+\frac{\eta Ng^2}{N^21}\frac{1}{2V}\frac{1}{q^2}\right)\delta _{\mu \nu }q_\mu q_\nu \left(\left(1\frac{1}{\alpha }\right)+\frac{\eta Ng^2}{N^21}\frac{1}{2V}\frac{1}{q^4}\right)\right)\delta ^{ab}.$$
(3.39)
Integrating over the gauge field, one has
$$𝒵_{\mathrm{quadr}}=𝒩\frac{d\eta }{2\pi i}e^{\eta \mathrm{log}\eta }\left(det𝒬_{\mu \nu }^{ab}\right)^{\frac{1}{2}}=𝒩^{}\frac{d\eta }{2\pi i}e^{f(\eta )},$$
(3.40)
where $`f(\eta )`$ is given by
$$f(\eta )=\eta \mathrm{log}\eta \frac{3}{2}(N^21)\underset{q}{}\mathrm{log}\left(q^4+\frac{\eta Ng^2}{N^21}\frac{1}{2V}\right).$$
(3.41)
Following , the expression $`\left(\text{3.40}\right)`$ can be now evaluated at the saddle point, namely
$$𝒵_{\mathrm{quadr}}e^{f(\eta _0)},$$
(3.42)
where $`\eta _0`$ is determined by the minimum condition
$$1\frac{1}{\eta _0}\frac{3}{4}\frac{Ng^2}{V}\underset{q}{}\frac{1}{\left(q^4+\frac{\eta _0Ng^2}{N^21}\frac{1}{2V}\right)}=0.$$
(3.43)
Taking the thermodynamic limit, $`V\mathrm{}`$, and introducing the Gribov parameter $`\gamma `$
$$\gamma ^4=\frac{\eta _0Ng^2}{N^21}\frac{1}{2V},V\mathrm{},$$
(3.44)
we get the gap equation
$$\frac{3}{4}Ng^2\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{q^4+\gamma ^4}=1,$$
(3.45)
where the term $`1/\eta _0`$ in $`\left(\text{3.43}\right)`$ has been neglected in the thermodynamic limit. To obtain the gauge propagator, we can now go back to the expression for $`𝒵_{\mathrm{quadr}}`$ which, after substituting the saddle point value $`\eta =\eta _0`$, becomes
$$𝒵_{\mathrm{quadr}}=𝒩DAe^{\frac{1}{2}_qA_\mu ^a(q)𝒬_{\mu \nu }^{ab}A_\nu ^b(q)},$$
(3.46)
with
$`𝒬_{\mu \nu }^{ab}`$ $`=`$ $`\left(\left(q^2+{\displaystyle \frac{\gamma ^4}{q^2}}\right)\delta _{\mu \nu }q_\mu q_\nu \left(\left(1{\displaystyle \frac{1}{\alpha }}\right)+{\displaystyle \frac{\gamma ^4}{q^4}}\right)\right)\delta ^{ab}`$ (3.47)
$`=`$ $`\left(\left(q^2+{\displaystyle \frac{\gamma ^4}{q^2}}\right)\left(\delta _{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right)+{\displaystyle \frac{q_\mu q_\nu }{q^2}}\left({\displaystyle \frac{q^2}{\alpha }}\right)\right)\delta ^{ab}.`$
Evaluating the inverse of $`𝒬_{\mu \nu }^{ab}`$ , we get the gluon propagator
$$A_\mu ^a(q)A_\nu ^b(q)=\delta ^{ab}\left(\frac{q^2}{q^4+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)+\frac{\alpha }{q^2}\frac{q_\mu q_\nu }{q^2}\right).$$
(3.48)
A few remarks are now in order.
Let us first observe that the gap equation $`\left(\text{3.45}\right)`$ defining the parameter $`\gamma `$ has the same form of that obtained by Gribov in the Landau gauge . This is an expected result, since the factor $`\rho (0,A^T)`$ in equation $`\left(\text{2.32}\right)`$ depends only on the transverse component $`A_\mu ^{Ta}`$.
As it is apparent from the expression $`\left(\text{3.48}\right)`$, the transverse component of the gluon propagator is suppressed in the infrared, while the longitudinal component is left unchanged. Moreover, as we shall see later, the behavior of the longitudinal part turns out to be modified once the gluon condensate $`A_\mu ^aA_\mu ^a`$ is taken into account.
Finally, in the limit $`\alpha 0`$, the results of the Landau gauge are recovered.
### 3.1 The infrared behavior of $`𝒢_{tr}(k)`$
Let us discuss now the infrared behavior of the Green function $`𝒢_{tr}(k)`$ of eq.$`\left(\text{1.5}\right)`$, which is obtained upon contraction of the gauge fields in the expression $`\left(\text{2.31}\right)`$, namely
$$𝒢_{tr}(k)\frac{1}{k^2}\frac{1}{1\rho (k)},$$
(3.49)
with
$$\rho (k)=g^2\frac{N}{N^21}\frac{k_\mu k_\nu }{k^2}\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{\left(kq\right)^2}A_\mu ^{Ta}(q)A_\nu ^{Ta}(q).$$
(3.50)
From the expression of the gluon propagator in eq.$`\left(\text{3.48}\right)`$, it follows
$$\rho (k)=g^2N\frac{k_\mu k_\nu }{k^2}\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{\left(kq\right)^2}\frac{q^2}{q^4+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right).$$
(3.51)
Furthermore, from the gap equation $`\left(\text{3.45}\right)`$, it turns out
$$Ng^2\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{q^4+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)=\delta _{\mu \nu },$$
(3.52)
so that
$$1\rho (k)=Ng^2\frac{k_\mu k_\nu }{k^2}\frac{d^4q}{\left(2\pi \right)^4}\frac{k^22qk}{\left(kq\right)^2}\frac{1}{q^4+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right).$$
(3.53)
Note that the integral in the right hand side of eq.$`\left(\text{3.53}\right)`$ is convergent and nonsingular at $`k=0`$. Therefore, for $`k0`$,
$$\left(1\rho (k)\right)_{k0}\frac{3Ng^2}{4}k^2,$$
(3.54)
where $``$ stands for the value of the integral
$$=\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{q^2(q^4+\gamma ^4)}=\frac{1}{32\pi }\frac{1}{\gamma ^2}.$$
(3.55)
Finally, for the infrared behavior of the Green function $`𝒢_{tr}(k)`$ we get
$$\left(𝒢_{tr}\right)_{k0}\frac{4}{3Ng^2}\frac{1}{k^4}.$$
(3.56)
One sees thus that $`𝒢_{tr}(k)`$ is enhanced in the infrared, obeying the Gribov-Zwanziger condition $`lim_{k0}\left(k^2𝒢_{tr}(k)\right)^1=0`$.
### 3.2 Analysis of the ghost propagator
For the sake of completeness, let us discuss here the infrared behavior of the ghost two-point function, $`𝒢_{gh}(k)`$, given by
$$𝒢_{gh}(k)=\frac{1}{N^21}\underset{ab}{}\delta ^{ab}\overline{c}^a(k)c^b(k)\frac{1}{k^2}\frac{1}{1\omega (k)},$$
(3.57)
with
$$\omega (k)=\frac{N}{N^21}\frac{g^2}{k^2}\frac{d^4q}{\left(2\pi \right)^4}\frac{k_\mu (kq)_\nu }{(kq)^2}A_\mu ^a(q)A_\nu ^a(q).$$
(3.58)
Making use of the expression for the gluon propagator in eq.$`\left(\text{3.48}\right)`$ and of the equation $`\left(\text{3.52}\right)`$, it follows that, in the infrared,
$$1\omega (k)\frac{3Ng^2}{128\pi }\frac{k^2}{\gamma ^2}\frac{\alpha Ng^2}{k^2}\frac{d^4q}{\left(2\pi \right)^4}\frac{k_\mu (kq)_\nu }{(kq)^2}\frac{q_\mu q_\nu }{q^4}.$$
(3.59)
The second term in the right hand-side of eq.$`\left(\text{3.59}\right)`$ can be easily evaluated by means of the dimensional regularization. Adopting the $`\overline{MS}`$ scheme, the final expression for the factor $`(1\omega (k))`$ is found
$$(1\omega (k))\frac{3Ng^2}{128\pi }\frac{k^2}{\gamma ^2}\frac{\alpha Ng^2}{64\pi ^2}log\frac{k^2}{\overline{\mu }^2}.$$
(3.60)
One sees that, in the present case, the Gribov-Zwanziger horizon condition is jeopardized by the term $`log\frac{k^2}{\overline{\mu }^2}`$ in expression $`\left(\text{3.60}\right)`$. As it is apparent from the presence of the gauge parameter $`\alpha `$, this term is due to the contribution of the longitudinal components of the gauge field. Note that the longitudinal modes do not contribute to the Green function $`𝒢_{tr}(k)`$ in eq.$`\left(\text{3.49}\right)`$.
## 4 Inclusion of the dimension two condensate $`A_\mu ^aA_\mu ^a`$
The dimension two gluon condensate $`A_\mu ^aA_\mu ^a`$ has received much attention in the last years . This condensate turns out to contribute to the gluon two-point function, as observed in within the operator product expansion. As such, it has to be taken into account when discussing the gluon propagator. A renormalizable effective potential for $`A_\mu ^aA_\mu ^a`$ in linear covariant gauges has been constructed and evaluated in analytic form in . The output of this investigation is that a nonvanishing value of the condensate $`A_\mu ^aA_\mu ^a`$ is favoured since it lowers the vacuum energy. As a consequence, a dynamical tree level gluon mass is generated in the gauge fixed Lagrangian . The inclusion of the condensate $`A_\mu ^aA_\mu ^a`$ in the present framework can be done along the lines outlined in , where the effects of the Gribov copies on the gluon and ghost propagators in the presence of $`A_\mu ^aA_\mu ^a`$ have been worked out in the Landau gauge. Let us begin by giving a brief account of the dynamical mass generation in linear covariant gauges. Following , the dynamical mass generation in these gauges is described by the following action
$$S(A,\sigma )=S_{YM}+S_{gf}+S_\sigma ,$$
(4.61)
where $`S_{YM}`$, $`S_{gf}`$ are the Yang-Mills and the gauge fixing terms, as given in eq.$`\left(\text{2.8}\right)`$. The term $`S_\sigma `$ in eq.$`\left(\text{4.61}\right)`$ contains the auxiliary scalar field $`\sigma `$ and reads
$$S_\sigma =d^4x\left(\frac{\sigma ^2}{2g^2\zeta }+\frac{1}{2}\frac{\sigma }{g\zeta }A_\mu ^aA_\mu ^a+\frac{1}{8\zeta }\left(A_\mu ^aA_\mu ^a\right)^2\right).$$
(4.62)
The introduction of the auxiliary field $`\sigma `$ allows us to study the condensation of the local operator $`A_\mu ^aA_\mu ^a`$. In fact, as shown in , the following relation holds
$$\sigma =\frac{g}{2}A_\mu ^aA_\mu ^a.$$
(4.63)
The dimensionless parameter $`\zeta `$ in expression $`\left(\text{4.62}\right)`$ is needed to account for the ultraviolet divergences present in the vacuum correlation function $`A^2(x)A^2(y)`$. For the details of the renormalizability properties of the local operator $`A_\mu ^aA_\mu ^a`$ in linear covariant gauges we refer to . The action $`S(A,\sigma )`$ is the starting point for evaluating the renormalizable effective potential $`V(\sigma )`$ for the auxiliary field $`\sigma `$, obeying the renormalization group equations. The minimum of $`V(\sigma )`$ occurs for a nonvanishing vacuum expectation value of the auxiliary field , i.e. $`\sigma 0`$. In particular, from expression $`\left(\text{4.61}\right)`$, the first order induced dynamical gluon mass is found to be
$$m^2=\frac{g\sigma }{\zeta _0},$$
(4.64)
where $`\zeta _0`$ is the first contribution to the parameter $`\zeta `$, given by
$`\zeta `$ $`=`$ $`{\displaystyle \frac{\zeta _0}{g^2}}+\zeta _1+O(g^2),`$
$`\zeta _0`$ $`=`$ $`{\displaystyle \frac{3\left(7826\alpha ^2+3\alpha ^3+18\alpha \mathrm{log}\alpha \right)}{2\left(3\alpha 13\right)^2}}{\displaystyle \frac{\left(N^21\right)}{N}}.`$ (4.65)
We remind here that in linear covariant gauges, the Faddeev-Popov ghosts $`\overline{c}^a,`$ $`c^a`$ remain massless, due to the absence of mixing between the composite operators $`A_\mu ^a(x)A_\mu ^a(x)`$ and $`\overline{c}^a(x)c^a(x)`$. This stems from additional Ward identities present in these gauges , which forbid the appearance of the term $`\overline{c}^a(x)c^a(x)`$.
### 4.1 Infrared behavior of the gluon propagator in the presence of $`A_\mu ^aA_\mu ^a`$
It is worth underlining that the action $`S(A,\sigma )`$ leads to a partition function which is still plagued by the Gribov copies. It might be useful to note in fact that the action $`\left(S_{YM}+S_\sigma \right)`$ is left invariant by the local gauge transformations
$`\delta A_\mu ^a`$ $`=`$ $`D_\mu ^{ab}\omega ^b,`$ (4.66)
$`\delta \sigma `$ $`=`$ $`gA_\mu ^a_\mu \omega ^a,`$
$$\delta \left(S_{YM}+S_\sigma \right)=0.$$
(4.67)
Therefore, implementing the restriction to the region $`\stackrel{~}{\mathrm{\Omega }}`$, for the partition function we obtain
$$𝒵=DAD\sigma det\left(_\mu \left(\delta ^{ab}_\mu +gf^{abc}A_\mu ^c\right)\right)e^{\left(\frac{1}{4}{\scriptscriptstyle d^4xF_{\mu \nu }^aF_{\mu \nu }^a}+\frac{1}{2\alpha }{\scriptscriptstyle d^4x\left(A^a\right)^2}+S_\sigma \right)}𝒱(\stackrel{~}{\mathrm{\Omega }}),$$
(4.68)
with the factor $`𝒱(\stackrel{~}{\mathrm{\Omega }})`$ given in eqs.$`\left(\text{2.32}\right)`$,$`\left(\text{2.33}\right)`$. To discuss the gluon propagator we proceed as before and retain only the quadratic terms in expression $`\left(\text{4.68}\right)`$ which contribute to the two-point correlation function $`A_\mu ^a(k)A_\nu ^b(k)`$. Expanding around the nonvanishing vacuum expectation value of the auxiliary field, $`\sigma 0`$, one easily get
$`𝒵_{\mathrm{quadr}}`$ $`=`$ $`𝒩{\displaystyle DA\frac{d\eta }{2\pi i\eta }e^{\eta (1\rho (0,A))}e^{(\frac{1}{4}{\scriptscriptstyle }d^4x((_\mu A_\nu ^a_\mu A_\nu ^a)^2+\frac{1}{2\alpha }{\scriptscriptstyle }d^4x(A^a)^2+\frac{1}{2}m^2{\scriptscriptstyle }d^4x\left(A_\mu ^aA_\mu ^a\right))}}`$ (4.69)
$`=`$ $`𝒩{\displaystyle DA\frac{d\eta }{2\pi i}e^{\eta \mathrm{log}\eta }e^{\frac{1}{2}_qA_\mu ^a(q)𝒬_{\mu \nu }^{ab}A_\nu ^b(q)}},`$
where the factor $`𝒬_{\mu \nu }^{ab}`$ is now given by
$$𝒬_{\mu \nu }^{ab}=\left(\left(q^2+m^2+\frac{\eta Ng^2}{N^21}\frac{1}{2V}\frac{1}{q^2}\right)\delta _{\mu \nu }q_\mu q_\nu \left(\left(1\frac{1}{\alpha }\right)+\frac{\eta Ng^2}{N^21}\frac{1}{2V}\frac{1}{q^4}\right)\right)\delta ^{ab}.$$
(4.70)
Integrating over the gauge field, one has
$$𝒵_{\mathrm{quadr}}=𝒩\frac{d\eta }{2\pi i}e^{\eta \mathrm{log}\eta }\left(det𝒬_{\mu \nu }^{ab}\right)^{\frac{1}{2}}=𝒩^{}\frac{d\eta }{2\pi i}e^{f(\eta )},$$
(4.71)
with
$$f(\eta )=\eta \mathrm{log}\eta \frac{3}{2}(N^21)\underset{q}{}\mathrm{log}\left(q^4+m^2q^2+\frac{\eta Ng^2}{N^21}\frac{1}{2V}\right).$$
(4.72)
Evaluating $`𝒵_{\mathrm{quadr}}`$ at the saddle point, yields
$$𝒵_{\mathrm{quadr}}e^{f(\eta _0)},$$
(4.73)
where $`\eta _0`$ is determined by the minimum condition
$$1\frac{1}{\eta _0}\frac{3}{4}\frac{Ng^2}{V}\underset{q}{}\frac{1}{\left(q^4+m^2q^2+\frac{\eta _0Ng^2}{N^21}\frac{1}{2V}\right)}=0.$$
(4.74)
Taking the thermodynamic limit, $`V\mathrm{}`$, and introducing the Gribov parameter
$$\gamma ^4=\frac{\eta _0Ng^2}{N^21}\frac{1}{2V},V\mathrm{},$$
(4.75)
we get the gap equation in the presence of the dynamical gluon mass, corresponding to a nonvanishing condensate $`A_\mu ^aA_\mu ^a`$, namely
$$\frac{3}{4}Ng^2\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{q^4+m^2q^2+\gamma ^4}=1.$$
(4.76)
Note that the dynamical mass $`m`$ appears now explicitly in the gap equation $`\left(\text{4.76}\right)`$. To obtain the gauge propagator, one goes back to the expression $`\left(\text{4.69}\right)`$ which, when evaluated at the saddle point value $`\eta =\eta _0`$, yields
$$𝒵_{\mathrm{quadr}}=𝒩DAe^{\frac{1}{2}_qA_\mu ^a(q)𝒬_{\mu \nu }^{ab}A_\nu ^b(q)},$$
(4.77)
with
$$𝒬_{\mu \nu }^{ab}=\left(\left(q^2+m^2+\frac{\gamma ^4}{q^2}\right)\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)+\frac{q_\mu q_\nu }{q^2}\left(\frac{q^2}{\alpha }+m^2\right)\right)\delta ^{ab}.$$
(4.78)
Thus, for the gauge propagator in the presence of the dynamical gluon mass $`m`$ we get
$$A_\mu ^a(q)A_\nu ^b(q)=\delta ^{ab}\left(\frac{q^2}{q^4+m^2q^2+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)+\frac{\alpha }{q^2+\alpha m^2}\frac{q_\mu q_\nu }{q^2}\right).$$
(4.79)
We note that, due to the presence of the mass $`m`$, the infrared suppression of the transverse component of the gluon propagator is enforced. Moreover, also the longitudinal component gets suppressed.
### 4.2 The infrared behavior of $`𝒢_{tr}(k)`$ in the presence of $`A_\mu ^aA_\mu ^a`$
It remains now to discuss the infrared behavior of the Green function $`𝒢_{tr}(k)`$ in the presence of $`A_\mu ^aA_\mu ^a`$. This can be easily worked out by repeating the analysis done in the previous sections. From the expression of the gluon propagator $`\left(\text{4.79}\right)`$, it follows that
$$𝒢_{tr}(k)\frac{1}{k^2}\frac{1}{1\rho (k)},$$
(4.80)
with
$`\rho (k)`$ $`=`$ $`g^2{\displaystyle \frac{N}{N^21}}{\displaystyle \frac{k_\mu k_\nu }{k^2}}{\displaystyle \frac{d^4q}{\left(2\pi \right)^4}\frac{1}{\left(kq\right)^2}A_\mu ^{Ta}(q)A_\nu ^{Ta}(q)}`$ (4.81)
$`=`$ $`g^2N{\displaystyle \frac{k_\mu k_\nu }{k^2}}{\displaystyle \frac{d^4q}{\left(2\pi \right)^4}\frac{1}{\left(kq\right)^2}\frac{q^2}{q^4+m^2q^2+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)}.`$
Also, from the gap equation $`\left(\text{4.76}\right)`$, one has
$$Ng^2\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{q^4+m^2q^2+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)=\delta _{\mu \nu },$$
(4.82)
so that
$`1\rho (k)`$ $`=`$ $`Ng^2{\displaystyle \frac{k_\mu k_\nu }{k^2}}{\displaystyle \frac{d^4q}{\left(2\pi \right)^4}\frac{k^22qk}{\left(kq\right)^2}\frac{1}{q^4+m^2q^2+\gamma ^4}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)}.`$
Thus, for $`k0`$,
$$\left(1\rho (k)\right)_{k0}\frac{3Ng^2𝒥}{4}k^2,$$
(4.84)
where $`𝒥`$ stands for the value of the integral
$$𝒥=\frac{d^4q}{\left(2\pi \right)^4}\frac{1}{q^2(q^4+m^2q^2+\gamma ^4)},$$
(4.85)
which is ultraviolet finite. Therefore, for the Green function $`𝒢_{tr}(k)`$, we get
$$\left(𝒢_{tr}(k)\right)_{k0}\frac{4}{3Ng^2𝒥}\frac{1}{k^4},$$
(4.86)
exhibiting the infrared enhancement which, thanks to the gap equation $`\left(\text{4.76}\right)`$, turns out to hold also in the presence of the gluon condensate $`A_\mu ^aA_\mu ^a`$.
## 5 Comparison with the results obtained from lattice simulations and from the Schwinger-Dyson equations
Having investigated the infrared behavior of the gluon propagator and of the Green function $`\left(𝒢_{tr}(k)\right)`$, as summarized by equations $`\left(\text{4.79}\right)`$ and $`\left(\text{4.86}\right)`$, it is useful to make a comparison with the results already available from lattice simulations and from the studies of the Schwinger-Dysons equations. Let us begin with the lattice data
### 5.1 Comparison with the lattice data
In a series of papers , Giusti et al. have managed to put the linear covariant gauges on the lattice. This has allowed for a numerical investigation of the transverse as well as of the longitudinal component of the gluon propagator. Following , let us introduce the transverse and longitudinal form factors $`D_T(q)`$ and $`D_L(q)`$ through
$$A_\mu ^a(q)A_\nu ^b(q)=\delta ^{ab}\left(\frac{D_T(q)}{q^2}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)+\frac{D_L(q)}{q^2}\frac{q_\mu q_\nu }{q^2}\right).$$
(5.87)
The results obtained in show that both $`D_T(q)`$ and $`D_L(q)`$ are suppressed in the low momentum region, see for instance Fig.3 and Fig.4 of . Our results are in qualitative agreement with the lattice data. Indeed, from the expression $`\left(\text{4.79}\right)`$, we obtain
$`D_T(q)`$ $`=`$ $`{\displaystyle \frac{q^4}{q^4+m^2q^2+\gamma ^4}},`$
$`D_L(q)`$ $`=`$ $`{\displaystyle \frac{\alpha q^2}{q^2+\alpha m^2}},`$ (5.88)
exhibiting infrared suppression. Note that, at least within the approximation considered in the present work, the suppression of the longitudinal form factor $`D_L(q)`$ in eq.$`\left(\text{5.88}\right)`$ is a consequence of the dynamical gluon mass, due to the gluon condensate $`A_\mu ^aA_\mu ^a`$, as already pointed out in . Concerning now the ghost propagator and the Green function $`𝒢_{tr}(k)`$, to our knowledge, no results from lattice data are available so far.
### 5.2 Comparison with the results obtained from the Schwinger-Dysons equations
The infrared behavior of the gluon and ghost propagator has been investigated within the Schwinger-Dyson framework in . Here, a power-law Ansatz for the transverse and longitudinal form factors of the gluon propagator as well as for the ghost form factor $`D_{gh}(q)`$ has been employed, according to
$`D_T(q)`$ $``$ $`\left(q^2\right)^\sigma ,`$
$`D_L(q)`$ $``$ $`\left(q^2\right)^\rho ,`$ (5.89)
and
$`\overline{c}^a(q)c^b(q)`$ $`=`$ $`\delta ^{ab}{\displaystyle \frac{D_{gh}(q)}{q^2}},`$
$`D_{gh}(q)`$ $``$ $`\left(q^2\right)^\beta ,`$ (5.90)
The results obtained for the infrared exponents $`(\sigma ,\rho ,\beta )`$ turn out to be similar to those of the Landau gauge, namelyThe explicit values of these infrared exponents as well as their dependence from the gauge parameter can be found in .
$`\sigma `$ $`>`$ $`0,`$
$`\rho `$ $`>`$ $`0,`$
$`\beta `$ $`=`$ $`{\displaystyle \frac{\sigma }{2}}={\displaystyle \frac{\rho }{2}},`$ (5.91)
indicating an infrared suppression of the transverse and longitudinal gluon form factors, and an infrared enhancement of the ghost propagator. Concerning the gluon propagator, these results are in qualitative agreement with our results as well as with the lattice data. However, concerning the ghost propagator, we have found that, instead of the ghost form factor, the quantity which is enhanced in the infrared is $`𝒢_{tr}(k)`$. For a better understanding of this point, it is worth reminding here that the result for the infrared exponents in eq.$`\left(\text{5.91}\right)`$ has been obtained by using a bare-vertex truncation scheme . This approximation has been proven successful in the Landau gauge . In particular, in the Landau gauge, no qualitative difference has been found if bare vertices are replaced by vertices dressed according to the Slanov-Taylor identities. This feature of the Landau gauge is believed to be deeply related to the nonrenormalization theorem of the ghost-antighost-gluon vertex, which holds to all orders of perturbation theory . Recently, the nonrenormalization theorem of the ghost-antighost-gluon vertex in the Landau gauge has been investigated through lattice simulations in , which have provided indications of its validity beyond perturbation theory. However, to our knowledge, no such a theorem is available in linear covariant gauges, for a nonvanishing value of the gauge parameter $`\alpha `$. Furthermore, according to the authors , it is yet an open question whether the values of the infrared exponents in eq.$`\left(\text{5.91}\right)`$ remain unchanged if bare vertices are replaced by dressed ones. Our results suggest that a different behavior might be expected when dressed vertices would be employed.
## 6 Conclusion
In this work we have attempted at analyzing the effects of the Gribov copies on the gluon propagator in linear covariant gauges. By considering small values of the gauge parameter $`\alpha `$, a few properties of the Gribov copies have been established, allowing us to investigate the infrared behavior of the gluon two-point function.
As in the case of the Landau gauge, it turns out that the transverse component of the gluon propagator is suppressed in the infrared. Moreover, the longitudinal part is left unchanged, as shown in eq.$`\left(\text{3.48}\right)`$. The infrared behavior of the gluon propagator has been investigated also in the presence of the gluon condensate $`A_\mu ^aA_\mu ^a`$. In this case, the infrared suppression of the transverse component is enforced. Furthermore, its longitudinal component turns out to be suppressed as well, as expressed by eq.$`\left(\text{4.79}\right)`$. These results are in qualitative agreement with those obtained from lattice simulations and from the analysis of the Schwinger-Dyson equations.
Concerning now the behavior of the ghost fields in linear covariant gauges, the output of our analysis is that, instead of the ghost propagator, the Green function which exhibits infrared enhancement is given by $`𝒢_{tr}(k)`$, as defined in eq.$`\left(\text{1.5}\right)`$. It should be remarked that $`𝒢_{tr}(k)`$ does not coincide with the ghost propagator for a generic value of the gauge parameter $`\alpha `$. However, $`𝒢_{tr}(k)`$ reduces to the ghost two-point function for vanishing $`\alpha `$, so that our results turn out to coincide with those of the Landau gauge in the limit $`\alpha 0`$.
Needless to say, many aspects of the covariant linear gauges remain still to be investigated. A partial list of them is:
As already pointed out, a suitable auxiliary functional corresponding to the linear covariant gauge fixing condition is not yet at our disposal. As in the case of the Landau gauge , this functional could be very helpful for a characterization of the properties of the Gribov copies not attainable by infinitesimal gauge transformations.
From the present analysis, it emerges that the Green function $`𝒢_{tr}(k)`$ has a special role, as it obeys the Gribov-Zwanziger horizon condition and reduces to the ghost two-point function in the Landau gauge. Although its dependence from the transverse component $`A_\mu ^{aT}`$ of the gauge field suggests that it might have a deeper meaning, it would be worth to have a better understanding of $`𝒢_{tr}(k)`$.
It would be useful to have a consistent framework to compute quantum corrections to the gluon propagator and to $`𝒢_{tr}(k)`$. This would amount to construct a local renormalizable action in linear covariant gauges incorporating the effects of the Gribov copies, as done by Zwanziger in the Landau gauge .
Finally, it would be interesting to have more data in the linear covariant gauges from lattice simulations on the gluon and ghost propagators as well as on the Green function $`𝒢_{tr}(k)`$.
## Acknowledgments.
The Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq-Brazil), the Faperj, Fundação de Amparo à Pesquisa do Estado do Rio de Janeiro, the SR2-UERJ and the Coordenação de Aperfeiçoamento de Pessoal de Nível Superior (CAPES) are gratefully acknowledged for financial support.
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# 1 Introduction
## 1 Introduction
The origin of the baryon asymmetry in our Universe (BAU) has always been one of the central topics in particle cosmology. Recently, the high-precision determination of many cosmological parameters, including the baryon-to-photon ratio of number densities, $`\eta _B6.1\times 10^{10}`$ , has given renewed momentum for extensive studies on this topic . The established BAU provides one of the strongest pieces of evidence towards physics beyond the Standard Model (SM). One interesting suggestion for explaining the BAU, known as leptogenesis , is linked with neutrinos. Although strictly massless in the SM, neutrinos can naturally acquire their small observed mass through the presence of superheavy partners and the so-called seesaw mechanism . These superheavy neutrinos are singlets under the SM gauge group and may therefore possess large Majorana masses that violate lepton number ($`L`$) conservation by two units. In an expanding Universe, these heavy Majorana neutrinos will in general decay out of equilibrium, potentially generating a net lepton asymmetry. The so-produced lepton asymmetry will eventually be converted into the observed BAU by means of in-thermal equilibrium $`(B+L)`$-violating sphaleron interactions .
One difficulty faced by ordinary seesaw models embedded in grand unified theories (GUTs) is associated with the natural mass scale of the heavy Majorana neutrinos. This is expected to be of order the GUT scale $`M_{\mathrm{GUT}}=10^{16}`$ GeV. On the other hand, inflationary supergravity models generically predict a reheating temperature $`T_{\mathrm{reh}}`$ of order $`10^9`$ GeV. In these models, a significant constraint on the upper bound for $`T_{\mathrm{reh}}`$ comes from the requirement that gravitinos are underabundant in the early Universe and so their late decays do not disrupt the nucleosynthesis of the light elements . However, the low $`T_{\mathrm{reh}}`$ gives rise to another constraint within the context of thermal leptogenesis. The heavy Majorana neutrino, whose $`L`$-violating decays are responsible for the BAU, has to be somewhat lighter than $`T_{\mathrm{reh}}10^9`$ GeV, so as to be abundantly produced in the early Universe. Such a mass for the heavy Majorana neutrino should be regarded as unnaturally low for GUT-scale thermal leptogenesis. Finally, further constraints on successful GUT-scale leptogenesis may be obtained from solar and atmospheric neutrino data .
The aforementioned problem with a low reheating temperature may be completely avoided in models that realize low-scale thermal leptogenesis . In particular, the lowering of the scale may rely on a dynamical mechanism, in which heavy-neutrino self-energy effects on the leptonic asymmetry become dominant and get resonantly enhanced , when a pair of heavy Majorana neutrinos has a mass difference comparable to the heavy neutrino decay widths. In , this dynamical mechanism was termed resonant leptogenesis (RL). As a consequence of RL, the heavy Majorana mass scale can be as low as $``$ 1 TeV in complete agreement with the solar and atmospheric neutrino data .
A crucial model-building aspect of RL models is that such models have to rely on a nearly degenerate heavy neutrino mass spectrum. Although, without any additional lepton-flavour symmetry, such a requirement would appear very fine-tuned, there is no theoretical or phenomenologically compelling reason that would prevent the singlet neutrino sector of the SM from possessing such a symmetry. Specifically, the RL model discussed in , which was motivated by E<sub>6</sub> unified theories , was based on a particular lepton symmetry in the heavy neutrino sector. This lepton symmetry was broken very approximately by GUT- and/or Planck-scale suppressed operators of dimension 5 and higher. In , another RL scenario was put forward based on the Froggatt–Nielsen (FN) mechanism , where two of the heavy neutrinos naturally had a mass difference comparable to their decay widths. Recently, several constructions of RL models appeared in the literature within the context of supersymmetric theories , or even embedded in SO(10) unified theories .
One of the great advantages of RL models is that their predictions for the BAU are almost independent of the primordial $`L`$-number, $`B`$-number and heavy neutrino abundances . This fact may be explained as follows: in RL scenarios, the $`L`$-violating decay widths of the heavy Majorana neutrinos can be significantly larger than the Hubble expansion rate $`H`$ of the Universe. As a consequence, the heavy Majorana neutrinos can rapidly thermalize from their decays, inverse decays and scatterings with the other SM particles in the plasma, even if there were no heavy Majorana neutrinos at high temperatures. Moreover, in this high temperature regime, any pre-existing lepton asymmetry will rapidly be driven to zero, due to the $`L`$-violating inverse decays and scattering processes which are almost in thermal equilibrium. As the Universe cools down, a net lepton asymmetry can be created at temperatures just below the heavy neutrino mass as a consequence of the aforementioned CP-violating resonant enhancement that occurs in RL models. This $`L`$ asymmetry will survive wash-out effects and will be converted by the $`(B+L)`$-violating sphalerons into the observed BAU.
In this paper we provide a detailed study of a new variant of RL where a given single lepton flavour asymmetry is resonantly produced by the quasi-in-equilibrium decays of heavy Majorana neutrinos of a particular family type. Such a variant of RL was first discussed in , and for the case of the $`\tau `$-lepton number this mechanism has been called resonant $`\tau `$-leptogenesis (R$`\tau `$L). This mechanism makes use of the property that, in addition to $`BL`$, sphalerons preserve the individual quantum numbers $`\frac{1}{3}BL_{e,\mu ,\tau }`$ . In a R$`\tau `$L model, the generated excess in the $`L_\tau `$ number will be converted into the observed BAU, provided the $`L_\tau `$-violating reactions are not strong enough to wash out such an excess.
Although our focus will be on minimal non-supersymmetric 3-generation RL models, supersymmetry could account for the origin of the electroweak-scale heavy Majorana neutrinos. In particular, one may tie the singlet Majorana neutrino mass scale $`m_N`$ to the $`\mu `$-parameter through the vacuum expectation value (VEV) of a chiral singlet superfield $`\widehat{S}`$ . The proposed model is a variant of the so-called Next-to-Minimal Supersymmetric Standard Model (NMSSM) and is described by the following superpotential (summation over repeated indices implied):
$$W=W_{\mathrm{MSSM}}(\mu =0)+h_{ij}^{\nu _R}\widehat{L}_i\widehat{H}_2\widehat{\nu }_{jR}+\lambda \widehat{S}\widehat{H}_1\widehat{H}_2+\frac{\rho }{2}\widehat{S}\widehat{\nu }_{iR}\widehat{\nu }_{iR}+\frac{\kappa }{3}\widehat{S}^3,$$
(1.1)
where $`W_{\mathrm{MSSM}}(\mu =0)`$ is the superpotential of the well-known Minimal Supersymmetric Standard Model (MSSM) without the $`\mu `$-term, and $`\widehat{H}_{1,2}`$, $`\widehat{L}_{1,2,3}`$ and $`\widehat{\nu }_{1,2,3R}`$ are the Higgs-doublet, lepton-doublet and right-handed neutrino superfields, respectively. Once the scalar component of $`\widehat{S}`$ develops a VEV $`v_S`$, then both the would-be $`\mu `$-parameter, $`\mu =\lambda v_S`$, and the SO(3)-symmetric singlet scale, $`m_N=\frac{1}{2}\rho v_S`$, are expected to be comparable in magnitude (asumming that $`\lambda \rho `$), thus providing a natural framework for the possible existence of 3 nearly degenerate electroweak-scale heavy Majorana neutrinos . In this minimal extension of the MSSM, the predictions for the BAU will depend on the size of the soft SUSY-breaking mass scale $`M_{\mathrm{SUSY}}`$. However, if $`M_{\mathrm{SUSY}}`$ is relatively larger than the singlet Majorana neutrino mass scale $`m_N`$, e.g. $`M_{\mathrm{SUSY}}\stackrel{>}{_{}}2m_N`$, the dominant source of leptogenesis will be the minimal non-supersymmetric sector that we are studying here, so our predictions will remain almost unaffected in this case.
As mentioned above, single lepton-flavour effects on the net $`L`$ and $`B`$ asymmetries play a key role in R$`\tau `$L models. To properly treat these as well as SM chemical potential effects, the relevant network of the Boltzmann equations (BEs) needs to be extended consistently. In particular, single lepton-flavour effects can have a dramatic impact on the predictions for the $`B`$ asymmetry. These predictions for the BAU can differ by many orders of magnitude with respect to those obtained in the conventional BE formalism, which is commonly used in the literature. Although our primary interest will be to analyze RL models, we should stress that single lepton-flavour effects could also significantly affect the predictions obtained in hierarchical leptogenesis scenarios. The improved set of BEs derived here will therefore be of general use.
Another important question we wish to address is whether the leptogenesis scale can be lowered to energies 100–250 GeV, very close to the critical temperature $`T_c`$, where the electroweak phase transition occurs. In this temperature region, freeze-out effects from sphaleron processes dropping out of equilibrium need to be considered, as they can significantly modify the predicted values for the final baryon asymmetry. Our treatment of these sphaleron freeze-out effects will be approximate and based on the calculations of . Our approximate treatment is motivated by the fact that, within the framework of RL models, the creation of a net lepton asymmetry at the electroweak scale does not require the electroweak phase transition to be strongly first order.
Most importantly, in models where the BAU is generated from an individual lepton-number asymmetry, a range of testable phenomenological implications may arise. The key aspect is that scenarios such as R$`\tau `$L can contain heavy Majorana neutrinos with appreciable Yukawa couplings to electrons and muons. The (normalized to the SM) $`W^\pm `$-boson couplings of $`e`$ and $`\mu `$ leptons to these heavy Majorana neutrinos could be as large as $`10^2`$. For electroweak-scale heavy neutrinos, such couplings would be sufficient to produce these particles at future $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ colliders. Furthermore, minimal (non-supersymmetric) 3-generation R$`\tau `$L models can predict $`\mu e\gamma `$ and $`\mu e`$ conversion in nuclei at rates that can be tested by the foreseeable experiments MEG at PSI and MECO at BNL , respectively. Finally, R$`\tau `$L models naturally realize an inverted hierarchy for the light neutrino spectrum and therefore also predict neutrinoless double beta ($`0\nu \beta \beta `$) decay with a sizeable effective neutrino mass $`|m|`$, as large as $`0.02`$ eV. This value falls within reach of proposals for future $`0\nu \beta \beta `$-decay experiments sensitive to $`|m|0.01`$–0.05 eV , e.g. CUORE ($`{}_{}{}^{130}\mathrm{Te}`$), GERDA ($`{}_{}{}^{76}\mathrm{Ge}`$), EXO ($`{}_{}{}^{136}\mathrm{Xe}`$), MOON ($`{}_{}{}^{100}\mathrm{Xe}`$), XMASS ($`{}_{}{}^{136}\mathrm{Xe}`$), CANDLES ($`{}_{}{}^{48}\mathrm{Ca}`$), SuperNEMO ($`{}_{}{}^{82}\mathrm{Se}`$) etc.
Our paper has been organized as follows: Section 2 presents a minimal model for resonant $`\tau `$-leptogenesis. In Section 3 we derive the BEs for single lepton flavours, by carefully taking into account SM chemical potential effects. Technical details pertinent to this derivation have been relegated to Appendix A. In Section 4 we review the calculation of out of equilibrium sphaleron effects at the electroweak phase transition and apply it to leptogenesis. In Section 5 we give several numerical examples of R$`\tau `$L models, focusing our attention on scenarios that can be tested at future $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ colliders and in low-energy experiments. In particular, in Section 6, we present predictions for lepton-flavour-violating (LFV) processes, such as $`\mu e\gamma `$, $`\mu eee`$ and $`\mu e`$ conversion in nuclei. Finally, we present our conclusions and future prospects in Section 7.
## 2 Minimal Model for Resonant $`𝝉`$-Leptogenesis
There have been several studies on RL models in the literature . Here, we will focus our attention on a variant of resonant leptogenesis where the BAU is generated by the production of an individual lepton number . For definiteness, we consider a minimal (non-supersymmetric) model for R$`\tau `$L.
Let us start our discussion by briefly reviewing the relevant low-energy structure of the SM symmetrically extended with one singlet neutrino $`\nu _{iR}`$ per $`i`$ family (with $`i=1,2,3`$). The leptonic Yukawa and Majorana sectors of such a model are given by the Lagrangian
$`_{\mathrm{M},\mathrm{Y}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=1}{\overset{3}{}}}\left((\overline{\nu }_{iR})^C(M_S)_{ij}\nu _{jR}+\text{h.c.}\right)`$ (2.1)
$`+{\displaystyle \underset{i=e,\mu ,\tau }{}}\left[\widehat{h}_{ii}^l\overline{L}_i\mathrm{\Phi }l_{iR}+\left({\displaystyle \underset{j=1}{\overset{3}{}}}h_{ij}^{\nu _R}\overline{L}_i\stackrel{~}{\mathrm{\Phi }}\nu _{jR}+\text{h.c.}\right)\right],`$
where $`L_i=(\nu _{iL},l_{iL})^T`$ are the left-handed lepton doublets <sup>1</sup><sup>1</sup>1Occasionally we will also denote the individual lepton numbers with $`L_{e,\mu ,\tau }`$. This apparent abuse of notation should cause no confusion to the reader, as the precise meaning of $`L_{e,\mu ,\tau }`$ can be easily inferred from the context., $`l_{iR}`$ are the right-handed leptons, and $`\stackrel{~}{\mathrm{\Phi }}`$ is the isospin conjugate of the Higgs doublet $`\mathrm{\Phi }`$.
In the Lagrangian (2.1), we have defined the individual lepton numbers $`L_{e,\mu ,\tau }`$ in the would-be charged-lepton mass basis, where the charged-lepton Yukawa matrix $`\widehat{h}^l`$ is positive and diagonal. In fact, without loss of generality, it can be shown that sphaleron transitions exhibit a U(3) flavour symmetry and so they can be rotated to become flavour diagonal in the same would-be mass basis. To prove this, one may write the operator $`O_{B+L}`$ responsible for $`B+L`$-violating sphaleron transitions as follows (group-invariant contraction of the colour and weak degrees of freedom implied) :
$$O_{B+L}=\underset{i=1}{\overset{3}{}}Q_i^{}Q_i^{}Q_i^{}L_i^{},$$
(2.2)
where $`Q_i^{}`$ and $`L_i^{}`$ denote the quark and lepton doublets defined in an arbitrary weak basis. The operator $`O_{B+L}`$ contains the fully antisymmetric operator combinations: $`Q_1^{}Q_2^{}Q_3^{}`$ and $`L_1^{}L_2^{}L_3^{}`$, which are invariant under U(3) flavour rotations . Thus, we can use this U(3)-rotational freedom to render the charged lepton and up-quark sectors flavour diagonal and positive.
To obtain a phenomenologically relevant model, at least 3 singlet heavy Majorana neutrinos $`\nu _{1,2,3R}`$ are needed and these have to be nearly degenerate in mass. To ensure the latter, we assume that to leading order, the heavy neutrino sector is SO(3) symmetric, i.e.
$$M_S=m_N\mathbf{\hspace{0.17em}1}_3+\mathrm{\Delta }M_S,$$
(2.3)
where $`\mathrm{𝟏}_3`$ is the $`3\times 3`$ identity matrix and $`\mathrm{\Delta }M_S`$ is a general SO(3)-breaking matrix. As we will discuss below, compatibility with the observed light neutrino masses and mixings requires that $`(\mathrm{\Delta }M_S)_{ij}/m_N\stackrel{<}{_{}}10^7`$, for electroweak-scale heavy Majorana neutrinos, i.e. for $`m_N0.1`$–1 TeV. One could imagine that the soft SO(3)-breaking matrix $`\mathrm{\Delta }M_S`$ originates from a sort of Froggatt–Nielsen mechanism .
In order to account for the smallness of the light neutrino masses, we require that the neutrino Yukawa sector possesses a leptonic U(1)<sub>l</sub> symmetry. This will explicitly break the imposed SO(3) symmetry of the heavy neutrino sector to a particular subgroup SO(2) $``$ U(1)<sub>l</sub>. For example, one possibility relevant to R$`\tau `$L is to couple all lepton doublets to a particular heavy neutrino combination: $`\frac{1}{\sqrt{2}}(\nu _{2R}+i\nu _{3R})`$. In detail, the U(1)<sub>l</sub> charges of the fields are
$$Q(L_i)=Q(l_{iR})=1,Q\left(\frac{\nu _{2R}+i\nu _{3R}}{\sqrt{2}}\right)=Q\left(\frac{\nu _{2R}i\nu _{3R}}{\sqrt{2}}\right)=1,Q(\nu _{1R})=0.$$
(2.4)
As a result of the U(1)<sub>l</sub> symmetry, the matrix for the neutrino Yukawa couplings reads:
$$h^{\nu _R}=(\begin{array}{ccc}0& ae^{i\pi /4}& ae^{i\pi /4}\\ 0& be^{i\pi /4}& be^{i\pi /4}\\ 0& ce^{i\pi /4}& ce^{i\pi /4}\end{array})+\delta h^{\nu _R}.$$
(2.5)
In the above, $`a,b`$ and $`c`$ are arbitrary complex parameters of the model. For electroweak-scale heavy neutrinos, the absolute value of these parameters has to be smaller than about $`10^2`$, for phenomenological reasons to be discussed below and in Section 6. In particular, the requirement that an excess in $`L_\tau `$ is protected from wash-out effects leads to the relatively stronger constraint $`|c|\stackrel{<}{_{}}10^5`$. In addition, $`\delta h^{\nu _R}`$ is a $`3\times 3`$ matrix that parameterizes possible violations of the U(1)<sub>l</sub> symmetry. It should be noted that the charged lepton sector and the leading SO(3)-invariant form of the heavy neutrino mass matrix are compatible with the U(1)<sub>l</sub> symmetry.
In this paper we shall not address the possible origin of the smallness of the SO(3)- and U(1)<sub>l</sub>-breaking parameters $`(\mathrm{\Delta }M_S)_{ij}`$ and $`\delta h_{ij}^{\nu _R}`$, as there are many different possibilities that could be considered for this purpose, e.g. the Froggatt–Nielsen mechanism , Planck- or GUT-scale lepton-number breaking . Instead, in our model-building we will require that the symmetry breaking terms do not induce radiative effects much larger than the tree-level contributions. This naturalness condition will be applied to the light and heavy neutrino mass matrices $`𝐦^\nu `$ and $`M_S`$, respectively.
We start by observing that the U(1)<sub>l</sub> symmetry is sufficient to ensure the vanishing of the light neutrino mass matrix $`𝐦^\nu `$ . In fact, if U(1)<sub>l</sub> is an exact symmetry of the theory, the light neutrino mass matrix will vanish to all orders in perturbation theory . To leading order in the U(1)<sub>l</sub>-breaking parameters $`\mathrm{\Delta }M_S`$, the tree-level light neutrino mass matrix $`𝐦^\nu `$ is given by
$$𝐦^\nu =\frac{v^2}{2}h^{\nu _R}M_S^1(h^{\nu _R})^T=\frac{v^2}{2m_N}\left(\frac{h^{\nu _R}\mathrm{\Delta }M_S(h^{\nu _R})^T}{m_N}h^{\nu _R}(h^{\nu _R})^T\right),$$
(2.6)
where $`v=2M_W/g_w=245`$ GeV is the vacuum expectation value of the SM Higgs field $`\mathrm{\Phi }`$. As a minimal departure from U(1)<sub>l</sub> in the neutrino Yukawa sector, we consider that this leptonic symmetry is broken only by $`\nu _{1R}`$, through
$$\delta h^{\nu _R}=(\begin{array}{ccc}\epsilon _e& 0& 0\\ \epsilon _\mu & 0& 0\\ \epsilon _\tau & 0& 0\end{array}).$$
(2.7)
In this case, the tree-level light neutrino mass matrix (2.6) takes on the form
$$𝐦^\nu =\frac{v^2}{2m_N}(\begin{array}{ccc}\frac{\mathrm{\Delta }m_N}{m_N}a^2\epsilon _e^2& \frac{\mathrm{\Delta }m_N}{m_N}ab\epsilon _e\epsilon _\mu & \frac{\mathrm{\Delta }m_N}{m_N}ac\epsilon _e\epsilon _\tau \\ \frac{\mathrm{\Delta }m_N}{m_N}ab\epsilon _e\epsilon _\mu & \frac{\mathrm{\Delta }m_N}{m_N}b^2\epsilon _\mu ^2& \frac{\mathrm{\Delta }m_N}{m_N}bc\epsilon _\mu \epsilon _\tau \\ \frac{\mathrm{\Delta }m_N}{m_N}ac\epsilon _e\epsilon _\tau & \frac{\mathrm{\Delta }m_N}{m_N}bc\epsilon _\mu \epsilon _\tau & \frac{\mathrm{\Delta }m_N}{m_N}c^2\epsilon _\tau ^2\end{array}),$$
(2.8)
where $`\mathrm{\Delta }m_N=2(\mathrm{\Delta }M_S)_{23}+i[(\mathrm{\Delta }M_S)_{33}(\mathrm{\Delta }M_S)_{22}]`$. It is interesting to notice that in this type of U(1)<sub>l</sub> breaking, the parameters $`\epsilon _{e,\mu ,\tau }`$ enter the tree-level light neutrino mass matrix $`𝐦^\nu `$ quadratically. As a consequence, one finds that for $`m_Nv`$, these U(1)<sub>l</sub>-breaking parameters need not be much smaller than the electron Yukawa coupling $`h_e10^6`$. Moreover, one should observe that only a particular combination of soft SO(3)- and U(1)<sub>l</sub>-breaking terms $`(\mathrm{\Delta }M_S)_{ij}`$ appears in $`𝐦^\nu `$ through $`\mathrm{\Delta }m_N`$. Nevertheless, for electroweak-scale heavy neutrinos with mass differences $`|\mathrm{\Delta }m_N|/m_N\stackrel{<}{_{}}10^7`$, one should have $`|a|,|b|\stackrel{<}{_{}}10^2`$ to avoid getting too large light neutrino masses much above 0.5 eV. As we will see more explicitly in Section 5, for the R$`\tau `$L scenario under study, the favoured solution will be an inverted hierarchical neutrino mass spectrum with large $`\nu _e\nu _\mu `$ and $`\nu _\mu \nu _\tau `$ mixings .
In addition to the tree level contributions given in (2.8), there are $`Z`$\- and Higgs-boson-mediated threshold contributions $`\delta 𝐦^\nu `$ to $`𝐦^\nu `$ . The contributing graphs are displayed in Fig. 1. In the heavy neutrino mass basis, where $`M_S\mathrm{diag}(m_{N_1},m_{N_2},m_{N_3})`$, with $`m_{N_1}m_{N_2}m_{N_3}`$, and $`h^{\nu _R}h^\nu `$, these finite radiative corrections may conveniently be expressed as follows :
$`(\delta 𝐦^\nu )_{ll^{}}`$ $`=`$ $`{\displaystyle \frac{\alpha _w}{32\pi }}{\displaystyle \underset{\alpha =1,2,3}{}}{\displaystyle \frac{h_{l\alpha }^\nu h_{l^{}\alpha }^\nu v^2}{m_{N_\alpha }}}[{\displaystyle \frac{3M_Z^2}{M_W^2}}(B_0(0,m_{N_\alpha }^2,M_Z^2)B_0(0,0,M_Z^2))`$ (2.9)
$`+{\displaystyle \frac{m_{N_\alpha }^2}{M_Z^2}}(B_0(0,m_{N_\alpha }^2,M_H^2)B_0(0,m_{N_\alpha }^2,m_{N_\alpha }^2))],`$
where $`\alpha _w=g_w^2/(4\pi )`$ and $`M_H`$ is the SM Higgs boson mass. In (2.9), $`B_0(0,m_1^2,m_2^2)`$ is the usual Pasarino–Veltman one-loop function , i.e.
$$B_0(0,m_1^2,m_2^2)=C_{\mathrm{UV}}+\mathrm{\hspace{0.25em}1}\mathrm{ln}\left(\frac{m_1m_2}{\mu ^2}\right)+\frac{m_1^2+m_2^2}{m_1^2m_2^2}\mathrm{ln}\left(\frac{m_2}{m_1}\right),$$
(2.10)
and $`C_{\mathrm{UV}}`$ is a UV infinite constant. Moreover, in writing (2.9), we have neglected terms of order $`[(h_{l\alpha }^\nu )^4v^3]/m_{N_\alpha }^2`$, which are suppressed by higher powers of the small Yukawa couplings. It can easily be verified that the radiative lepton-number-violating contribution $`\delta 𝐦^\nu `$ to the light neutrino mass matrix is UV finite and $`\mu `$-scale independent. For $`m_{N_\alpha }^2M_H^2`$ and $`(m_{N_\alpha }m_{N_1})/m_{N_1}1`$, the expression (2.9) evaluated in the original weak basis simplifies to
$$\delta 𝐦^\nu =\frac{\alpha _w}{16\pi }\frac{M_H^2+3M_Z^2}{M_W^2}\frac{v^2}{m_N}\frac{h^{\nu _R}\mathrm{\Delta }M_S(h^{\nu _R})^T}{m_N}.$$
(2.11)
For electroweak-scale heavy Majorana neutrinos $`m_{N_\alpha }v`$ and $`M_H=120`$–200 GeV, one may estimate that for $`(m_{N_\alpha }m_{N_1})/m_{N_1}\stackrel{<}{_{}}10^7`$ and $`|a|,|b|\stackrel{<}{_{}}10^2`$, the finite radiative effects $`\delta 𝐦^\nu `$ stay well below 0.01 eV. In fact, up to an overall coupling-suppressed constant, these corrections have the same analytic form as the first term on the RHS of (2.6). They can be absorbed by appropriately rescaling $`\mathrm{\Delta }m_N`$ defined after (2.8). As a consequence, these finite radiative effects do not modify the parametric dependence of the tree-level light neutrino mass matrix given in (2.8).
We now turn our attention to the heavy Majorana neutrino sector. In this case, renormalization-group (RG) running effects become very significant. These effects explicitly break the SO(3)-symmetric form of the heavy neutrino mass matrix, $`M_S(M_X)=m_N\mathbf{\hspace{0.17em}1}_3`$, imposed at some high energy scale $`M_X`$, e.g. at the GUT scale. A fairly good quantitative estimate of these SO(3)-breaking effects can be obtained by solving the RG equation for the heavy neutrino mass matrix $`M_S`$:
$$\frac{dM_S}{dt}=\frac{1}{16\pi ^2}\left\{\left[(h^{\nu _R})^{}h^{\nu _R}\right]M_S+M_S\left[(h^{\nu _R})^T(h^{\nu _R})^{}\right]\right\},$$
(2.12)
with $`t=\mathrm{ln}(M_X/\mu )`$. Considering that $`h^{\nu _R}`$ has only a mild RG-scale dependence and assuming that $`M_S(M_X)=m_N\mathbf{\hspace{0.17em}1}_3`$ at some high scale $`M_X`$, we may calculate the RG effects by running from $`M_X`$ to the low-energy scale $`m_Nv`$ through the relation
$`M_S(m_N)`$ $`=`$ $`M_S(M_X){\displaystyle \frac{m_N}{8\pi ^2}}\mathrm{Re}\left[(h^{\nu _R})^{}h^{\nu _R}\right]\mathrm{ln}\left({\displaystyle \frac{M_X}{m_N}}\right)`$ (2.13)
$`=M_S(M_X){\displaystyle \frac{|a|^2+|b|^2}{8\pi ^2}}m_N\mathrm{ln}\left({\displaystyle \frac{M_X}{m_N}}\right)\left[\mathrm{diag}(0,1,1)+𝒪\left({\displaystyle \frac{|\epsilon _{e,\mu ,\tau }|}{(|a|^2+|b|^2)^{1/2}}}\right)\right].`$
If the scale $`M_X`$ of the SO(3) symmetry imposed on $`M_S(M_X)`$ is to be naturally associated with the scale $`M_{\mathrm{GUT}}10^{16}`$ GeV relevant to GUT dynamics, it can be estimated from (2.13) that the mass splittings $`|m_{N_2}m_{N_1}|/m_N`$ and $`|m_{N_3}m_{N_1}|/m_N`$ should be larger than $`10^5`$ for $`|a|,|b|10^2`$ ($`|c|,|\epsilon _{e,\mu ,\tau }|\stackrel{<}{_{}}10^5`$). Instead, the mass difference $`|m_{N_3}m_{N_2}|/m_N`$ should be comparatively much smaller, as it is protected by an approximate U(1)<sub>l</sub> symmetry. In particular, we find that $`|m_{N_3}m_{N_2}|/m_N=𝒪(|\epsilon _{e,\mu ,\tau }a|,|\epsilon _{e,\mu ,\tau }b|)\stackrel{<}{_{}}10^7`$. At this point we should stress that in the scenarios we are considering, RG effects predominantly modify the entries $`(\mathrm{\Delta }M_S)_{1i}`$ (with $`i=1,2,3`$) in (2.3) and so they do not affect the light neutrino mass matrix (2.8). However, these effects may affect the single lepton flavour asymmetries and the flavour-dependent wash-out factors that are discussed in the next section.
In addition to RG effects, one might worry that thermal effects could significantly modify the heavy neutrino mass spectrum. However, thermal effects respect the underlying symmetries of the theory, such as global, chiral and gauge symmetries . Hence, their impact on the heavy neutrino mass spectrum is controlled by the size of the SO(3)- and U(1)<sub>l</sub>-breaking parameters in the Yukawa neutrino sector. In the hard thermal loop (HTL) approximation , thermal corrections give rise to an effective heavy neutrino mass matrix $`M_S(T)`$, which differs from the one evaluated at $`T=0`$ by an amount
$$M_S(T)M_S(0)\frac{1}{16}\mathrm{Re}\left[(h^{\nu _R})^{}h^{\nu _R}\right]\frac{T^2}{m_N}.$$
(2.14)
By comparing (2.14) with (2.13), we notice that thermal corrections have a parametric dependence very similar to the RG effects and are opposite in sign. Nevertheless, if the SO(3)-breaking scale $`M_X`$ is identified with $`M_{\mathrm{GUT}}`$, RG effects become larger than thermal effects by at least a factor 3, for the temperature regime relevant to leptogenesis $`T\stackrel{<}{_{}}m_N`$.
In Section 5 we will present numerical estimates of the BAU for electroweak-scale RL models that are motivated by the naturalness of the light and heavy neutrino sectors. As we mentioned above, this condition provides a potential link between these models and GUT-scale physics.
## 3 Boltzmann Equations for Single Lepton Flavours
In this section we derive a set of coupled BEs for the abundances of heavy Majorana neutrinos and each lepton flavour. We follow a procedure analogous to the one presented in , where a number of controllable approximations were made. In particular, we assume Maxwell-Boltzmann statistics for the heavy Majorana neutrinos. For the SM particles, we instead consider the proper Bose–Einstein and Fermi–Dirac statistics, but ignore condensate effects . The above simplifications are expected to introduce an error no larger than 20%. Furthermore, we neglect thermal effects on the collision terms, which become less significant in the temperature domain $`T\stackrel{<}{_{}}m_{N_1}`$ relevant to RL. As we will see more explicitly in Section 5, the latter approximation may partially be justified by the observation that the resulting BAU predicted in RL models is independent of the initial abundances of the heavy neutrinos and any initial baryon or lepton asymmetry.
Various definitions and notations that will be useful in deriving the BEs are introduced in Appendix A. Adopting the formalism of , the evolution of the number density, $`n_a`$, of all particle species $`a`$ can be modelled by a set of BEs. These are coupled first order differential equations and may be generically written down as<sup>2</sup><sup>2</sup>2This formalism neglects coherent time-oscillatory terms describing particle oscillations in terms of number densities, as well as off-diagonal number densities $`n_{a\overline{b}}`$, for the destruction of a particle species $`b`$ and the correlated creation of a particle species $`a`$, where $`a`$ and $`b`$ could represent the 3 lepton flavours or the 3 heavy neutrinos $`N_{1,2,3}`$. Although these effects can be modelled as well , their impact on the BAU is expected to be negligible. Specifically, coherent time-oscillatory terms between heavy Majorana neutrinos will rapidly undergo strong damping, as a consequence of the quasi-in-thermal equilibrium dynamics governing RL models. This results from the fact that the decay widths $`\mathrm{\Gamma }_{N_{1,2,3}}`$ of the heavy neutrinos are much larger than the expansion rate of the Universe. Additionally, the correlated off-diagonal number densities $`n_{a\overline{b}}`$ will be Yukawa-coupling suppressed $`𝒪((h^\nu )^2)`$ with respect to the diagonal ones $`n_{a,b}`$, if the heavy neutrinos and the charged leptons are defined in the diagonal mass basis. In particular, the contribution of $`n_{a\overline{b}}`$ to $`n_{a,b}`$ will be further suppressed $`𝒪((h^\nu )^4)`$. We will therefore neglect the effects of the coherent time-oscillatory terms and the off-diagonal number densities $`n_{a\overline{b}}`$ on the BEs.
$$\frac{dn_a}{dt}+\mathrm{\hspace{0.25em}3}Hn_a=\underset{aX^{}Y}{}\left[\frac{n_an_X^{}}{n_a^{\mathrm{eq}}n_X^{}^{\mathrm{eq}}}\gamma (aX^{}Y)\frac{n_Y}{n_Y^{\mathrm{eq}}}\gamma (YaX^{})\right],$$
(3.1)
where all possible reactions of the form $`aX^{}Y`$ or $`YaX^{}`$, in which $`a`$ can be created or annihilated are summed over. If $`a`$ is unstable, it could occur as a real intermediate state (RIS) in a resonant process like $`XaY`$. In this case, special treatment is required to avoid overcounting processes.
In principle, there is a large number of coupled BEs, one for each particle degree of freedom. This number can be drastically reduced by noting that rapidly occurring SM processes hold most of the different particle degrees of freedom and particle species in thermal equilibrium. The non-zero chemical potentials of the particle species other than heavy Majorana neutrinos and leptons produce effects of $`𝒪(1)`$ on the final baryon asymmetry . These effects will be consistently included in the BEs for the heavy Majorana neutrinos $`N_{1,2,3}`$ and the lepton doublets $`L_{e,\mu ,\tau }`$.
Although an infinite series of collision terms could be added to each BE, only a few will have a significant contribution. Following the procedure in , terms of order $`\overline{h}_\pm ^{\nu \mathrm{\hspace{0.17em}4}}h_u^2`$ and higher will be neglected, where $`\overline{h}_\pm ^\nu h^\nu `$ are the one loop resummed effective Yukawa couplings introduced in . Also neglected are terms of order $`\overline{h}_\pm ^{\nu \mathrm{\hspace{0.17em}4}}`$ for $`22`$ scatterings with two external heavy Majorana neutrinos. This leaves $`12`$ decays and inverse decays of heavy Majorana neutrinos $`𝒪(\overline{h}_\pm ^{\nu \mathrm{\hspace{0.17em}2}})`$ and $`22`$ scatterings between heavy Majorana neutrinos, lepton doublets, gauge bosons, quarks and the Higgs field, which are formally of order $`\overline{h}_\pm ^{\nu \mathrm{\hspace{0.17em}2}}g^2`$, $`\overline{h}_\pm ^{\nu \mathrm{\hspace{0.17em}2}}g^{\mathrm{\hspace{0.17em}2}}`$ and $`(\overline{h}_\pm ^\nu )^2h_u^2`$.
An important step in the following derivation is the proper subtraction of RISs. For example, the process $`L_j\mathrm{\Phi }L_k^C\mathrm{\Phi }^{}`$ will contain real intermediate heavy Majorana neutrino states. Their inverse decay and subsequent decay have already been accounted for in the BEs and must be subtracted to ensure that unitarity and CPT are respected .
In analogy to $`22`$ scatterings, $`23`$ processes, such as $`L_jQ^CL_k^C\mathrm{\Phi }^{}u^C`$, may also contain the heavy neutrinos $`N_\alpha `$ as RISs. The resonant part of such a process consists of the reaction $`L_jQ^CN_\alpha u^C`$, followed by the decay $`N_\alpha L_k^C\mathrm{\Phi }^{}`$. As before, to avoid double counting, we subtract the RISs from such a $`23`$ process. Although the off-shell $`23`$ process is a higher order effect than those we are considering, the subtracted resonant part contributes terms of order $`(\overline{h}_\pm ^\nu )^2h_u^2`$ and must be consistently included within the given approximations for the BEs. Specifically, the following relations among the collision terms are derived:
$`\gamma ^{}(L_k^C\mathrm{\Phi }^{}L_j\mathrm{\Phi })\gamma ^{}(L_k\mathrm{\Phi }L_j^C\mathrm{\Phi }^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^k\delta _{N_\alpha }^j+B_{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_l\mathrm{\Phi }}^{N_\alpha },`$
$`\gamma ^{}(L_k\mathrm{\Phi }L_j\mathrm{\Phi })\gamma ^{}(L_k^C\mathrm{\Phi }^{}L_j^C\mathrm{\Phi }^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^k\delta _{N_\alpha }^jB_{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_l\mathrm{\Phi }}^{N_\alpha },`$
$`\gamma ^{}(Qu^CL_jL_k\mathrm{\Phi })\gamma ^{}(Q^CuL_j^CL_k^C\mathrm{\Phi }^{})`$ $`=`$ $`S_{jk}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^k\delta _{N_\alpha }^j+B_{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{Qu^C}^{N_\alpha L_l},`$
$`\gamma ^{}(Qu^CL_jL_k^C\mathrm{\Phi }^{})\gamma ^{}(Q^CuL_j^CL_k\mathrm{\Phi })`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^k\delta _{N_\alpha }^jB_{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{Qu^C}^{N_\alpha L_l},`$
$`\gamma ^{}(L_jQ^Cu^C\mathrm{\Phi }^{}L_k^C)\gamma ^{}(L_j^CQu\mathrm{\Phi }L_k)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^k\delta _{N_\alpha }^j+B_{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_lQ^C}^{N_\alpha u^C},`$
$`\gamma ^{}(L_jQ^Cu^C\mathrm{\Phi }L_k)\gamma ^{}(L_j^CQu\mathrm{\Phi }^{}L_k^C)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^k\delta _{N_\alpha }^jB_{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_lQ^C}^{N_\alpha u^C},`$
where a prime denotes subtraction of RISs, the indices $`j,k=e,\mu ,\tau `$ label lepton flavour, and $`S_{jk}=(1+\delta _{jk})^1`$ is a statistical factor that corrects for the production or annihilation of identical lepton flavours. In addition, we have defined the individual lepton-flavour asymmetries and branching ratios as
$`\delta _{N_\alpha }^l`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(N_\alpha L_l\mathrm{\Phi })\mathrm{\Gamma }(N_\alpha L_l^C\mathrm{\Phi }^{})}{_k\left[\mathrm{\Gamma }(N_\alpha L_k\mathrm{\Phi })+\mathrm{\Gamma }(N_\alpha L_k^C\mathrm{\Phi }^{})\right]}},`$
$`B_{N_\alpha }^l`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(N_\alpha L_l\mathrm{\Phi })+\mathrm{\Gamma }(N_\alpha L_l^C\mathrm{\Phi }^{})}{_k\left[\mathrm{\Gamma }(N_\alpha L_k\mathrm{\Phi })+\mathrm{\Gamma }(N_\alpha L_k^C\mathrm{\Phi }^{})\right]}}.`$ (3.3)
As CP violation in these processes is predominantly caused by the resonant exchange of heavy Majorana neutrinos, the CP-violating collision terms have been approximated in terms of the CP-conserving ones as follows:
$`\delta \gamma _{L_j\mathrm{\Phi }}^{N_\alpha }`$ $`=`$ $`\delta _{N_\alpha }^j{\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_l\mathrm{\Phi }}^{N_\alpha },\delta \gamma _{L_jQ^C}^{N_\alpha u^C}=\delta _{N_\alpha }^j{\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_lQ^C}^{N_\alpha u^C},`$
$`\delta \gamma _{Qu^C}^{N_\alpha L_j}`$ $`=`$ $`\delta _{N_\alpha }^j{\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{Qu^C}^{N_\alpha L_l}\mathrm{etc}.`$ (3.4)
Unlike the $`23`$ reactions, $`32`$ processes are treated differently. Although $`32`$ processes could contain real intermediate $`N_\alpha `$ states, collision terms for their associated annihilation processes have not been included before. For example, in the process $`L_jL_k\mathrm{\Phi }Qu^C`$, a real intermediate $`N_\alpha `$ state could be coherently created from $`L`$ and $`\mathrm{\Phi }`$ states. This coherent RIS would then interact with another $`L`$ state producing $`Q`$ and $`u^C`$. Previously, the process $`N_\alpha LQu^C`$ has only been considered for heavy $`N_\alpha `$ neutrinos in a thermally incoherent state. This implies that $`32`$ processes containing $`N_\alpha `$ as RISs have not yet been accounted for and should not be subtracted. With the help of CPT and unitarity, one may therefore obtain the following relations for the $`32`$ processes:
$`\gamma (L_jL_k\mathrm{\Phi }Qu^C)\gamma (L_j^CL_k^C\mathrm{\Phi }^{}Q^Cu)`$ $`=`$ $`𝒪(h^{\nu \mathrm{\hspace{0.17em}4}}h_u^2),`$
$`\gamma (L_jL_k^C\mathrm{\Phi }^{}Qu^C)\gamma (L_j^CL_k\mathrm{\Phi }Q^Cu)`$ $`=`$ $`𝒪(h^{\nu \mathrm{\hspace{0.17em}4}}h_u^2),`$
$`\gamma (L_j\mathrm{\Phi }uL_k^CQ)\gamma (L_j^C\mathrm{\Phi }^{}u^CL_kQ^C)`$ $`=`$ $`𝒪(h^{\nu \mathrm{\hspace{0.17em}4}}h_u^2),`$
$`\gamma (L_j\mathrm{\Phi }u^CL_kQ^C)\gamma (L_j^C\mathrm{\Phi }^{}uL_k^CQ)`$ $`=`$ $`𝒪(h^{\nu \mathrm{\hspace{0.17em}4}}h_u^2).`$ (3.5)
As a consequence of this, $`32`$ processes will contribute extra CP-conserving $`22`$ collision terms, through the resonant exchange of real intermediate $`N_\alpha `$ states. Applying the narrow width approximation, we find
$`\gamma (L_jL_k\mathrm{\Phi }Qu^C)+\gamma (L_j^CL_k^C\mathrm{\Phi }^{}Q^Cu)`$ $`=`$ $`S_{jk}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^jB_{N_\alpha }^k+\delta _{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{Qu^C}^{N_\alpha L_l},`$
$`\gamma (L_jL_k^C\mathrm{\Phi }^{}Qu^C)+\gamma (L_j^CL_k\mathrm{\Phi }Q^Cu)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^jB_{N_\alpha }^k\delta _{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{Qu^C}^{N_\alpha L_l},`$
$`\gamma (L_j\mathrm{\Phi }uL_k^CQ)+\gamma (L_j^C\mathrm{\Phi }^{}u^CL_kQ^C)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^jB_{N_\alpha }^k+\delta _{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_lQ^C}^{N_\alpha u^C},`$
$`\gamma (L_j\mathrm{\Phi }u^CL_kQ^C)+\gamma (L_j^C\mathrm{\Phi }^{}uL_k^CQ)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\left(B_{N_\alpha }^jB_{N_\alpha }^k\delta _{N_\alpha }^j\delta _{N_\alpha }^k\right){\displaystyle \underset{l=e,\mu ,\tau }{}}\gamma _{L_lQ^C}^{N_\alpha u^C}.`$
We may now employ (3.1) and write down the BEs in terms of the number densities of heavy Majorana neutrinos $`n_{N_\alpha }`$ and the lepton-doublet asymmetries $`n_{\mathrm{\Delta }L_{e,\mu ,\tau }}`$,
$`{\displaystyle \frac{dn_{N_\alpha }}{dt}}+\mathrm{\hspace{0.25em}3}Hn_{N_\alpha }`$ $`=`$ $`(1{\displaystyle \frac{n_{N_\alpha }}{n_{N_\alpha }^{\mathrm{eq}}}}){\displaystyle \underset{k=e,\mu ,\tau }{}}(\gamma _{L_k\mathrm{\Phi }}^{N_\alpha }+\gamma _{Qu^C}^{N_\alpha L_k}+\gamma _{L_kQ^C}^{N_\alpha u^C}+\gamma _{L_ku}^{N_\alpha Q}`$ (3.7)
$`+\gamma _{L_k\mathrm{\Phi }}^{N_\alpha V_\mu }+\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_k}+\gamma _{L_kV_\mu }^{N_\alpha \mathrm{\Phi }^{}})`$
$`{\displaystyle \underset{k=e,\mu ,\tau }{}}{\displaystyle \frac{n_{\mathrm{\Delta }L_k}}{2n_{l_k}^{\mathrm{eq}}}}[\delta \gamma _{L_k\mathrm{\Phi }}^{N_\alpha }+\delta \gamma _{L_kQ^C}^{N_\alpha u^C}+\delta \gamma _{L_ku}^{N_\alpha Q}+\delta \gamma _{L_k\mathrm{\Phi }}^{N_\alpha V_\mu }+\delta \gamma _{L_kV_\mu }^{N_\alpha \mathrm{\Phi }^{}}`$
$`+{\displaystyle \frac{n_{N_\alpha }}{n_{N_\alpha }^{\mathrm{eq}}}}(\delta \gamma _{Qu^C}^{N_\alpha L_k}+\delta \gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_k})],`$
$`{\displaystyle \frac{dn_{\mathrm{\Delta }L_j}}{dt}}+\mathrm{\hspace{0.25em}3}Hn_{\mathrm{\Delta }L_j}`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{3}{}}}({\displaystyle \frac{n_{N_\alpha }}{n_{N_\alpha }^{\mathrm{eq}}}}\mathrm{\hspace{0.25em}1})(\delta \gamma _{L_j\mathrm{\Phi }}^{N_\alpha }\delta \gamma _{Qu^C}^{N_\alpha L_j}+\delta \gamma _{L_jQ^C}^{N_\alpha u^C}+\delta \gamma _{L_ju}^{N_\alpha Q}`$ (3.8)
$`+\delta \gamma _{L_j\mathrm{\Phi }}^{N_\alpha V_\mu }\delta \gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_j}+\delta \gamma _{L_jV_\mu }^{N_\alpha \mathrm{\Phi }^{}})`$
$`{\displaystyle \frac{n_{\mathrm{\Delta }L_j}}{2n_{l_j}^{\mathrm{eq}}}}[{\displaystyle \underset{\alpha =1}{\overset{3}{}}}(\gamma _{L_j\mathrm{\Phi }}^{N_\alpha }+\mathrm{\hspace{0.25em}2}\gamma _{L_jQ^C}^{N_\alpha u^C}+\mathrm{\hspace{0.25em}2}\gamma _{L_ju}^{N_\alpha Q}+\mathrm{\hspace{0.25em}2}\gamma _{L_j\mathrm{\Phi }}^{N_\alpha V_\mu }+\mathrm{\hspace{0.25em}2}\gamma _{L_jV_\mu }^{N_\alpha \mathrm{\Phi }^{}}`$
$`+\mathrm{\hspace{0.25em}2}\gamma _{Qu^C}^{N_\alpha L_j}+\mathrm{\hspace{0.25em}2}\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_j}+{\displaystyle \frac{n_{N_\alpha }}{n_{N_\alpha }^{\mathrm{eq}}}}(\gamma _{Qu^C}^{N_\alpha L_j}+\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_j}))`$
$`+{\displaystyle \underset{k=e,\mu ,\tau }{}}(\gamma _{L_k^C\mathrm{\Phi }^{}}^{L_j\mathrm{\Phi }}+\gamma _{\mathrm{\Phi }^{}\mathrm{\Phi }^{}}^{L_jL_k}+\gamma _{L_k\mathrm{\Phi }}^{L_j\mathrm{\Phi }}+\gamma _{\mathrm{\Phi }\mathrm{\Phi }^{}}^{L_jL_k^C})]`$
$`{\displaystyle \underset{k=e,\mu ,\tau }{}}{\displaystyle \frac{n_{\mathrm{\Delta }L_k}}{2n_{l_k}^{\mathrm{eq}}}}[\gamma _{L_j^C\mathrm{\Phi }^{}}^{L_k\mathrm{\Phi }}+\gamma _{\mathrm{\Phi }^{}\mathrm{\Phi }^{}}^{L_kL_j}\gamma _{L_j\mathrm{\Phi }}^{L_k\mathrm{\Phi }}\gamma _{\mathrm{\Phi }\mathrm{\Phi }^{}}^{L_kL_j^C}`$
$`+{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\delta _{N_\alpha }^j\delta _{N_\alpha }^k{\displaystyle \underset{l=e,\mu ,\tau }{}}(\gamma _{L_lQ^C}^{N_\alpha u^C}+\gamma _{L_lu}^{N_\alpha Q}+\gamma _{L_l\mathrm{\Phi }}^{N_\alpha V_\mu }+\gamma _{L_lV_\mu }^{N_\alpha \mathrm{\Phi }^{}}`$
$`+\mathrm{\hspace{0.25em}2}\gamma _{Qu^C}^{N_\alpha L_l}+\mathrm{\hspace{0.25em}2}\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_l})].`$
In the above set of BEs, we have only kept terms to leading order in $`n_{\mathrm{\Delta }L_j}/n_{l_j}^{\mathrm{eq}}`$, and implemented the relations given in (3)–(3).
All SM species in the thermal bath, including the lepton doublets $`L_{e,\mu ,\tau }`$, possess non-zero chemical potentials. These chemical potentials can be expressed in terms of the lepton-doublet chemical potentials only, under the assumption that SM processes are in full thermal equilibrium . This analysis yields the following relations:
$`\mu _V`$ $`=`$ $`0,\mu _\mathrm{\Phi }={\displaystyle \frac{4}{21}}{\displaystyle \underset{l=e,\mu ,\tau }{}}\mu _{L_l},\mu _Q={\displaystyle \frac{1}{9}}{\displaystyle \underset{l=e,\mu ,\tau }{}}\mu _{L_l},\mu _u={\displaystyle \frac{5}{63}}{\displaystyle \underset{l=e,\mu ,\tau }{}}\mu _{L_l},`$
$`\mu _{e_l}`$ $`=`$ $`\mu _{L_l}{\displaystyle \frac{4}{21}}{\displaystyle \underset{l=e,\mu ,\tau }{}}\mu _{L_l},`$ (3.9)
where $`\mu _x`$ denotes the chemical potential of a particle species $`x`$. The relations (3) can be used to implement the effects of the SM chemical potentials in the BEs. They result in corrections to the so-called wash-out terms in both the lepton and heavy neutrino BEs. At this point we should also note that the BEs in their present form are most accurate above $`T_c`$. As $`T`$ approaches $`T_c`$, the assumption that the sphaleron processes are in thermal equilibrium becomes less valid. This will result in $`𝒪(v/T)`$ corrections to the relations in (3). The inclusion of the bulk of these corrections will be considered in the next section.
To numerically solve the BEs, it proves convenient to introduce a number of new variables. In the radiation dominated epoch of the Universe relevant to baryogenesis, the cosmic time $`t`$ is related to the temperature $`T`$ through
$$t=\frac{z^2}{2H(z=1)},$$
(3.10)
where
$$z=\frac{m_{N_1}}{T},H(z)17.2\times \frac{m_{N_1}^2}{z^2M_{\mathrm{Planck}}},$$
(3.11)
with $`M_{\mathrm{Planck}}=1.2\times 10^{19}`$ GeV. Also, we normalize the number density of a particle species, $`n_a`$, to the number density of photons, $`n_\gamma `$, thereby defining the new parameter $`\eta _a`$,
$$\eta _a(z)=\frac{n_a(z)}{n_\gamma (z)},$$
(3.12)
with
$$n_\gamma (z)=\frac{2T^3}{\pi ^2}=\frac{2m_{N_1}^3}{\pi ^2}\frac{1}{z^3}.$$
(3.13)
To allow the BEs to be written in a slightly more compact form, we will use the summation conventions
$$\gamma _{LY}^{N_\alpha X}=\underset{l=e,\mu ,\tau }{}\gamma _{L_lY}^{N_\alpha X},\eta _{\mathrm{\Delta }L}=\underset{l=e,\mu ,\tau }{}\eta _{\mathrm{\Delta }L_l},$$
(3.14)
where $`X`$ and $`Y`$ stand for any particle state other than $`L_l`$ and $`N_\alpha `$.
Using (3.7)–(3.14) and incorporating corrections due to the SM chemical potentials, the BEs for heavy Majorana neutrinos and lepton doublets are written down
$`{\displaystyle \frac{d\eta _{N_\alpha }}{dz}}`$ $`=`$ $`{\displaystyle \frac{z}{H(z=1)}}[(1{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}){\displaystyle \underset{k=e,\mu ,\tau }{}}(\mathrm{\Gamma }^{D(\alpha k)}+\mathrm{\Gamma }_{\mathrm{Yukawa}}^{S(\alpha k)}+\mathrm{\Gamma }_{\mathrm{Gauge}}^{S(\alpha k)})`$ (3.15)
$`{\displaystyle \frac{2}{3}}{\displaystyle \underset{k=e,\mu ,\tau }{}}\eta _{\mathrm{\Delta }L_k}\delta _{N_\alpha }^k(\widehat{\mathrm{\Gamma }}^{D(\alpha k)}+\widehat{\mathrm{\Gamma }}_{\mathrm{Yukawa}}^{S(\alpha k)}+\widehat{\mathrm{\Gamma }}_{\mathrm{Gauge}}^{S(\alpha k)})],`$
$`{\displaystyle \frac{d\eta _{\mathrm{\Delta }L_j}}{dz}}`$ $`=`$ $`{\displaystyle \frac{z}{H(z=1)}}\{{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\delta _{N_\alpha }^j({\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}\mathrm{\hspace{0.25em}1}){\displaystyle \underset{k=e,\mu ,\tau }{}}(\mathrm{\Gamma }^{D(\alpha k)}+\mathrm{\Gamma }_{\mathrm{Yukawa}}^{S(\alpha k)}+\mathrm{\Gamma }_{\mathrm{Gauge}}^{S(\alpha k)})`$ (3.16)
$`{\displaystyle \frac{2}{3}}\eta _{\mathrm{\Delta }L_j}[{\displaystyle \underset{\alpha =1}{\overset{3}{}}}B_{N_\alpha }^j(\stackrel{~}{\mathrm{\Gamma }}^{D(\alpha j)}+\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{Yukawa}}^{S(\alpha j)}+\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{Gauge}}^{S(\alpha j)}+\mathrm{\Gamma }_{\mathrm{Yukawa}}^{W(\alpha j)}+\mathrm{\Gamma }_{\mathrm{Gauge}}^{W(\alpha j)})`$
$`+{\displaystyle \underset{k=e,\mu ,\tau }{}}(\mathrm{\Gamma }_{\mathrm{Yukawa}}^{\mathrm{\Delta }L=2(jk)}+\mathrm{\Gamma }_{\mathrm{Yukawa}}^{\mathrm{\Delta }L=0(jk)})]`$
$`{\displaystyle \frac{2}{3}}{\displaystyle \underset{k=e,\mu ,\tau }{}}\eta _{\mathrm{\Delta }L_k}[{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\delta _{N_\alpha }^j\delta _{N_\alpha }^k(\mathrm{\Gamma }_{\mathrm{Yukawa}}^{W(\alpha k)}+\mathrm{\Gamma }_{\mathrm{Gauge}}^{W(\alpha k)})`$
$`+\mathrm{\Gamma }_{\mathrm{Yukawa}}^{\mathrm{\Delta }L=2(kj)}\mathrm{\Gamma }_{\mathrm{Yukawa}}^{\mathrm{\Delta }L=0(kj)}]\},`$
where
$`\mathrm{\Gamma }^{D(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}\gamma _{L_l\mathrm{\Phi }}^{N_\alpha },`$ (3.17)
$`\widehat{\mathrm{\Gamma }}^{D(\alpha l)}`$ $`=`$ $`\stackrel{~}{\mathrm{\Gamma }}^{D(\alpha l)}={\displaystyle \frac{1}{n_\gamma }}\left(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}}\right)\gamma _{L\mathrm{\Phi }}^{N_\alpha },`$ (3.18)
$`\mathrm{\Gamma }_{\mathrm{Yukawa}}^{S(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}\left(\gamma _{Qu^C}^{N_\alpha L_l}+\gamma _{L_lQ^C}^{N_\alpha u^C}+\gamma _{L_lu}^{N_\alpha Q}\right),`$ (3.19)
$`\widehat{\mathrm{\Gamma }}_{\mathrm{Yukawa}}^{S(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}[({\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{Qu^C}^{N_\alpha L}+(1+{\displaystyle \frac{1}{9}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}}{\displaystyle \frac{5}{63}}{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{LQ^C}^{N_\alpha u^C}`$ (3.20)
$`+(1+{\displaystyle \frac{5}{63}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}}{\displaystyle \frac{1}{9}}{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{Lu}^{N_\alpha Q}],`$
$`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{Yukawa}}^{S(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}[({\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{Qu^C}^{N_\alpha L}+(1+{\displaystyle \frac{1}{9}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}}+{\displaystyle \frac{5}{63}}{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{LQ^C}^{N_\alpha u^C}`$ (3.21)
$`+(1+{\displaystyle \frac{5}{63}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}}+{\displaystyle \frac{1}{9}}{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{Lu}^{N_\alpha Q}],`$
$`\mathrm{\Gamma }_{\mathrm{Gauge}}^{S(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}\left(\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L_l}+\gamma _{L_l\mathrm{\Phi }}^{N_\alpha V_\mu }+\gamma _{L_lV_\mu }^{N_\alpha \mathrm{\Phi }^{}}\right),`$ (3.22)
$`\widehat{\mathrm{\Gamma }}_{\mathrm{Gauge}}^{S(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}[({\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L}+(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{L\mathrm{\Phi }}^{N_\alpha V_\mu }`$ (3.23)
$`+(1{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{LV_\mu }^{N_\alpha \mathrm{\Phi }^{}}],`$
$`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{Gauge}}^{S(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}[({\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L}+(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{L\mathrm{\Phi }}^{N_\alpha V_\mu }`$ (3.24)
$`+(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{N_\alpha }}{\eta _{N_\alpha }^{\mathrm{eq}}}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{LV_\mu }^{N_\alpha \mathrm{\Phi }^{}}],`$
$`\mathrm{\Gamma }_{\mathrm{Yukawa}}^{W(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}[(2+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{Qu^C}^{N_\alpha L}+(1+{\displaystyle \frac{17}{63}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{LQ^C}^{N_\alpha u^C}`$ (3.25)
$`+(1+{\displaystyle \frac{19}{63}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{Lu}^{N_\alpha Q}],`$
$`\mathrm{\Gamma }_{\mathrm{Gauge}}^{W(\alpha l)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}[(2+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{\mathrm{\Phi }^{}V_\mu }^{N_\alpha L}+(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{L\mathrm{\Phi }}^{N_\alpha V_\mu }`$ (3.26)
$`+(1+{\displaystyle \frac{8}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_l}}})\gamma _{LV_\mu }^{N_\alpha \mathrm{\Phi }^{}}],`$
$`\mathrm{\Gamma }_{\mathrm{Yukawa}}^{\mathrm{\Delta }L=2(jk)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}\left[\left(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_j}}}\right)\left(\gamma _{L_k^C\mathrm{\Phi }^{}}^{L_j\mathrm{\Phi }}+\gamma _{\mathrm{\Phi }^{}\mathrm{\Phi }^{}}^{L_jL_k}\right)\right],`$ (3.27)
$`\mathrm{\Gamma }_{\mathrm{Yukawa}}^{\mathrm{\Delta }L=0(jk)}`$ $`=`$ $`{\displaystyle \frac{1}{n_\gamma }}\left[\left(1+{\displaystyle \frac{4}{21}}{\displaystyle \frac{\eta _{\mathrm{\Delta }L}}{\eta _{\mathrm{\Delta }L_j}}}\right)\gamma _{L_k\mathrm{\Phi }}^{L_j\mathrm{\Phi }}+\gamma _{L_k\mathrm{\Phi }^{}}^{L_j\mathrm{\Phi }^{}}+\gamma _{\mathrm{\Phi }\mathrm{\Phi }^{}}^{L_jL_k^C}\right].`$ (3.28)
Notice that the would-be singularities in the limit of a vanishing $`\eta _{\mathrm{\Delta }L_l}`$ in (3.18)–(3.21) and (3.23)–(3.28) are exactly cancelled by corresponding factors $`\eta _{\mathrm{\Delta }L_l}`$ that multiply the collision terms in the BEs (3.15) and (3.16). We should also note that the flavour-diagonal $`\mathrm{\Delta }L=0`$ processes, with $`k=j`$, do not contribute to the BEs, as it can be explicitly checked in (3.16). Finally, it is worth commenting on the earlier form of the BEs, obtained in . This can be recovered from (3.15)–(3.28), after summing over the three lepton-doublet BEs, with the assumption that $`n_{\mathrm{\Delta }L_i}=\frac{1}{3}n_{\mathrm{\Delta }L}`$, and after neglecting SM chemical potential corrections.
## 4 Out of Equilibrium Sphaleron Effects
In the SM, the combination of the baryon and lepton numbers, $`B+L`$, is anomalous . Although at low energies this $`B+L`$ violation is unobservably small, at temperatures close to and above the electroweak phase transition, e.g. for $`T\stackrel{>}{_{}}150`$ GeV, thermal fluctuations more efficiently overcome the so-called sphaleron barrier allowing rapid $`B+L`$ violation in the SM .
The temperature dependence of the rate of $`B+L`$ violation is of particular interest in models of low-scale leptogenesis. Any lepton asymmetry produced after the $`(B+L)`$-violating interactions drop out of thermal equilibrium will not be converted into a baryon asymmetry. Therefore, in electroweak-scale leptogenesis scenarios, it is important to consider the rate of $`B+L`$ violation in the BEs, in order to offer a more reliable estimate of the final baryon asymmetry.
The rate of $`(B+L)`$-violating transitions has been studied in detail in for temperatures satisfying the double inequality
$$M_W(T)T\frac{M_W(T)}{\alpha _w},$$
(4.1)
where $`\alpha _w=g^2/4\pi `$ is the SU(2)<sub>L</sub> fine structure constant, $`M_W(T)=gv(T)/2`$ is the $`T`$-dependent $`W`$-boson mass and
$$v(T)=v(0)\left(\mathrm{\hspace{0.17em}1}\frac{T^2}{T_c^2}\right)^{\frac{1}{2}}$$
(4.2)
is the $`T`$-dependent VEV of the Higgs field. The critical temperature of the electroweak phase transition, $`T_c`$, is given by
$$T_c=v(0)\left(\frac{1}{2}+\frac{3g^2}{16\lambda }+\frac{g^{\mathrm{\hspace{0.17em}2}}}{16\lambda }+\frac{h_t}{4\lambda }\right)^{\frac{1}{2}},$$
(4.3)
where $`\lambda `$ is the quartic Higgs self-coupling, $`g^{}`$ is the U(1)<sub>Y</sub> gauge coupling and $`h_t`$ is the top-quark Yukawa coupling.
The rate of $`B+L`$ violation per unit volume is
$$\gamma _{\mathrm{\Delta }(B+L)}\frac{\mathrm{\Gamma }}{V}=\frac{\omega _{}}{2\pi }𝒩_{\mathrm{tr}}(𝒩V)_{\mathrm{rot}}\left(\frac{\alpha _WT}{4\pi }\right)^3\alpha _3^6e^{E_{\mathrm{sp}}/T}\kappa .$$
(4.4)
According to (4.1), this expression is valid for temperatures $`T\stackrel{<}{_{}}T_c`$. The various quantities in (4.4) are related to the sphaleron dynamics and are discussed in . Following the notation of , the parameters $`\omega _{}`$, $`𝒩_{\mathrm{tr}}`$ and $`𝒩_{\mathrm{rot}}`$ are functions of $`\lambda /g^2`$, $`V_{\mathrm{rot}}=8\pi ^2`$ and $`\alpha _3=g_3^2/4\pi `$, where
$$g_3^2=\frac{g^2T}{2M_W(T)}.$$
(4.5)
$`E_{\mathrm{sp}}`$ is the energy of the sphaleron and is given by
$$E_{\mathrm{sp}}=A\frac{2M_W(T)}{\alpha _W},$$
(4.6)
where $`A`$ is a function of $`\lambda /g^2`$ and is of order 1 for all phenomenologically relevant values of $`\lambda /g^2`$. The dependence of the parameter $`\kappa `$ on $`\lambda /g^2`$ has been calculated in , using various techniques. The results of those studies are summarized in Table 1, where the values of $`\kappa `$ and the other sphaleron-related parameters are exhibited for $`\lambda /g^2=0.278`$, which corresponds to a SM Higgs-boson mass $`M_H`$ of 120 GeV.
Given the present experimental limits on the SM Higgs-boson mass, $`M_H\stackrel{>}{_{}}115`$ GeV, it can be shown that the electroweak phase transition in the SM will either be a weakly first order one, or even a second or higher order one, without bubble nucleation and the formation of large spatial inhomogeneities in particle densities. Therefore, the use of a formalism describing the $`(B+L)`$-violating sphaleron dynamics in terms of spatially independent $`B`$\- and $`L`$-number densities $`n_B`$ and $`n_{L_j}`$ may be justified. Further refinements to this approach will be presented elsewhere.
We should bear in mind that heavy Majorana neutrino decays, sphaleron effects and other processes considered in the BEs (3.15) and (3.16) act directly on the number densities of SU(2)<sub>L</sub> lepton doublets, $`n_{\mathrm{\Delta }L_{e,\mu ,\tau }}`$. However, the quantity usually referred to as lepton number, $`L`$, has a contribution from the right-handed charged leptons $`l_{e,\mu ,\tau R}`$ as well. In thermal equilibrium, one may relate the asymmetry in right-handed charged leptons to the asymmetry in SU(2)<sub>L</sub> lepton doublets by virtue of (3), leading to the result
$$\eta _{\mathrm{\Delta }l_{iR}}=\frac{1}{2}\eta _{\mathrm{\Delta }L_i}\frac{2}{21}\eta _{\mathrm{\Delta }L}.$$
(4.7)
The change in lepton flavour can be thought of as having two components, one component termed leptogenesis due to lepton-number-violating processes considered in Section 3, and another due to the $`(B+L)`$-violating sphalerons:
$$\frac{d\eta _{L_i}}{dz}=\frac{d\eta _{L_i}}{dz}|_{\mathrm{leptogenesis}}+\frac{d\eta _{L_i}}{dz}|_{\mathrm{sphaleron}},$$
(4.8)
where, up to $`𝒪(v/T)`$ corrections,
$$\frac{d\eta _{L_i}}{dz}|_{\mathrm{leptogenesis}}=\frac{3}{2}\frac{d\eta _{\mathrm{\Delta }L_i}}{dz}\frac{2}{21}\frac{d\eta _{\mathrm{\Delta }L}}{dz}.$$
(4.9)
The BEs determining the leptogenesis component of (4.8) have been discussed in Section 3. We shall now discuss the BEs determining the sphaleron component of (4.8), and the generation of a net $`B`$-number asymmetry.
Within the context of the above formalism, the sphaleron components of the BEs for $`n_B`$ and $`n_{L_j}`$ are given by
$`{\displaystyle \frac{dn_B}{dt}}+\mathrm{\hspace{0.25em}3}Hn_B`$ $`=`$ $`n_G(e^{\beta (\mu _BQ_B+\mu _LQ_L)}e^{\beta (\mu _BQ_B^{}+\mu _LQ_L^{})}e^{\beta (\mu _BQ_{\overline{B}}+\mu _LQ_{\overline{L}})}`$
$`+e^{\beta (\mu _BQ_{\overline{B}}^{}+\mu _LQ_{\overline{L}}^{})})\gamma _{\mathrm{\Delta }(B+L)},`$
$`{\displaystyle \frac{dn_{L_j}}{dt}}+\mathrm{\hspace{0.25em}3}Hn_{L_j}`$ $`=`$ $`(e^{\beta (\mu _BQ_B+\mu _LQ_L)}e^{\beta (\mu _BQ_B^{}+\mu _LQ_L^{})}`$ (4.10)
$`e^{\beta (\mu _BQ_{\overline{B}}+\mu _LQ_{\overline{L}})}+e^{\beta (\mu _BQ_{\overline{B}}^{}+\mu _LQ_{\overline{L}}^{})})\gamma _{\mathrm{\Delta }(B+L)},`$
where $`n_G`$ is the number of generations and $`\beta =1/T`$. Furthermore, $`Q_{B(L)}`$ is the baryonic (or leptonic) charge of the system before the $`(B+L)`$-violating sphaleron transition and $`Q_{B(L)}^{}`$ is the charge after the transition. Klinkhamer and Manton showed that a sphaleron carries a baryon (and lepton) number of $`n_G/2`$, therefore $`Q_BQ_B^{}=n_G/2`$ and $`Q_{L_i}Q_{L_i}^{}=1/2`$. Finally, assuming that the baryon and lepton chemical potentials are small with respect to the temperature, the BEs (4) may be approximated by
$`{\displaystyle \frac{dn_B}{dt}}+\mathrm{\hspace{0.25em}3}Hn_B`$ $`=`$ $`n_G\beta \left(n_G\mu _B+{\displaystyle \underset{i}{}}\mu _{L_i}\right)\gamma _{\mathrm{\Delta }(B+L)},`$
$`{\displaystyle \frac{dn_{L_j}}{dt}}+\mathrm{\hspace{0.25em}3}Hn_{L_j}`$ $`=`$ $`\beta \left(n_G\mu _B+{\displaystyle \underset{i}{}}\mu _{L_i}\right)\gamma _{\mathrm{\Delta }(B+L)}.`$ (4.11)
Notice that the BEs (4) are linear in the chemical potentials, which is a very useful approximation for our numerical estimates.
We now need to determine the relation between the baryon and lepton chemical potentials and their respective number densities. These relations can be found by considering the effective potential, $`V`$, of the Higgs and the SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge fields. They have been computed in at finite temperatures, for small chemical potentials, $`\mu _B,\mu _LT`$ and when $`v(T)\stackrel{<}{_{}}(\text{a few})\times T`$. In this framework, the neutrality of the system with respect to gauge charges can be accounted for by minimizing the potential with respect to the temporal components of the SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge fields, $`W_0^a`$ ($`a=1,2,3`$) and $`B_0`$, respectively. The baryon and lepton number densities are then given by
$$n_B=\frac{V}{\mu _B},n_{L_i}=\frac{V}{\mu _{L_i}}.$$
(4.12)
For the SM with 3 generations and 1 Higgs doublet, we obtain
$`\mu _B`$ $`=`$ $`3n_B{\displaystyle \frac{77T^2+27v^2(T)}{132T^4+51T^2v^2(T)}}2{\displaystyle \frac{22T^2+3v^2(T)}{132T^4+51T^2v^2(T)}}{\displaystyle \underset{j=e,\mu ,\tau }{}}n_{L_j},`$
$`\mu _{L_i}`$ $`=`$ $`{\displaystyle \frac{2}{51T^2}}\left(51n_{L_i}\mathrm{\hspace{0.25em}3}n_B+\mathrm{\hspace{0.25em}4}{\displaystyle \underset{j=e,\mu ,\tau }{}}n_{L_j}\right)`$ (4.13)
$`{\displaystyle \frac{484}{153\left(44T^2+17v^2(T)\right)}}\left(\mathrm{\hspace{0.17em}3}n_B\mathrm{\hspace{0.25em}4}{\displaystyle \underset{j=e,\mu ,\tau }{}}n_{L_j}\right).`$
Employing the relations (4), we may now extend the system of BEs (3.15) and (3.16), by explicitly taking the $`(B+L)`$-violating sphaleron transitions into account,
$`{\displaystyle \frac{d\eta _B}{dz}}`$ $`=`$ $`{\displaystyle \frac{z}{H(z=1)}}`$
$`\times \left[\eta _B+{\displaystyle \frac{28}{51}}{\displaystyle \underset{j=e,\mu ,\tau }{}}\eta _{L_j}+{\displaystyle \frac{225}{561}}{\displaystyle \frac{v^2(T)}{T^2}}\left(\eta _B+{\displaystyle \frac{108}{225}}{\displaystyle \underset{j=e,\mu ,\tau }{}}\eta _{L_j}\right)\right]\mathrm{\Gamma }_{\mathrm{\Delta }(B+L)},`$
$`{\displaystyle \frac{d\eta _{L_i}}{dz}}`$ $`=`$ $`{\displaystyle \frac{d\eta _{L_i}}{dz}}|_{\mathrm{leptogenesis}}+{\displaystyle \frac{1}{3}}{\displaystyle \frac{d\eta _B}{dz}},`$ (4.15)
with
$$\mathrm{\Gamma }_{\mathrm{\Delta }(B+L)}=\frac{(3366/\pi ^2)T^2}{132T^2+\mathrm{\hspace{0.25em}51}v^2(T)}\frac{\gamma _{\mathrm{\Delta }(B+L)}}{n_\gamma }.$$
(4.16)
The leptogenesis component of (4.15) may be determined using relation (4.9), along with the BEs (3.15) and (3.16).
Observe that in the limit of infinite sphaleron transition rate, $`\mathrm{\Gamma }_{\mathrm{\Delta }(B+L)}/H(z=1)\mathrm{}`$, and at high temperatures $`Tv(T)`$, the conversion of lepton-to-baryon number densities is given by the known relation:
$$\eta _B=\frac{28}{51}\underset{j=e,\mu ,\tau }{}\eta _{L_j}.$$
(4.17)
To account for the $`T`$-dependent $`(B+L)`$-violating sphaleron effects, our numerical estimates given in the next section will be based on the BEs (3.15), (3.16), (4.9), (4) and (4.15).
## 5 Numerical Examples
We shall now analyze R$`\tau `$L models that comply with the constraints obtained from the existing low-energy neutrino data and provide successful baryogenesis. As was discussed in Section 2, our specific choice of model parameters will be motivated by the naturalness of the light and heavy neutrino sectors.
Phenomenologically relevant R$`\tau `$L models can be constructed for an SO(3) invariant heavy neutrino mass of the size of the electroweak scale, e.g. $`m_N=250`$ GeV \[cf. (2.3)\], if $`|a||b||c|`$ and $`|a|,|b|10^2`$. To protect the $`\tau `$-lepton number from wash-out effects, we also require that the small U(1)<sub>l</sub>-breaking parameters $`|\epsilon _{e,\mu ,\tau }|`$ be no larger than about $`10^6`$ and $`|c|\stackrel{<}{_{}}10^5`$. For definiteness, the model parameters determining the light neutrino mass spectrum are chosen to be (in arbitrary complex units)
$$\frac{\mathrm{\Delta }m_N}{m_N}a^2=4,\epsilon _e=2+\frac{21}{250},\epsilon _\mu =\frac{13}{50},\epsilon _\tau =\frac{49}{128},$$
(5.1)
where the ratio $`b/a`$ is kept fixed:
$$\frac{b}{a}=\frac{19}{50}.$$
(5.2)
The actual values selected for the relevant parameters $`a`$, $`(\mathrm{\Delta }M_S)_{22}`$, $`(\mathrm{\Delta }M_S)_{33}`$ and $`(\mathrm{\Delta }M_S)_{23}`$ vary with the SO(3) invariant mass $`m_N`$. As we will see in more detail below, Table 2 illustrates choices of these parameters consistent with the light neutrino data. For $`m_N`$ in the range 100–1000 GeV, the chosen parameters are consistent with the naturalness condition mentioned in Section 2, whilst giving rise to phenomenologically rich models.
In our numerical analysis, we will focus on 4 examples, with $`m_N=100`$, 250, 500, and 1000 GeV. Clearly, the model parameters selected in (5.1), (5.2) and Table 2 imply that all the scenarios have the same tree-level light neutrino mass matrix:
$$𝐦^\nu (\begin{array}{ccc}1.27& 3.63& 2.96\\ 3.63& 1.89& 0.370\\ 2.96& 0.370& 0.544\end{array})\times 10^2\mathrm{eV}.$$
(5.3)
This leads to an inverted hierarchy of light neutrino masses, with mass differences and mixings compatible with the current $`3\sigma `$ bounds . Adopting the convention $`m_{\nu _3}<m_{\nu _1}<m_{\nu _2}`$, we find the mass squared differences and mixing angles
$`m_{\nu _2}^2m_{\nu _1}^2`$ $`=`$ $`7.54\times 10^5\mathrm{eV}^2,m_{\nu _1}^2m_{\nu _3}^2=2.45\times 10^3\mathrm{eV}^2,`$
$`\mathrm{sin}^2\theta _{12}`$ $`=`$ $`0.362,\mathrm{sin}^2\theta _{23}=0.341,\mathrm{sin}^2\theta _{13}=0.047.`$ (5.4)
Since the mass matrix (2.8) is rank 2, one light neutrino will be massless at the tree level ($`m_{\nu _3}=0`$), thus fixing the absolute scale of the light neutrino hierarchy.
The remaining soft SO(3)-breaking parameters, $`(\mathrm{\Delta }M_S)_{11}`$, $`(\mathrm{\Delta }M_S)_{12}`$, $`(\mathrm{\Delta }M_S)_{13}`$, do not affect the light neutrino mass spectrum. These together with the parameter $`c`$ play a key role in obtaining the correct BAU and are exhibited in Table 3. We choose $`(\mathrm{\Delta }M_S)_{11}`$ to be relatively large, $`(\mathrm{\Delta }M_S)_{11}10^5m_N`$, providing large mass differences $`|m_{N_2}m_{N_1}|/m_N`$ and $`|m_{N_3}m_{N_1}|/m_N10^5`$. Such a choice is consistent with thermal and RG effects running from the GUT scale $`10^{16}`$ GeV to the electroweak scale $`m_N`$ (see also the discussion in Section 2). The other two soft SO(3)-breaking parameters, $`(\mathrm{\Delta }M_S)_{12}`$ and $`(\mathrm{\Delta }M_S)_{13}`$, are selected so as to give the observed BAU.
To assess the degree of cancellation between tree-level and RG contributions to $`\mathrm{\Delta }M_S`$, we introduce the parameter $`r_C`$ defined as
$$r_C\underset{(i,j)}{}\frac{|(\mathrm{\Delta }M_S^{\mathrm{RG}})_{ij}|}{|(\mathrm{\Delta }M_S)_{ij}|}.$$
(5.5)
In (5.5), the product $`(i,j)`$ is taken over contributions where $`|(\mathrm{\Delta }M_S^{\mathrm{RG}})_{ij}|>|(\mathrm{\Delta }M_S)_{ij}|`$. The parameter $`r_C`$ is always greater than 1 and represents that the degree of cancellation is 1 part in $`r_C`$. From the values of $`r_C`$ displayed in Table 3, we observe that electroweak-scale heavy Majorana neutrinos are favoured by naturalness.
The baryon asymmetry predicted for each model can be determined by solving the BEs (3.15), (3.16), (4.9), (4) and (4.15), and using the collision terms derived in Appendix A and . These collision terms are calculated in the basis where the charged-lepton and heavy-Majorana mass matrices are positive and diagonal. They have been appropriately expressed in terms of the one-loop resummed effective Yukawa couplings derived in . It is worth noting that all SM reactions, including those involving the $`e`$-Yukawa coupling, are in full thermal equilibrium for the temperatures relevant to our scenarios, $`T\stackrel{<}{_{}}10`$ TeV . Moreover, since heavy Majorana neutrino decays are thermally blocked at temperatures $`T\stackrel{>}{_{}}3m_{N_\alpha }`$ , we will only display numerical estimates of the evolution of lepton and baryon asymmetries, for $`z=m_{N_1}/T\stackrel{>}{_{}}0.1`$. Nevertheless, as we will see below, the predictions for the final BAU are relatively robust in RL models, because of the near or complete independence on the primordial baryon and lepton number abundances.
Some of the Yukawa and gauge-mediated collision terms contain IR divergences, which are usually regulated in thermal field theory by considering the thermal masses of the exchanged particles . To assess the theoretical errors introduced by the choice of a universal thermal mass regulator (see the discussion in Appendix A), we have estimated the response of the final baryon asymmetry under variations of the IR mass regulator $`m_{\mathrm{IR}}`$. We find that the predicted BAU only varies by $`\pm 7\%`$, for a variation of $`m_{\mathrm{IR}}`$ by $`\pm 25\%`$.
The BEs are solved numerically, using the Fortran code LeptoGen<sup>2</sup><sup>2</sup>2LeptoGen may be obtained from http://hep.man.ac.uk/u/thomasu/leptogen. Fig. 2 shows the predicted evolution of the baryon and individual lepton asymmetries, $`\eta _B`$ and $`\eta _{L_l}`$, as functions of the $`T`$-related parameter $`z=m_{N_1}/T`$, for each of the 4 examples, with $`m_N=`$ 100, 250, 500 and 1000 GeV. The specific model parameters are given in (5.1), (5.2), and Tables 2 and 3. Each scenario had an initially thermal heavy Majorana neutrino abundance and zero initial baryon and lepton asymmetries, i.e. $`\eta _{N_\alpha }^{\mathrm{in}}=1`$ and $`\eta _B^{\mathrm{in}}=\eta _{L_l}^{\mathrm{in}}=0`$. The 4 panels show that the large $`L_\tau `$ asymmetry is slightly reduced by less significant, but opposite sign $`L_e`$ and $`L_\mu `$ asymmetries. Clearly visible in each scenario is the effect of the rapidly decreasing rate of $`B+L`$ violation; the lepton and baryon asymmetries quickly decouple at $`TT_c`$. This decoupling is particularly pronounced in the $`m_N=100`$ GeV scenario where the baryon asymmetry freezes out exactly when the lepton asymmetry is maximal. In particular, the rapid decoupling of $`\eta _B`$ from $`\eta _{L_l}`$ at temperatures $`T`$ close to $`T_c`$ has the virtue that, unlike $`\eta _{L_l}`$, $`\eta _B`$ remains almost unaffected from ordinary SM mass effects due to a non-zero VEV $`v(T)`$ \[cf. (4.2)\], since it is $`v(TT_c)v(T=0)`$.
Fig. 3 shows the evolution of the baryon asymmetry for varying initial lepton, baryon and heavy neutrino abundances. For the 250 GeV scenario, Fig. 3(a) illustrates the near independence of the resultant baryon asymmetry on the initial conditions. Even for the most extreme initial conditions $`\eta _{L_l}^{\mathrm{in}}=\mathrm{\hspace{0.17em}0.1}`$ and $`\eta _B^{\mathrm{in}}=\pm \mathrm{\hspace{0.17em}0.1}`$, the variation in the final baryon asymmetry is only $`\pm \mathrm{\hspace{0.17em}38}\%`$.
For heavy neutrino masses $`m_N\stackrel{<}{_{}}250`$ GeV, the dependence on initial conditions becomes stronger. In the $`m_N=100`$ GeV scenario, Fig. 3(b) shows the dependence of the final BAU on the initial lepton and baryon asymmetries in a R$`\tau `$L scenario with $`m_N=100`$ GeV. It is interesting to observe that the final $`B`$ asymmetry will remain almost unaffected, even if the primordial baryon asymmetry $`\eta _B^{\mathrm{in}}`$ at $`T\stackrel{>}{_{}}10m_N`$ is as large as $`10^6`$, namely two orders of magnitude larger than the one required to agree with observational data.
In R$`\tau `$L scenarios with $`m_N>250`$ GeV, the final baryon asymmetry is completely independent of the initial conditions. This is illustrated in Fig. 4 for the R$`\tau `$L scenario with $`m_N=`$ 500 GeV. In this numerical example, it is most striking to notice that the prediction for the final BAU remains unchanged, even if the initial conditions are set at temperatures below the heavy neutrino mass scale $`m_N`$, e.g. at $`T0.5m_N`$.
Some insight into the independence on initial conditions is provided by Fig. 5. The ratios of various collision terms to the Hubble parameter are plotted for the $`m_N=`$ 250 GeV scenario. These ratios show that RL can take place almost completely in thermal equilibrium; in certain cases, the reaction rates are many orders of magnitude above the Hubble parameter $`H(z=1)`$. In spite of this fact, RL (R$`\tau `$L) can successfully generate the required excess in $`L`$ ($`L_\tau `$), because of the resonantly enhanced CP asymmetry.
To allow for a simple quantitative understanding of the baryon asymmetry in R$`\tau `$L (and similar) scenarios, we need to introduce the individual lepton flavour $`K`$-factors
$$K_{N_\alpha }^l=\frac{\mathrm{\Gamma }(N_\alpha L_l\mathrm{\Phi })+\mathrm{\Gamma }(N_\alpha L_l^C\mathrm{\Phi }^{})}{H(z=1)}.$$
(5.6)
Note that the decay widths are calculated in terms of the one-loop resummed effective Yukawa couplings .
Table 4 shows the various components of $`K_{N_\alpha }^l`$ for the $`m_N=`$ 250 GeV scenario. This explicitly demonstrates how the texture provided by (2.5) and (2.7) allows for a heavy Majorana neutrino to decay relatively out of equilibrium, whilst simultaneously protecting the $`\tau `$-lepton number from being washed-out, even though large $`e`$\- and $`\mu `$-Yukawa couplings to $`N_{1,2}`$ exist. Bear in mind that we use the convention $`m_{N_1}<m_{N_2}<m_{N_3}`$ upon diagonalization of the heavy Majorana neutrino mass matrix $`M_S`$. As can be seen from Table 4, $`K`$-factors $`K_{N_3}^{e,\mu ,\tau }10`$–100 and a CP-asymmetry $`\delta _{N_3}^\tau 10^6`$ are sufficient to generate a large $`\tau `$-lepton asymmetry. Although the $`K`$-factors $`K_{N_{1,2}}^{e,\mu }`$ associated with $`N_{1,2}`$ and the $`e`$ and $`\mu `$ leptons are enormous of order $`10^9`$$`10^{10}`$, these turn out to be harmless to the $`\tau `$-lepton asymmetry, as the latter is protected by the low $`\tau `$-lepton $`K`$-factors $`K_{N_{1,2,3}}^\tau 10`$.
An order of magnitude estimate of the final baryon asymmetry, including single lepton flavour effects, may be obtained using
$$\eta _B\mathrm{\hspace{0.17em}10}^2\times \underset{l=1}{\overset{3}{}}\underset{N_\alpha }{}e^{(m_{N_\alpha }m_{N_1})/m_{N_1}}\delta _{N_\alpha }^l\frac{K_{N_\alpha }^l}{K_lK_{N_\alpha }}.$$
(5.7)
The above estimate for $`\eta _B`$ is also consistent with the one stated earlier in . In (5.7), the $`K`$-factors are summed in the following way:
$`K_{N_\alpha }`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{3}{}}}K_{N_\alpha }^l,K_l={\displaystyle \underset{N_\alpha }{}}e^{(m_{N_\alpha }m_{N_1})/m_{N_1}}K_{N_\alpha }^l.`$ (5.8)
Notice that all $`K`$-factors are evaluated at $`T=m_{N_1}`$ (i.e. $`z=m_{N_1}/T=1`$), where $`m_{N_1}`$ is the lightest of the heavy Majorana neutrinos. The intuitive estimate (5.7) is applicable for all leptogenesis scenarios satisfying the approximate inequality
$$K_{lN_\alpha }\stackrel{>}{_{}}1,$$
(5.9)
for each of the lepton flavours $`l`$ and the heavy Majorana neutrinos $`N_{1,2,3}`$. The inequality (5.9) ensures that the energy scale $`m_{N_1}`$ can be identified as the true scale of leptogenesis.
In RL scenarios, such as R$`\tau `$L, the importance of taking individual lepton flavour effects into account can be demonstrated by comparing (5.7) with the naive estimate, in which lepton flavour effects are treated indiscriminately in a universal manner,
$$\eta _B^{\mathrm{univ}.}\mathrm{\hspace{0.17em}10}^2\times \underset{N_\alpha }{}e^{(m_{N_\alpha }m_{N_1})/m_{N_1}}\frac{\delta _{N_\alpha }}{K},$$
(5.10)
where $`K=_{l=e,\mu \tau }K_l`$. In the R$`\tau `$L scenario with $`m_N=250`$ GeV, the dominant contribution to this estimate will come from $`N_3`$, with a total CP asymmetry $`\delta _{N_3}10^3`$. Taking the ratio of the two estimates yields
$$\frac{\eta _B^{\mathrm{univ}.}}{\eta _B}\frac{\delta _{N_3}}{\delta _{N_3}^\tau }\frac{K_{N_3}K_\tau }{K_{N_3}^\tau K}\frac{\delta _{N_3}}{\delta _{N_3}^\tau }\frac{|c|^2}{|a|^2+|b|^2}\mathrm{\hspace{0.17em}10}^6.$$
(5.11)
Thus, without considering single lepton flavour effects in this particular R$`\tau `$L model, one obtains an erroneous prediction for the BAU, which is suppressed by 6 orders of magnitude and has the wrong sign. These estimates are confirmed by solving the total lepton number BEs presented in .
In a hierarchical scenario, the number densities of the heavier neutrinos $`N_{2,3}`$ at $`T=m_{N_1}`$ will be Boltzmann suppressed. To account for this phenomenon, we have included the Boltzmann factors $`e^{(m_{N_\alpha }m_{N_1})/m_{N_1}}`$ in the estimates (5.7), (5.10) and in the definition of $`K_l`$. Clearly, in RL models with each heavy neutrino nearly degenerate in mass, this last factor can be set to 1.
Flavour effects can also play a significant role in mildly hierarchical scenarios. Figure 6 shows the predicted evolution of the lepton asymmetry in a scenario where $`m_{N_3}=3m_{N_1}`$, $`m_{N_2}=2m_{N_1}`$ and $`m_{N_1}=10^{10}`$ GeV. The Yukawa texture was chosen to be consistent with light neutrino data and a normal hierarchical light neutrino spectrum. In this example, neglecting individual lepton flavour effects introduces an $`𝒪(10)`$ suppression of the final lepton and baryon asymmetry.
In fully hierarchical scenarios satisfying (5.9), it can be seen that the estimates (5.7) and (5.10) are completely equivalent. A large hierarchy in the heavy neutrino spectrum, combined with the condition (5.9), implies that the final lepton asymmetry is determined entirely by the decay of the lightest heavy Majorana neutrino $`N_1`$. This fact makes it impossible for a single lepton flavour to be protected from wash-out, whilst the neutrino decays out of equilibrium.
Likewise, in flavour universal scenarios, where $`\eta _{L_{e,\mu ,\tau }}=\frac{1}{3}\eta _L`$, the estimates (5.7) and (5.10) are completely equivalent for both nearly degenerate and hierarchical leptogenesis scenarios.
Our numerical analysis presented in this section has explicitly demonstrated that models of R$`\tau `$L can provide a viable explanation for the observed BAU, in accordance with the current light neutrino data. In the next section, we will see how the scenarios considered here have far reaching phenomenological implications for low-energy observables of lepton flavour/number violation and collider experiments.
## 6 Phenomenological Implications
RL models, and especially R$`\tau `$L models, can give rise to a number of phenomenologically testable signatures. In particular, we will analyze the generic predictions of R$`\tau `$L models for the $`0\nu \beta \beta `$ decay, and for the LFV processes: $`\mu e\gamma `$, $`\mu eee`$ and $`\mu e`$ conversion in nuclei. Finally, we will present simple and realistic numerical estimates of production cross sections of heavy Majorana neutrinos at future $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ colliders, and apply these results to the R$`\tau `$L models.
### 6.1 $`\mathrm{𝟎}𝝂𝜷𝜷`$ Decay
Neutrinoless double beta decay ($`0\nu \beta \beta `$) corresponds to a process in which two single $`\beta `$ decays occur simultaneously in one nucleus. As a consequence of this, a nucleus $`(Z,A)`$ gets converted into a nucleus $`(Z+2,A)`$, i.e.
$${}_{Z}{}^{A}X_{Z+2}^AX+\mathrm{\hspace{0.25em}2}e^{}.$$
Evidently, this process violates $`L`$-number by two units and can naturally take place in minimal RL models, in which the observed light neutrinos are Majorana particles. The observation of such a process would provide further information on the structure of the light neutrino mass matrix $`𝐦^\nu `$.
To a very good approximation, the half life for a $`0\nu \beta \beta `$ decay mediated by light Majorana neutrinos is given by
$$[T_{1/2}^{0\nu \beta \beta }]^1=\frac{|m|^2}{m_e^2}|_{0\nu \beta \beta }|^2G_{01},$$
(6.1)
where $`m`$ denotes the effective Majorana neutrino mass, $`m_e`$ is the electron mass and $`_{0\nu \beta \beta }`$ and $`G_{01}`$ denote the appropriate nuclear matrix element and the phase space factor, respectively. More details regarding the calculation of $`T_{1/2}^{0\nu \beta \beta }`$ may be found in .
In models of interest to us, the effective neutrino mass can be related to the entry $`\{11\}(\{ee\})`$ of the light neutrino mass matrix $`𝐦^\nu `$ in (2.8), i.e.
$$|m|=|(𝐦^\nu )_{ee}|=\frac{v^2}{2m_N}\left|\frac{\mathrm{\Delta }m_N}{m_N}a^2\epsilon _e^2\right|.$$
(6.2)
As has been explicitly demonstrated in the previous section, R$`\tau `$L models realize a light neutrino mass spectrum with an inverted hierarchy , thus giving rise to a sizeable effective neutrino mass. The prediction for $`|m|`$ in these models is
$$|m|=|(𝐦^\nu )_{ee}|0.013\mathrm{eV}.$$
(6.3)
Such a prediction lies at the very low end of the value of the effective Majorana neutrino mass, reported recently by the Heidelberg–Moscow collaboration . There are proposals for future $`0\nu \beta \beta `$-decay experiments that will be sensitive to values of $`|m|`$ of order $`10^2`$ , significantly increasing the constraints on this parameter.
### 6.2 $`𝝁\mathbf{}𝒆𝜸`$
As shown in Fig. 7, heavy Majorana neutrinos may induce LFV couplings to the photon ($`\gamma `$) and the $`Z`$ boson via loop effects. These couplings give rise to LFV decays, such as $`\mu e\gamma `$ and $`\mu eee`$ . Our discussion and notation closely follows the extensive studies . Related phenomenological analyses of LFV effects in the SM with singlet neutrinos may be found in .
To properly describe LFV in low-energy observables, we first introduce the so-called Langacker–London (LL) parameters :
$$(s_L^{\nu _l})^2=1\underset{l^{}=e,\mu ,\tau }{}|B_{l\nu _l^{}}|^2\left(m_D^{}M_S^1M_S^1m_D^T\right)_{ll},$$
(6.4)
where $`m_D=\frac{1}{\sqrt{2}}h^{\nu _R}v`$ and $`B_{l\nu _l^{}}`$ are mixing-matrix factors close to 1 that multiply the SM tree-level $`Wl\nu _l^{}`$ vertices . The LL parameters $`(s_L^{\nu _l})^2`$ quantify the deviation of the actual squared $`Wl\nu _l^{}`$ couplings (summed over all light neutrinos) from the corresponding sum of squared couplings in the SM. The parameters $`(s_L^{\nu _{e,\mu ,\tau }})^2`$ are constrained by LEP and low-energy electroweak observables . Independent constraints on these parameters typically give: $`(s_L^{\nu _{e,\mu ,\tau }})^2\stackrel{<}{_{}}10^2`$. As we will see in a moment, however, LFV observables impose much more severe constraints on products of the LL parameters, and especially on $`s_L^{\nu _e}s_L^{\nu _\mu }`$.
As a first LFV observable, we consider the decay $`\mu e\gamma `$. The branching fraction for the decay $`\mu e\gamma `$ is given by
$$B(\mu e\gamma )=\frac{\alpha _w^3s_w^2}{256\pi ^2}\frac{m_\mu ^4}{M_W^4}\frac{m_\mu }{\mathrm{\Gamma }_\mu }|G_\gamma ^{\mu e}|^2\frac{\alpha _w^3s_w^2}{1024\pi ^2}\frac{m_\mu ^4}{M_W^4}\frac{m_\mu }{\mathrm{\Gamma }_\mu }(s_L^{\nu _\mu })^2(s_L^{\nu _e})^2,$$
(6.5)
where $`\mathrm{\Gamma }_\mu =2.997\times 10^{19}`$ GeV is the experimentally measured muon decay width, and $`G_\gamma ^{\mu e}`$ is a composite form-factor defined in . In arriving at the last equality in (6.5), we have assumed that $`m_N^2M_W^2`$, for a model with two nearly degenerate heavy Majorana neutrinos. In this case, one finds that $`G_\gamma ^{\mu e}\frac{e^{i\varphi }}{2}s_L^{\nu _\mu }s_L^{\nu _e}`$, where $`\varphi `$ is an unobservable model-dependent phase. Confronting the theoretical prediction (6.5) with the experimental upper limit
$$B_{\mathrm{exp}}(\mu e\gamma )<1.2\times 10^{11},$$
(6.6)
we obtain the following constraint:
$$s_L^{\nu _e}s_L^{\nu _\mu }<1.2\times 10^4.$$
(6.7)
This last constraint is stronger by one to two orders of magnitude with respect to those derived on $`(s_L^{\nu _e})^2`$ and $`(s_L^{\nu _\mu })^2`$ individually.
In R$`\tau `$L models, only two of the right-handed neutrinos, $`\nu _{2R}`$ and $`\nu _{3R}`$, which have appreciable $`e`$\- and $`\mu `$-Yukawa couplings, will be relevant to LFV effects. In this case, the LL parameters $`(s_L^{\nu _e})^2`$ and $`(s_L^{\nu _\mu })^2`$ are, to a very good approximation, given by
$$(s_L^{\nu _e})^2=\frac{|a|^2v^2}{m_N^2},(s_L^{\nu _\mu })^2=\frac{|b|^2v^2}{m_N^2}.$$
(6.8)
Then, the following theoretical prediction is obtained:
$$B(\mu e\gamma )=910^4\times \frac{|a|^2|b|^2v^4}{m_N^4}.$$
(6.9)
For the particular scenarios considered in Section 5, we find $`B(\mu e\gamma )10^{12}`$. These values are well within reach of the MEG collaboration, which will be sensitive to $`B(\mu e\gamma )10^{14}`$ .
### 6.3 $`𝝁\mathbf{}𝒆𝒆𝒆`$
As illustrated in Fig. 8, quantum effects mediated by heavy Majorana neutrinos may also give rise to the 3-body LFV decay mode $`\mu ^{}e^{}e^+e^{}`$. The branching ratio for this LFV decay may conveniently be expressed as
$`B(\mu eee)`$ $`=`$ $`{\displaystyle \frac{\alpha _w^4}{24576\pi ^3}}{\displaystyle \frac{m_\mu ^4}{M_W^4}}{\displaystyle \frac{m_\mu }{\mathrm{\Gamma }_\mu }}\{\mathrm{\hspace{0.17em}2}|\frac{1}{2}F_{\mathrm{box}}^{\mu eee}+F_Z^{\mu e}2s_w^2(F_Z^{\mu e}F_\gamma ^{\mu e})|^2`$ (6.10)
$`+\mathrm{\hspace{0.17em}4}s_w^4|F_Z^{\mu e}F_\gamma ^{\mu e}|^2+\mathrm{\hspace{0.17em}16}s_w^2\mathrm{Re}\left[(F_Z^{\mu e}+\frac{1}{2}F_{\mathrm{box}}^{\mu eee})G_\gamma ^{\mu e}\right]`$
$`\mathrm{\hspace{0.17em}48}s_w^4\mathrm{Re}\left[(F_Z^{\mu e}F_\gamma ^{\mu e})G_\gamma ^{\mu e}\right]+\mathrm{\hspace{0.17em}32}s_w^4|G_\gamma ^{\mu e}|^2(\mathrm{ln}{\displaystyle \frac{m_\mu ^2}{m_e^2}}{\displaystyle \frac{11}{4}})\}.`$
The expressions $`F_\gamma ^{\mu e}`$, $`F_Z^{\mu e}`$ and $`F_{\mathrm{box}}^{\mu eee}`$ are composite form-factors, defined and computed in . In the limit $`m_N^2M_W^2`$ and up to an overall physically irrelevant phase factor $`e^{i\varphi }`$, these composite form-factors simplify to
$`F_\gamma ^{\mu e}`$ $``$ $`{\displaystyle \frac{7}{12}}s_L^{\nu _\mu }s_L^{\nu _e}{\displaystyle \frac{1}{6}}s_L^{\nu _\mu }s_L^{\nu _e}\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right),`$ (6.11)
$`F_Z^{\mu e}`$ $``$ $`\left[{\displaystyle \frac{5}{2}}{\displaystyle \frac{3}{2}}\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right)\right]s_L^{\nu _\mu }s_L^{\nu _e}{\displaystyle \frac{1}{2}}s_L^{\nu _\mu }s_L^{\nu _e}{\displaystyle \underset{k=e,\mu ,\tau }{}}(s_L^{\nu _k})^2{\displaystyle \frac{m_N^2}{M_W^2}},`$ (6.12)
$`F_{\mathrm{box}}^{\mu eee}`$ $``$ $`2s_L^{\nu _\mu }s_L^{\nu _e}+{\displaystyle \frac{1}{2}}s_L^{\nu _\mu }s_L^{\nu _e}(s_L^{\nu _e})^2{\displaystyle \frac{m_N^2}{M_W^2}}.`$ (6.13)
Correspondingly, the analytic result (6.10) in the same limit may be cast into the form:
$`B(\mu eee)`$ $``$ $`{\displaystyle \frac{\alpha _w^4}{294912\pi ^3}}{\displaystyle \frac{m_\mu ^4}{M_W^4}}{\displaystyle \frac{m_\mu }{\mathrm{\Gamma }_\mu }}(s_L^{\nu _\mu })^2(s_L^{\nu _e})^2\{\mathrm{\hspace{0.17em}54}300s_w^2+217s_w^4+96s_w^4\mathrm{ln}\left({\displaystyle \frac{m_\mu ^2}{m_e^2}}\right)`$ (6.14)
$`\left(108492s_w^2+800s_w^4\right)\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right)+\left(54192s_w^2+256s_w^4\right)\mathrm{ln}^2\left({\displaystyle \frac{m_N^2}{M_W^2}}\right)`$
$`+{\displaystyle \frac{m_N^2}{M_W^2}}[(1850s_w^2(1832s_w^2)\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right))(s_L^{\nu _e})^2`$
$`(36172s_w^2+300s_w^4(36136s_w^2+192s_w^4)\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right)){\displaystyle \underset{l=e,\mu ,\tau }{}}(s_L^{\nu _l})^2]`$
$`+{\displaystyle \frac{m_N^4}{M_W^4}}[{\displaystyle \frac{3}{2}}(s_L^{\nu _e})^46(12s_w^2)(s_L^{\nu _e})^2{\displaystyle \underset{l=e,\mu ,\tau }{}}(s_L^{\nu _l})^2`$
$`+\mathrm{\hspace{0.25em}6}(14s_w^2+6s_w^4)\left({\displaystyle \underset{l=e,\mu ,\tau }{}}(s_L^{\nu _l})^2\right)^2]\}.`$
It can be seen from (6.14) that the so-called non-decoupling terms proportional to $`m_N^4/M_W^4`$ are always multiplied with higher powers of the LL parameters. In general, these terms do not decouple and become very significant , for large heavy neutrino masses $`m_N`$ and fixed values of $`s_L^{\nu _l}`$, which amounts to scenarios with large neutrino Yukawa couplings $`|h_{ij}^{\nu _R}|\stackrel{>}{_{}}0.5`$ . However, these non-decoupling terms are negligible, as long as $`s_L^{\nu _l}m_N/M_W1`$. This is actually the case for the R$`\tau `$L models discussed in Section 5. Neglecting terms proportional to $`m_N^2/M_W^2`$ and $`m_N^4/M_W^4`$, we may relate $`B(\mu eee)`$ to $`B(\mu e\gamma )`$ through:
$$B(\mu eee)8.210^3\times \left[\mathrm{\hspace{0.17em}1}0.8\mathrm{ln}\left(\frac{m_N^2}{M_W^2}\right)+0.5\mathrm{ln}^2\left(\frac{m_N^2}{M_W^2}\right)\right]B(\mu e\gamma ).$$
(6.15)
For example, for an R$`\tau `$L model with $`m_N=250`$ GeV, (6.15) implies
$$B(\mu eee)1.410^2\times B(\mu e\gamma )1.410^{14}$$
(6.16)
This value is a factor $`70`$ below the present experimental bound : $`B_{\mathrm{exp}}(\mu eee)<1.0\times 10^{12}`$. In this respect, it would be very encouraging, if higher sensitivity experiments could be designed to probe this observable.
### 6.4 Coherent $`𝝁\mathbf{}𝒆`$ Conversion in Nuclei
One of the most sensitive experiments to LFV is the coherent conversion of $`\mu e`$ in nuclei, e.g. $`\mu ^{}{}_{22}{}^{48}\mathrm{Ti}e^{}{}_{22}{}^{48}\mathrm{Ti}`$ . The Feynman graphs responsible for such a process are displayed in Fig. 8.
Our calculation of $`\mu e`$ conversion in nuclei closely follows . We consider the kinematic approximations: $`q^2m_\mu ^2`$ and $`p_e^0|\stackrel{}{p}_e|m_\mu `$, which are valid for $`\mu e`$ conversion. Given the above approximation, the $`\mu e`$ conversion rate in a nucleus with nucleon numbers $`(N,Z)`$, is given by
$$B_{\mu e}(N,Z)\frac{\mathrm{\Gamma }[\mu (N,Z)e(N,Z)]}{\mathrm{\Gamma }[\mu (N,Z)\mathrm{capture}]}\frac{\alpha _{\mathrm{em}}^3\alpha _w^4m_\mu ^5}{32\pi ^2M_W^4\mathrm{\Gamma }_{\mathrm{capt}.}}\frac{Z_{\mathrm{eff}}^4}{Z}|F(m_\mu ^2)|^2|Q_W|^2,$$
(6.17)
where $`\alpha _{\mathrm{em}}=1/137`$ is the electromagnetic fine structure constant, $`Z_{\mathrm{eff}}`$ is the effective atomic number of coherence and $`\mathrm{\Gamma }_{\mathrm{capt}.}`$ is the muon nuclear capture rate. For $`{}_{22}{}^{48}\mathrm{Ti}`$, experimental measurements give $`Z_{\mathrm{eff}}17.6`$ for $`{}_{22}{}^{48}\mathrm{Ti}`$ and $`\mathrm{\Gamma }[\mu {}_{22}{}^{48}\mathrm{Ti}\mathrm{capture}]1.705\times 10^{18}`$ GeV . Moreover, $`|F(m_\mu ^2)|0.54`$ is the nuclear form factor . Finally, $`Q_W=V_u(2Z+N)+V_d(Z+2N)`$ is the coherent charge of the nucleus, which is associated with the vector current. Its explicit form is given by
$`V_u`$ $`=`$ $`{\displaystyle \frac{2}{3}}s_w^2\left(F_\gamma ^{\mu e}G_\gamma ^{\mu e}F_Z^{\mu e}\right)+{\displaystyle \frac{1}{4}}\left(F_Z^{\mu e}F_{\mathrm{box}}^{\mu euu}\right),`$ (6.18)
$`V_d`$ $`=`$ $`{\displaystyle \frac{1}{3}}s_w^2\left(F_\gamma ^{\mu e}G_\gamma ^{\mu e}F_Z^{\mu e}\right){\displaystyle \frac{1}{4}}\left(F_Z^{\mu e}+F_{\mathrm{box}}^{\mu edd}\right).`$ (6.19)
The composite form-factors $`F_{\mathrm{box}}^{\mu euu}`$ and $`F_{\mathrm{box}}^{\mu edd}`$ are defined in . In the SM with two nearly degenerate heavy Majorana neutrinos and in the limit $`m_N^2/M_W^21`$, these form-factors can be written down in the simplified forms:
$$F_{\mathrm{box}}^{\mu euu}F_{\mathrm{box}}^{\mu edd}s_L^{\nu _\mu }s_L^{\nu _e}.$$
(6.20)
In the same limit $`m_N^2/M_W^21`$, $`B_{\mu e}(N,Z)`$ is given by
$`B_{\mu e}(N,Z)`$ $``$ $`{\displaystyle \frac{\alpha _{\mathrm{em}}^3\alpha _w^4m_\mu ^5}{18432\pi ^2M_W^4\mathrm{\Gamma }_{\mathrm{capt}.}}}{\displaystyle \frac{Z_{\mathrm{eff}}^4}{Z}}|F(m_\mu ^2)|^2(s_L^{\nu _\mu })^2(s_L^{\nu _e})^2`$ (6.21)
$`\times \{[\mathrm{\hspace{0.17em}3}N+(3386s_w^2)Z+(9N(932s_w^2)Z)\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right)]^2`$
$`+6{\displaystyle \frac{m_N^2}{M_W^2}}\left[\mathrm{\hspace{0.17em}3}N+(3386s_w^2)Z+\left(9N(932s_w^2)Z\right)\mathrm{ln}\left({\displaystyle \frac{m_N^2}{M_W^2}}\right)\right]`$
$`\times \left(N(14s_w^2)Z\right){\displaystyle \underset{l=e,\mu ,\tau }{}}(s_L^{\nu _l})^2`$
$`+9{\displaystyle \frac{m_N^4}{M_W^4}}(N(14s_w^2)Z)^2\left({\displaystyle \underset{l=e,\mu ,\tau }{}}(s_L^{\nu _l})^2\right)^2\}.`$
For the $`{}_{22}{}^{48}\mathrm{Ti}`$ case, $`B_{\mu e}(26,22)`$ is related to $`B(\mu e\gamma )`$ through
$$B_{\mu e}(26,22)0.1\times \left[\mathrm{\hspace{0.17em}1}+0.5\mathrm{ln}\left(\frac{m_N^2}{M_W^2}\right)\right]^2B(\mu e\gamma ).$$
(6.22)
On the experimental side, the strongest upper bound on $`B_{\mu e}(N,Z)`$ is obtained from experimental data on $`\mu e`$ conversion in $`{}_{22}{}^{48}\mathrm{Ti}`$ :
$$B_{\mu e}^{\mathrm{exp}}(26,22)<4.3\times 10^{12},$$
(6.23)
at the 90% CL. However, the proposed experiment by the MECO collaboration will be sensitive to conversion rates of order $`10^{16}`$.
In the R$`\tau `$L model with $`m_N=250`$ GeV, one obtains, on the basis of (6.22), the prediction for $`\mu e`$ conversion in $`{}_{22}{}^{48}\mathrm{Ti}`$:
$$B_{\mu e}(26,22)0.46\times B(\mu e\gamma )4.5\times 10^{13}.$$
(6.24)
The above prediction falls well within reach of the sensitivity proposed by the MECO collaboration.
In Table 5, we summarize our results for the branching ratios of the 3 LFV processes: $`\mu e\gamma `$, $`\mu eee`$ and coherent $`\mu e`$ conversion in $`{}_{22}{}^{48}\mathrm{Ti}`$ nuclei, for each R$`\tau `$L model considered in Section 5.
As a final general remark, we should mention that R$`\tau `$L models, and leptogenesis models in general, do not suffer from too large contributions to the electron electric dipole moment (EDM) , which first arises at two loops. The reason is that EDM effects are suppressed either by higher powers of small Yukawa couplings of order $`10^4`$ and less, or by very small factors, such as $`(m_{N_1}m_{N_{2,3}})/m_N10^7`$. The latter is the case in R$`\tau `$L models, which leads to unobservably small EDM effects of order $`10^{37}e`$ cm, namely 10 orders of magnitude smaller than the present experimental limits .
### 6.5 Collider Heavy Majorana Neutrino Production
If heavy Majorana neutrinos have electroweak-scale masses and appreciable couplings to electrons and muons they can be copiously produced at future $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ colliders. As shown in Fig. 9, this is exactly the kinematic situation for the heavy Majorana neutrinos $`N_{2,3}`$ described by the R$`\tau `$L models. The heavy Majorana neutrino $`N_1`$ has a very small coupling to leptons and it would be very difficult to produce this state directly.
For collider c.m.s. energies $`\sqrt{s}m_N`$, the $`t`$-channel $`W^{}`$-boson exchange graphs will dominate over the $`Z`$-boson exchange graph, which is $`s`$-channel propagator suppressed (see Fig. 9). In this high-energy limit, the production cross section for heavy Majorana neutrinos approaches a constant , i.e.
$$\sigma \left[e^+e^{}(\mu ^+\mu ^{})N_{2,3}\nu \right]\frac{\pi \alpha _w^2}{4M_W^2}(s_L^{\nu _{e(\mu )}})^210\mathrm{fb}\times \left(\frac{s_L^{\nu _{e(\mu )}}}{10^2}\right)^2.$$
(6.25)
Since $`s_L^{\nu _\tau }0`$ in R$`\tau `$L models, the produced heavy Majorana neutrinos $`N_{2,3}`$ will have the characteristic signature that they will predominantly decay into electrons and muons, but not into $`\tau `$ leptons. Assuming that $`m_N\stackrel{>}{_{}}M_H`$, the branching fraction of $`N_{2,3}`$ decays into charged leptons and into $`W^\pm `$ bosons decaying hadronically is
$$B\left(N_{2,3}e^{},\mu ^{}W^\pm (\mathrm{jets})\right)\frac{1}{2}\times \frac{2}{3}=\frac{1}{3}.$$
(6.26)
Given (6.25), (6.26) and an integrated luminosity of 100 fb<sup>-1</sup>, we expect to be able to analyze about 100 signal events for $`(s_L^{\nu _{e(\mu )}})^2=10^2`$ and $`m_N\stackrel{<}{_{}}300`$ GeV, at future $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ colliders with c.m.s. energy $`\sqrt{s}=0.5`$$`1\mathrm{TeV}`$.
These simple estimates are supported by a recent analysis, where competitive background reactions to the signal have been considered . This analysis showed that the inclusion of background processes reduces the number of signal events by a factor of 10. The authors in find that an $`e^+e^{}`$ linear collider with c.m.s. energy $`\sqrt{s}=0.5`$ TeV will be sensitive to values of $`s_L^{\nu _e}=|a|v/m_N0.7\times 10^2`$. This amounts to the same level of sensitivity to the parameter $`|a|`$, for R$`\tau `$L scenarios with $`m_N=250`$ GeV. The sensitivity to $`s_L^{\nu _e}`$ could be improved by a factor of 3, i.e. $`s_L^{\nu _e}0.2\times 10^2`$, in proposed upgraded $`e^+e^{}`$ accelerators such as CLIC.
A similar analysis should be envisaged to hold for future $`\mu ^+\mu ^{}`$ colliders, leading to similar findings for $`s_L^{\nu _\mu }=|b|v/m_N`$. In general, we expect that the ratio of the two production cross sections of $`N_{2,3}`$ at the two colliders under identical conditions of c.m.s. energy and luminosity will give a direct measure of the ratio of $`|a|^2/|b|^2`$. This information, together with that obtained from low-energy LFV observables, $`0\nu \beta \beta `$-decay experiments, and neutrino data, will significantly constrain the parameters of the R$`\tau `$L models. Finally, since the heavy Majorana neutrinos $`N_{2,3}`$ play an important synergetic role in resonantly enhancing $`\delta _{N_1}^\tau `$, potentially large CP asymmetries in their decays will determine the theoretical parameters of these models further. Evidently, more detailed studies are needed before one could reach a definite conclusion concerning the exciting possibility that electroweak-scale R$`\tau `$L models may naturally constitute a laboratory testable solution to the cosmological problem of the BAU.
## 7 Conclusions
We have studied a novel variant of RL, which may take place at the electroweak phase transition. This RL variant gives rise to a number of phenomenologically testable signatures for low-energy experiments and future high-energy colliders. The new RL scenario under study makes use of the property that, in addition to $`BL`$ number, sphalerons preserve the individual quantum numbers $`\frac{1}{3}BL_{e,\mu ,\tau }`$ . The observed BAU can be produced by lepton-to-baryon conversion of an individual lepton number. For the case of the $`\tau `$-lepton number this mechanism has been called resonant $`\tau `$-leptogenesis .
In studying leptogenesis, we have extended previous analyses of the relevant network of BEs. More explicitly, we have consistently taken into account SM chemical potential effects, as well as effects from out of equilibrium sphalerons and single lepton flavours. In particular, we have found that single lepton flavour effects become very important in R$`\tau `$L models. In this case, the difference between our improved formalism of BEs and the usual formalism followed in the literature could be dramatic. The predictions of the usual formalism could lead to an erroneous result which is suppressed by many orders of magnitude. The suppression factor could be enormous of order $`10^6`$ for the R$`\tau `$L scenarios considered in Section 5. Even within leptogenesis models with a mild hierarchy between the heavy neutrino masses, the usual formalism turns out to be inadequate to properly treat single lepton flavour effects; its predictions may differ even up to one order of magnitude with respect to those obtained with our improved formalism.
One generic feature of R$`\tau `$L models is that their predictions for the final baryon asymmetry are almost independent of the initial values for the primordial $`B`$-number, $`L`$-number and heavy Majorana neutrino abundances. Specifically, we have investigated the dependence of the BAU on the initial conditions, as a function of the heavy neutrino mass scale $`m_N`$. We have found that for $`m_N\stackrel{>}{_{}}250`$ GeV, the dependence of the BAU is always less than 7%, even if the initial baryon asymmetry is as large as $`\eta _B^{\mathrm{in}}=10^2`$ at $`z=m_N/T=0.1`$. For smaller values of $`m_N`$, this dependence starts getting larger. Thus, for $`m_N=100`$ GeV, the dependence of the final baryon asymmetry on the initial conditions is stronger, unless the primordial baryon asymmetry is smaller than $`10^6`$ at $`z=0.1`$.
In order to have successful leptogenesis in the R$`\tau `$L models under study, the heavy Majorana neutrinos are required to be nearly degenerate. This nearly degenerate heavy neutrino mass spectrum may be obtained by enforcing an SO(3) symmetry, which is explicitly broken by the Yukawa interactions to a particular SO(2) sub-group isomorphic to a lepton-type group U(1)<sub>l</sub>. The approximate breaking of U(1)<sub>l</sub>, which could result from a FN mechanism, leads to a Yukawa texture that accounts for the existing neutrino oscillation data, except those from the LSND experiment . Our choice of the breaking parameters was motivated by the naturalness of the light and heavy neutrino sectors. To obtain natural R$`\tau `$L models, we have followed the principle that there should be no excessive cancellations between tree-level and radiative or thermal effects. In this way, we have found that R$`\tau `$L models strongly favour a light neutrino mass spectrum with an inverted hierarchy. Moreover, when the same naturalness condition is applied to the heavy neutrino sector, a particular hierarchy for the mass differences of the heavy Majorana neutrinos is obtained. In particular, the mass difference of one pair of heavy Majorana neutrinos is much smaller than the other two possible pairs.
R$`\tau `$L models offer a number of testable phenomenological signatures for low-energy experiments and future high-energy colliders. These models contain electroweak-scale heavy Majorana neutrinos with appreciable couplings to electrons and muons, e.g. $`N_{1,2}`$. Specifically, the (normalized to the SM) $`W^\pm `$-boson couplings of electrons and muons to the heavy Majorana neutrinos $`N_{1,2}`$ could be as large as 0.01, for $`m_{N_{1,2}}=100`$–300 GeV. As a consequence, these heavy Majorana particles can be produced at future $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ colliders, operating with a c.m.s. energy $`\sqrt{s}=0.5`$–1 TeV. Another feature of R$`\tau `$L models is that thanks to the inverted hierarchic structure of the light neutrino mass spectrum, they can account for sizeable $`0\nu \beta \beta `$ decay. The predicted effective neutrino mass $`|(𝐦^\nu )_{ee}|`$ can be as large as 0.02 eV, which is within the sensitivity of the proposed next round of $`0\nu \beta \beta `$ decay experiments. The most striking phenomenological feature of 3-generation (non-supersymmetric) R$`\tau `$L models is that they can predict $`e`$\- and $`\mu `$-number-violating processes, such as the decay $`\mu e\gamma `$ and $`\mu e`$ conversion in nuclei, with observable rates. In particular, these LFV effects could be as large as $`10^{12}`$ for $`B(\mu e\gamma )`$ and as large as $`5\times 10^{13}`$ for a $`\mu e`$ conversion rate in $`{}_{22}{}^{}{}_{}{}^{48}`$Ti, normalized to the $`\mu `$ capture rate. The above predicted values are within reach of the experiments proposed by the MEG and MECO collaborations.
Although the present study improves previous analyses of the BEs related to leptogenesis models, there are still some additional smaller but relevant effects that would require special treatment. The first obvious improvement would be to calculate the thermal effects on the collision terms, beyond the HTL approximation. These corrections would eliminate some of the uncertainties pertinent to the actual choice of the IR regulator in some of the collision terms. These effects limit the accuracy of our predictions and introduce an estimated theoretical uncertainty of 30% for leptogenesis models operating well above the electroweak phase transition, with relatively large $`K`$ factors, i.e. $`K_{lN_\alpha }\stackrel{>}{_{}}5`$. For models at the electroweak phase transition, the IR problem is less serious, but larger uncertainties may enter due to the lack of a satisfactorily accurate quantitative framework for sphaleron dynamics. Although the implementation of the sphaleron dynamics in our BEs for RL models was based on the calculations of , particular treatment would be needed, if the electroweak phase transition was a strong first-order one. In this case, the dynamics of the expanding bubbles during the electroweak phase transition becomes relevant . This possibility may emerge in supersymmetric versions of RL models. Nevertheless, the inclusion of the aforementioned additional effects is expected not to modify the main results of the present analysis drastically and will be the subject of a future communication.
### Acknowledgements
We thank Mikko Laine, Costas Panagiotakopoulos, Graham Ross, Kiriakos Tamvakis and Carlos Wagner for useful discussions and comments. The work of AP has been supported in part by the PPARC research grants: PPA/G/O/2002/00471 and PP/C504286/1. The work of TU has been funded by the PPARC studentship PPA/S/S/2002/03469.
## Appendix A Collision Terms
### A.1 Useful Notation and Definitions
The following notation and definitions are used in the derivation of the BEs. The number density, $`n_a`$, of a particle species, $`a`$, with $`g_a`$ internal degrees of freedom is given by
$`n_a(T)`$ $`=`$ $`g_a{\displaystyle \frac{d^3𝐩}{(2\pi )^3}\mathrm{exp}\left[\left(\sqrt{𝐩^2+m_a^2}\mu _a(T)\right)/T\right]}`$ (A.1)
$`=`$ $`{\displaystyle \frac{g_am_a^2Te^{\mu _a(T)/T}}{2\pi ^2}}K_2\left({\displaystyle \frac{m_a}{T}}\right),`$
where $`\mu _a`$ is the $`T`$-dependent chemical potential and $`K_n(x)`$ is the $`n`$th-order modified Bessel function . In our minimal leptogenesis model, the factors $`g_a`$ are: $`g_{W^a}=3g_B=6`$ and $`g_\mathrm{\Phi }=g_\mathrm{\Phi }^{}=2`$, and for the $`i`$th family: $`g_{N_\alpha }=2`$, $`g_{L_i}=g_{L_i^C}=4`$, $`g_{Q_i}=g_{Q_i^C}=12`$, and $`g_{u_i}=g_{u_i^C}=6`$. Using the same formalism as the CP-conserving collision term for a generic process $`XY`$ and its CP-conjugate $`\overline{X}\overline{Y}`$ is defined as
$$\gamma _Y^X\gamma (XY)+\gamma (\overline{X}\overline{Y}),$$
(A.2)
with
$$\gamma (XY)=𝑑\pi _X𝑑\pi _Y(2\pi )^4\delta ^{(4)}(p_Xp_Y)e^{p_X^0/T}|(XY)|^2.$$
(A.3)
In the above, $`|(XY)|^2`$ is the squared matrix element which is summed but not averaged over the internal degrees of freedom of the initial and final multiparticle states $`X`$ and $`Y`$. Moreover, $`d\pi _X`$ represents the phase space factor of a multiparticle state $`X`$,
$$d\pi _X=\frac{1}{S_X}\underset{i=1}{\overset{n_X}{}}\frac{d^4p_i}{(2\pi )^3}\delta (p_i^2m_i^2)\theta (p_i^0),$$
(A.4)
where $`S_X=n_{\mathrm{id}}!`$ is a symmetry factor depending on the number of identical particles, $`n_{\mathrm{id}}`$, contained in $`X`$.
As CPT is preserved, the CP-conserving collision term $`\gamma _Y^X`$ obeys the relation
$$\gamma _Y^X=\gamma _X^Y.$$
(A.5)
Analogously, it is possible to define a CP-violating collision term $`\delta \gamma _Y^X`$ as
$$\delta \gamma _Y^X\gamma (XY)\gamma (\overline{X}\overline{Y})=\delta \gamma _X^Y,$$
(A.6)
where the last equality follows from CPT invariance.
### A.2 CP-Conserving Collision Terms
In numerically solving the BEs, we introduce the dimensionless parameters:
$$z=\frac{m_{N_1}}{T},x=\frac{s}{m_{N_1}^2},a_\alpha =\left(\frac{m_{N_\alpha }}{m_{N_1}}\right)^2,a_r=\left(\frac{m_{\mathrm{IR}}}{m_{N_1}}\right)^2,c_\alpha ^l=\left(\frac{\mathrm{\Gamma }_{N_\alpha }^l}{m_{N_1}}\right)^2,$$
(A.7)
where $`\alpha =1,2,3`$ labels the heavy Majorana neutrino states, $`s`$ is the usual Mandelstam variable and $`m_{\mathrm{IR}}`$ is an infra-red (IR) mass regulator which is discussed below.
In terms of the resummed effective Yukawa couplings $`(\overline{h}_\pm ^\nu )_{l\alpha }`$ introduced in , the radiatively corrected decay width $`\mathrm{\Gamma }_{N_\alpha }^l`$ of a heavy Majorana neutrino $`N_\alpha `$ into a lepton flavour $`l`$ is given by
$$\mathrm{\Gamma }_{N_\alpha }^l=\frac{m_{N_\alpha }}{16\pi }\left[(\overline{h}_+^\nu )_{l\alpha }^{}(\overline{h}_+^\nu )_{l\alpha }+(\overline{h}_{}^\nu )_{l\alpha }^{}(\overline{h}_{}^\nu )_{l\alpha }\right].$$
(A.8)
By means of (A.3), the $`12`$ CP-conserving collision term $`\gamma _{L_l\mathrm{\Phi }}^{N_\alpha }`$ is found to be
$`\gamma _{L_l\mathrm{\Phi }}^{N_\alpha }=\gamma (N_\alpha L_l\mathrm{\Phi })+\gamma (N_\alpha L_l^C\mathrm{\Phi }^{})`$ $`=`$ $`\mathrm{\Gamma }_{N_\alpha }^lg_{N_\alpha }{\displaystyle \frac{d^3𝐩_{N_\alpha }}{(2\pi )^3}\frac{m_{N_\alpha }}{E_{N_\alpha }(𝐩)}e^{E_{N_\alpha }(𝐩)/T}}`$ (A.9)
$`=`$ $`{\displaystyle \frac{m_{N_1}^4a_i\sqrt{c_i^l}}{\pi ^2z}}K_1(z\sqrt{a_i}),`$
where $`E_{N_\alpha }(𝐩)=\sqrt{𝐩^2+m_{N_\alpha }^2}`$ and $`g_{N_\alpha }=2`$ is the number of internal degrees of freedom of $`N_\alpha `$. Upon summation over lepton flavours $`l`$, this collision term reduces to the corresponding one given in (B.4) of .
For $`22`$ processes, one can make use of the reduced cross section $`\widehat{\sigma }(s)`$ defined as
$$\widehat{\sigma }(s)8\pi \mathrm{\Phi }(s)𝑑\pi _Y(2\pi )^4\delta ^{(4)}(qp_Y)\left|(XY)\right|^2,$$
(A.10)
where $`s=q^2`$ and the initial phase space integral is given by
$$\mathrm{\Phi }(s)𝑑\pi _X(2\pi )^4\delta ^{(4)}(p_Xq).$$
(A.11)
These expressions simplify to give
$$\widehat{\sigma }(s)=\frac{1}{8\pi s}\underset{t_{}}{\overset{t_+}{}}𝑑t\left|(XY)\right|^2,$$
(A.12)
where $`t`$ is the usual Mandelstam variable, and the phase-space integration limits $`t_\pm `$ will be specified below.
In processes, such as $`N_\alpha V_\mu L_l\mathrm{\Phi }`$, the exchanged particles (e.g. $`L`$ and $`\mathrm{\Phi }`$) occurring in the $`t`$ and $`u`$ channels are massless. These collision terms possess IR divergences at the phase-space integration limits $`t_\pm `$ in (A.12). Within a more appropriate framework, such as finite temperature field theory, these IR singularities would have been regulated by the thermal masses of the exchanged particles. In our $`T=0`$ field theory calculation, we have regulated the IR divergences by cutting off the phase-space integration limits $`t_\pm `$ using a universal thermal regulator $`m_{\mathrm{IR}}`$ related to the expected thermal masses of the exchanged particles. This procedure preserves chirality and gauge invariance, as would be expected within the framework of a finite temperature field theory .
Thermal masses for the Higgs and leptons are predominantly generated by gauge and top-quark Yukawa interactions. In the HTL approximation, they are given by
$`{\displaystyle \frac{m_L^2(T)}{T^2}}`$ $`=`$ $`{\displaystyle \frac{1}{32}}\left(3g^2+g^{\mathrm{\hspace{0.17em}2}}\right),`$
$`{\displaystyle \frac{m_\mathrm{\Phi }^2(T)}{T^2}}`$ $`=`$ $`2d\left(1{\displaystyle \frac{T_c^2}{T^2}}\right),`$ (A.13)
where $`d=(8M_W^2+M_Z^2+2m_t^2+M_H^2)/(8v^2)`$. In our numerical estimates, we choose the regulator $`m_{\mathrm{IR}}`$ to vary between the lepton and Higgs thermal masses, evaluated at $`Tm_N`$. The resulting variation in the predicted baryon asymmetry can be taken as a contribution to the theoretical uncertainties in our zero temperature calculation.
For reduced cross-sections with an apparent singularity at the upper limit $`t_+`$, the following upper and lower limits are used:
$$t_+=m_{\mathrm{IR}}^2,t_{}=m_{N_\alpha }^2s.$$
(A.14)
Likewise, for reduced cross-sections with apparent singularities at both the upper and lower limits $`t_\pm `$, the following limits are employed:
$$t_+=m_{\mathrm{IR}}^2,t_{}=m_{N_\alpha }^2+m_{\mathrm{IR}}^2s.$$
(A.15)
It is important to remark here that the collision terms do not suffer from IR singularities at $`T\stackrel{<}{_{}}T_c`$, because the leptons, $`W`$ and $`Z`$ bosons receive $`v(T)`$-dependent masses during the electroweak phase transition. The full implementation of such effects will be given elsewhere.
Substituting (A.10) and (A.11) into (A.3), one obtains
$$\gamma _Y^X=\frac{m_{N_1}^4}{64\pi ^4z}\underset{x_{\mathrm{thr}}}{\overset{\mathrm{}}{}}𝑑x\sqrt{x}K_1(z\sqrt{x})\widehat{\sigma }_Y^X(x),$$
(A.16)
where $`x_{\mathrm{thr}}`$ is the kinematic threshold for a given $`22`$ process.
For $`22`$ $`\mathrm{\Delta }L=1`$ processes, one can repeat the procedure in (Appendix B), without summing over lepton flavours. Each $`\mathrm{\Delta }L=1`$ process has an identical factor dependent on $`\overline{h}_\pm ^\nu `$. To produce the $`\mathrm{\Delta }L=1`$ collision terms for each lepton flavour, this factor needs to be replaced with its un-summed equivalent,
$$(\overline{h}_+^\nu )_{l\alpha }^{}(\overline{h}_+^\nu )_{l\alpha }+(\overline{h}_{}^\nu )_{l\alpha }^{}(\overline{h}_{}^\nu )_{l\alpha },$$
(A.17)
exactly as was done in (A.8). The remainder of the analytic expression for each of these terms is presented in .
In addition to the Higgs and gauge mediated $`\mathrm{\Delta }L=1`$ terms, there are also $`22`$ $`\mathrm{\Delta }L=2`$ processes. As before, these processes are $`L_k\mathrm{\Phi }L_l^C\mathrm{\Phi }^{}`$ and $`L_kL_l\mathrm{\Phi }^{}\mathrm{\Phi }^{}`$ where the former has its real intermediate states subtracted. The analytic forms of these collision terms are identical to the total lepton number case but lepton flavour is not summed over. The reduced cross sections are given by
$`\widehat{\sigma }_{L_l^C\mathrm{\Phi }^{}}^{L_k\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle \underset{\alpha ,\beta =1}{\overset{3}{}}}\mathrm{Re}\{[(\overline{h}_+^\nu )_{k\alpha }^{}(\overline{h}_+^\nu )_{k\beta }(\overline{h}_+^\nu )_{l\alpha }^{}(\overline{h}_+^\nu )_{l\beta }+(\overline{h}_{}^\nu )_{k\alpha }^{}(\overline{h}_{}^\nu )_{k\beta }(\overline{h}_{}^\nu )_{l\alpha }^{}(\overline{h}_{}^\nu )_{l\beta }]𝒜_{\alpha \beta }^{(ss)}`$ (A.18)
$`+\mathrm{\hspace{0.25em}2}\left[(\overline{h}_+^\nu )_{l\alpha }^{}h_{l\beta }^\nu (\overline{h}_+^\nu )_{k\alpha }^{}h_{k\beta }^\nu +(\overline{h}_{}^\nu )_{l\alpha }^{}h_{l\beta }^\nu (\overline{h}_{}^\nu )_{k\alpha }^{}h_{k\beta }^\nu \right]𝒜_{\alpha \beta }^{(st)}`$
$`+\mathrm{\hspace{0.25em}2}\left(h_{k\alpha }^\nu h_{k\beta }^\nu h_{l\alpha }^\nu h_{l\beta }^\nu \right)𝒜_{\alpha \beta }^{(tt)}\},`$
and
$$\widehat{\sigma }_{\mathrm{\Phi }^{}\mathrm{\Phi }^{}}^{L_kL_l}=\underset{\alpha ,\beta =1}{\overset{3}{}}\mathrm{Re}\left(h_{k\alpha }^\nu h_{k\beta }^\nu h_{l\alpha }^\nu h_{l\beta }^\nu \right)_{\alpha \beta },$$
(A.19)
where the $`𝒜`$ and $``$ factors are presented in .
As we now consider lepton flavours separately, it is necessary to include $`\mathrm{\Delta }L=0`$, but lepton flavour violating interactions. The three lowest order $`22`$ processes are shown diagrammatically in Figure 10: $`L_k\mathrm{\Phi }L_l\mathrm{\Phi }`$, $`L_k\mathrm{\Phi }^{}L_l\mathrm{\Phi }^{}`$ and $`L_kL_l^C\mathrm{\Phi }^{}\mathrm{\Phi }`$ (note that $`kl`$). The first of these reactions contains heavy Majorana neutrinos as RISs. These need be removed using the procedure outlined in . The reduced cross section for each of these processes is
$$\widehat{\sigma }_{L_l\mathrm{\Phi }}^{L_k\mathrm{\Phi }}=\underset{\alpha ,\beta =1}{\overset{3}{}}\left[(\overline{h}_+^\nu )_{l\alpha }^{}(\overline{h}_{}^\nu )_{k\alpha }^{}(\overline{h}_+^\nu )_{l\beta }(\overline{h}_{}^\nu )_{k\beta }+(\overline{h}_{}^\nu )_{l\alpha }^{}(\overline{h}_+^\nu )_{k\alpha }^{}(\overline{h}_{}^\nu )_{l\beta }(\overline{h}_+^\nu )_{k\beta }\right]𝒞_{\alpha \beta }$$
(A.20)
with
$$𝒞_{\alpha \beta }=\{\begin{array}{cc}\frac{xa_\alpha }{4\pi \left|D_\alpha ^2\right|}0\hfill & (\alpha =\beta )\\ \frac{x\sqrt{a_\alpha a_\beta }}{4\pi P_\alpha ^{}P_\beta }\hfill & (\alpha \beta )\end{array}$$
(A.21)
In (A.21), $`P_\alpha ^1(x)`$ is the Breit–Wigner $`s`$-channel propagator
$$P_\alpha ^1(x)=\frac{1}{xa_\alpha +i\sqrt{a_\alpha c_\alpha }}.$$
(A.22)
Therefore, following the procedure in , the RIS-subtracted propagator is determined by
$$|D_\alpha ^1(x)|^2=|P_\alpha ^1(x)|^2\frac{\pi }{\sqrt{a_\alpha c_\alpha }}\delta (xa_\alpha )0.$$
(A.23)
Processes (b) and (c) in Fig. 10 do not contain RISs and have the following reduced cross sections:
$`\widehat{\sigma }_{L_l\mathrm{\Phi }^{}}^{L_k\mathrm{\Phi }^{}}`$ $`=`$ $`{\displaystyle \underset{\alpha ,\beta =1}{\overset{3}{}}}\mathrm{Re}\left(h_{l\alpha }^\nu h_{k\alpha }^\nu h_{l\beta }^\nu h_{k\beta }^\nu \right)𝒟_{\alpha \beta },`$ (A.24)
$`\widehat{\sigma }_{\mathrm{\Phi }^{}\mathrm{\Phi }}^{L_kL_l^C}`$ $`=`$ $`{\displaystyle \underset{\alpha ,\beta =1}{\overset{3}{}}}\mathrm{Re}\left(h_{l\alpha }^\nu h_{k\alpha }^\nu h_{l\beta }^\nu h_{k\beta }^\nu \right)_{\alpha \beta },`$ (A.25)
where for $`\alpha \beta `$,
$`𝒟_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{\sqrt{a_\alpha a_\beta }}{\pi x(a_\alpha a_\beta )}}\left[(x+a_\beta )\mathrm{ln}\left({\displaystyle \frac{x+a_\beta }{a_\beta }}\right)(x+a_\alpha )\mathrm{ln}\left({\displaystyle \frac{x+a_\alpha }{a_\alpha }}\right)\right],`$ (A.26)
$`_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{\sqrt{a_\alpha a_\beta }}{\pi (a_\alpha a_\beta )}}\mathrm{ln}\left({\displaystyle \frac{a_\alpha (x+a_\beta )}{a_\beta (x+a_\alpha )}}\right),`$ (A.27)
and for $`\alpha =\beta `$,
$`𝒟_{\alpha \alpha }`$ $`=`$ $`{\displaystyle \frac{a_\alpha }{\pi x}}\left[{\displaystyle \frac{x}{a_\alpha }}\mathrm{ln}\left({\displaystyle \frac{x+a_\alpha }{a_\alpha }}\right)\right],`$ (A.28)
$`_{\alpha \alpha }`$ $`=`$ $`{\displaystyle \frac{x}{\pi (x+a_\alpha )}}.`$ (A.29)
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# A fresh view on Henize 2-10 with VLT/NAOS-CONICA1footnote 11footnote 1based on observations made at the ESO VLT under program 71.B-0492, and 270.B-5011
## 1 Introduction
Henize 2-10 (He 2-10 (catalog He2-10)) is a blue compact galaxy with quite interesting properties. It is a nearby starburst galaxy with a heliocentric distance of $`9\pm 5`$ Mpc (Johansson, 1987; Tully, 1988). In this paper we adopt a distance of 9 $`h^1`$Mpc, with $`H_0=75`$ km/s, which yields a scale of $``$45pc/arcsec. The galaxy shows spectroscopic Wolf-Rayet features which indicates the presence of a very young starburst region (Hutsemekers & Surdej, 1984; Vacca & Conti, 1992), and has a slightly sub-solar metallicity (Kobulnicky et al., 1999). Figure 1 shows a RGB-composite image (in color in electronic version only) of the central $`18\mathrm{}\times 27\mathrm{}`$ ($`800\times 1200`$ parsecs) of He 2-10, where most of the current star formation occur. The blue and green channels are 0.3″-convolved HST WFPC archive in $`F555W`$ and $`F814W`$. The red channel is our $`0.3\mathrm{}`$-seeing $`K_S`$ taken with VLT/ISAAC. The bright central nucleus, generally referred to as region A, is an arc of UV-bright super-star clusters (Vacca & Conti, 1992). It is surrounded by two presumably older star-forming regions. Region B on the east shows a mixed population of blue and red clusters (only detected in $`K_S`$ and longer wavelength bands). Region C, on the north-west side, has a long tail containing bright red clusters as well. Dust is clearly apparent in this image, as shown by the red filaments observed to the south-east side between A and B. These appear to be blown away from the central region, rapidly dissolving in He 2-10’s interstellar medium. Assuming for the dust clouds a velocity equal to the typical sound speed in the ISM, i.e. 10 $`10`$ km$``$s<sup>-1</sup>, hence a growing rate of $`10`$pc$`/`$Myr, the present radii of curvature of the filaments $`50100`$ pc yield dynamical ages of $`510`$Myr. Region A is also flanked by two compact red sources that are only visible at $`K_S`$ and longer wavelengths. These two sources get brighter with longer wavelengths, as shown in Figure 3. The various colors of the sources hint at a highly heterogeneous dust content, and possibly age differences among the cluster population.
Recent observations in the optical, IR, and radio, have brought new exciting facts.
First, the youth of the starburst event in the center was confirmed by STIS analysis of the brightest UV/optical knots by Chandar et al. (2003) which yielded a coeval formation age of 4-5 Myr for all optical clusters.
Second, the presence and importance of dust in the central region was confirmed by high-resolution mid-infrared observations (Sauvage et al., 1997; Beck et al., 2001; Vacca et al., 2002): the majority (60%) of the MIR emission is confined to a $`5\mathrm{"}`$ region, compatible in size with the location of the observed starburst. However the intrinsically large uncertainties in MIR astrometry prevented a clear identification of the MIR sources.
Finally, the radio observations of Kobulnicky & Johnson (1999, 2000) evidenced 5 compact radio sources (hereafter called the radio knots) characterized by mostly thermal spectra. The striking morphological resemblance between the radio knots and the MIR emission allowed to tie down the location of the MIR sources precisely (Beck et al., 2001). Furthermore, comparing with HST images, Kobulnicky & Johnson (1999, 2000) argued that most of these radio-MIR sources were off-centered, in the dusty area between region A and B. They thus attributed this emission to young ($`1<`$Myr) ultra dense (UD) HII regions with ongoing star formation hidden in dense molecular clouds. Following this interpretation, Beck et al. (2001), from $`11.7\mu `$m observations, derived that up to $`10^4`$ O7V stars (equivalent to $`10^{49}`$ Lyman photons$``$s<sup>-1</sup>) must be hidden in these dense cocoons, and Vacca et al. (2002), from $`10.8\mu `$m observations, computed that $`10^7M_{}`$ of dust and gas must be surrounding the UD HII regions and estimated the bolometric luminosity of the brightest MIR region overlapping with radio knot 4 to be as much as $`2\times 10^9L_{}`$. Johnson & Kobulnicky (2003) extended the radio observations and refined the measurement of physical properties of the UD HII regions, in particular an anomalous low mass HII content of 2-8$`\times 10^3M_{}`$ was found, attributed to the extreme youth of the objects.
Therefore, as of today, He 2-10 appears as a spectacular case of a starburst galaxy where a large fraction of its most current star formation activity lies completely buried in dust, and has absolutely no visible counterpart.
We present here high-resolution observations in $`K_S`$ ($`2.2\mu `$m) with VLT/ISAAC (under ESO program 270.B-5011(A) (catalog 270.B-5011(A))), $`L^{}`$ (3.8$`\mu `$m), and $`M^{}`$ (4.8$`\mu `$m) bands with the Adaptive Optics VLT / NAOS - CONICA (under ESO program 71.B-0492(A) (catalog 71.B-0492(A))) that give the highest resolution to date of the nucleus of He 2-10 in the NIR, a wavelength regime adequately located between the stellar optical regime and the dust thermal regime. The high quality of the observations allows the identification, for the first time, of bright $`L^{}`$ regions which correlate with radio knots (Kobulnicky & Johnson, 1999), and $`K_S`$ regions that correlate with the optically bright cluster, thus bridging the existing wavelength gap.
## 2 Observations and data reduction
We observed He 2-10 in $`K_S`$ using ISAAC at the ESO VLT/Antu under 0.3$`\mathrm{}`$ seeing (Figure 1), and in $`L^{}`$ and $`M^{}`$ with the adaptive optics system NAOS-CONICA at the ESO VLT/Yepun, reaching a corrected PSF full width at half maximum (FWHM) of 0.12$`\mathrm{}`$. The observations were obtained in service mode in 2003-2004. We used standard observing strategies. In $`L^{}`$, we alternated exposures between the object and the sky in 15 ABBA sequences with a throw of 40″. Object cubes (respectively sky cubes) consisted of 220 (110) random 0.175 s jitter for a total exposure time of 1155 s (577.5 s). In $`M^{}`$, as NACO was still in early operation, we used the only allowed chopping mode with pre-defined exposure time and a 30″chopping throw. Typical strehl ratios of validated exposures were in the range 20-25% for both filters.
The ABBA sequences were combined/substracted using custom iraf scripts. We performed accurate astrometry reconstruction and photometry as described in the following sections.
### 2.1 Astrometry
In order to compare the astrometry of our NAOS-CONICA images with previous multiwaveband observations, we performed an accurate calibration of the NAOS-CONICA to the USNO-B based J2000 epoch. We first calibrated the astrometry of the $`K_S`$ image using 23 field stars catalogued in the USNO-B (Monet et al., 1998), and applied general transformations using the NOAO iraf routine geomap. The positional accuracies of the transformations are $`\mathrm{\Delta }\alpha _{rms}=0.008`$ sec and $`\mathrm{\Delta }\delta _{rms}=0.08\mathrm{}`$. We then derived the transformations between the $`K_S`$ USNO-calibrated image and $`L`$ image using four common sources in the central region (two field stars and two bright clusters). The number of available common sources is too small to derive a robust general transformation but the clear morphological similarities between the $`K_S`$ and $`L`$ images (Figure 3) make us confident that the reconstructed $`L`$ astrometry has an accuracy equivalent to that of the $`K_S`$ image. The positional accuracies of the transformations are $`\mathrm{\Delta }\alpha _{rms}=0.001`$ sec and $`\mathrm{\Delta }\delta _{rms}=0.012\mathrm{}`$. The resulting astrometry, computed for all $`L^{}`$ sources, is shown Table 1.
### 2.2 Photometry
We extracted $`L^{}`$\- and $`M^{}`$ fluxes of all detected sources of the inner core of He 2-10 in aperture sizes of 0.5$`\mathrm{}`$ (the finding chart in $`L^{}`$ is given Figure 2, left panel), the good match between VLT-NACO and HST resolutions allows us to compute consistent aperture photometry over the entire optical and near-infrared (NIR) range. We measured the $`F555W`$, H$`\alpha `$, $`F814W`$ fluxes from the HST calibrated archive images, on the same positions and aperture sizes, including a systematic centroid position mismatch in the error computation. All photometric measurements were then calibrated to Vega magnitudes using IR standards provided by ESO service observing team. All ancillary data may be found in http://archive.eso.org/eso/eso\_archive\_main.html, under Program ID 71.B-0492(A). Typically, $`L^{}`$ errors on zeropoints are ca. 10%, because of a highly variable background. Additional uncertainties are due to the anisoplanetism of the FOV of view and the different strehl ratio between the photometric standard and the science observations. Both effect are difficult to quantify, but we assume that, because of our large apertures ($`5\times `$FWHM), the strehl difference will not dominate the systematics and most of the error will actually come from background subtraction. Unfortunately NACO’s response was not fully known when the $`M^{}`$ data were taken and the standard stars were observed with high ADU counts. They possibly reach the non-linear regime of the detector and the derived zero-points from different nights disagree to $`\pm 0.25`$ mag. Hence the $`M^{}`$ data are mostly useful for morphology and only deliver a crude photometry.
HST $`F555W`$, $`F658N`$/H$`\alpha `$, and $`F814W`$ images were analysed following standard recipes detailed by Holtzman et al. (1995). We extracted the sources of the nucleus with the full-resolution images in order to compare with $`L^{}`$ data, and we convolved the HST data with a $`0.3\mathrm{}`$ gaussian in order to compare with $`K_S`$ data. We derived H$`\alpha `$ equivalent width interpolating the continuum between $`F555W`$ and $`F814W`$ following Johnson et al. (2000). We recover the original magnitudes and similar H$`\alpha `$ equivalent widths of Johnson et al. (2000) and the V magnitude of knot 4 (Chandar et al., 2003) using different strategies with smaller apertures within photon and aperture position errors.
All Vega calibrated magnitudes were corrected for the galactic extinction of $`E(BV)=0.11`$ (Schlegel et al., 1998) yielding $`A_V=0.369`$, $`A_{\text{H}\alpha }=0.298`$ $`A_I=0.216`$, $`A_K=0.04`$, $`A_L=0.017`$, and $`A_M=0.0`$.
We also extracted all sources detected in $`K_S`$, $`F814W`$ and $`F555W`$ in the outer 1′ of He 2-10 (0.5$`\mathrm{}`$ aperture size, see Figure 2 right panel). A special care was taken to check the effect of the larger PSF in $`K_S`$ and HST convolved resolution, compared to HST original resolution. We computed the total flux lost between unconvolved and convolved point sources for 0.5$`\mathrm{}`$ apertures. The loss amounts to $`0.31\pm 0.05`$ mag in $`F814W`$ and $`F555W`$. Although substantial in absolute, the correction is negligible in color, hence the color-color diagrams are not affected by the correction applied to the sources of the outer 1′.
Finally, another photometric extraction was done on the brightest $`L^{}`$ and $`M^{}`$ sources in order to build the full spectral energy distributions of these regions. We convolved $`L^{}`$ and $`M^{}`$ images with 0.3″ and extracted fluxes around the brightest radio knots in 1″ apertures.
The largest source of uncertainty comes from the diffuse background in the inner core of He 2-10. We adopted different strategies to estimate the background in the different bands, in order to account for the variable resolutions and field crowding of the data. The HST images were sky subtracted using sky annuli of radius 0.3-0.5$`\mathrm{}`$ around the objects located in the inner 10$`\mathrm{}`$ of the nucleus, and using annuli 0.4-0.6 pix for outer sources (also for $`K_S`$). In $`K_S`$, $`L^{}`$ and $`M^{}`$, we computed the diffuse emission background in annuli of radius 0.3-1$`\mathrm{}`$ around the objects located in the inner 10$`\mathrm{}`$ of the nucleus. We additionally computed 2 sky annuli (0.3-0.5$`\mathrm{}`$, and 0.3-2$`\mathrm{}`$) in order to estimate the systematic errors resulting from the measured flux differences. Sytematic errors of $`0.15`$ mag were typical. The $`M^{}`$ magnitudes have higher systematics due the uncertainty on the zeropoints. Finally the background around the 1″ apertures was measured in annuli of 0.5-1.5″.
## 3 Results
### 3.1 Optical to radio identification of the central sources
Figure 3 shows the entire series of high-resolution observations of the nucleus of Henize 2-10, in $`F555W`$, $`F658N/H\alpha `$, $`F814W`$, $`K_S`$, $`L^{}`$, $`M^{}`$ (shifted to USNO-B absolute astrometry, see section 2.1), $`N`$ and $`2`$ cm. The $`M^{}`$ image is much shallower than the $`L^{}`$ image, and only the three reddest sources are robustly detected. As is evident from this figure, and also from Table 1, the $`L^{}`$ and $`M^{}`$ sources are clear counterparts to some of the radio knots, and because of the strong morphological and physical association between the radio and MIR emissions, to the MIR sources as well. Figure 4 shows a precise comparison of the $`L^{}`$ images (gray scale) with $`N`$ GEMINI/OSCIR $`N`$ contours (left, the $`N`$ source associations to the radio knots are according to Vacca et al., 2002), and with VLA 2-cm contours (right, with the radio knots identification of Kobulnicky & Johnson, 1999). It is striking to see, on Fig 4, how well sources L1 and L5 correlate with the N-band sources. Furthermore Figure 3 also demonstrates the association of L4a with the brightest source in F814W, and of L6, L3a and L3b with the string of three clusters detected west of the brightest source in F814W. Therefore, in constrast to the statement of Kobulnicky & Johnson (1999, 2000), we see a clear connection between the structures observed from 0.5 $`\mu `$m to 2 cm.
At this stage it is important to stress that the positional association from the optical to the radio that we present here is not a new proposition that we are making, but is a fact imposed by the astrometry. The HST, $`K_S`$, $`L^{}`$, and $`M^{}`$ images are now tied to the same astrometric reference, which is the USNO-B astrometry. This is accurate to $`\pm 0.3\mathrm{}`$ which will be the dominant positional uncertainty in this wavelength regime (for instance the relative astrometric accuracy of NAOS-CONICA is $`\pm 0.02\mathrm{}`$). The absolute astrometry of the VLA is accurate to $`\pm 0.1\mathrm{}`$. With these accuracies, it becomes impossible to accept the relative positions of the HST and MIR/radio sources as presented in Kobulnicky & Johnson (1999, 2000); Vacca et al. (2002), i.e. a shift of 1.2″ in $`\alpha `$, even taking into account the beam size of the VLA, which is 0.4″ in $`\alpha `$. In fact, the main source of astrometric error in previous papers was an incorrect astrometry in the HST data used at the time. This is now corrected in the archived versions of the same data. Therefore Figure 3 now reveals for the first time the correct evolution of the central region of He 2-10 from the optical to the radio. This allows now to present a much firmer identification of the different components across the spectrum.
L1 is clearly detected in $`M^{}`$ and can now be associated to the western N band source. It has a rather faint counterpart in $`K_S`$ and is no longer detected in the optical bands, and shows a clear anti-correlation with H$`\alpha `$ (Fig. 5). Figure 4 shows that L1 appears to fall between the radio knots 1 and 2, however this time the shift is too small to be significant. Since L1 is accompanied by a faint source L2, we propose to associate the two sources L1 and L2 with knot 1 and knot 2 in the radio map. L2 has no counterpart in the optical bands, an anti-correlation with H$`\alpha `$ (Fig. 5), much like L1, but is much fainter. Source L3a and L3b have clear optical and H$`\alpha `$ counterparts west of the brightest optical source (cf. Fig. 5). They fall in a region of diffuse N-band emission and their association with knot 3 is unclear. The brightest $`L^{}`$ source L4a is the counterpart of the main source in the F814W image. It is fainter in the H$`\alpha `$ image but reappears strongly in the F555W image. Again L4a appears displaced from its possible radio and MIR counterpart knot 4. In fact it is quite likely that knot 4 and the associated brightest MIR source are the counterparts of the group L4b, c and d. These sources are easily detectable in the optical wide bands and are quite strong in H$`\alpha `$. L5 is another bright source that has no optical or H$`\alpha `$ counterpart (Fig. 5) but is cleary detected in $`K_S`$, N and radio and can be associated with knot 5. L6 has no clear radio or MIR counterpart but is associated to the optical cluster just west of the brightest F814W source. Finally L7 and L8 have no definite counterparts but are associated with diffuse emission in all bands.
We emphazise that the fact that we detect spatially correlated sources accross the spectrum doesn’t mean that we actually detect the *same* sources. An error of 0.2″ translates to 9 pc at He 2-10’s distance, while the compact sources are not expected to have sizes much larger than 1-2pc, hence the $`L^{}`$ sources and radio knots could belong to contiguous yet distinct star forming regions. A robust identification will have to wait for deeper $`L^{}`$ observations and higher resolution radio observations. Nevertheless, it is interesting to follow the working hypothesis that association between radio and NIR sources is indeed a physical one.
### 3.2 Near-infrared/optical colors of He 2-10 nucleus
We selected the sources likely to be clusters at 9 Mpc with a magnitude threshold of $`M_V<8.5`$ (Whitmore et al., 1999; Johnson et al., 2000). Table 2 shows the corresponding Vega magnitudes of the 12 central sources (Fig 2). All 12 sources seen in $`L^{}`$ would still be too bright to be supergiant stars if Henize 2-10 was at a distance of 6 Mpc, but 2 sources out of the 21 outer sources detected in $`V`$, $`I`$ and $`K_S`$ would fall below the threshold. We plot $`VI`$ vs $`VK`$, $`VI`$ vs $`IL^{}`$ and $`VI`$ vs $`KL^{}`$ colors in Figure 6 along with the dust free STARBURST99 (Leitherer et al., 1999, SB99) tracks of solar metallicity instantaneous bursts, with Salpeter IMF, from 0 (low $`VI`$) to 1 Gyr (high $`VI`$). The sharp turnoff seen in $`VI`$ vs $`KL^{}`$ correspond to 6.3 Myr, after which nebular emission ceases to be important. The 12 sources of the nucleus are plotted as squares the size of which relate to the total errors. The source of the outer 1′are plotted as filled triangles for the sources belonging to the regions B and C of the north-east and west/north-west (Fig. 1), and as circles for sources falling at least 1′ away from the nucleus (all circles happen to be in the southern part of He 2-10, visible as bright blue \[probably foreground\] clusters in Fig. 1). Photometric errors are typically contained within the size of the symbol, but systematics related to the diffuse background subtraction can be of the order of 15%.
The solid arrows show the direction and the amplitude of a screen extinction of $`A_V=1`$ mag (assuming a Calzetti form instead of the canonical galactic extinction curve does not change the values by more than 3% in the range 0.55 $`\mu `$m-3.8 $`\mu `$m). All diagrams show the well-known age-extinction degeneracy of optical colors (here $`VI`$), as well as the gradual decrease of the extinction impact on colors as one goes to the infrared. All show that most of the data points cannot be explained by reddenning any point of the models by any amount of screen extinction, especially for the diagrams including the $`L^{}`$ photometry. This is a rather common observation in starburst regions (Johnson et al., 2004; Vanzi & Sauvage, 2004; Cresci et al., 2005).
On the top left panel, the squares and the triangles seem to occupy the same locus whereas the circles are shifted up by $`VI0.3\pm 0.1`$ mag. This shift can be due to a systematic effect in the background subtraction although unlikely because the colors of triangles and squares were derived with independent methods (cf section 2.2) and still lies on top of each other, hence do not seem to be dominated by systematics. The shift more likely comes from an infrared excess affecting the central sources of He 2-10. This conclusion finds some support in the large excesses of $`IL^{}`$ and $`KL^{}`$ colors observed for the 12 central sources in the other two color-color diagrams. These infrared color excesses can have a number of explanations that will be explored later. But let us first characterize the magnitude of the extinction affecting the central sources. This will help us setting constraints on clusters ages as well as on their actual infrared excess from their location in the color-color diagrams.
As expected in a complex region such as revealed by Figure 1, extinctions measured at different wavelengths do not agree well with one another. First Allen et al. (1976) measure $`A_V`$=2.3 from optical observations while Johansson (1987) correcting for the contribution of stellar absorption features obtain $`A_V=0.86`$. Then we can use the observation of Vanzi & Rieke (1997) to derive the extinction of the central part of He 2-10 from the H$`\alpha `$ to Br$`\gamma `$ ratio. Vanzi & Rieke measure a Br$`\gamma `$ flux of 6.3$`\times 10^{14}`$ erg$``$s$`{}_{}{}^{1}`$cm<sup>-2</sup> over an aperture of 2.4x15.6″centered on the K bright nucleus. From the HST H$`\alpha `$ image we measure the H$`\alpha `$ flux on the same aperture and obtain a value of 1.86e-12 erg/s/cm2. Assuming an intrisic H$`\alpha `$/Br$`\gamma `$ ratio of 155.4 corresponding to Te=7500 K, ne=1000cm<sup>-3</sup> (Storey & Hummer, 1995) we find $`A_V=1.25`$. Finally, from the ratio Br$`\gamma `$/Br10 of Vanzi & Rieke we obtain $`A_V`$=10.5 which is in full agreement with the extinction derived using the Br$`\alpha `$, Br$`\gamma `$ fluxes of Kawara et al. (1989). We thus observe a clear trend of increasing extinction from the optical to the IR which is typical of a system where the absorbing material, gas and dust is mixed with emitting sources. It is worthwhile to note here that correcting the clusters for extinction will bring them closer to the youngest part of the SB99.
But this correction cannot explain the significant red excess observed in $`VK`$, $`IL^{}`$ and $`KL^{}`$ for all the central sources. This red excess was also observed in a similar diagram built for NGC 5253 (Vanzi & Sauvage, 2004; Cresci et al., 2005) and for Haro 3 (Johnson et al., 2004), and most likely has the same origin: a contribution of hot dust in the NIR bands. Indeed as the earliest evolutionary phase of star clusters occurs deep in molecular clouds, we can expect to find dust and molecular clouds in the immediate vicinity of these clusters. Dust close to its sublimation temperature would be able to contaminate the $`L^{}`$ band. Using radiative transfer models, Vanzi & Sauvage (2004) showed that the location of the reddest clusters in the ($`VI`$,$`KL^{}`$) diagram of NGC 5253, could be explained by a combination of extinction, scattering and emission by dust in a shell around the cluster. This is likely what we observe here as well. In that respect, it is noteworthy that the three reddest sources in $`KL^{}`$ are L1, L2 and L5, the sources with no visible counterparts and associated with thermal radio knots. If we apply to the 12 central sources of He 2-10 color corrections such as those derived from the modeling of the red cluster in NGC 5253, this will bring all the sources to the bottom left of the diagrams, indicating ages of less than 6.3 Myr.
### 3.3 Physical properties
We study now the physical properties of the sources which dominates the $`L^{}`$ emission of the nucleus (Fig. 2). The bright $`L^{}`$ sources L1+L2, L4 and L5 are not or barely resolved with radii $`0.1\mathrm{}`$ or 4.5 pc, L5 is not isolated, there is a faint diffuse emission extending as far as 3 FWHMs (10 pc) (the apertures shown on Fig. 2 are 22.5pc wide). These sources are not resolved in $`M^{}`$ either (radii$`0.12\mathrm{}`$). Fig.6 bottom panel shows the histogram of the H$`\alpha `$ equivalent width Log(EW\[H$`\alpha `$\]) of the 12 central sources in units of Log\[Å\] (from Table 2) and the associated ages according to Leitherer & Heckman (1995). The histogram peaks at ages $`6`$ Myr in agreement with recent works (Johnson et al., 2000; Chandar et al., 2003) and with age trends inferred from previous section color analysis. We emphasize that the H$`\alpha `$ equivalent width measurements are prone to strong systematic effects. We substracted out the \[NII\] $`\lambda 6584`$ measured to be 31% of H$`\alpha `$ in an unpublished high-resolution spectrum covering the 1.6″ around source $`L4a`$. But the main source of error is a poorly-known continuum level. Our method to derive the continuum by interpolation of the contiguous broadband filters continuum for lack of a better estimator is crude and the quoted Log(EW\[H$`\alpha `$\]) should be carefully used.
The measured Log(EW\[H$`\alpha `$\]) centered on the sources L1, L2 and L5 are 2.60, 2.66 and 2.37, yielding ages$`5`$ Myr. But the H$`\alpha `$ emission shows remarkable anti-correlations with the source L1+2 and L5 and no emission in $`F555W`$, $`F658N`$ and $`F814W`$ can be clearly associated with source L1, L2 and L5. Hence the measured Log(EW\[H$`\alpha `$\]) is probably more related to He 2-10 foreground medium rather than the $`L^{}`$ emitting clusters. An interesting check on the age derived from Log(EW\[H$`\alpha `$\]) comes from the sources L3a, L3b, L4a and L6 because all three have unambiguous $`F555W`$, $`F658N`$ and $`F814W`$ counterparts. Log(EW\[H$`\alpha `$\])$`<1.5`$ and the brightest source in $`L^{}`$ (L4a of Fig. 2) have Log(EW\[H$`\alpha `$\])$`<1.5`$ (Table 2), consistent with the measurements of Johnson & Kobulnicky (2003, Fig. 9), yielding ages$`>`$6 Myr, whereas Chandar et al. (2003) derived ages of 4-5 Myr fitting single-burst SB99 templates to a de-reddened UV/STIS spectrum of the same sources. Thus there appears to be an age dichotomy, with source L1, L2 and L5 being younger than 5 Myr and the other $`L^{}`$ source being older.
On the basis of the new identification of the components accross the spectrum (see Section 3.1), it is important to review previous assumptions on the nature of He 2-10 radio knots emission. The presence of UD HII in He 2-10 is primarily inferred from observed turnovers from flat radio spectra ($`F_\nu =\nu ^\alpha `$ with $`\alpha =0`$, characteristic of an optically thin ionized gas), to thermal spectra ($`\alpha =2`$, signature of a free-free optically thick gas) (Kobulnicky & Johnson, 1999). This feature is common in galactic HII regions. These UD HII regions were presumed to be extremely young, e.g. less than one million year, mostly because of their compactness, their high electron densities, and of their lack of counterpart in the NIR and below. Indeed only HII regions buried extremely deep in dust, i.e. right at the start of their expansion, would remain invisible in the NIR. Since most of the radio sources now have counterparts, the interpretation of the radio knots must be updated.
Among the 5 radio sources, we would only classify knot 1+2 and knot 5, as bona fide UD HII regions. They have counterparts in the NIR but not in the visible. This implies a significant optical depth (typically $``$ 10 if we follow existing models of embedded super-star clusters), and thus a young age although possibly not as young as previously postulated.
The radio knot 3 was already known to be a different type of source than other radio knots from a strong non-thermal signature Johnson & Kobulnicky (2003), (very likely a supernova remnant). We also propose that the radio knot 4, associated with diffuse $`L^{}`$ emission (L4b,c,d) and showing strong correlations with bright optical and H$`\alpha `$ sources, is not an UD HII region but a complex mix of normal HII regions and supernova remnants. Knot 4 does show a non-thermal signature with a slightly negative spectral index $`\alpha `$. The proposed re-classification of the radio sources would lower the total mass of hidden $`O`$ stars (as deduced from the radio) by a factor of $``$2 Johnson & Kobulnicky (2003).
Detailed modelling of the spectral energy distributions is essential to have a deeper understanding of the star formation process ongoing in He 2-10, especially of sources L1+L2 and L5 which represent the best candidate for very young dust-enshrouded super-star clusters. This is however a long endeavour that we differ to a later paper. Nevertheless the present observations outline the importance of high-resolution multiwavelength datasets to disentangle the intrinsically complex nature of star forming regions.
## 4 Conclusions
We detected compact sources in the nucleus of Henize 2-10 with $`K_S`$, $`L^{}`$ and $`M^{}`$ observations using ISAAC and NAOS-CONICA on the VLT. The sources are compact ($`<4.5`$ pc), highly correlated with radio and mid-infrared ultradense HII regions, previously thought to be optically thick.
The color-color magnitude diagrams show the presence of strong red-excess in $`K_S`$ and $`L^{}`$. Such red excesses point at highly heterogeneous dust distribution and at the presence of a hot dust component emitting and scattering down to $`L^{}`$ and $`M^{}`$.
We tentatively review the previous classification of the radio knots by identifying two bona fide UD HII (knot 1+2 and knot 5) and propose that knots 3 and 4 are non-thermal radio sources, akin to supernova remnants.
These new high-resolution data uncover a complex structure in infrared. We suggest that to understand He 2-10 star forming history, a detailed model of the radiative processes is needed. This model should include all known components of the galaxy in a consistent way in order to fit the spectral energy distribution from radio to UV.
RAC wish to acknowledge an ESO grant for visiting scientist in Santiago, and thank Chip Kobulnicky for sharing his VLA 2-cm data and Kelsey Johnson for useful comment on the first version of the paper.
Facilities: VLT(ISAAC) VLT(NAOS-CONICA)
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# The physical origin of the Fresnel drag of light by a moving dielectric medium
## I Introduction
It is usual to consider the famous experiment of Fizeau (1851)Fizeau on the drag of light by a uniformly moving medium as one of the crucial experiments which, just as the Michelson-Morley experiment, cannot be correctly understood without profound modification of Newtonian space-time concepts (for a review of Einstein’s relativity as well as a discussion of several experiments the reader is invited to consult Jackson ; Bladel ; Smith ). The result of this experiment which was predicted by FresnelFresnel , in the context of elastic theory, is indeed completely justified by well known arguments due to von Laue (1907) Laue . He deduced the Fresnel-Fizeau result for the light velocity $`v`$ in a medium, corresponding to a relativistic first order expansion of the Einstein velocity transformation formula:
$$v=\frac{\frac{c}{n_0}+v_e}{1+\frac{v_e}{cn_0}}\frac{c}{n_0}+v_e\left(1\frac{1}{n_0^2}\right)+O\left(\frac{v^2}{c^2}\right).$$
(1)
Here, $`n_0`$ represents the optical index of refraction of the dielectric medium in its proper frame, and we suppose that the uniform medium motion with velocity $`v_e`$ is parallel to the path of the light and oriented in the same direction of propagation. In the context of electromagnetic theory Minkowsky ; Pauli all derivations of this effect are finally based on the invariance property of the wave operator $`_\mu ^\mu [\mathrm{}]`$ in a Lorentz transformation. It is easy to write the wave equation $`[_𝐱^{}^2\left(n_0^2/c^2\right)_t^{}^2]\psi =0`$ in the co-moving frame R’$`(𝐱^{},t^{})`$ of the medium in covariant formLeonhardt $`[_\mu ^\mu +\left(n_0^21\right)\left(v^\nu _\nu \right)^2]\psi =0`$ which is valid in all inertial frames and which for a plane waves, implies the result of Eq. 1. In this calculation we obtain the result $`v=c/n_0+v_e`$ if we use the Galilean transformation which proves the insufficiency of Newtonian dynamics.
However, the question of the physical meaning of this phenomenon is not completely clear. This fact is in part due to the existence of a derivation made by Lorentz (1895) Lorentz based on the mixing between the macroscopic Maxwell’s equations and a microscopic electronic oscillator model which is classical in the sense of the Newtonian dynamics. In his derivation Lorentz did not use the relativistic transformation between the two coordinate frames: laboratory and moving medium. Consequently, the relativistic nature of the reasoning does not appear explicitly. Following the point of view of Einstein (1915) Einstein the Lorentz demonstration must contain an implicit hypothesis of relativistic nature, however, this point has not been studied in the literature. Recent developments in optics of moving media Leonhardt ; Wilkens ; Artoni ; Yu allows us to consider this question as an important one to understand the relation between optics, relativity and newtonian dynamics. This constitutes the subject of the present paper. Here, we want to analyze the physical origin of the Fresnel-Fizeau effect. In particular we want to show that this phenomenon is, in its major part, independent of relativistic dynamics.
The paper is organized as follows. In section II we present the generalized Lorentz “microscopic-macroscopic” derivation of the Fresnel formula and the principal defect of this treatment. In section III we show how to derive the Fresnel result in a perturbation approach based on the Lorentz oscillator model and finally in IV we justify this effect independently from all physical assumptions concerning the electronic structure of matter.
## II The Lorentz electronic model and its generalization
In this part, we are going to describe the essential contents of the Lorentz model and of its relativistic extension. Let $`𝝃(𝒙,t)`$ be the displacement of an electron from its equilibrium position at rest, written as an explicit function of the atomic position $`𝒙`$ and of the time $`t`$. In the continuum approximation we can write the equation of motion for the oscillator as $`_t^2𝝃(𝒙,t)+\omega _0^2𝝃(𝒙,t)\frac{e}{m}𝐄_0e^{i\left(\omega t𝒌𝒙\right)}`$ where the supposed harmonic electric incident field appears and where the assumption of small velocity allows us to neglect the magnetic force term. In the case of a non relativistic uniformly moving medium we have
$$\left(_t+𝒗_e\mathbf{}\right)^2𝝃+\omega _0^2𝝃\frac{e}{m}\left(𝐄_0+\frac{𝒗_e}{c}\times 𝑩_0\right)e^{i\left(\omega t𝒌𝒙\right)}$$
(2)
which includes the magnetic field $`𝑩=c𝒌\times 𝑬/\omega `$ of the plane wave and the associated force due to the uniform motion with velocity $`𝒗_e`$. The equation of propagation of the electromagnetic wave in the moving medium has an elementary solution when the velocity of the light and of the medium are parallel. If we refer to a cartesian frame $`𝒌=k\widehat{𝒆_x}`$, $`𝒗_e=v_e\widehat{𝒆_x}`$ we have in this case $`𝑬=E_0e^{i\left(\omega tkx\right)}\widehat{𝒆}_y`$, $`𝑩=c\frac{k}{\omega }E_0e^{i\left(\omega tkx\right)}\widehat{𝒆}_z`$ for the electromagnetic field and
$`𝝃={\displaystyle \frac{e}{m}}{\displaystyle \frac{𝑬_0\left(1\frac{kv_e}{\omega }\right)}{\omega _0^2\left(\omega kv_e\right)^2}}e^{i\left(\omega tkx\right)}`$ (3)
for the displacement vector parallel to the $`y`$ axis. The relativistic extension of this model can be obtained directly putting $`v_e=0`$ in Eq. 2 or 3 and using a Lorentz transformation between the moving frame and the laboratory one. We deduce the displacement
$`𝝃={\displaystyle \frac{e}{m}}\gamma _e{\displaystyle \frac{𝑬_0\left(1\frac{kv_e}{\omega }\right)}{\omega _0^2\gamma _e^2\left(\omega kv_e\right)^2}}e^{i\left(\omega tkx\right)}`$ (4)
where $`\gamma _e=1/\sqrt{\left(1v_e^2/c^2\right)}`$. We could alternatively obtain the same result considering the generalization of the Newton dynamics i. e. by doing the substitutions $`mm\gamma _e`$ and $`\omega _0\omega _0\gamma _e^1`$ in Eq. 2. The dispersion relation is then completely fixed by the Maxwell equation $`\frac{^2}{x^2}E\frac{1}{c^2}\frac{^2}{t^2}E=\frac{4\pi }{c^2}\frac{}{t}J`$, where the current density $`J`$ is given by the formula $`J=eN\left(_t+v_e_x\right)\xi `$ depending on the local number of atoms per unit volume $`N`$ supposed to be constant. Using $`J`$ and Eq. 3 or 4 we obtain a dispersion relation $`k^2=n^2\left(\omega \right)\omega ^2/c^2`$ where the effective refractive index $`n\left(\omega \right)`$ depends on the angular frequency $`\omega `$ and on the velocity $`v_e`$. The more general index obtained using Eq. 4 is defined by the implicit relation
$$n^2\left(\omega \right)=1+\gamma _e^2[n_0^2\left(\omega ^{}\right)1][1\frac{n\left(\omega \right)v_e}{c}]^2.$$
(5)
Here $`\omega ^{}=\omega \left(1\frac{nv_e}{c}\right)\gamma _e`$, and $`n_0^2\left(\omega \right)=1+4\pi N_0e^2/\left(\omega _0^2\omega ^2\right)/m`$ is the classical Lorentz index (also called Drude index) which contains the local proper density which is defined in the frame where the medium is immobile by $`N_0=N\gamma _e^1`$. These relativistic equations imply directly the correct relativistic formula for the velocity of light in the medium: Writing $`n_0^21=(n^21)(1v_e^2/c^2)/(1nv_e/c)^2=(nv_e/c)^2/(1nv_e/c)^21`$ we deduce
$$\frac{c}{n_0}=\frac{c/nv_e}{1\frac{v_e}{cn}}.$$
(6)
which can be easily transformed into
$$v=\frac{c}{n}=\frac{c/n_0\left(\omega ^{}\right)+v_e}{1+\frac{v_e}{cn_0\left(\omega ^{}\right)}}.$$
(7)
It can be added that by combining these expressions we deduce the explicit formula
$$n^2\left(\omega \right)=1+\gamma _e^2\frac{[n_0^2\left(\omega ^{}\right)1]}{[1+\frac{n_0\left(\omega ^{}\right)v_e}{c}]^2}.$$
(8)
The non relativistic case can be obtained directly from Eq. 3 or by writting $`\gamma _e=1`$ in Eqs. 5,7. This limit
$$v=\frac{c}{n}\frac{c}{n_0}+v_e[1\frac{1}{n_0^2}+\omega \frac{d\mathrm{ln}n_0}{d\omega }]+O\left(\frac{v_e^2}{c^2}\right)$$
(9)
is the Fresnel-Fizeau formula corrected by a “frequency-dispersion” term due to LorentzLorentz . For our purpose, it is important to note that in the non-relativistic limit of Eq. 5 we can always write the equality
$$\frac{c}{n}=\frac{cv_e}{n^{}}+v_e$$
(10)
where $`n^{}=n\left(1v_e/c\right)/\left(1nv_e/c\right)`$ is the index of refraction defined relatively to the moving medium. We then can see directly that the association of Maxwell’s equation with Newtonian dynamics implies a modification of the intuitive assumption “$`c/n_0+v_e`$” used in the old theory of emission. In fact, the problem can be understood in the Newtonian mechanics using the absolute time $`t=t^{}`$ and the transformation $`x=x^{}+v_et^{}`$. In the laboratory frame the speed of light, which in vacuum is $`c`$, becomes $`c/n_0`$ in a medium at rest. In the moving frame the speed of light in vacuum is now $`cv_e`$Jackson . However, due to invariance of acceleration and resultant force in a galilean transformation we can interpret the presence of the magnetic term in Eq. 2 as a correction to the electric field in the moving frame. This effective electric field affecting the oscillator in the moving frame is then transformed into $`E\left(1nv_e/c\right)`$. It is this term which essentially implies the existence of the effective optical index $`n^{}n_0`$ and the light speed $`\left(cv_e\right)/n^{}`$ in the moving frame. It can be observed that naturally Maxwell’s equations are not invariant in a Galilean transformation. The interpretation of $`E\left(1nv_e/c\right)`$ as an effective electric field is in the context of Newtonian dynamics only formal: This field is introduced as an analogy with the case $`v_e=0`$ only in order to show that $`n^{}`$ must be different from $`n_0`$.
## III Perturbation approach and optical theorem
The difficulty of the preceding model is that the Lorentz derivation does not clarify the meaning of the Fresnel-Fizeau phenomenon. Indeed we justify Eq. 1 using a microscopical model which is in perfect agreement with the principle of relativity. However we observe that at the limit $`v_ec`$ the use of the non relativistic dynamics of Newton (see Eq. 2) gives the same result. More precisely one can see from Eq. 2 that the introduction of the magnetic force $`e𝒗_e\times 𝑩/c`$ in addition to the electric force is already sufficient to account for the Fresnel-Fizeau effect and this even if the classical force formula $`𝐅=m\ddot{𝐱}(t)`$ is conserved. Since the electromagnetic force contains the ratio $`v_e/c`$ and originates from Maxwell’s equations this is already a term of relativistic nature (Einstein used indeed this fact to modify the dynamical laws of NewtonEinstein2 ). The derivation of Lorentz is then based on Newton as well as on Einstein dynamics. It is well know in counterpart that the Doppler-Fizeau effect, which includes the same factor $`1v_e/c`$, can be understood without introducing Einstein’s relativity. Indeed this effect is just a consequence of the invariance of the phase associated with a plane wave when we apply a Galilean transformation (see Jackson , Chap. 11) as well as a Lorentz transformation. We must then analyze further in detail the interaction of a plane wave with a moving dipole in order to see if the Fresnel phenomenon can be understood independently of the specific Lorentz dynamics.
We consider in this part a different calculation based on a perturbation method and inspired by a derivation of the optical theorem by FeynmanFeynman ; Hulst . Consider a thin slab of thickness $`L`$ perpendicular to the $`x`$ axis. Let this slab move along the positive $`x`$ direction with the constant velocity $`v_e\widehat{𝒆}_x`$. Let in addition $`𝐄_0e^{i\omega \left(tx/c\right)}`$ be the incident electric field of a plane wave which pursues the moving slab (see Fig. 1). Therefore, the electric field after the slab can be formally written as
$$𝐄_{\text{after}}=𝐄_0e^{i\omega \left(t\delta t_{v_e}x/c\right)},$$
(11)
where $`\delta t_{v_e}`$ appear as a retardation time produced by the interaction of light with the slab and where all reflections are neglected ($`𝐄_{\text{after}}=𝐄_{\text{before}}`$). For a “motionless” slab (i. e., the case considered by Feynman) we can write the travel time of the light through the slab as $`\mathrm{\Delta }\tau _0=L_0/c+\delta t_0=n_0L_0/c`$ and therefore $`\delta t_0=\left(n_01\right)L_0/c`$ where $`L_0`$ defines the proper length of the slab in the frame where it is at rest. For the general case of a moving slab of reduced length $`L=L_0\gamma _e^1`$ we find for the travel time:
$$\mathrm{\Delta }\tau _e=\frac{\left(L+v_e\mathrm{\Delta }\tau _e\right)}{c}+\delta t_{v_e}=n\frac{\left(L+v_e\mathrm{\Delta }\tau _e\right)}{c}$$
(12)
and therefore the perturbation time is
$`\delta t_{v_e}={\displaystyle \frac{\left(n1\right)L}{\left(cnv_e\right)}}.`$ (13)
We can obtain this result more rigorously by using Maxwell’s boundary conditions at the two moving interfaces separating the matter of the slab and the air (see Appendix A).
In order to evaluate the diffracted field which is $`𝐄_{\text{after}}𝐄_{\text{before}}`$ we can limit our calculation to a first order approximation. Thereby, each dipole of the Lorentz model as discussed above can be considered as being excited directly by the incident electromagnetic wave and where we can neglect all phenomena implying multiple interactions between light and matter. In this limit Eq. 11 reduces to
$`𝐄_{\text{after}}=𝐄_0e^{i\omega \left(tx/c\right)}e^{+i\omega \left(n1\right)\frac{L}{cnv_e}}`$
$`𝐄_0e^{i\omega \left(tx/c\right)}\left(1+i\omega \left(n1\right){\displaystyle \frac{L}{cnv_e}}\right).`$ (14)
If the distance between the slab and an observation point is much larger than $`L`$ we can consider the slab as a 2D continuous distribution of radiating point dipoles. The vector potential $`𝐀_{\text{rad}}`$ radiated by a relativistically moving point charge $`e`$ is in according with the Lienard-Wiechert’s formula given by Jackson :
$`𝐀_{\text{rad}}(𝐱,t)=e{\displaystyle \frac{𝐯/c}{\left(1\widehat{𝐑}\beta \right)R}}|_{ret}.`$ (15)
Here $`R=𝐱𝐱_0\left(t\right)`$ is the distance separating the observation point $`𝐱`$ (denoted by P) and the point charge position located at $`𝐱_0\left(t\right)`$ at the time $`t`$; additionally $`𝒗\left(t\right)=\dot{𝒙}_0\left(t\right)`$ is the velocity of the point charge and $`\widehat{𝑹}\left(t\right)`$ is the unit vector $`\left(𝐱𝐱_0\left(t\right)\right)/R\left(t\right)`$. In this formula, in agreement with causality, all point charge variables are evaluated at the retarded time $`t_{ret}=tR\left(t_{ret}\right)/c`$.
In the present case the motion of the point charge can be decomposed into a uniform longitudinal component $`𝐯_et`$ oriented along the positive $`x`$ direction and into a transversal oscillating part $`\xi \left(t\right)=\xi _0e^{i\omega \left(1\frac{v_e}{c}\right)t}`$ obeying the condition $`\dot{\xi }\left(t\right)/c1`$. Owing to this condition we can identify $`\left(1\widehat{𝐑}𝐯/c\right)`$ with $`\left(1\widehat{𝐑}𝐯_𝐞/c\right)`$. Consequently, in the far-field the contribution of the electron uniform velocity is cancelled by the similar but opposite contribution associated with the nucleus of the atomic dipole: Only the vibrating contribution of the electron survives at a long distance from the diffraction source. If we add the contribution of each dipole of the slab acting on the observation point P at the time $`t`$ we obtain then the total diffracted vector potential $`𝐀_{\text{diff}}`$ produced by the moving medium:
$`𝐀_{\text{diff}}(𝐱,t)2\pi \gamma _eN_0Li\omega {\displaystyle \frac{e}{c}}\left(1\beta _e\right)\xi _0`$
$`{\displaystyle _0^+\mathrm{}}\rho 𝑑\rho {\displaystyle \frac{e^{i\omega \left(1\beta _e\right)\left(tR\left(t_{ret}\right)/c\right)}}{\left(1\widehat{𝑹}\left(t_{ret}\right)𝒗_e/c\right)R\left(t_{ret}\right)}},`$ (16)
Here $`\rho `$ is the radial coordinate in a cylindrical coordinate system using the direction $`x`$ as a revolution axis, and the quantity $`\gamma _eN_0L2\pi \rho d\rho `$ is the number of dipoles contained in the cylindrical volume of length $`L`$ and of radius varying between $`\rho `$ and $`\rho +d\rho `$ if we consider a local dipole density given by $`\gamma _eN_0`$. In this formula the retarded distance $`R\left(t_{ret}\right)`$ is a function of $`\rho `$ and we have (see the textbook of Jackson Jackson )
$$R\left(t_{ret}\right)=\gamma _e^1\left(1\widehat{𝐑}\left(t_{ret}\right)𝐯_e/c\right)^1\sqrt{\rho ^2+\gamma _e^2\left(xv_et\right)^2}.$$
(17)
This expression shows that the minimum $`R_{min}`$ is obtained for a point charge on the $`x`$ axis, and that:
$$R_{min}=\left(xv_et\right)/\left(1\beta _e\right).$$
(18)
In order to evaluate the integral in Eq. 16 we must use in addition the following relation (see Appendix C):
$$R\left(t_{ret}\right)=\gamma _e^2\beta _e\left(xv_et\right)+\gamma _e\sqrt{\rho ^2+\gamma _e^2\left(xv_et\right)^2}$$
(19)
Hence, we obtain the following integral :
$`𝐀_{\text{diff}}(𝐱,t)2\pi i\omega \gamma _e^2N_0L{\displaystyle \frac{e}{c}}\left(1\beta _e\right)\xi _0`$
$`e^{i\omega \left(1\beta _e\right)\left(t\gamma _e^2\beta _e\left(xv_et\right)\right)}`$
$`{\displaystyle _0^+\mathrm{}}\rho d\rho {\displaystyle \frac{e^{i\frac{\omega }{c}(1\beta _e)\gamma _e\sqrt{\rho ^2+\gamma _e^2\left(xv_et\right)^2}\}}}{\sqrt{\rho ^2+\gamma _e^2\left(xv_et\right)^2}}},`$ (20)
where we have used the relations Eq. 17, Eq. 19 in the denominator and in the exponential argument of the right hand side of Eq. 16, respectively. The diffracted field is therefore directly calculable by using the variable $`u=\sqrt{\rho ^2+\gamma _e^2\left(xv_et\right)^2}`$. We obtain the result
$`𝐀_{\text{diff}}2\pi \gamma _eLN_0{\displaystyle \frac{e}{c}}\xi _0e^{i\omega \left(tx/c\right)}.`$ (21)
The total diffracted electric field $`𝐄_{\text{diff}}`$ is obtained using Maxwell’s formula $`𝐄=\left(1/c\right)_t𝐀`$, which gives:
$`𝐄_{\text{diff}}2\pi i\gamma _eLN_0\omega {\displaystyle \frac{e}{c}}\xi _0e^{i\omega \left(tx/c\right)}.`$ (22)
The final result is given substituting Eq. 4 in Eq. 22 and implies by comparison with Eq. 14
$`n1+2\pi N_0\gamma _e^2{\displaystyle \frac{e^2}{m}}{\displaystyle \frac{\left(1\frac{v_e}{c}\right)^2}{\omega _0^2\gamma _e^2\omega ^2\left(1\frac{v_e}{c}\right)^2}}.`$ (23)
This equation constitutes the explicit limit $`N_00`$ of Eq. 5 and implies the correct velocity formula Eq. 7 when we neglect terms of $`\text{O}[N_0^2]`$. It can again be observed that the present calculation can be reproduced in the non relativistic case by neglecting all terms of order $`(v_e/c)^2`$.
## IV Physical meaning and discussion
The central fact in this reasoning is “the travel condition” given by Eqs. 12,13. Indeed, of the same order in power of $`N_0`$ we can deduce the relation
$$\delta t_{v_e}=\gamma _e\delta t_0\left(1v_e/c\right)$$
(24)
and consequently the condition Eq. 12 reads
$$\mathrm{\Delta }\tau _e=\frac{L+v_e\left(\mathrm{\Delta }\tau _e\delta t_0\gamma _e\right)}{c}+\delta t_0\gamma _e=n\frac{\left(L+v_e\mathrm{\Delta }\tau _e\right)}{c}.$$
(25)
If we call $`\delta t_0`$ the time during which the energy contained in a plane of light moving in the positive x direction is absorbed by the slab at rest in the laboratory, $`\delta t_0\gamma _e`$ is evidently the enlarged time for the moving case. During the period where this plane of light is absorbed by the slab its energy moves at the velocity $`v_e`$. This fact can be directly deduced of the energy and momentum conservation laws. Indeed, let $`M\gamma _ev_e`$ be the momentum of the slab of mass $`M`$ before the collision and $`ϵ`$ the energy of the plane of light, then during the interaction the slab is in a excited state and its energy is now $`E^{}=ϵ+M\gamma _ec^2`$ and its momentum $`P^{}=ϵ/c+M\gamma _ev_e`$. The velocity of the excited slab is defined by $`w=c^2P^{}/E^{}`$ and we can see that in the approximation $`M\mathrm{}`$ used here $`wv_e`$ (we neglect the recoil of the slab). During $`\delta t_0\gamma _e`$ the slab moves along a path length equal to $`v_e\delta t_0\gamma _e`$ and thus the travel condition of the plane of energy in the moving slab can be written
$$c\left(\mathrm{\Delta }\tau _e\gamma _e\delta t_0\right)=L+v_e\left(\mathrm{\Delta }\tau _e\gamma _e\delta t_0\right),$$
(26)
which is an other form for Eq. 25. Now eliminating directly $`\mathrm{\Delta }\tau _e`$ in Eq. 26 give us the velocity $`v`$ of the wave:
$$v=v_e+\frac{cv_e}{1+\frac{c\delta t_0}{L}\left(1\beta _e\right)\gamma _e},$$
(27)
i. e.
$$v=v_e+\frac{cv_e}{1+\left(n_01\right)\left(1\beta _e\right)\gamma _e^2},$$
(28)
which depends on the optical index $`n_0=1+c\delta t_0/L_0`$. After straightforward manipulations this formula becomes
$$v=\frac{c/n_0+v_e}{1+\frac{v_e}{cn_0}}$$
(29)
which is the Einstein formula containing the Fresnel result as the limit behavior for small $`v_e`$.
It can be observed that this reasoning is even more natural if we think in term of particles. A photon moving along the axis $`x`$ and pursuing an atom moving at the velocity $`v_e`$ constitutes a good analogy to understand the Fresnel phenomenon. This analogy is evidently not limited to the special case of the plane wave $`e^{i\omega \left(tx/c\right)}`$. If for example we consider a small wave packet which before the interaction with the slab has the form
$$E_{\text{before}}(x,t)=_{\mathrm{\Delta }\omega }𝑑\omega a_\omega e^{i\left(kx\omega t\right)},$$
(30)
where $`\mathrm{\Delta }\omega `$ is a small interval centered on $`\omega _m`$, then after the interaction we must have:
$$E_{\text{after}}(x,t)=_{\mathrm{\Delta }\omega }𝑑\omega a_\omega e^{i(kx\omega [t\delta t\left(\omega \right))},$$
(31)
where $`\delta t\left(\omega \right)`$ is given by Eq. 13. After some manipulation we can write these two wave packets in the usual approximative form:
$`E_{\text{before}}e^{i\left(k_mx\omega _mt\right)}{\displaystyle _{\mathrm{\Delta }\omega }}𝑑\omega a_\omega e^{i\left(\omega \omega _m\right)[tk/\omega _mx]}`$
$`=e^{i\left(k_mx\omega _mt\right)}F\left(tx/v_g\right)`$
$`E_{\text{after}}e^{i\left(k_mx\omega _m[t\delta t\left(\omega _m\right)]\right)}F\left(tx/v_g\delta t_g\right).`$ (32)
Here, $`v_g=\omega _m/k_m=c`$ is the group velocity of the pulse in vacuum and $`\delta t_g=\left(\omega _m\delta t\left(\omega _m\right)\right)/\omega _m`$ is the perturbation time associated with this group motion. This equation for $`F`$ possesses the same form as Eq. 11 and then the same analogy which implies Eq. 25 is possible. This can be seen from the fact that we have
$$\delta t\left(\omega \right)=\gamma _e\delta t_0(\omega ^{})\left(1v_e/c\right)$$
(33)
with $`\omega ^{}=\gamma _e\omega (1v_e/c)`$. We deduce indeed
$$\delta t_g=\gamma _e\delta t_{0g}(\omega _m^{})\left(1v_e/c\right),$$
(34)
where we have $`\delta t_{0g}(\omega _m^{})=\left(\omega _m\delta t_0\left(\omega _m^{}\right)\right)/\omega _m`$ i. e. $`\delta t_{0g}(\omega _m^{})=\left(\omega _m^{}\delta t_0\left(\omega _m^{}\right)\right)/\omega _m^{}`$. Since Eq. 33 and Eq. 34 have the same form the Fresnel law must be true for the group velocity.
It is important to remark that all this reasoning conserves its validity if we put $`\gamma _e=1`$ and if we think only in the context of Newtonian dynamics. Since the reasoning with the travel time does not explicitly use the structure of the medium involved (and no more the magnetic force $`e𝒗_e\times 𝑩/c`$) it must be very general and applicable in other topics of physics concerning for example elasticity or sound.
Consider as an illustration the case of a cylindrical wave guide with revolution axis $`x`$ and of constant length $`L`$ pursued by a wave packet of sound. We suppose that the scalar wave $`\psi `$ obeys the equation $`[c^2^2/𝐫^2^2/t^2]\psi =0`$ where $`c`$ is the constant sound velocity. The propagative modes in the cylinder considered at rest in the laboratory are characterized by the classical dispersion relation
$$\omega ^2/c^2=\gamma _{n,m}^2+k_x^2$$
(35)
where the cut off wave vector $`\gamma _{n,m}`$ depend only of the two “quantum” numbers $`n,m`$ and of the cross section area $`A`$ of the guide ($`\gamma ^21/A`$). The group velocity $`\omega /k_x`$ of the wave in the guide is defined by $`v_g=\left(c^2/\omega \right)\sqrt{\omega ^2/c^2\gamma ^2}c[1\frac{1}{2}c^2\gamma ^2/\omega ^2]`$ and the travel time $`\mathrm{\Delta }\tau `$ by $`L/v_gL[1+\frac{1}{2}c^2\gamma ^2/\omega ^2]/c`$ which implies $`\delta t_0=\frac{1}{2}Lc\gamma ^2/\omega ^2`$. In the moving case where the cylinder possesses the velocity $`v_e`$ we can directly obtain the condition given by Eq. 25 (with $`\gamma _e=1`$) and then we can deduce the group velocity of the sound in the guide with the formula
$`v=v_e+{\displaystyle \frac{cv_e}{1+\frac{c\delta t_0}{L}\left(1\beta _e\right)}}.`$ (36)
This last equation give us the Fresnel result if we put the effective sound index $`n_0=1+c\delta t_0/L`$. We can control the self consistency of this calculation by observing that the dispersion relation Eq. 35 allows the definition of a phase index $`n_{\text{phase}}=ck/\omega 1c^2\gamma _{n,m}^2/(2\omega ^2)`$ which is equivalent to Eq. 25 when $`\omega _0=0`$ and $`2\pi N_0e^2/m=c^2\gamma _{n,m}^2/2`$. This reveals a perfect analogy between the sound wave propagating in a moving cylinder and the light wave propagating in a moving slab. It is then not surprising that the Fresnel result is correct in the two cases.
The principal limitation of our deduction is contained in the assumption expressed above for the slab example: $`𝐄_{\text{after}}𝐄_{\text{before}}`$ i. e. the condition of no reflection supposing the perturbation on the motion of the wave to be small. Nevertheless, the principal origin of the Fresnel effect is justified in our scheme without the use of the Einstein relativity principle.
We can naturally ask if the simple analogy proposed can not be extended to a dense medium i. e. without the approximation of a weak density $`N_0`$ or of a low reflectivity. In order to see that it is indeed true we return to the electromagnetic theory and we suppose an infinite moving medium like the one considered in the second section. In the rest frame of the medium we can define a slab of length $`L_0`$. The unique difference with the section 3 is that now this slab is not bounded by two interfaces separating the atoms from the vacuum but is surrounded by a continuous medium having the same properties and moving at the same velocity $`v_e`$. In the laboratory frame the length of the moving slab is $`L=L_0\gamma _e^1`$. We can write the time $`\mathrm{\Delta }\tau _e`$ taken by a signal like a wave packet, a wave front or a plane of constant phase to travel through the moving slab:
$$c\mathrm{\Delta }\tau _e=n\left(L+v_e\mathrm{\Delta }\tau _e\right).$$
(37)
The optical index $`n`$ can be the one defined in section 2 for the case of the Drude model but the result is very general. We can now introduce a time $`\delta t_0`$ such that Eq. 25, and consequently Eq. 27, are true *by definition*. We conclude that this last equation Eq. 27 is equivalent to the relativistic Eq. 29 if, and only if, we define the time $`\delta t_0`$ by the formula
$$\delta t_0=\left(n_01\right)L_0/c.$$
(38)
In other terms we can always use the analogy with a photon pursuing an atom since the general formula Eq. 29 is true whatever the microscopic and Electrodynamics model considered. In this model - based on a retardation effect- the absorbtion time $`\delta t_0`$ is always given by Eq. 38.
This opens new perspectives when we consider the problem of a sound wave propagating in an effective moving medium. Indeed there are several situations where we can develop a deep analogy between the propagation of sound and the propagation of light. This implies that the conclusions obtained for the Fresnel effect for light must to a large part be valid for sound as well. This is in particular true if we consider an effective meta material like the one that is going to be described now:
We consider a system of mirrors as represented in Fig. 2A, at rest in the laboratory. A beam of light propagates along the zigzag trajectory $`A_0,B_0,A_1,\mathrm{},A_n,B_n,\mathrm{}`$. The length $`A_nB_n`$ is given by $`\sqrt{(L_0^2+D^2)}`$ where the distance $`L_0`$ and $`D`$ are represented on the figure. The time $`\mathrm{\Delta }\tau _0`$ spent by a particle of light to move along $`A_nB_n`$ is then $`\sqrt{(L_0^2+D^2)}/c`$. We can equivalently define an effective optical index $`n_0`$ such that we have
$$\frac{(L_0^2+D^2)}{c^2}=\mathrm{\Delta }\tau _0^2=\frac{L_0^2n_0^2}{c^2}.$$
(39)
This implies
$$n_0^2=1+\frac{D_0^2}{L_0^2}.$$
(40)
We consider now the same problem for a system of mirrors moving with the velocity $`v_e`$. In order to be consistent with relativity we introduce the reduced length $`L=L_0\gamma _e^1`$. The beam propagating along the path $`A_0,B_0,A_1,\mathrm{},A_n,B_n,\mathrm{}`$ must pursue the set of mirrors. We then define the travel time $`\mathrm{\Delta }\tau _e`$ along an elementary path $`A_nB_n`$ by
$$\frac{((L+v_e\mathrm{\Delta }\tau _e)^2+D^2)}{c^2}=\mathrm{\Delta }\tau _e^2=\frac{((L+v_e\mathrm{\Delta }\tau _e)^2n^2}{c^2},$$
(41)
where $`n`$ is the effective optical index for the moving medium. From this equation we deduce first $`\mathrm{\Delta }\tau _e=(L/c)n/(1v_en/c)`$ and then
$$n^21=(n_0^21)\frac{(1\frac{v_en}{c})^2}{(1(\frac{v_e}{c})^2)}$$
(42)
which finally give us the formula
$$\frac{c}{n}=\frac{c/n_0+v_e}{1+\frac{v_e}{cn_0}}$$
(43)
We can again justify the Fresnel formula at the limit $`v_e/c1`$.
The simplicity of this model is such that it does not depend on the physical properties of atoms, electrons and photons but only on geometrical parameters. Clearly we can make the same reasoning for a sound wave by putting $`\gamma _e=1`$. This still gives us the Fresnel formula when we neglect terms equal or smaller than $`O\left(v_e^2/c^2\right)`$. In addition this model allows us to conclude that the essential element justifying the Fresnel-Fizeau result is the emergence of a delay time – a retardation effect– when we consider the propagation of the signal at a microscopic or internal level. The index $`n`$ which characterizes the macroscopic or external approach is then just a way to define an effective velocity without looking for a causal explanation of the retardation.
The essential message of our analysis is that by taking explicitly into account the physical origin of the delay we can justify the essence of the Fresnel-Fizeau effect in a non relativistic way. The Fresnel-Fizeau effect is then a very general phenomenon. It is a consequence of the conservation of energy and momentum and of the constant value of the wave velocity in vacuum or in the considered medium. The so called travel condition (Eq. 26) which is a combination of these two points can be compared to the usual demonstration for the Doppler effect. In these two cases of light pulses pursuing a moving particle the perturbation time $`\delta t_{v_e}\delta t_0\left(1v_e/c\right)`$ is a manifestation of the Doppler phenomenon. It should be emphasized that the analogy between sound and electromagnetic waves discussed in this article could be compared to the similarities between sound wave and gravitational waves discussed in particular by Unruh. On this subject and some connected discussions concerning the acoustic Aharonov-Bohm effect (that is related to the optical Aharonov-Bohm effect that follows from the Fizeau effect) the reader should consult Unruh1 ; Unruh2 .
## V Summary
We have obtain the Fresnel-Fizeau formula using a perturbation method based on the optical theorem and in a more general way by considering the physical origin of the refractive index. The modification of the speed of light in the medium appears then as a result of a retardation effect due to the duration of the interaction or absorbtion of light by the medium, and the Fresnel-Fizeau effect, as a direct consequence of the medium’s flight in front of the light. These facts rely on the same origin as the Doppler-Fizeau effect. We finally have shown that it is not correct to assume, as frequently done in the past, that a coherent and “Newtonian interpretation” of these phenomena would be impossible. On the contrary, the results do not invalidate the derivation of the Fresnel-Fizeau effect based on the principle of relativity but clarify it. We observe indeed than all reasoning is in perfect agreement with the principle of relativity. We must emphasize that even if the Fizeau/Fresnel effect is conceptually divorced from relativity it strongly motivated Einstein’s work (more even than the Michelson and Morley result). The fact that the Fizeau as well as the Michelson-Morley experiment can be justified so easily with special relativity clearly show the advantages of Einstein’s principle to obtain quickly the correct results. Nevertheless, if we look from a dynamical point of view, as it is the case here, this principle plays a role only for effects of order $`𝐯_e^2/c^2`$ which however are not necessary to justify the Fresnel formula.
###### Acknowledgements.
The author acknowledges S. Huant, M. Arndt, J. Krenn, D. Jankowska as well as the two anonymous referees for interesting and fruitful discussions during the redaction process.
## Appendix A
Maxwell’s equations impose the continuity of the electric field on each interface of the slab. More precisely these boundary relations impose: $`𝐄_{\text{medium A}}|_S=𝐄_{\text{medium B}}|_S`$ where $`S`$ is one of the two moving interfaces separating vacuum and matter. Hence we obtain an equality condition between the two phases $`\varphi _{\text{medium A}}`$ and $`\varphi _{\text{medium B}}`$ valid for all times at the interface. Let $`\mathrm{\Phi }_1=i\omega \left(tx/c\right)`$ be the phase of the plane wave before the slab. In a similar way let $`\mathrm{\Phi }_2=i\omega _2\left(tn_{\omega _2,\beta _e}x/c\delta _2\right)`$ and $`\mathrm{\Phi }_3=i\omega _3\left(tx/c\delta _3\right)`$ be the phases in the slab and in vacuum after traversing the slab, respectively. In these expressions there appear two retardation constants, $`\delta _{2,3}`$ and the optical index of the slab. On the first interface denoted by (I-II) we have $`x=c\beta _et`$ and consequently
$`\omega \left(1\beta _e\right)t=\omega _2\left(1n(\omega ,\beta _e)\beta _e\right)`$
$`\left(t{\displaystyle \frac{\delta _2}{1n(\omega _2,\beta _e)\beta _e}}\right),`$ (44)
which is valid for each time and possesses the unique solution:
$`\omega _2=\omega {\displaystyle \frac{1\beta _e}{1n(\omega ,\beta _e)\beta _e}}`$ $`,\delta _2=0.`$ (45)
Considering the second interface (II-III) in a similar way we obtain the following conditions
$`\omega _3=\omega _2{\displaystyle \frac{1n(\omega ,\beta _e)\beta _e}{1\beta _e}}=\omega `$
$`\delta _3=L{\displaystyle \frac{n(\omega ,\beta _e)1}{cn(\omega ,\beta _e)c\beta _e}}`$ (46)
where the $`2^{nd}`$ equality is Eq. 13.
## Appendix B
Using geometrical considerations (see Fig. 3) we can deduce the relation
$$R\left(t_{ret}\right)^2=\rho ^2+\left(\beta _eR_{ret}+xv_et\right)^2,$$
(47)
which is equivalent after manipulations to the other:
$$\rho ^2+\gamma _e^2\left(xv_et\right)^2=\left(1\beta _e^2\right)\left(R_{ret}\beta _e\gamma _e^2\left(xv_et\right)\right)^2.$$
(48)
We can in a second step rewrite this equality as follows:
$$R_{ret}=\gamma _e^2\beta _e\left(xv_et\right)+\gamma _e\sqrt{\rho ^2+\gamma _e^2\left(xv_et\right)^2}$$
(49)
which is Eq. 19.
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# Norms of Schur Multipliers
## 1. Background
If $`A=[a_{ij}]_{i,jS}`$ is a finite or infinite matrix, the Schur (a.k.a. Hadamard) multiplier is the operator $`S_A`$ on $`(l^2(S))`$ that acts on an operator $`T=[t_{ij}]`$ by pointwise multiplication: $`S_A(T)=[a_{ij}t_{ij}]`$. To distinguish from the norm on bounded operators, we will write $`A_m`$ for the norm of a Schur multiplier. In general it is very difficult to compute the norm of a Schur multiplier. Nevertheless, much is known in a theoretical sense about the norm. In this section, we will quickly review some of the most important results.
The following classical result owes most credit to Grothendieck. For a proof, see the books by Pisier \[17, Theorem 5.1\] and Paulsen \[14, Theorem 8.7\].
###### Theorem 1.1.
For $`X`$ an arbitrary set, let $`S=[s_{ij}]`$ be an $`|X|\times |X|`$ matrix with bounded entries considered as a Schur multiplier on $`(l^2(X))`$. Then the following are equivalent:
1. $`S_m1`$.
2. $`S_{cb}1`$.
3. There are contractions $`V`$ and $`W`$ from $`l^2(X)`$ to $`l^2(X)l^2(Y)`$ such that $`S(A)=W^{}(AI)V`$.
4. There are unit vectors $`x_i`$ and $`y_j`$ in $`l^2(Y)`$ so that $`s_{ij}=x_i^{}y_j`$.
5. $`\gamma _2(S)1`$ where $`\gamma _2(S)=inf_{S=AB}A_{2,\mathrm{}}B_{1,2}`$.
6. There are $`|X|\times |X|`$ matrices $`A=[a_{ij}]`$ and $`B=[b_{ij}]`$ with $`a_{ii}=b_{ii}=1`$ so that $`\left[\begin{array}{cc}A& S\\ S^{}& B\end{array}\right]`$ is positive semidefinite.
Recall that the complete bound norm of $`S`$ is the norm of the inflation of $`S`$ acting on operators with operator entries. The most elegant proof of (1) implies (2) is due to Smith . The converse is trivial. The equivalence of (2) and ($`2^{}`$) is Wittstock’s Theorem for representing completely bounded maps.
The equivalence of (1), (3) and (4) is due to Grothendieck. (3) follows from ($`2^{}`$) by taking $`y_j=(E_{1j}I)Ve_j`$ and $`x_i=(E_{1i}I)We_i`$. Conversely, (3) implies ($`2^{}`$) by taking $`Ve_j=e_jy_j`$ and $`We_i=e_ix_i`$. This condition was rediscovered by Haagerup, and became well-known as his observation. So we shall refer to these as the Grothendieck–Haagerup vectors for the Schur multiplier.
The $`\gamma _2`$ norm is the optimal factorization through Hilbert space of $`S`$ considered as a map from $`l^1`$ to $`l^{\mathrm{}}`$. The norm $`A_{2,\mathrm{}}`$ is the maximum of the 2-norm of the rows; while $`B_{1,2}`$ is the maximum of the 2-norm of the columns. Thus (3) implies (4) follows from $`A=_ie_ix_i^{}`$ and $`B=_jy_je_j^{}`$. And this implication is reversible.
The equivalence of (5) is due to Paulsen, Power and Smith . This follows from (3) by taking $`a_{ij}=x_i^{}x_j`$ and $`b_{ij}=y_i^{}y_j`$. Conversely, assume first that $`X`$ is finite. Then the positive matrix $`P`$ decomposes as a sum of positive rank one matrices, and thus have the form $`[\overline{z}_iz_j]`$ which can be seen to be a scalar version of (3). Indeed it is completely positive. Hence the sum $`S_P`$ is also a completely positive Schur multiplier. Consequently $`S_P_{cb}=S_P(I)=\mathrm{max}\{a_{ii},b_{ii}\}=1`$. So (2) holds. The case of general $`X`$ is a routine limit argument.
The $`\gamma _2`$ norm is equivalent to the norm in the Haagerup tensor product $`\mathrm{}^{\mathrm{}}(X)_h\mathrm{}^{\mathrm{}}(X)`$, where we identify an elementary tensor $`ab`$ with the matrix $`[a_ib_j]`$. The Haagerup norm of a tensor $`\tau `$ is given by taking the infimum over all representations $`\tau =_ka_kb_k`$ of
$$\underset{k}{}a_ka_k^{}^{1/2}\underset{k}{}b_k^{}b_k^{1/2}.$$
See \[14, Chapter 17\]. Of course, since $`\mathrm{}^{\mathrm{}}`$ is abelian, the order of the adjoints is irrelevant. One can see the equivalence by taking a factorization $`S=AB`$ from (4). Consider $`A`$ as a matrix with columns $`a_k\mathrm{}^{\mathrm{}}(X)`$ and $`B`$ as a matrix with rows $`b_k\mathrm{}^{\mathrm{}}(X)`$. Identify the product with the tensor $`_ka_kb_k`$. The norm $`_ka_ka_k^{}^{1/2}`$ can be seen to be $`A_{2,\mathrm{}}`$ and the norm of $`_kb_k^{}b_k^{1/2}`$ to be $`B_{1,2}`$.
Generally, it is difficult to compute the norm of a Schur multiplier. The exception occurs when the matrix $`S`$ is positive definite. Then it is a classical fact that $`S`$ is a completely positive map. Consequently, $`S_{cb}=S(I)=sup_{iX}s_{ii}`$.
Grothendieck proved another remarkable result about Schur multipliers. Recall that the projective tensor product $`\mathrm{}^{\mathrm{}}(X)\widehat{}\mathrm{}^{\mathrm{}}(X)`$ norms a tensor $`\tau `$ as the infimum over representations $`\tau =_ka_kb_k`$ of the quantity $`_ka_kb_k`$. It is a surprising fact that this norm and the $`\gamma _2`$ or Haagerup norm are equivalent. We will need this connection to understand the relevance of work of Varopoulos. For the moment, we state this result in a way that makes a stronger connection to Schur multipliers. An elementary tensor $`ab`$ yields a rank one matrix $`[a_ib_j]`$. Thus Grothendieck’s result is equivalent to:
###### Theorem 1.2 (Grothendieck).
The convex hull of the rank one Schur multipliers of norm one contains the ball of all Schur multipliers of norm at most $`K_G^1`$, where $`K_G`$ is a universal constant.
In terms of the projective tensor product norm for a tensor $`\tau `$ and the corresponding Schur multiplier $`S_\tau `$, this result says that
$$K_G^1\tau _{\mathrm{}^{\mathrm{}}(S)\widehat{}\mathrm{}^{\mathrm{}}(S)}S_\tau _m\tau _{\mathrm{}^{\mathrm{}}(S)\widehat{}\mathrm{}^{\mathrm{}}(S)}$$
The constant $`K_G`$ is not known exactly. In the complex case Haagerup showed that $`1.338<K_G<1.405`$; and in the real case Krivine obtained the range $`[1.676,1.783]`$ and conjectured the correct answer to be $`\frac{\pi }{2\mathrm{log}(1+\sqrt{2})}`$.
We turn to the results of Varopoulos and Pisier which relate to our work. The paper of Varopoulos is famous for showing that three commuting contractions need not satisfy the von Neumann inequality. Proofs of this, including the one in the appendix of Varopoulos’s paper, are generally constructive. But the argument in the main part of his paper instead establishes a result about $`\mathrm{}^{\mathrm{}}(X)\widehat{}\mathrm{}^{\mathrm{}}(X)`$. He does not establish precise information about constants. This result was extended and sharpened by Pisier, who casts it in the language of Schur multipliers, to deal with multipliers and lacunary sets on nonamenable groups.
Consider $`\{\pm 1\}^{X\times X}`$ to be the space of functions from $`X\times X`$ to $`\{1,1\}`$ with the product measure $`\mu `$ obtained from $`p(1)=p(1)=.5`$.
###### Theorem 1.3 (Varopoulos–Pisier).
Let $`S=[s_{ij}]`$. The following are equivalent.
1. For all $`\epsilon \{\pm 1\}^{X\times X}`$, $`[\epsilon _{ij}s_{ij}]_m<\mathrm{}`$.
2. For almost all $`\epsilon \{\pm 1\}^{X\times X}`$, $`[\epsilon _{ij}s_{ij}]_m<\mathrm{}`$.
3. $`S=A+B`$ and there is a constant $`M`$ so that
$$\underset{i}{sup}\underset{j}{}|a_{ij}|^2M^2\text{and}\underset{j}{sup}\underset{i}{}|b_{ij}|^2M^2.$$
4. There is a constant $`M`$ so that for every pair of finite subsets $`R,CX`$, $`_{iR,jC}|s_{ij}|^2M^2\mathrm{max}\{|R|,|C|\}`$.
Pisier shows that if the average Schur multiplier norm
$$[\epsilon _{ij}s_{ij}]_m𝑑\mu (\epsilon )1,$$
then one can take $`M=1`$ in (3). Our results are not quite so sharp, as we require a constant (Lemma 2.9) of approximately 1/4. The constant $`M`$ in the two conditions (3) and (4) are not the same. The correct relationship replaces $`\mathrm{max}\{|R|,|C|\}`$ by $`|R|+|C|`$ (see Lemma 2.7); but they are related within a constant. If $`M`$ is the bound in (3), it is not difficult to obtain a bound of $`2M`$ for (1) (see Corollary 2.6). Thus one obtains that the average Schur norm is within a factor of 2 of the maximum.
## 2. Schur Bounded Patterns
A pattern $`𝒫`$ is a subset of $`\times `$. An infinite matrix $`A=[a_{ij}]`$ is supported on $`𝒫`$ if $`\{(i,j):a_{ij}0\}`$ is contained in $`𝒫`$. We let $`𝒮(𝒫)`$ denote the set of Schur multipliers supported on $`𝒫`$ with matrix entries $`|s_{ij}|1`$.
More generally, we will also consider Schur multipliers dominated by a given infinite matrix $`A=[a_{ij}]`$ with nonnegative entries. Let $`𝒮(A)`$ denote the set of all Schur multipliers with matrix entries $`|s_{ij}|a_{ij}`$.
###### Definition 2.1.
Say that a pattern $`𝒫\times `$ is Schur bounded if every $`X𝒮(𝒫)`$ yields bounded Schur multiplier $`S_X`$. The Schur bound of $`𝒫`$ is defined as $`𝔰(𝒫):=sup_{X𝒮(𝒫)}X_m.`$ Similarly, for a matrix $`A`$ with nonnegative entries, define $`𝔰(𝒮(A))=sup_{X𝒮(A)}X_m`$; and say that $`𝒮(A)`$ is Schur bounded if this value is finite.
It is easy to see that if $`𝒮(𝒫)`$ is Schur bounded, then $`𝔰(𝒫)`$ is finite. Note that if $`A_𝒫`$ is the matrix with 1s on the entries of $`𝒫`$ and 0s elsewhere, then $`𝒮(A_𝒫)=𝒮(𝒫)`$. We will maintain a distinction because we will require integral decompositions when working with a pattern $`𝒫`$.
Certain patterns are easily seen to be Schur bounded and this is the key to our result. The following two definitions of row bounded for patterns and matrices are not parallel, as the row bound of $`A_𝒫`$ is actually the square root of the row bound of $`𝒫`$. Each definition seems natural for its context, so we content ourselves with this warning.
###### Definition 2.2.
A pattern is row bounded by $`k`$ if there are at most $`k`$ entries in each row; and row finite if it is row bounded by $`k`$ for some $`k`$. Similarly we define column bounded by $`k`$ and column finite.
A nonnegative matrix $`A=[a_{ij}]`$ is row bounded by $`L`$ if the rows of $`A`$ are bounded by $`L`$ in the $`l^2`$-norm: $`sup_{i1}_{j1}|a_{ij}|^2L^2<\mathrm{}`$. Similarly we define column bounded by $`L`$.
The main result of this section is:
###### Theorem 2.3.
For a pattern $`𝒫`$, the following are equivalent:
1. $`𝒫`$ is Schur bounded.
2. $`𝒫`$ is the union of a row finite set and a column finite set.
3. $`\underset{R,C\text{ finite}}{sup}{\displaystyle \frac{|𝒫(R\times C)|}{|R|+|C|}}<\mathrm{}`$.
Moreover, the optimal bound $`m`$ on the size of the row and column finite sets in $`(2)`$ coincides with the least integer dominating the supremum in $`(3)`$; and the Schur bound satisfies
$$\sqrt{m}/4𝔰(𝒫)2\sqrt{m}.$$
This theorem has a direct parallel for nonnegative matrices.
###### Theorem 2.4.
For a nonnegative infinite matrix $`A=[a_{ij}]`$, the following are equivalent:
1. $`𝒮(A)`$ is Schur bounded.
2. $`A=B+C`$ where $`B`$ is row bounded and $`C`$ is column bounded.
3. $`\underset{R,C\text{ finite}}{sup}{\displaystyle \frac{_{iR,jC}a_{ij}^2}{|R|+|C|}}<\mathrm{}`$.
Moreover, the optimal bound $`M`$ on the row and column bounds in $`(2)`$ coincides with the square root of the supremum $`M^2`$ in $`(3)`$; and the Schur bound satisfies
$$M/4𝔰(𝒫)2M.$$
###### Lemma 2.5.
If $`𝒫`$ is row $`(`$or column$`)`$ bounded by $`n`$, then $`𝔰(𝒫)\sqrt{n}`$.
Likewise if $`A`$ is row $`(`$or column$`)`$ bounded by $`L`$, then $`𝔰(𝒮(A))L`$.
Proof. The pattern case follows from the row bounded case for the nonnegative matrix $`A=A_𝒫`$ with $`L=\sqrt{n}`$. Suppose that $`𝒮(A)`$ is row bounded by $`L`$. Consider any $`S𝒮(A)`$. Then $`sup_{i1}_{j1}|s_{ij}|^2L^2`$. Define vectors $`x_i=_{j1}s_{i,j}e_j`$ for $`i1`$. Then $`sup_{i1}x_iL`$; and $`x_i,e_j=s_{ij}`$. So by the Grothendieck–Haagerup condition,
$$S_m\underset{i,j}{sup}x_ie_jL.$$
Thus $`𝔰(𝒮(A))L`$.
###### Corollary 2.6.
If $`𝒫`$ is the union of a set row bounded by $`n`$ and a set column bounded by $`m`$, then $`𝒫`$ is Schur bounded with bound $`\sqrt{n}+\sqrt{m}`$.
Likewise, if $`A=B+C`$ such that $`B`$ is row bounded by $`L`$ and $`C`$ is column bounded by $`M`$, then $`𝔰(𝒮(A))L+M`$.
We require a combinatorial characterization of sets which are the union of an $`n`$-row bounded set and an $`m`$-column bounded set. This will be a consequence of the min-cut-max-flow theorem (see , for example). This is an elementary result in combinatorial optimization that has many surprising consequences. For example, it has been used by Richard Haydon to give a short proof of the reflexivity of commutative subspace lattices . It should be more widely known.
###### Lemma 2.7.
A pattern $`𝒫`$ is the union of a set $`𝒫_r`$ row bounded by $`m`$ and a set $`𝒫_c`$ column bounded by $`n`$ if and only if for every pair of finite subsets $`R,C`$,
$$|𝒫R\times C|m|R|+n|C|.$$
Similarly, a matrix $`A=[a_{ij}]`$ with nonnegative entries decomposes as a sum $`A=A_r+A_c`$ where $`A_r`$ is row bounded by $`M^{1/2}`$ and $`A_c`$ is column bounded by $`N^{1/2}`$ if and only if for every pair of finite subsets $`R,C`$,
$$\underset{iR}{}\underset{jC}{}a_{ij}^2M|R|+N|C|.$$
Proof. The two proofs are essentially identical. However the decomposition of $`𝒫`$ must be into two disjoint subsets. This means that the decomposition $`A_𝒫=A_{𝒫_1}+A_{𝒫_2}`$ is a split into $`0,1`$ matrices. We will work with $`A`$, but will explain the differences in the pattern version when it arises.
The condition is clearly necessary.
For the converse, we first show that it suffices to solve the finite version of the problem. For $`p`$, let $`A_p`$ be the restriction of $`A`$ to the first $`p`$ rows and columns. Suppose that we can decompose $`A_p=A_{r,p}+A_{c,p}`$ where $`A_{r,p}`$ is row bounded by $`M^{1/2}`$ and $`A_{c,p}`$ column bounded by $`N^{1/2}`$ for each $`p`$. Fix $`k`$ so that $`A_k0`$. For each $`pk`$, the set of such decompositions for $`A_p`$ is a compact subset of $`𝔐_p\times 𝔐_p`$. In the pattern case, we consider only $`0,1`$ decompositions. The restriction to the $`k\times k`$ corner is also a compact set, say $`𝒳_{k,p}`$. Observe that this is a decreasing sequence of nonempty compact sets. Thus $`_{pk}𝒳_{k,p}=𝒳_k`$ is nonempty. Therefore there is a consistent choice of a decomposition $`A=A_r+A_c`$ so that the restriction to each $`k\times k`$ corner lies in $`𝒳_k`$ for each $`k1`$. In the pattern case, the entries are all zeros and ones.
So now we may assume that $`A=[a_{ij}]`$ is a matrix supported on $`R_0\times C_0`$, where $`R_0`$ and $`C_0`$ are finite. We may also suppose that the $`l^2`$-norm of each row is greater than $`M^{1/2}`$ and the $`l^2`$-norm of each column is greater than $`N^{1/2}`$. For otherwise, we assign all of those entries in the row to $`A_r`$ (or all entries in the column to $`A_c`$) and delete the row (column). Solving the reduced problem will suffice. If after repeated use of this procedure, the matrix is empty, we are done. Otherwise, we reach a reduced situation in which the $`l^2`$-norm of each row is greater than $`M^{1/2}`$ and the $`l^2`$-norm of each column is greater than $`N^{1/2}`$.
Define a graph $`𝒢`$ with vertices $`\alpha `$, $`r_i`$ for $`iR_0`$, $`c_j`$ for $`jC_0`$, and $`\omega `$. Put edges from each $`r_iR_0`$ to each $`c_jC_0`$, from $`\alpha `$ to $`r_i`$, $`iR_0`$, and from $`c_j`$ to $`\omega `$, $`jC_0`$. Consider a network flow on the graph in which the edge from $`r_i`$ to $`c_j`$ may carry $`a_{ij}`$ units; edges leading out of $`\alpha `$ can carry up to $`M`$ units; and the edge from $`c_j`$ to $`\omega `$ can carry $`v_jN`$ units, where $`v_j=_{iR_0}a_{ij}^2`$. In the pattern case, these constraints are integers.
The min-cut-max-flow theorem states that the maximal possible flow from $`\alpha `$ to $`\omega `$ across this network equals the minimum flow across any cut that separates $`\alpha `$ from $`\omega `$. Moreover, when the data is integral, the maximal flow comes from an integral solution. A cut $`𝒳`$ is just a partition of the vertices into two disjoint sets $`\{\alpha \}R_1C_1`$ and $`\{\omega \}R_2C_2`$. The flow across the cut is the total of allowable flows on each edge between the two sets.
The flow across the cut $`𝒳`$ is
$`f(𝒳)`$ $`={\displaystyle \underset{iR_1}{}}{\displaystyle \underset{jC_2}{}}a_{ij}^2+M|R_2|+{\displaystyle \underset{jC_1}{}}(v_jN)`$
$`={\displaystyle \underset{iR_1}{}}{\displaystyle \underset{jC_2}{}}a_{ij}^2+M|R_2|N|C_1|+{\displaystyle \underset{iR_0}{}}{\displaystyle \underset{jC_1}{}}a_{ij}^2`$
$`={\displaystyle \underset{iR_0}{}}{\displaystyle \underset{jC_0}{}}a_{ij}^2{\displaystyle \underset{iR_2}{}}{\displaystyle \underset{jC_2}{}}a_{ij}^2+M|R_2|+N|C_2|N|C_0|`$
$`{\displaystyle \underset{iR_0}{}}{\displaystyle \underset{jC_0}{}}a_{ij}^2N|C_0|`$
The last inequality uses the hypothesis on $`A`$ with $`R=R_2`$ and $`C=C_2`$. On the other hand, the cut separating $`\omega `$ from the rest has flow exactly
$$\underset{jC_0}{}(v_jN)=\underset{iR_0}{}\underset{jC_0}{}a_{ij}^2N|C_0|.$$
Therefore there is a network flow that achieves this maximum. In the pattern case, the solution is integral. Necessarily this will involve a flow of exactly $`v_jN`$ from each $`jC_0`$ to $`\omega `$. Let $`b_{ij}`$ be the optimal flow from $`r_i`$ to $`c_j`$. So $`0b_{ij}a_{ij}`$. The flow out of each $`r_i`$ equals the flow into $`r_i`$ from $`\alpha `$, whence $`_{jC_0}b_{ij}M`$.
Define the matrix $`A_r=\left[\sqrt{b_{ij}}\right]`$ and $`A_c=\left[\sqrt{a_{ij}b_{ij}}\right]`$. In the pattern case, these entries are 0 or 1. Then the rows of $`A_r`$ are bounded by $`M^{1/2}`$. The $`j`$th column of $`A_c`$ has norm squared equal to
$$\underset{iR_0}{}a_{ij}b_{ij}=v_j(v_jN)=N.$$
This is the desired decomposition and it is integral for patterns.
To construct large norm Schur multipliers on certain patterns, we will make use of the following remarkable result by Françoise Lust-Piquard \[11, Theorem 2\]. While the method of proof is unexpected, it is both short and elementary.
###### Theorem 2.8 (Lust-Piquard).
Given any $`(`$finite or infinite$`)`$ nonnegative matrix $`X=[x_{ij}]`$ satisfying
$$\underset{i}{\mathrm{max}}\underset{j}{}x_{ij}^21\text{and}\underset{j}{\mathrm{max}}\underset{i}{}x_{ij}^21\text{for all}i,j,$$
there is an operator $`Y=[y_{ij}]`$ so that
$$Y\sqrt{6}\text{and}|y_{ij}|x_{ij}\text{ for all }i,j.$$
The constant of $`\sqrt{6}`$ is optimal, as shown in an addendum to .
###### Lemma 2.9.
Let $`A=[a_{ij}]`$ be a nonnegative $`m\times m`$ matrix such that $`_{i=1}^m_{j=1}^ma_{ij}^2=m\alpha `$. Then there is a Schur multiplier $`S𝒮(A)`$ such that $`S_m\frac{1}{2}\sqrt{\frac{\alpha }{3}}`$.
Proof. We may assume that there are no nonzero rows or columns. Let
$$r_i=\underset{j=1}{\overset{m}{}}a_{ij}^2\text{and}c_j=\underset{i=1}{\overset{m}{}}a_{ij}^2.$$
Define
$$x_{ij}=\frac{a_{ij}}{\sqrt{r_i+c_j}}.$$
Let $`X=[x_{ij}]`$. The row norms of $`X`$ satisfy
$$\underset{j=1}{\overset{m}{}}x_{ij}^2\underset{j=1}{\overset{m}{}}\frac{a_{ij}^2}{r_i}=1;$$
and similarly the column norms are bounded by 1.
By Theorem 2.8, there is a matrix $`Y`$ such that
$$Y\sqrt{6}\text{and}|y_{ij}|x_{ij}\text{for all}i,j.$$
Define $`s_{ij}=a_{ij}x_{ij}/y_{ij}`$ (where $`0/0:=0`$). Then $`S=[s_{ij}]`$ belongs to $`𝒮(A)`$. Observe that
$$S(Y)=Z:=[a_{ij}x_{ij}]=\left[\begin{array}{c}\frac{a_{ij}^2}{\sqrt{r_i+c_j}}\end{array}\right].$$
Hence $`S_mZ/K`$.
Define vectors $`u=(u_i)`$ and $`v=(v_j)`$ by
$$u_i=\left(\frac{r_i}{m\alpha }\right)^{1/2}\text{and}v_j=\left(\frac{c_j}{m\alpha }\right)^{1/2}.$$
Then $`u_2^2={\displaystyle \frac{1}{m\alpha }}{\displaystyle \underset{i=1}{\overset{m}{}}}r_i=1`$ and similarly $`v_2=1`$. Compute
$$Zu^{}Zv=\frac{1}{m\alpha }\underset{i=1}{\overset{m}{}}\underset{j=1}{\overset{m}{}}a_{ij}^2\sqrt{\frac{r_ic_j}{r_i+c_j}}.$$
Observe that $`\sqrt{{\displaystyle \frac{r_ic_j}{r_i+c_j}}}=\left({\displaystyle \frac{1}{r_i}}+{\displaystyle \frac{1}{c_j}}\right)^{1/2}`$. Also
$`{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \underset{j=1}{\overset{m}{}}}a_{ij}^2\left({\displaystyle \frac{1}{r_i}}+{\displaystyle \frac{1}{c_j}}\right)`$ $`={\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{a_{ij}^2}{r_i}}+{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{m}{}}}{\displaystyle \frac{a_{ij}^2}{c_j}}`$
$`={\displaystyle \underset{i=1}{\overset{m}{}}}1+{\displaystyle \underset{j=1}{\overset{m}{}}}1=2m.`$
A routine Lagrange multiplier argument shows that if $`\alpha _k0`$ are constants, $`t_k>0`$ are variables, and $`_{k=1}^{m^2}\alpha _kt_k=2m`$, then $`_{k=1}^{m^2}\alpha _kt_k^{1/2}`$ is minimized when all $`t_k`$ are equal. Hence if $`_{k=1}^{m^2}\alpha _k=m\alpha `$,
$$\underset{k=1}{\overset{m^2}{}}\alpha _kt_k^{1/2}m\alpha \left(\frac{2m}{m\alpha }\right)^{1/2}=m\alpha \sqrt{\frac{\alpha }{2}}$$
Applying this to the numbers $`\frac{1}{r_i}+\frac{1}{c_j}`$ yields
$$Z\frac{1}{m\alpha }\underset{i=1}{\overset{m}{}}\underset{j=1}{\overset{m}{}}a_{ij}^2\left(\frac{1}{r_i}+\frac{1}{c_j}\right)^{1/2}\sqrt{\frac{\alpha }{2}}.$$
Thus $`S_m\frac{\sqrt{\alpha }}{\sqrt{6}\sqrt{2}}=\frac{1}{2}\sqrt{\frac{\alpha }{3}}`$.
Proof of Theorem 2.3 and Theorem 2.4.
Statements (2) and (3) are equivalent by Lemma 2.7, taking $`m=n`$ and $`M=N`$.
Assuming (2), Corollary 2.6 shows that $`𝒫`$ or $`𝒮(A)`$ is Schur bounded by $`2\sqrt{m}`$ or $`2M`$. Assuming (3) in the pattern case, the supremum exceeds $`m1`$; so Lemma 2.9 shows that
$$𝔰(𝒫)\frac{\sqrt{m1}}{2\sqrt{3}}\frac{\sqrt{m}}{4}$$
for $`m4`$. For $`m16`$, $`\sqrt{m}/41`$; and $`1`$ is also a lower bound for any pattern. For the matrix case, we use the exact supremum in Lemma 2.9, so we obtain a lower bound of $`M/4`$.
Conversely, if the supremum in (3) is infinite, the same argument shows that the Schur bound is infinite. In fact it is easy to see that this implies that $`𝒮(𝒫)`$ or $`𝒮(A)`$ contains unbounded Schur multipliers. It is not difficult to produce disjoint finite rectangles $`R_n\times C_n`$ on which the ratio in (3) exceeds $`n^2`$. So by Lemma 2.9, we construct a Schur multiplier $`S_n`$ in $`𝒮(𝒫)`$ or $`𝒮(A)`$ supported on $`R_n\times C_n`$ with Schur norm at least $`n/4`$. Take $`S`$ to be defined on each rectangle as $`S_n`$ and zero elsewhere. Then $`S`$ is an unbounded Schur multiplier in this class.
###### Remark 2.10.
One might suspect, from the $`\sqrt{n}`$ arising in Lemma 2.5, that if two matrices are supported on pairwise disjoint patterns, there might be an $`L^2`$ estimate on the Schur norm of the sum. This is not the case, as the following example shows.
Let $`\mathrm{𝟏}=(1,1,1,1)^t^4`$ and $`A=\mathrm{𝟏𝟏}^{}I`$. If $`U=\mathrm{diag}(1,i,1,i)`$, then the diagonal expectation is
$$\mathrm{\Delta }(X)=S_I(X)=\frac{1}{4}\underset{k=0}{\overset{3}{}}U^kXU^k.$$
We use a device due to Bhatia–Choi–Davis . Observe that
$`S_{A+tI}(X)`$ $`=X+(t1)\mathrm{\Delta }(X)`$
$`=(1+{\displaystyle \frac{t1}{4}})X+{\displaystyle \frac{t1}{4}}{\displaystyle \underset{k=1}{\overset{3}{}}}U^kXU^k.`$
Therefore
$`S_{A+tI}_m`$ $`\left|1+{\displaystyle \frac{t1}{4}}\right|+{\displaystyle \frac{3|t1|}{4}}`$
$`=\{\begin{array}{cc}|t|\hfill & \text{if}t1\text{ or }t3\hfill \\ \frac{1}{2}|3t|\hfill & \text{if}3t1\hfill \end{array}.`$
On the other hand, $`S_{A+tI}(I)=tI`$; so $`S_{A+tI}_m|t|`$. Observe that $`\frac{1}{4}\mathrm{𝟏𝟏}^{}`$ is a projection. Hence $`A+tI=\mathrm{𝟏𝟏}^{}+(t1)I`$ has spectrum $`\{t1,t+3\}`$; and thus
$$A+tI=\mathrm{max}\{|t1|,|t+3|\}.$$
So $`AI=2`$. If $`3t1`$, then $`S_{A+tI}(AI)=AtI`$ has norm $`|3t|`$ and so $`S_{A+tI}_m|3t|/2`$.
In particular, $`S_A_m=\frac{3}{2}`$ and $`S_I_m=1`$, but
$$S_{AI}_m=2>\left(S_A_m^2+S_I_m^2\right)^{1/2}.$$
###### Remark 2.11.
In , Bennett, Goodman and Newman show that if $`A`$ is an $`n\times n`$ matrix with entries taking the values $`\pm 1`$ with probability .5, then on average the norm of $`A`$ is bounded by $`K\sqrt{n}`$, where $`K`$ is a universal constant. This is best possible as each row and column has norm $`\sqrt{n}`$. This minimum can be achieved in certain cases, for example by tensoring copies of $`\left[\begin{array}{cc}1& 1\\ 1& 1\end{array}\right]`$ together. The maximum norm occurs for the matrix $`\mathrm{𝟏𝟏}^{}`$ for which all entries are 1, in which case the norm is $`n`$. So we see that the average norm is within a constant of the minimum.
This can be used to show that, on average, the Schur norm $`A_m`$ is near the maximum $`\sqrt{n}`$. Indeed, $`S_A(A)=\mathrm{𝟏𝟏}^{}`$. So
$$A_m\frac{n}{A}K^1\sqrt{n}$$
on average.
## 3. Hankel and Toeplitz Patterns
A Hankel pattern is a set of the form
$$(S)=\{(i,j):i,j,i+jS\}\text{for}S.$$
A Toeplitz pattern is a set of the form
$$𝒯(S)=\{(i,j):i,j_0,ijS\}\text{for}S.$$
Recall that a set $`S=\{s_1<s_2<\mathrm{}\}`$ is lacunary if there is a constant $`q>1`$ so that $`s_{i+1}/s_i>q`$ for all $`i1`$.
Nikolskaya and Farforovskaya show that a Hankel pattern is Schur bounded if and only if it is a finite union of lacunary sets \[13, Theorem 3.8\], by considering Fejér kernels and Toeplitz extensions. We give an elementary proof based on Theorem 2.3.
###### Proposition 3.1.
Consider a Hankel pattern $`(S)`$ of a set $`S`$. Then the following are equivalent:
1. $`(S)`$ is Schur bounded.
2. $`(S)`$ is the union of a row finite and a column finite set.
3. $`sup_{k0}|S(2^{k1},2^k]|<\mathrm{}`$.
4. $`S`$ is the union of finitely many lacunary sets.
Proof. By Theorem 2.3, (1) and (2) are equivalent.
Let $`a_k=|S(2^{k1},2^k]|`$ for $`k0`$. If (3) holds, $`\mathrm{max}_{k0}a_k=L<\mathrm{}`$. So $`S`$ splits into $`2L`$ subsets with at most one element in every second interval $`(2^{k1},2^k]`$; which are therefore lacunary with ratio at least 2. Conversely, suppose that $`S`$ is the union of finitely many lacunary sets. A lacunary set with ratio $`q`$ may be split into $`d`$ lacunary sets of ratio 2 provided that $`q^d2`$. So suppose that there are $`L`$ lacunary sets of ratio 2. Then each of these sets intersects $`(2^{k1},2^k]`$ in at most one element. Hence $`\mathrm{max}_{k0}a_kL<\mathrm{}`$. Thus (3) and (4) are equivalent.
Suppose that $`S`$ is the union of $`L`$ sets $`S_i`$ which are each lacunary with constant $`2`$. Split each $`(S_i)`$ into the subsets $`R_i`$ on or below the diagonal and $`C_i`$ above the diagonal. Observe that $`R_i`$ is row bounded by 1, and $`C_i`$ is column bounded by 1. Hence (4) implies (2).
Consider the subset of $`(S)`$ in the first $`2^k`$ rows and columns $`R_k\times C_k`$. This square will contain at least $`2^{k1}a_k`$ entries corresponding to the backward diagonals for $`S(2^{k1},2^k]`$, which all have more than $`2^{k1}`$ entries. Thus
$$\underset{k0}{sup}\frac{|(S)(R_k\times C_k)|}{|R_k|+|C_k|}\underset{k0}{sup}\frac{2^{k1}a_k}{2^k+2^k}=\underset{k0}{sup}\frac{a_k}{4}.$$
Hence if (3) fails, this supremum if infinite. Thus $`(S)`$ is not the union of a row finite and a column finite set. So (2) fails.
The situation for Toeplitz patterns is quite different. It follows from classical results, as we explain below, and Nikolskaya and Farforovskaya outline a related proof \[13, Remark 3.9\]. But first we show how it follows from our theorem.
###### Proposition 3.2.
The Toeplitz pattern $`𝒯(S)`$ of any infinite set $`S`$ is not Schur bounded. Further,
$$\frac{1}{4}|S|^{1/2}𝔰(𝒯(S))|S|^{1/2}.$$
Proof. Since $`𝒯(S)`$ is clearly row bounded by $`|S|`$, the upper bound follows from Lemma 2.5.
Suppose that $`S=\{s_1<s_2<\mathrm{}<s_n\}`$. Consider the $`m\times m`$ square matrix with upper left hand corner equal to $`(s_1,0)`$ if $`s_10`$ or $`(0,s_1)`$ if $`s_1<0`$. Then beginning with row $`m(s_ns_1)`$, there will be $`n`$ entries of $`𝒯(S)`$ in each row. Thus the total number of entries is at least $`n(m(s_ns_1))`$. For $`m`$ sufficiently large, this exceeds $`(n1)m`$. Hence by Lemma 2.9,
$$𝔰(𝒯(S))\frac{\sqrt{n1}}{2\sqrt{3}}\frac{\sqrt{n}}{4}$$
provided $`n4`$. The trivial lower bound of 1 yields the lower bound for $`n<4`$.
To see how this is done classically, we recall the following \[2, Theorem 8.1\]. Here, $`𝒯`$ denotes the space of Toeplitz operators.
###### Theorem 3.3 (Bennett).
A Toeplitz matrix $`A=[a_{ij}]`$ determines a bounded Schur multiplier if and only if there is a finite complex Borel measure $`\mu `$ on the unit circle $`𝕋`$ so that $`\widehat{\mu }(n)=a_n`$, $`n`$. Moreover
$$A_m=S_A|_𝒯=\mu .$$
We combine this with estimates obtained from the Khintchine inequalities.
###### Theorem 3.4.
Let $`(a_k)_k`$ be an $`l^2`$ sequence and let $`A=[a_{ij}]`$. Then
$$\frac{1}{\sqrt{2}}(a_k)_2𝔰(A)(a_k)_2.$$
Proof. Suppose $`S𝒮(A)`$, that is, $`S=[s_{ij}]`$ with $`|s_{ij}|a_{ij}`$. Then each row of $`S`$ has norm bounded by $`(a_k)_2`$. Hence by Lemma 2.5, $`S_m(a_k)_2`$. So $`𝔰(A)(a_k)_2`$.
Conversely, let $`X:=\{1,1\}^{}`$. Put the measure $`\mu `$ on $`X`$ which is the product of measures on $`\{1,1\}`$ assigning measure $`1/2`$ to both $`\pm 1`$. For $`\epsilon =(\epsilon _k)_k`$ in $`X`$, define $`f_\epsilon (\theta )=_k\epsilon _ka_ke^{ik\theta }`$. Then $`f_\epsilon L^2(𝕋)L^1(𝕋)`$. Hence $`S_\epsilon :=S_{T_{f_\epsilon }}`$ defines a bounded Schur multiplier with
$$S_\epsilon _m=f_\epsilon _1f_\epsilon _2=(a_k)_2.$$
Then we make use of the Khintchine inequality :
$$\frac{1}{\sqrt{2}}(a_k)_2_Xf_\epsilon _1𝑑\mu (\epsilon )(a_k)_2.$$
It follows that on average, most $`f_\epsilon `$ have $`L^1`$-norm comparable to the $`L^2`$-norm. In particular, there is some choice of $`\epsilon `$ with $`f_\epsilon _1\frac{1}{\sqrt{2}}(a_k)_2`$. Thus $`𝔰(A)S_\epsilon _m\frac{1}{\sqrt{2}}(a_k)_2.`$
###### Remark 3.5.
In the case of a finite Toeplitz pattern $`𝒯(S)`$, say $`S=\{s_1<s_2<\mathrm{}<s_n\}`$, $`f_\epsilon =_{k=1}^n\epsilon _ke^{is_k\theta }`$. We can use the Khintchine inequality for $`L^{\mathrm{}}`$:
$$(a_k)_2_Xf_\epsilon _{\mathrm{}}𝑑\mu (\epsilon )\sqrt{2}(a_k)_2.$$
Thus there will be choices of $`\epsilon `$ so that $`f_\epsilon _{\mathrm{}}\sqrt{2n}`$. Then note that $`S_{T_{f_\epsilon }}(T_{f_\epsilon })=T_{f_\mathrm{𝟏}}`$, where $`f_\mathrm{𝟏}=_{k=1}^ne^{is_k\theta }`$. Clearly $`f_\mathrm{𝟏}_{\mathrm{}}=f_\mathrm{𝟏}(0)=n`$. Thus $`S_{T_{f_\epsilon }}|_{𝒯(S)}\sqrt{n/2}`$.
## 4. Patterns with a Symmetry Group
Consider a finite group $`G`$ acting transitively on a finite set $`X`$. Think of this as a matrix representation on the Hilbert space $`_X`$ with orthonormal basis $`\{e_x:xX\}`$. Let $`\pi `$ denote the representation of $`G`$ on $`_X`$ and $`𝒯`$ the commutant of $`\pi (G)`$. The purpose of this section is to compute the norm of $`S_T`$ for $`T𝒯`$.
Decompose $`X^2`$ into G-orbits $`X_i`$ for $`0in`$, beginning with the diagonal $`X_0=\{(x,x):xX\}`$. Let $`T_i(_X)`$ denote the matrix with $`1`$s on the entries of $`X_i`$ and $`0`$ elsewhere. Then it is easy and well-known that $`𝒯`$ is $`\mathrm{span}\{T_i:0in\}`$. In particular, $`𝒯`$ is a C\*-algebra. Also observe that every element of $`𝒯`$ is constant on the main diagonal.
Since $`G`$ acts transitively on $`X`$, $`r_i:=|\{yX:(x,y)X_i\}|`$ is independent of the choice of $`xX`$. Thus the vector $`\mathrm{𝟏}`$ of all ones is a common eigenvector for each $`T_i`$, and hence for all elements of $`𝒯`$, corresponding to a one-dimensional reducing subspace on which $`G`$ acts via the trivial representation.
First we establish an easy, general upper bound for $`T_m`$ where $`T𝒯`$. As usual, $`\mathrm{\Delta }`$ is the expectation onto the diagonal.
###### Proposition 4.1.
For a matrix $`T`$,
$$T_m\mathrm{\Delta }(|T^{}|)^{1/2}\mathrm{\Delta }(|T|)^{1/2}=|T^{}|_m^{1/2}|T|_m^{1/2}.$$
Proof. Use polar decomposition to factor $`T=U|T|`$. Define vectors $`x_i=|T|^{1/2}e_i`$ and $`y_j=|T|^{1/2}U^{}e_j`$. Then
$$x_i,y_j=|T|^{1/2}e_i,|T|^{1/2}U^{}e_j=Te_i,e_j.$$
This yields a Grothendieck–Haagerup form for $`S_T`$. Now
$$x_i^2=|T|^{1/2}e_i,|T|^{1/2}e_i=|T|e_i,e_i.$$
Hence $`\mathrm{max}_ix_i=\mathrm{\Delta }(|T|)^{1/2}`$. Similarly, since $`|T|^{1/2}U^{}=U^{}|T^{}|^{1/2}`$
$$y_j^2=U^{}|T^{}|^{1/2}e_j,U^{}|T^{}|^{1/2}e_j=|T^{}|e_j,e_j.$$
So $`\mathrm{max}_jy_j=\mathrm{\Delta }(|T^{}|)^{1/2}`$. Therefore
$$T_m\underset{i,j}{\mathrm{max}}x_iy_j=\mathrm{\Delta }(|T^{}|)^{1/2}\mathrm{\Delta }(|T|)^{1/2}.$$
Since $`|T|`$ and $`|T^{}|`$ are positive, the Schur norm is just the sup of the diagonal entries.
###### Corollary 4.2.
If $`T=T^{}`$, then $`T_m\mathrm{\Delta }(|T|)`$.
###### Remark 4.3.
In general this is a strict inequality. If $`T=\left[\begin{array}{cc}4& 3\\ 3& 1\end{array}\right]`$, then $`|T|=\left[\begin{array}{cc}2\sqrt{5}& \sqrt{5}\\ \sqrt{5}& \sqrt{5}\end{array}\right]`$. But $`S_T_m=4<2\sqrt{5}`$. Indeed, take $`x_1=y_1=2e_1`$ and $`x_2=\frac{3}{2}e_1+\frac{\sqrt{5}}{2}e_2`$ and $`y_2=\frac{3}{2}e_1\frac{\sqrt{5}}{2}e_2`$.
The main result of this section is:
###### Theorem 4.4.
Let $`X`$ be a finite set with a transitive action by a finite group $`G`$. If $`T`$ belongs to $`𝒯`$, the commutant of the action of $`G`$, then for any $`x_0X`$,
$$T_m=S_T|_𝒯=|X|^1\mathrm{Tr}(|T|)=|T|e_{x_0},e_{x_0}.$$
This result is a special case of a nice result of Mathias . As far as we know, the application of Mathias’ result to the case of matrices invariant under group actions has not been exploited. As Mathias’s argument is short and elegant, we include it.
###### Theorem 4.5 (Mathias).
If $`T`$ is an $`n\times n`$ matrix with $`\mathrm{\Delta }(|T^{}|)`$ and $`\mathrm{\Delta }(|T|)`$ scalar, then
$$T_m=\frac{1}{n}\mathrm{Tr}(|T|).$$
Proof. For an upper bound, Proposition 4.1 shows that
$`T_m`$ $`\mathrm{\Delta }(|T^{}|)^{1/2}\mathrm{\Delta }(|T|)^{1/2}`$
$`=\left(\frac{1}{n}\mathrm{Tr}(|T^{}|)\right)^{1/2}\left(\frac{1}{n}\mathrm{Tr}(|T|)\right)^{1/2}=\frac{1}{n}\mathrm{Tr}(|T|),`$
because $`|T|`$ and $`|T^{}|`$ are constant on the main diagonal, and $`|T^{}|`$ is unitarily equivalent to $`|T|`$, and so has the same trace.
For the lower bound, use the polar decomposition $`T=W|T|`$. Let $`\overline{W}`$ have matrix entries which are the complex conjugates of the matrix entries of $`W`$. Write $`T=[t_{ij}]`$ and $`W=[w_{ij}]`$ as $`n\times n`$ matrices in the given basis. Set $`\mathrm{𝟏}`$ to be the vector with $`n`$ 1’s. Then
$`T_m`$ $`S_T(\overline{W}){\displaystyle \frac{1}{n}}S_T(\overline{W})\mathrm{𝟏},\mathrm{𝟏}`$
$`={\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{n}{}}}\overline{w}_{ij}t_{ij}={\displaystyle \frac{1}{n}}{\displaystyle \underset{j=1}{\overset{n}{}}}W^{}Te_j,e_j={\displaystyle \frac{1}{n}}\mathrm{Tr}(|T|)`$
Thus $`T_m=\frac{1}{n}\mathrm{Tr}(|T|)`$.
Proof of Theorem 4.4. We have already observed that elements of $`𝒯`$ are constant on the diagonal. Thus $`T_m=\frac{1}{n}\mathrm{Tr}(|T|)=|T|e_{x_0},e_{x_0}`$. For the rest, observe that $`W`$ belongs to $`\mathrm{C}^{}(T)`$. Hence so does $`\overline{W}`$ because the basis $`T_i`$ of $`𝒯`$ has real entries.
We will provide an interesting example in the next section. For now we provide a couple of more accessible ones.
###### Example 4.6.
Consider the action of the symmetric group $`𝔖_n`$ acting on a set $`X`$ with $`n`$ elements in the canonical way. Then the orbits in $`X^2`$ are just the diagonal $`X_0`$ and its complement $`X_1`$. So $`S_{X_1}`$ is the projection onto the off-diagonal part of the matrix.
Observe that $`X_1=\mathrm{𝟏𝟏}^{}I`$, where $`\mathrm{𝟏}`$ is the vector of $`n`$ ones. Since $`\mathrm{𝟏𝟏}^{}=nP`$, where $`P`$ is the projection onto $`\mathrm{𝟏}`$, $`X_1=(n1)PP^{}`$. Therefore we obtain a formula due to Bhatia, Choi and Davis
$`X_1_m`$ $`={\displaystyle \frac{1}{n}}\mathrm{Tr}(|X_1|)={\displaystyle \frac{1}{n}}\mathrm{Tr}\left((n1)P+P^{}\right)`$
$`={\displaystyle \frac{1}{n}}(n1+n1)=2{\displaystyle \frac{2}{n}}.`$
###### Example 4.7.
Consider the cyclic group $`C_n`$ acting on an $`n`$-element set, $`n3`$. Let $`U`$ be the unitary operator given by $`Ue_k=e_{k+1}`$ for $`1kn`$, working modulo $`n`$. The powers of $`U`$ yields a basis for the commutant of the group action.
Consider $`T=U+I`$. The spectrum of $`U`$ is just $`\{\omega ^k:0kn1\}`$ where $`\omega =e^{2\pi i/n}`$. Thus the spectrum of $`|T|`$ consists of the points
$$|1+\omega ^k|=2|\mathrm{cos}(\frac{k\pi }{n})|\text{for}0kn1.$$
Hence
$$T_m=\frac{1}{n}\mathrm{Tr}(|T|)=\frac{2}{n}\underset{k=0}{\overset{n1}{}}|\mathrm{cos}(\frac{k\pi }{n})|=\{\begin{array}{cc}\frac{2\mathrm{cos}(\frac{\pi }{2n})}{n\mathrm{sin}(\frac{\pi }{2n})}\hfill & n\text{ even}\hfill \\ \frac{2}{n\mathrm{sin}(\frac{\pi }{2n})}\hfill & n\text{ odd}\hfill \end{array}$$
Thus the limit as $`n`$ tends to infinity is $`{\displaystyle \frac{4}{\pi }}`$. The multiplier norms for the odd cycles decrease to $`\frac{4}{\pi }`$, while the even cycles increase to the same limit.
###### Example 4.8.
Mathias considers polynomials in the circulant matrices $`C_z`$ given by $`C_ze_k=e_{k+1}`$ for $`1k<n`$ and $`C_ze_n=ze_1`$, where $`|z|=1`$. This falls into our rubric because there is a diagonal unitary $`D`$ so that $`DC_zD^{}=wU`$ where $`U`$ is the cycle in the previous example and $`w`$ is any $`n`$th root of $`z`$. It is easy to see that conjugation by a diagonal unitary has no effect on the Schur norm. Thus any polynomial in $`C_z`$ is unitarily equivalent to an element of $`\mathrm{C}^{}(U)`$ via the diagonal $`D`$. Hence the Schur norm equals the normalized trace of the absolute value.
The most interesting example of this was obtained with $`z=1`$ and $`S_n=_{k=0}^{n1}C_1^k`$ which is the matrix with entries $`\mathrm{sgn}(ij)`$. So the Schur multiplier defined by $`S_n`$ is a finite Hilbert transform. Mathias shows that
$$S_n_m=\frac{2}{n}\underset{j=1}{\overset{n/2}{}}\mathrm{cot}\frac{(2j1)\pi }{2n}.$$
From this, he obtains sharper estimates on the norm of triangular truncation than are obtained in .
## 5. Kneser and Johnson Graph Patterns
In this section, we consider an interesting family of symmetric patterns which arise commonly in graph theory and combinatorial codes. The Johnson graphs $`J(v,n,i)`$ have $`\left(\genfrac{}{}{0pt}{}{v}{n}\right)`$ vertices indexed by $`n`$ element subsets of a $`v`$ element set, and edges between $`A`$ and $`B`$ if $`|AB|=i`$. Thus $`0in`$. We consider only $`1nv/2`$ since, if $`n>v/2`$, one obtains the same graphs by considering the complementary sets of cardinality $`vn`$. We will explicitly carry out the calculation for the Kneser graphs $`K(v,n)=J(v,n,0)`$, and in particular, for $`K(2n+1,n)`$. For more on Johnson and Kneser graphs, see .
We obtained certain Kneser graphs from Toeplitz patterns. Take a finite subset $`S=\{s_1<s_2<\mathrm{}<s_{2n+1}\}`$ and consider the Toeplitz pattern $`𝒫`$ with diagonals in $`S`$, namely $`𝒫=\{(i,j):jiS\}`$. Consider $`R`$ to be the set of all sums of $`n`$ elements from $`S`$ and $`C`$ to be the set of all sums of $`n+1`$ elements from $`S`$. Index $`R`$ by the corresponding subset $`A`$ of $`\{1,2,\mathrm{},2n+1\}`$ of cardinality $`n`$; and likewise index each element of $`C`$ by a subset $`B`$ of cardinality $`n+1`$. Then for each entry $`A`$ in $`R`$, there are exactly $`n+1`$ elements of $`C`$ which contain it. The difference of the sums is an element of $`S`$. It is convenient to re-index $`C`$ by sets of cardinality $`n`$, replacing $`B`$ by its complement $`\{1,2,\mathrm{},2n+1\}B`$. Then the pattern can be seen to be the Kneser graph $`K(2n+1,n)`$ with $`\left(\genfrac{}{}{0pt}{}{2n+1}{n}\right)`$ vertices indexed by $`n`$ element subsets of a $`2n+1`$ element set, with an edge between vertices $`A`$ and $`B`$ if $`AB=\mathrm{}`$. In general, unfortunately, $`𝒫(R\times C)`$ will contain more than just these entries, because two subsets of $`S`$ of size $`n+1`$ can have the same sum.
The adjacency matrix of a graph $`𝒢`$ is a $`v\times v`$ matrix with a 1 in each entry $`(i,j)`$ corresponding to an edge from vertex $`i`$ to vertex $`j`$, and 0’s elsewhere. This is a symmetric matrix and its spectral theory is available in the graph theory literature; see, for example, . We prove the simple facts we need.
Fix $`(v,n)`$ with $`nv`$ and let $`X`$ denote the set of $`n`$ element subsets of $`\{1,\mathrm{},v\}`$. Define a Hilbert space $`=_X`$ as in the previous section but write the basis as $`\{e_A:AX\}`$. Observe that there is a natural action $`\pi `$ of the symmetric group $`𝔖_v`$ on $`X`$. The orbits in $`X^2`$ are
$$X_i=\{(A,B):A,BX,|AB|=i\}\text{for}0in.$$
The matrix $`T_i`$ is just the adjacency matrix of the Johnson graph $`J(v,n,i)`$ and, in particular, $`T_n=I`$.
This action has additional structure that does not hold for arbitrary transitive actions.
###### Lemma 5.1.
The commutant $`𝒯=\mathrm{span}\{T_i:0in\}`$ of $`\pi (𝔖_v)`$ is abelian. Thus $`\pi `$ decomposes into a direct sum of $`n+1`$ distinct irreducible representations.
Proof. Equality with the span was observed in the last section. To see that the algebra $`𝒯`$ is abelian, observe that $`T_iT_j=_{k=0}^na_{ijk}T_k`$ where we can find the coefficients $`a_{ijk}`$ by fixing any two sets $`A,BV`$ of size $`n`$ with $`|AB|=k`$ and computing
$$a_{ijk}=|\{CV:|C|=n,|AC|=i,|CB|=j\}.$$
This is clearly independent of the order of $`i`$ and $`j`$. As $`𝒯`$ is abelian and $`n+1`$ dimensional, the representation $`\pi `$ decomposes into a direct sum of $`n+1`$ distinct irreducible representations.
###### Corollary 5.2.
$`T_i=\left(\genfrac{}{}{0pt}{}{n}{i}\right)\left(\genfrac{}{}{0pt}{}{vn}{ni}\right)`$ and this is an eigenvalue of multiplicity one. The spectrum of $`T_i`$ contains at most $`n+1`$ points.
Proof. Observe that if $`|A|=n`$, then the number of subsets $`BX`$ with $`|AB|=i`$ is $`\left(\genfrac{}{}{0pt}{}{n}{i}\right)\left(\genfrac{}{}{0pt}{}{vn}{ni}\right)`$. Thus $`T_i`$ has this many $`1`$’s in each row. Hence
$$T_i\mathrm{𝟏}=\left(\genfrac{}{}{0pt}{}{n}{i}\right)\left(\genfrac{}{}{0pt}{}{vn}{ni}\right)\mathrm{𝟏}.$$
Clearly $`T_i`$ has nonnegative entries and is indecomposable (except for $`i=n`$, the identity matrix). So by the Perron–Frobenius Theorem, $`\left(\genfrac{}{}{0pt}{}{n}{i}\right)\left(\genfrac{}{}{0pt}{}{vn}{ni}\right)`$ is the spectral radius and $`\mathrm{𝟏}`$ is the unique eigenvector; and there are no other eigenvalues on the circle of this radius. Since $`T=T^{}`$, the norm equals spectral radius. As $`𝒯`$ is $`n+1`$ dimensional, the spectrum can have at most $`n+1`$ points.
We need to identify the invariant subspaces of $`𝔖_v`$ as they are the eigenspaces of $`T_i`$. The space $`V_0=\mathrm{𝟏}`$ yields the trivial representation. Define vectors associated to sets $`C\{1,\mathrm{},v\}`$ of cardinality at most $`n`$, including the empty set, by
$$v_C:=\underset{|A|=n,AC=\mathrm{}}{}e_A.$$
Then define subspaces $`V_i=\mathrm{span}\{v_C:|C|=i\}`$ for $`0in.`$ It is obvious that each $`V_i`$ is invariant for $`𝔖_v`$. Given $`C`$ with $`|C|=i`$, we have
$$\underset{CD,|D|=i+1}{}v_D=(vni)v_C,$$
as the coefficient of $`e_A`$ counts the number of choices for the $`(i+1)`$st element of $`D`$ disjoint from an $`A`$ already disjoint from $`C`$. Therefore
$$\mathrm{𝟏}=V_0V_1V_2\mathrm{}V_n.$$
So the $`n+1`$ subspaces $`W_i=V_iV_{i1}`$ are invariant for $`𝔖_v`$.
Let $`E_i`$ denote the idempotent in $`𝒯`$ projecting onto $`W_i`$. Observe that $`𝒯=\mathrm{span}\{E_i:0in\}`$. We need to know the dimension of these subspaces.
###### Lemma 5.3.
The vectors $`\{v_C:|C|=i\}`$ are linearly independent. Hence $`\mathrm{dim}W_i=\left(\genfrac{}{}{0pt}{}{v}{i}\right)\left(\genfrac{}{}{0pt}{}{v}{i1}\right)`$.
Proof. Suppose that $`v_{C_0}+_{|C|=i,CC_0}\gamma _Cv_C=0`$. By averaging over the subgroup of $`𝔖_v`$ which fixes $`C_0`$, namely $`𝔖_i\times 𝔖_{vi}`$, we may assume that the coefficients are invariant under this action. Hence $`\gamma _C=\alpha _j`$ where $`j=|CC_0|`$. So with $`w_j:=_{|C|=i,|CC_0|=j}v_C`$, we have $`_{j=0}^i\alpha _jw_j=0`$ where $`\alpha _i=1`$. We also define vectors $`x_k=_{|AC_0|=k}e_A`$, which are clearly linearly independent for $`0ki`$. Compute for $`0ji`$ (here $`A`$ implicitly has $`|A|=n`$)
$$w_j=\underset{\begin{array}{c}|C|=i\\ |CC_0|=j\end{array}}{}\underset{AC=\mathrm{}}{}e_A=\underset{k=0}{\overset{ij}{}}b_{jk}x_k$$
where the coefficients are obtained by counting, for a fixed set $`A`$ with $`|C_0A|=k`$ and $`kij`$:
$$b_{jk}=|\{C:|C|=i,|CC_0|=j,AC=\mathrm{}\}|=\left(\genfrac{}{}{0pt}{}{ik}{j}\right)\left(\genfrac{}{}{0pt}{}{v+kni}{ij}\right).$$
It is evident by induction that
$$\mathrm{span}\{w_j:ikji\}=\mathrm{span}\{x_j:0jk\}.$$
So $`\{v_C:|C|=i\}`$ are linearly independent.
We write $`T_i=_{j=0}^n\lambda _{ij}E_j`$ be the spectral decomposition of each $`T_i`$. The discussion above shows that if $`|C|=j`$, then $`v_C`$ is contained in $`V_j`$ but not $`V_{j1}`$. Thus $`\lambda _{ij}`$ is the unique scalar so that $`(T_i\lambda _{ij}I)v_CV_{j1}`$. This idea can be used to compute the eigenvalues, but the computations are nontrivial. We refer to \[6, Theorem 9.4.3\] for the Kneser graph $`K(2n+1,n)`$ which is the only one we work out in detail.
###### Lemma 5.4.
The adjacency matrix for the Kneser graph $`K(2n+1,n)`$ has eigenvalues are $`(1)^i(n+1i)`$ with eigenspaces $`W_i`$ for $`0in`$.
###### Theorem 5.5.
If $`T`$ is the adjacency matrix of $`K(2n+1,n)`$, then
$$T_m=S_T|_𝒯=\frac{2^{2n}}{\left(\genfrac{}{}{0pt}{}{2n+1}{n}\right)}=\frac{(4)(6)\mathrm{}(2n+2)}{(3)(5)\mathrm{}(2n+1)}>\frac{1}{2}\mathrm{log}(2n+3).$$
Proof. By Theorem 4.4 and Lemma 5.4,
$`T_m`$ $`=\mathrm{\Delta }(|T|)=\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n}}\right)^1{\displaystyle \underset{i=0}{\overset{n}{}}}(n+1i)\mathrm{Tr}(E_i)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n}}\right)^1{\displaystyle \underset{i=0}{\overset{n}{}}}(n+1i)\left(\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{i1}}\right)\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n}}\right)^1{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{i}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n}}\right)^1{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=0}{\overset{2n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{i}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n}}\right)^12^{2n}={\displaystyle \frac{2^{2n}n!(n+1)!}{(2n+1)!}}`$
$`={\displaystyle \frac{24\mathrm{}(2n)\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}4\mathrm{}(2n)(2n+2)}{24\mathrm{}(2n)\mathrm{\hspace{0.17em}\hspace{0.17em}1}3\mathrm{}(2n1)(2n+1)}}`$
$`={\displaystyle \frac{24\mathrm{}(2n)(2n+2)}{13\mathrm{}(2n1)(2n+1)}}`$
$`={\displaystyle \underset{i=0}{\overset{n}{}}}\left(1+{\displaystyle \frac{1}{2i+1}}\right)>{\displaystyle \frac{1}{2}}\mathrm{log}(2n+3).`$
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# GMOS-IFU spectroscopy of 167-317 (LV2) proplyd in OrionBased on observations obtained at the Gemini Observatory, which is operated by the Association of Universities for Research in Astronomy, Inc., under a cooperative agreement with the NSF on behalf of the Gemini partnership: the National Science Foundation (United States), the Particle Physics and Astronomy Research Council (United Kingdom), the National Research Council (Canada), CONICYT (Chile), the Australian Research Council (Australia), CNPq (Brazil) and CONICET (Argentina).
## 1 Introduction
The Orion Nebula (M42) is the most active site of star formation of the Orion Molecular Cloud and contains many young stars, among them the Trapezium cluster. Its high mass young stellar members generate an intense ultraviolet radiation field that photodissociates and photoionizes the nearby material. These mechanisms promote the appearance of some complex structures as the so-called proplyds (O’Dell & Wen, 1994). Proplyds are low mass YSOs that are being exposed to an intense ultraviolet radiation field which renders them visible. In Orion, there are approximately 160 proplyds (Bally et al., 2000) which are being photoionized mainly by $`\theta ^1`$ Ori C, an O6 spectral type star. There are several objects that have been identified as being proplyds not only near the Trapezium cluster (Bally et al., 2000) but also in other star forming regions (like NGC 3372, Smith et al., 2003). Most of them share the same features: a bow-shaped head that faces the ionization source, a tail, that is primordially directed away from the source, a young star, that may (or not) be visible, and a disk, sometimes seen in silhouette against the HII region (e.g., O’Dell et al., 1993; O’Dell, 1998; Bally et al., 2000; Smith et al., 2005). Early studies of these objects, however, were only able to determine their apparent ubiquity close to the $`\theta ^1`$ Ori C star, their high ionization level and, in later studies, the presence of high velocity structures (outflows) associated with these condensations (e.g., Laques & Vidal, 1979; Meaburn, 1988; Meaburn et al., 1993).
The proplyds are presently explained by a set of models which include a photoevaporated wind, an ionization front and a photoionized wind (Johnstone et al., 1998; Störzer & Hollenbach, 1998; Henney & Arthur, 1998; Bally et al., 1998; O’Dell, 1998; Richling & Yorke, 2000). Johnstone et al. (1998) proposed models in which far ultraviolet radiation (FUV) and extreme ultraviolet radiation (EUV) are responsible for the mass loss rates of the proplyds depending on the distance of the proplyd to the ionization source. For the proplyds situated at intermediate distances from the ionizing star, FUV photons penetrate the ionization front and photodissociate and photoevaporate material from the accretion disk surface, generating a supersonic neutral wind that passes through a shock front before reaching the ionization front as a neutral, subsonic wind. At the ionization front, EUV photons ionize the wind, and the material is then reacelerated to supersonic velocities. The interaction of this supersonic ionized wind with the star wind generates bow shocks seen in H$`\alpha `$ and in \[O III\]$`\lambda `$5007 in some proplyds (Bally et al., 1998). For proplyds closer to the ionizing star, the ionization front reaches the disk surface, and the resulting wind is initially subsonic, becoming supersonic further out.
In general, the models cited above are able to explain the main observational features of the Orion proplyds. The region close to the ionization front is responsible for most of the emission of these objects. Störzer & Hollenbach (1998) have shown that heating and dissociation of H<sub>2</sub> can explain the observed \[O I\]$`\lambda `$6300 emission at the disk surface (Bally et al., 1998). Bally et al. (2000) proposed that neutral gas compressed between the shock front and the ionization front can explain the \[O III\]$`\lambda `$5007 emission seen near the ionization front in some objects. Henney & Arthur (1998) were able to reproduce the observed H$`\alpha `$ intensity profile of the proplyds in Orion in terms of accelerating photoevaporated flows. Numerical 2D HD simulations with a treatment of the radiative transfer (Richling & Yorke, 2000) also reproduce the proplyd morphology and emission, although in this work the treatment of the diffuse radiation is somewhat simplified.
There are several open issues about these objects. One of them is related with the calculation of the mass loss rate of the proplyds, that is strongly model dependent and which may pose severe constraints on the age of the Orion proplyds (as well as for $`\theta ^1`$ Ori C). As the flux of FUV photons dissociates the molecules in the disk of the proplyds, the sub- or trans-sonic *wind* is accelerated at the proplyd’s ionization front (IF) (see Henney et al. 2002 for a discussion), leading to derived mass-loss rates of $`8\times 10^7M_{}`$ yr<sup>-1</sup> which imply a short lifetime for these systems (Churchwell et al., 1987). It is hard to reconcile such short lifetimes with the fact that the region of the Trapezium cluster is populated by several proplyds. Improvements in both models and observational techniques have slightly reduced the calculated mass-loss rates for these systems, and Henney et al. (2002) find $`\dot{M}=8.2\times 10^7M_{}`$ yr<sup>-1</sup> for the 167-317<sup>1</sup><sup>1</sup>1The proplyds mentioned in the present paper will be denoted following the O’Dell & Wen (1994) notation, that is based on the coordinate of the object. (LV2) proplyd, implying an age of less than $`10^5`$ yr for $`\theta ^1`$ Ori C, and $`\dot{M}1.5\times 10^6M_{}`$ yr<sup>-1</sup> for the 170-377, 177-341, 182-413, 244-440 proplyds (e.g., Henney & O’Dell, 1999).
It is presently known that many proplyds show, besides protostellar features such as accretion disks and a young low mass protostar, the presence of jets. Bally et al. (2000), in a survey carried out using the WFPC2 camera on the HST, find 23 objects which appear to have collimated outflows seen as one-sided jets or bipolar chains of bow shocks. Meaburn et al. (2002) also found kinematical traces of jets associated with LV5 (158-323) and GMR 15 (161-307). The proplyd 167-317 shows evidence of the existence of a collimated outflow, detected spectroscopically (Meaburn, 1988; Meaburn et al., 1993; Henney et al., 2002). Until now, this outflow was detected as a one-sided jet, with a P.A. of $`120^{}`$ and a propagation velocity of about 100 km s<sup>-1</sup>. Recently, Smith et al. (2005) and Bally et al. (2005), using the Wide Field Camera of the Advanced Camera for Surveys (ACS/WFC) on HST, found jets associated with silhouette disks in the outer regions of Orion.
In this work, we present the first Integral Field Unity (IFU) observation of a proplyd. The observed object is 167-317, one of the brightest proplyds of the Trapezium. We discuss one of the very first results from the Gemini Multi-Object Spectrograph, in its IFU mode (hereafter, GMOS-IFU), at the Gemini South Telescope. In §2, we present the observations and the steps to obtain a calibrated, clean datacube. In §3 we present the spectral analysis for the observed region of the 167-317 proplyd, as well as the observed lines and intensity maps. We use the observed line profiles to identify the outflows in the system. In §4, we present the discussion and the conclusions.
## 2 The observations and data reduction
The data were taken during the System Verification run of the GMOS-IFU at the Gemini South Telescope (GST), under the GS-2003B-SV-212 program, on 2004 February 26<sup>th</sup> and 27<sup>th</sup>. The science field of view (FOV) is $`3\stackrel{}{\mathrm{.}}5\times 5\mathrm{}`$ with an array of 1000 lenslets of $`0\stackrel{}{\mathrm{.}}2`$. The sky is sampled with 500 lenses which are located $`1\mathrm{}`$ from the science FOV. We used the R831\_G5322 grating in single slit mode, giving a sampling of 0.34Å per pixel ($`15`$ km s<sup>-1</sup> per pixel at H$`\alpha `$) and a spectral coverage between $`5515`$ Å to $`7630`$Å. The spectral resolution (i.e., the instrumental profile) is 47 km s<sup>-1</sup> $``$ FWHM $`63`$ km s<sup>-1</sup> (the FWHM decreasing with increasing wavelength). A field centered on the 167-317 proplyd has been observed with an exposure time of 300s. A 60 s exposure of the standard star the Hiltner 600 has also been taken in order to derive the sensitivity function.
The data were reduced using the standard Gemini IRAF v1.6 <sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. routines. An average bias image has been prepared using the GBIAS task. The extraction has been worked out using the GCAL flat and the response using the twilight flat. Spectra extraction has been performed and wavelength calibration done using an arc taken during the run. The 167-317 field and Hiltner 600 star have been processed in the same way. The sensitivity function derived from the standard star has been applied to the observed field. Finally, data cubes have been built using the GFCUBE routine. We maintain the IFU original resolution of $`0\stackrel{}{\mathrm{.}}2`$ px<sup>-1</sup> although an interpolation was performed in order to turn the IFU hexagonal lens shape into squared pixels.
Because some of the emission lines are very intense, it was not possible to eliminate cosmic rays using the standard GSCRREJ routine of the IRAF Gemini v1.6 package. In order to remove cosmic rays we have then programmed an IDL routine. This routine works directly on the data cube by taking out impacts that are above a certain convenient level calculated based on the mean intensity of nearby pixels. The cosmic ray is then eliminated doing a linear interpolation in the wavelength direction. With this procedure, it is possible to keep very intense lines and eliminate lower level cosmic ray impacts that could be misinterpreted as low emission lines.
Further manipulation of the cube, such as the production of channel maps, emission line maps, dispersion and velocity maps, have been done using Starlink <sup>3</sup><sup>3</sup>3See documentation at this website: http://star-www.rl.ac.uk/, IDL routines and Fortran programs developed specially for this purpose.
The subtraction of the background emission from the proplyd spectrum is a very challenging task (Henney & O’Dell, 1999). This problem arises since the background nebular emission is strongly inhomogeneous on an arc-second scale. In this work, we have used standard $`\chi ^2`$-method in order to fit an intensity versus position plane in a semi-rectangular field located $`1\mathrm{}1\stackrel{}{\mathrm{.}}5`$ away from the peak proplyd emission in each of the individual velocity channel maps. We assure that we only take samples of the background avoiding the region of the proplyd. These planar fits to the nebular emission are then subtracted from the corresponding velocity channel maps. In Figure 1, we show the observed field, the object orientation in the plane of the sky, and also the region that defines our background (region 4, labeled as R4 in Figure 1). Region 4 (see Figure 1) contains 76 spectra which are used to define the background. This Figure also shows an internal box (labeled as R1, or region 1), for which we have defined an arbitrary $`xy`$-coordinate system limited by $`1x11`$ (in the East direction) and $`1y17`$ (in the North direction; see Figure 1). In the following sections, we show maps of spatially distributed physical variables that will be constrained to the limits defined by this box. The spatial positions inside this domain will be defined using the $`xy`$-coordinate system described above. Regions 2 and 3 (labelled as R2 and R3 in Figure 1) define, respectively, a region near the center of the 167-317 proplyd and a region which has a high-velocity, redshifted feature (see below).
## 3 Observational results: spectral line identification and high velocity features
### 3.1 The emission lines intensities and ratios
The lines that we identify in the spectra have already been reported in previous papers of the Orion Nebula (e.g,. Baldwin et al., 2000). We select 4 different regions to measure the line intensities, namely, regions R1, R2, R3 and R4 (see Fig. 1 and §2 for the coordinate system definition).
Four spectra are then obtained by co-adding the spectra of the pixels included in each region. The spectrum of region 1 is the result of the sum of 187 pixels, for region 2 the sum of 6 pixels, region 3 represents the spectrum of a single pixel and region 4 the sum of 76 pixels. A set of 38 lines has been selected from a manual search for emission features in the spectrum integrated over the selected fields. The intensities of the emission lines are determined by fitting a flat continuum and integrating over the whole emission feature <sup>4</sup><sup>4</sup>4The radial velocity range that defines the integration limits in each intensity line determination varies from line to line, since the FWHM changes with the wavelength (see §2).. Then, a mean flux was obtained dividing the integrated flux by the total number of pixels of each region. Table 1 gives the list of the observed lines together with the mean flux of each region (after subtraction of the emission from the background, region R4) and the ratio to H$`\alpha `$ (normalized to H$`\alpha =100`$). The fourth and eigth columns show the background emission and ratios to H$`\alpha `$, respectively, for each line. The mean fluxes are given in units of $`10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> px<sup>-1</sup>, and are not corrected for reddening. The absolute error of the mean flux is given in parentheses. This error was calculated taking into account the variations of the local continuum. The errors can be very large for the weaker lines. Also, the detection of spectral lines in region R3 is affected by the low signal to noise ratio. There are some lines which show stronger mean fluxes in the background, namely, the lines of Si III $`\lambda 5740`$, Si II $`\lambda 5979`$, \[Ni II\]$`\lambda 7378`$ and \[Fe II\]$`\lambda 7155`$ (see Table 1). These lines have low ionization potentials and are expected to appear mainly in regions of lower ionization degree. The lower line intensities within the proplyd indicate that these lines are mainly emitted by the background nebula, and that they are absorbed at least partially by the dust in the proplyd.
In Figure 2, we show the full extracted spectra from the data cube, for regions R1 to R4 as defined in Figure 1. In this figure, the spectra for regions R1, R2, R3 and R4 are depicted from bottom to top in each diagram. The data shown here are not background subtracted. Most of the detected 38 lines can be clearly seen in the spectra. It is also evident that the S/N is higher in regions R1 and R4 (which is due to the increase of the pixel number of these regions). Here, a third order cubic spline was fitted to the continuum and then subtracted from the data. Because of the high order spline polynome subtraction, some minor variations are still present near the H$`\alpha `$ line.
Figure 3 depicts (for region R1), 2D-intensity maps of the lines \[N II\]$`\lambda `$5755, \[N II\]$`\lambda `$6548, H$`\alpha `$, \[N II\]$`\lambda `$6583, \[He I\]$`\lambda `$7065 and \[Ar III\]$`\lambda `$7135, superimposed by line profiles. The data shown here are background subtracted, as described in §2. The line profiles can be better seen in Figure 4, in which we show the mean line profile for each line and for each of the regions R1, R2, R3 and R4, in the same diagram. As could be anticipated, the flux is more intense in region R2 (dotted line), where the proplyd is located, except for the \[N II\]$`\lambda 6548`$ line, which shows a stronger flux in region R3. For the H$`\alpha `$ and \[N II\]$`\lambda 6583`$ lines the emission of region R3 is more intense than the emission of region R1. For lines with good S/N, the presence of a redshifted feature starts to become clear. An analisys of the profile components will be given in more detail below.
In order to see how the spectra of the different lines change spatially, in Figure 5 we present an intensity versus position plot for the same lines shown in Figures 3 and 4. We show how the intensity of these lines changes along a diagonal line, that crosses region R1 passing through the pixel (3,4) (region R3) and through the centre of region R2. We note that there are 2 intensity maxima for the H$`\alpha `$, \[N II\]$`\lambda 6583`$ and \[N II\]$`\lambda 6548`$ lines, clearly showing in which lines region R3 is clearly visible (also see Figure 3).
### 3.2 High velocity features
We associate region R3 with the redshifted jet of the proplyd, although in Figure 9 we will show that the emission of the jet is not restricted only to this pixel. A redshifted jet associated with the 167-317 proplyd was previously detected, with a propagation velocity of $``$ 100 km s<sup>-1</sup> (e.g., Henney et al., 2002). We have also detected this high velocity feature in our data. In order to see this, we examine the behaviour of the moments of the radial velocity distribution. These are the flux-weighted mean radial velocity $`<v>`$, the flux-weighted rms width of the line $`\mathrm{\Delta }v^2`$ and the skewness $`\mathrm{\Delta }v^3`$, that are given, respectively, by:
$$<v>=\frac{1}{I}I_vv𝑑v$$
(1)
$$\mathrm{\Delta }v^2=\frac{1}{I}I_v(v<v>)^2𝑑v$$
(2)
$$\mathrm{\Delta }v^3=\frac{1}{I}I_v(v<v>)^3𝑑v$$
(3)
where
$$I=I_v𝑑v$$
(4)
The integrated intensity (see eq. 4 and Figure 3) shows that if we move away from the proplyd position (R2), the emission for several lines (for example, H$`\alpha `$ and \[Ar III\]) rapidly drops to low values (close to the background value). We also note that the maps for the flux-weighted mean radial velocity <sup>5</sup><sup>5</sup>5Unless explicitly mentioned, the radial velocities presented here are not corrected for the systemic radial velocity, that is of the order of $`v_{}26`$ km s<sup>-1</sup> (e.g., Meaburn et al., 1993). The error in the radial velocity measurements is of the order of 2.5 km s<sup>-1</sup>. (not shown here; see eq. 1) show that $`<v>=50\pm 20`$ km s<sup>-1</sup>.
The flux-weighted rms width of the line $`(\mathrm{\Delta }v^2)^{1/2}`$ and the skewness $`(\mathrm{\Delta }v^3)^{1/3}`$ moment of the radial velocity distribution (see equations 2 and 3, respectively) are depicted in Figure 6, for the H$`\alpha `$ (left) and \[Ar III\] (right) lines. Both of them were computed after carrying out both background and continuum subtractions. From Figure 6, we see that there is a clear enhancement of both of these moments from the proplyd position (R2 region) towards the R3 region (i.e., towards the SE direction, at a PA $`135^{}`$). This behaviour of the moments leads us to infer that: 1) there is an increase in the line width as we go from the proplyd position to the SE direction; and 2) that there are redshifted wings in the line profiles. These results indicate the presence of a redshifted outflow in this region. We note that the same behaviour is seen in the He I $`\lambda `$6678.15, \[N II\]$`\lambda `$6583.46 lines (not shown here).
In order to confirm the presence of high velocity components in the observed field, as well as to see the behaviour of the photoevaporated proplyd flow and the background emission, we have computed three-component Gaussian minimum $`\chi ^2`$ fits for each position-dependent line profile. The profiles of all emission lines have a major peak at $`v_{rad}=4060`$ km s<sup>-1</sup>, the exact radial velocity of this peak changing with spatial position. In some spectral lines, there is an evident second peak at redshifted velocities ranging from 100 to 150 km s<sup>-1</sup>, corresponding to the jet associated with the proplyd. Finally, we have been able to detect a blueshifted component, which is fainter than the redshifted components. As an example of our three-component Gaussian fit, in Figure 7 we show the data (full line) and the fit (crosses), for the H$`\alpha `$ and \[Ar III\]$`\lambda `$7135 lines in the ($`x`$,$`y`$)=(6,7) position. In this figure, for each emission line, the top-left panel represents the data and the three Gaussian fits, the top-right panel shows the main, low velocity component (and the Gaussian fit for this component), the bottom-left panel depicts the data minus the main component together with the fits for the blue- and red-shifted components and, finally, the bottom-right panel shows the residual, obtained by subtracting the three-component fit from the observed line profile. The fits depicted in this figure show the presence of redshifted and blueshifted components. The redshifted component is present in several lines in pixels around the SE direction, as already mentioned before. On the other hand, blueshifted emission can be found in several pixels to the NW of the proplyd peak emission as can be seen in the left bottom panel of Figure 7 and this could represent the first spectroscopic determination of the presence of a blueshifted counter-jet, with systemic velocities of -60 km s<sup>-1</sup> $`v_{rad}90`$ km s<sup>-1</sup>. This blueshifted component is also suggested by the fitted profiles of HeI $`\lambda `$7065 (not shown here).
Figure 8 depicts the spatially distributed intensity over the R1 region (see Figure 1), of the main peak for the \[N II\]$`\lambda `$5754 (top-left), \[N II\]$`\lambda `$6548 (top-middle), H$`\alpha `$ (top-right), \[N II\]$`\lambda `$6583 (bottom-left), HeI $`\lambda `$7035 (bottom-middle) and \[Ar III\]$`\lambda `$7135 (bottom-right) lines. The values shown in this Figure were obtained from three-component Gaussian fits to the observed line profile. All of the lines peak at the proplyd position, at ($`x`$,$`y`$) = (6,7), inside the R2 region, and three of them (namely \[N II\]$`\lambda `$6548, top-middle; H$`\alpha `$, top-right; and \[N II\]$`\lambda `$6583, bottom-left), have a secondary, less intense peak at ($`x`$,$`y`$)=(3,4) (i.e., in the R3 region; see Figure 1 for the definition of the coordinate system). This secondary peak is probably related to the jet, since, as we have seen before, the jet propagates from the R2 region towards the SE direction. It is also interesting to note that there is a tail of faint emission (compared with the maximum in a given map) that extends in the NE direction, and that can be seen in both of the \[N II\] lines that bracket H$`\alpha `$ (top-middle and bottom-left maps in Figure 8).
In Figure 9, we show the spatially dependent intensity (Fig. 9a) and the central velocity (Fig. 9b) of the redshifted component of the fitted H$`\alpha `$ profile. We limit the maps to the spatial pixels which have a skewness $`|(\mathrm{\Delta }v^3)^{1/3}|>30`$ km s<sup>-1</sup>. The pixels which satisfy this criterium have a well defined high velocity, redshifted component. These figures show that the high velocity component has an intensity maximum near the center of the proplyd, with an extension towards the SE direction (along the jet axis), surrounded by a region of decreasing fluxes. The high intensity spike seen in H$`\alpha `$, \[N II\]$`\lambda 6548`$ and \[N II\]$`\lambda 6583`$ lines (Figure 3) can be seen here. The jet shows a spread in velocity values, although most of the pixels present velocities around 120-140 km s<sup>-1</sup>. This figure also shows that the jet velocity decreases with increasing distance from the proplyd. This trend can be seen in the Figure 9c, where we show the radial velocity of the redshifted component as a function of distance from the proplyd position, ($`x`$,$`y`$) = (6,7) in region R1. To build this figure, we assume that the jet is propagating in a PA $`135^{}`$ position angle. We define a box with one axis aligned with the jet axis, and the second, perpendicular dimension extending two pixels to each side of the jet axis. We then take averages (of the mean velocity of the redshifted component) perpendicular to the jet direction and plot the resulting velocity as a function of distance from the proplyd (see the bottom frame of Figure 8). There is an indication that the jet velocity is slowly diminishing as a function of distance from the source.
### 3.3 The spatial distribution of the nitrogen ratio
In Figure 10a, we show a map for the R1 region of the line ratio:
$$\frac{[\mathrm{N}\mathrm{II}]\lambda 6548+[\mathrm{N}\mathrm{II}]\lambda 6583}{[\mathrm{N}\mathrm{II}]\lambda 5754}=\frac{A_{6548}\sigma _{6548}+A_{6583}\sigma _{6583}}{A_{5754}\sigma _{5754}},$$
(5)
where $`A`$ and $`\sigma `$ are the height and the dispersion of the fitted Gaussian profile (of the main, low velocity component). The \[NII\](6548+6583)/5754 ratio is a classical electron temperature diagnostic of low/medium density nebulae ($`N_e<10^5`$ cm<sup>-3</sup>). At high densities, due to the collisional desexcitations, this ratio cannot be used for the determination of the electron temperature, but it turns out to be a quite good electron density diagnostic. This can be seen, for example, from the Osterbrock (1989) equation:
$$[NII](6548+6583)/5754=\frac{6.91\mathrm{exp}(2.5\times 10^4/T_e)}{1+2.5\times 10^3(N_e/T_e^{1/2})}$$
(6)
Figure 10b shows the lines corresponding to \[NII\](6548+6583)/5754 = 9, 10, 11, 80, 90, 100 in a $`T_eN_e`$ plane. We can see that for the values obtained in our case (line ratios close to 10), the \[NII\](6548+6583)/5754 ratio is strongly dependent on the electron density, and that the dependence on the electron temperature is small.
In Figures 10c and 10d, we show the densities computed from the \[NII\](6548+6583)/5754 ratio for $`T=10^4`$ K and $`T=1.5\times 10^4`$ K, respectively. Although presenting a very complicated pattern, these figures show that the region coinciding with the proplyd center, and extending towards the SE direction (the jet propagation direction), shows the lowest line ratio, indicating electron densities ($`n_e`$) of the order of $`2\times 10^6`$ cm<sup>-3</sup>. This value is at least in qualitatively agreement with the $`n_0=(3.0\pm 0.5)\times 10^6`$ cm<sup>-3</sup> obtained by Henney et al. (2002) considering a $`T=1.2\times 10^4`$ K temperature for the ionization front (that is inside the R2 region). We can also note that both figures are similar, presenting few differences in the density values.
The extended structure from the R2 region towards the SE direction is surrounded by a region in which the nitrogen line ratio reaches values of up to $`30`$. In this limit, $`n_e<10^5`$ cm<sup>-3</sup>. It is interesting to note that in the clump associated with the redshifted jet, the nitrogen line ratio increases substantially, indicating a decrease in the electron density when compared with the same values near the center of the proplyd. We have also obtained the \[S II\]$`\lambda 6716,\lambda 6730`$ emission maps. Unfortunately, we are not able to use these lines as a diagnostic because neither the ratio of this doublet is a good electron density indicator for this case (since it is constant for densities higher than 10<sup>5</sup> cm<sup>-3</sup>) nor do we have a good enough signal to noise ratio for these lines. Furthermore, the background subtracted spectra for these two lines give us negative fluxes in the proplyd region (see Table 1). The presence of these negative features in the line profiles is probably related with the presence of dust, as discussed by Henney & O’Dell (1999).
## 4 Discussion and conclusions
We have presented in this work the first Integral Field Unit (IFU) spectroscopic observations of a photoevaporating disk immersed in a HII region. In particular, we have taken advantage of the System Verification Run of the Gemini South Telescope Multi Object Spectrograph (GMOS) to obtain spectra of the 167-317 proplyd in the Orion nebula. The 167-317 proplyd, also known as LV2 (from the pioneering work from Laques & Vidal, 1979), is one of the brightest and best studied proplyds. In a single exposure, we took 400 spectra with a spatial resolution of $`0\stackrel{}{\mathrm{.}}2`$. These spectra have been combined in order to optimize this instrumental feature, and, before discussing of our results, we want to make a few comments concerning the potential use of such IFU data-cubes for these objects. As previously discussed in the literature (Henney et al., 2002), the subtraction of the background emission is a very challenging task. This problem arises since the background near these proplyds (which are immersed in a highly non-homogeneous, photoionized ambient medium) is very complex and changes in an arc-second spatial scale. A careful computation of the background contribution to the spectra of these objects is one of the most important tasks in order to have reliable radial velocity measurements for both the proplyd photoevaporated flow and the high velocity features (jets) present in these systems. Here, we have used the 76 spectra, defined in Figure 1 as region R4, in order to obtain a planar fit for the background emission. We think that such a definition of the background emission is more precise, since we use several regions to calculate a better function to describe the background. One more thing to point out is that, with spatially distributed spectra, we are able to confidently separate the contribution of the background from the contribution of the object itself.
The data from the IFU observation of the 167-317 proplyd also allow us to investigate the spatially dependent properties of the outflows associated with this object. In particular, the 167-317 proplyd is known to have a redshifted, collimated jet that propagates with heliocentric velocities of $`100`$ km s<sup>-1</sup> towards the SE direction, and with a spatial extension of $`2\mathrm{}`$ (e.g, Meaburn, 1988; Meaburn et al., 1993; Massey & Meaburn, 1995; Henney, 2000; Henney et al., 2002). We find that a prominent, high velocity redshifted component can be detected in some emission lines, particularly in H$`\alpha `$, in this SE direction. The redshifted jet has a trend in its radial velocity (as a function of distance from the proplyd) with higher velocities close to the proplyd and lower velocities at increasing distances. We note that Henney et al. (2002) had previously suggested a variation in the jet velocity as a function of distance from the source.
There is a subtle peak (in several emission lines, particularly in \[N II\]$`\lambda `$5754 and H$`\alpha `$) that is associated with the SE jet emission, and located at $`0\stackrel{}{\mathrm{.}}67`$ SE of the center of the proplyd \[at (x,y)=(3,4); the R3 region in Figure 1\]. For this intensity peak, and also using the H$`\alpha `$ profiles of the neighbouring pixels, we find a mean heliocentric velocity $`v_{red}=116\pm 10`$ km s<sup>-1</sup>. It is interesting to note that previous HST images and spectroscopic analyses (Bally et al., 2000; Henney et al., 2002) also reveal a spike in the jet emission at $`0\stackrel{}{\mathrm{.}}4`$ from the proplyd cusp. If we associate these two emission regions as belonging to the same condensation, we can estimate a proper motion of $`140`$ km s<sup>-1</sup>, which combined with the $`116`$ km s<sup>-1</sup> radial velocity give a full jet velocity of $`180\pm 90`$ km s<sup>-1</sup> (where the error was inferred from the spatial resolution of our data). This velocity is of the same order as the values inferred for Herbig-Haro jets associated with T Tauri stars. However, the exact nature of this spike is unknown. We could relate this spike with a bow shock of the jet, since the line profiles for both are similar (Hartigan et al., 1987; Beck et al., 2004), showing a low velocity, more intense peak together with a high velocity, less intense peak (similar to the line profiles seen in Figures 3, 4 and 7).
We have found evidence for a blueshifted component in several emission lines around the northwest tip of the proplyd position. In particular, the H$`\alpha `$ profiles reveal the presence of a blue-shifted component with systemic velocity $`v_{blue}=(75\pm 15)`$ km s<sup>-1</sup>, which indicates that its velocity is similar though smaller than the one of the redshifted jet. However, the intensities of the blue-shifted features are much lower than the intensity of the redshifted components associated with the SE jet: the counter-jet is at least four times less intense than the redshifted jet (in H$`\alpha `$). The presence of such a blue-shifted component very close to the LV2 proplyd, together with previously reported evidence that (at several arcseconds to the NW of this proplyd; see Massey & Meaburn 1995) there is a blueshifted component (in the \[O III\]$`\lambda `$5007 line) strongly suggests the existence of a faint counterjet.
It is interesting to estimate the mass loss rate for the proplyd and the associated redshifted jet. In order to do this, we consider a simple model that assumes that the proplyd flow arises from a hemispherical, constant velocity wind that originates at the proplyd ionization front, at a radius $`r_0`$. Using the H$`\alpha `$ luminosity $`L_{H\alpha }`$, we can then obtain the particle density $`n_0`$ at this point and the mass loss rate, from:
$$n_0=\left(\frac{2L_{H\alpha }\mathrm{e}^{\tau _{H\alpha }}}{4\pi \alpha _{H\alpha }h\nu _{H\alpha }r_0^3}\right)^{1/2},$$
(7)
$$\dot{M}=4\pi r_0^2n_0\mu m_Hc_{II},$$
where $`\mathrm{e}^{\tau _{H\alpha }}`$ accounts for the extintion correction, $`\alpha _{H\alpha }=5.83\times 10^{14}`$ cm<sup>3</sup> is the effective recombination coefficient for H$`\alpha `$, $`h\nu _{H\alpha }`$ is the energy of the H$`\alpha `$ transition, $`\mu =1.3`$ is the mean molecular weight and $`c_{II}=10`$ km s<sup>-1</sup> is a typical value for the sound velocity in an HII region. The extintion correction is obtained from the base 10 logarithm of the extintion at $`H\beta `$ ($`c_{H\beta }`$) which for the proplyd 167-317 is equal to 0.83 (O’Dell, 1998) through the relation, $`c_{H\beta }=K\tau _{H\alpha }`$ (O’Dell et al., 1992), where here we take $`K`$ = 0.56. Assuming for the ionization front radius a value $`r_0=(7.9\pm 0.2)\times 10^{14}`$ cm (Henney et al., 2002) we obtain $`n_0=(2.3\pm 0.6)\times 10^6`$ cm<sup>-3</sup> and $`\dot{M}=(6.2\pm 0.6)\times 10^7`$ M year<sup>-1</sup>. Henney et al. (2002) have obtained a $`\dot{M}=8.2\times 10^7`$ M year<sup>-1</sup> $`\pm 10\%`$, and the small differences between both results may arise from the different observational techniques employed in both cases.
For the redshifted jet, we assume that the outflow arises from the blob located in region R3 (or pixel 3,4). In this case, the particle density, mass and mass loss rate of the jet are given by,
$$n=\left(\frac{3L_{H\alpha }}{4\pi \alpha _{H\alpha }h\nu _{H\alpha }r_B^3}\right)^{1/2},$$
(8)
$$M=\frac{4\pi }{3}r_B^3n\mu m_H,$$
(9)
$$\dot{M}M\frac{v_j}{L_j},$$
(10)
where $`r_B=6.7\times 10^{14}`$ cm is the radius of the blob, $`v_j=180`$ km s<sup>-1</sup> and $`L_j=4.5\times 10^{15}`$ cm are the propagation velocity and length of the jet, respectively, derived above. With these values, we obtain for the jet a mass loss rate equal to $`(2.0\pm 0.7)\times 10^8`$ M year<sup>-1</sup>, which is similar to the mass loss rate of typical HH jets from T Tauri stars.
We have been able to construct a (\[NII\]$`\lambda 6548`$ \+ \[NII\]$`\lambda 6583`$)/\[NII\]$`\lambda 5754`$ line ratio map. This ratio indicates densities higher than $`10^5`$ cm<sup>-3</sup> for the proplyd emission region, consistent with values previously reported in the literature. We find a subtle enhancement of this ratio in the region of the redshifted jet. Further observations of electron temperature and density diagnostic lines would allow a more accurate determination of the spatial dependence of these important parameters.
We thank the anonymous referee for his/her comments and suggestions, which have improved significatively the presentation of the paper. We are greatful to W. Henney for enlightning discussions. M.J.V. and A.H.C. would like to thank the staff of the Instituto de Ciencias Nucleares, UNAM, México, for their kind hospitality, as well as for partial financial support during our visit. We would like to thank Bryan Miller, Tracy Beck and Rodrigo Carrasco from the Gemini Observatory, and Bruno Castilho from LNA/CNPq, for help us in the early stages of IFU-data processing. The work of M.J.V., A.H.C. and H.P. was partially supported by the Bahia state funding agency (FAPESB), by the Universidade Estadual de Santa Cruz (UESC) and by the Milenium Institute (CNPq Project n. 62.0053/01-1-PADCT III/Milênio). This work was partially funded by the CONACYT grants 41320-F and 43103-F and the DGAPA (UNAM) grants IN112602 and IN111803.
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# Entanglement generation in harmonic chains: tagging by squeezing
## I Introduction
The development of reliable strategies for quantum communication and information transfer has gained, in these recent years, an increasing importance in the quantum information processing (QIP) panorama. The reliable implementation of a quantum channel for the exchange and distribution of information is indeed central in many potential QIP applications varie0 ; varie1 . Intuitively, one could think about a scenario in which the quantum channel and the processing device are two different entities which have to be interfaced at the right time of a given protocol. This implies the ability of switching on and off the interfacing interaction with sufficient degree of accuracy and a reliable control at the single-qubit level, which are very demanding requirements in general. On the other hand, the idea of exploiting collective interactions of intrinsically multipartite systems, governed by external potentials which globally address the entire register, has very recently encountered the interest of the QIP community sougato . A global addressing scheme offers advantages in terms of controllability of the device and protection from the decoherence channels unavoidably opened by any sort of local external intervention.
Inspired by the progresses performed in the design, coupling and management of bosonic nanostructures, which can behave quantum mechanically buks , important efforts have been produced in order to better understand the role that multipartite systems of coupled bosons have in the transfer and propagation of quantum information plenio ; iohelenmyung . The application of global addressing techniques to systems of continuous variable (CV) bosonic systems is appealing for a less demanding implementation of CV quantum information processing.
In this work, we re-consider the issue of entanglement generation in a chain of harmonic oscillators coupled through nearest-neighbor spring-like forces induced by an external potential which addressed the whole system. One of the points of interests in our analysis is the role played by counter-rotating terms (present in the interaction Hamiltonian) in the entanglement generation process. This point, anticipated by the studies in plenio , is analyzed here by a change of perspective. Instead of solving by brute force the dynamical equations ruling the evolution of the bosonic register, we look for a formal decomposition of the time evolution operator in terms of linear optics elements, following the successful route initiated in iohelenmyung . We believe that this alternative approach clarifies the entanglement dynamics within the register and provides a more transparent picture of the role of the counter-rotating terms in such a process. Entanglement is found to be always present if the counter-rotating terms are included in the interaction Hamiltonian. However, we find the degree of bipartite entanglement between the first and last oscillator to be very small (a feature which is evident, despite it has not been stressed, in the analyses in plenio ). In order to quantitatively improve the entanglement settled between the ends of an open chain, we design a strategy based on proper initialization of the register (performed by locally acting on the state of the extremal oscillators only) and global addressing, following the same lines depicted in quantum state transfer protocols qst . We show how, physically, this improvement is possible because of the symmetry properties of the bosonic system.
The reminder of the manuscript is organized as follows. In Section II we introduce the interaction model here at hand. We discuss the technical tools used in order to derive effective decompositions of the time-evolution operator into linear optics operations. Effective all-optical setups can be introduced, which provide a visual picture of the evolution of an $`N`$-element register and we give an explicit example for a simple case. In Section III, the entanglement generated in an open chain is quantified by means of the corresponding equivalent decompositions. We show that, as long as only the quantum correlations generated by the counter-rotating terms alone are considered, end-to-end entanglement in the chain is not favoured. Strong quantum correlations, which never disappear, are found between the first and the second oscillator in the chain. On the other hand, the entanglement between the first and the last oscillator is always very weak. A transparent physical interpretation of the time delay with which entanglement appears in the first-last oscillators subsystem is possible through the analysis of the corresponding all-optical setup. Section IV addresses a way to improve the results discussed in Section III. By considering the physical system as a fictitious two-terminal quantum black box, we show that simple local pre-squeezing of the first and last element of the channel allows us to obtain several interesting effects. The end-to-end degree of entanglement can be quantitatively improved and any other bipartite $`1j`$ quantum correlation ($`j=2,..,N1`$) can be correspondingly suppressed.
## II The model and the effective decomposition
We consider $`N`$ oscillators labelled by $`j[1,N]`$ and arranged in an open linear chain. The coupling between the oscillators is provided by a nearest-neighbor spring-like force settled by an external potential. By including the free dynamics of each harmonic oscillator, the corresponding Hamiltonian reads
$$\widehat{H}_{chain}=\frac{\omega }{2}\underset{j=1}{\overset{N1}{}}\left(\widehat{q}_j^2+\widehat{p}_j^2\right)+\kappa \underset{i=1}{\overset{N1}{}}\widehat{q}_j\widehat{q}_{j+1},(\mathrm{}=1)$$
(1)
with $`\widehat{q}_j=(\widehat{b}_j+\widehat{b}_j^{})/\sqrt{2}`$ and $`\widehat{p}_j=i(\widehat{b}_j^{}\widehat{b}_j)/\sqrt{2}`$ the position and momentum quadrature operators of the $`j`$-th oscillator respectively and $`\widehat{b}_j`$ ($`\widehat{b}_j^{}`$) the corresponding annihilation (creation) operator barnett . The coupling rates $`\kappa `$’s are real and time-independent. A sketch of the interaction configuration is provided in Fig. 1.
The form of the coupling terms deserves some comments as it is straightforward to see that each $`\kappa \widehat{q}_j\widehat{q}_{j+1}`$ in Eq. (1), expressed by means of the annihilation and creation operators, reads
$$\kappa \widehat{q}_j\widehat{q}_{j+1}=\frac{\kappa }{2}(\widehat{b}_j\widehat{b}_{j+1}+\widehat{b}_j\widehat{b}_{j+1}^{}+\widehat{b}_j^{}\widehat{b}_{j+1}+\widehat{b}_j^{}\widehat{b}_{j+1}^{}).$$
(2)
Eq. (2) includes co-rotating terms ($`\widehat{b}_j\widehat{b}_{j+1}^{}+h.c.`$) as well as counter-rotating (CR) terms ($`\widehat{b}_j\widehat{b}_{j+1}+h.c.`$barnett .
In this paper, we treat the CR terms on the same footage as the co-rotating ones, analyzing their relevance in entanglement generation processes in interacting bosonic systems. In order to analyze the time-evolution of the chain, we look for an effective decomposition of the propagator $`\widehat{U}(t)=e^{i\widehat{H}_{chain}t}`$ in terms of linear-optics elements. We order the quadrature-operators as $`\widehat{𝐱}=(\widehat{q}_1,\widehat{q}_2,..,\widehat{q}_N,\widehat{p}_1,\widehat{p}_2,..,\widehat{p}_N)^T`$ and divide $`\widehat{H}_{chain}`$ as $`\widehat{H}_{chain}=\widehat{H}_{chain}^p+\widehat{H}_{chain}^q`$, where $`\widehat{H}_{chain}^q`$ ($`\widehat{H}_{chain}^p`$) involves only the $`q`$-part ($`p`$-part) of $`\widehat{𝐱}`$. In matrix form
$$\widehat{H}_{chain}^q=\frac{1}{2}\left(\begin{array}{ccccccc}\omega & \kappa & 0& 0& \mathrm{}& 0& 0\\ \kappa & \omega & \kappa & 0& \mathrm{}& 0& 0\\ 0& \kappa & \omega & \kappa & \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \omega & \kappa \\ 0& 0& \mathrm{}& \mathrm{}& \mathrm{}& \kappa & \omega \end{array}\right).$$
(3)
This is a tridiagonal matrix whose dimension depends on the number of elements in the register. The formal diagonalization of $`\widehat{H}_{chain}^q`$ guides us in expressing the $`q`$-part of Eq. (1) in a picture defined by eigen-operators which are linear superpositions of $`\widehat{q}_j`$’s. On the other hand, $`\widehat{H}_{chain}^p`$ is already diagonal in the $`\widehat{𝐱}`$ basis and its form is not changed by orthogonal transformations. Therefore, we discard this part of Eq. (1) from our explicit analysis and will include it only when necessary.
The simple form of Eq. (3) allows for an efficient diagonalization tridiagonal , which helps us in identifying a proper pattern of coupling operations for the decomposition of $`\widehat{U}`$. In order to keep our analysis general, we will refer to the well-known beam-splitter (BS) operator $`BS_{jk}(\theta )=\mathrm{exp}[i\theta (\widehat{q}_j\widehat{p}_k\widehat{p}_j\widehat{q}_k)]`$ of its reflectivity $`\mathrm{sin}^2\theta `$ as a coupler operator because this term can be used for both optical fields and mechanical oscillators. In the eigen-operator basis we write
$$\widehat{H}_{chain}^{q,N}=\underset{j=1}{\overset{N}{}}E_j^N(\widehat{O}_j^N)^2,$$
(4)
where $`E_j^N`$’s are the eigen-frequencies of Eq. (3) and $`\widehat{O}_j^N=_{k=1}^N\alpha _{jk}^N\widehat{q}_k`$ ($`j=1,..,N`$) are the corresponding eigen-operators, expressed as normalized superpositions of the $`\widehat{q}_k`$ quadratures with coefficient $`\alpha _{jk}^N`$. The set $`\{E_j^N,\widehat{O}_j^N\}`$ is parameterized by the dimension $`N`$ of the chain. As an explicit example, we consider the first non-trivial case represented by an open chain of $`N=3`$ where we have $`\{E_1^3,E_2^3,E_3^3\}=\{\omega /2,(\omega +\sqrt{2}\kappa )/2,(\omega \sqrt{2}\kappa )/2\}`$. We introduce the matrix of coefficients $`\alpha _{jk}^3`$
$$\alpha ^3=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}1& 0& 1\\ \frac{1}{\sqrt{2}}& +1& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}& 1& \frac{1}{\sqrt{2}}\end{array}\right).$$
(5)
Thus, the spectrum of $`\widehat{H}_{chain}^{q,3}`$ is a ladder, symmetric with respect to the bare eigen-frequency $`\omega /2`$. This is a general result for an odd number of oscillators: by increasing the dimension of the register, the spectrum of $`\widehat{H}_{chain}^{q,2l+1}`$ ($`l`$) is never degenerate ans consists of $`l`$ different pairs of frequencies symmetrically shifted with respect to $`\omega /2`$.
Coming back to our example, the structure of the eigen-operators $`\widehat{O}_j^3`$ ($`j[1,3]`$) turns out to be very informative in the research for a set of operations which can be used, starting from Eq. (3), in order to get the diagonal form (4). Indeed, Eq. (5), suggests that a $`50:50`$ coupler operator $`BS_{13}`$ simplifies the structure of the coupling terms leaving us with oscillator $`1`$ being decoupled from the dynamics of the rest of the register. Then, a $`50:50`$ $`\widehat{B}_{23}`$ operation will complete the diagonalization of $`\widehat{H}_{chain}^q`$:
$$\widehat{H}_{chain}^q\stackrel{\widehat{B}_{23}\widehat{B}_{13}}{}\frac{\omega }{2}\widehat{q}_1^{{}_{}{}^{\prime \prime }2}+\frac{\omega \sqrt{2}\kappa }{2}\widehat{q}_3^{{}_{}{}^{\prime \prime }2}+\frac{\omega +\sqrt{2}\kappa }{2}\widehat{q}_2^{{}_{}{}^{\prime \prime }2}$$
(6)
with $`\widehat{q}_j^{^{\prime \prime }}`$ which are the new quadratures after BS’s to be put in correspondence with $`O_j^3`$’s. The matching with Eq. (4) is evident. Thus, the explicit diagonalization procedure of the $`q`$-part of the chain’s Hamiltonian gives us information about the pattern of BS operations which have to be applied to the bare expression Eq. (3) in order to get Eq. (4).
In Fig. 2 we provide the sequence of BS’s to apply in order to diagonalize the interaction Hamiltonian for $`N=3,\mathrm{\hspace{0.17em}4}`$ and $`5`$ (panels $`(𝐚)`$, $`(𝐛)`$ and $`(𝐜)`$ respectively). Straight lines represent $`50:50`$ BS’s, while curved ones stand for unbalanced BS’s. Some comments are in order. First, there is a striking difference between the even and odd number of oscillators. For the odd number case, up to $`N=7`$, a single unbalanced BS is required in order to diagonalize $`H_{chain}^q`$. On the other hand, the BS pattern for the even number case appears to be more complicated already for $`N=4`$, which is the first non-trivial configuration with even $`N`$, involving two unbalanced BS’s between the pairs of oscillators $`(1,2)`$ and $`(3,4)`$. A second important difference between chains of opposite parity will be highlighted later on. We stress that there might be other, inequivalent ordering of coupler operations which diagonalize the interaction Hamiltonian. The choices presented in Figs. 2 $`(𝐚)`$ and $`(𝐜)`$ allow us to write the interaction part of the chain Hamiltonian in the special diagonal form corresponding to the association between oscillators and eigen-frequencies shown in the bottom part of Fig. 2. There, the set $`\{E_j^N\}`$ is written as $`\{\omega /2,_1^N,..,_{N1}^N\}`$ with the subset $`\{_j^N\}`$ arranged in increasing order of frequencies (with $`\kappa >0`$commento .
After the action of the collective coupler operator $`\widehat{B}_{coll}`$, which collects the pattern (for a given $`N`$) discussed above, the total Hamiltonian of the chain reads $`\widehat{H}_{chain}^q=(\omega /2)[(\widehat{O}_1^N)^2+_{j=1}^N(\widehat{P}_j^N)^2]+_{j=1}^{N1}_j^N(\widehat{O}_j^N)^2`$. Here, $`\{\widehat{P}_j^N\}`$ is the new set of momentum quadrature operators determined by the application of the coupler operations to $`\{\widehat{p}_j\}`$. By specializing the discussion to the odd number of oscillators, the next step in order to find the decomposition of $`\widehat{U}(t)`$ is the introduction of proper operations (acting on the elements of the register, oscillator $`1`$ exluded) which balance the differences between $`\omega /2`$ and $`_j^N`$. In the case of even $`N`$, these operations would involve the entire set of oscillators, without exclusions, as the eigen-spectrum of $`\widehat{H}_{chain}^{q,2l}`$ ($`l`$) does not include the bare frequency $`\omega /2`$. Conceptually, this balancing is an important step as it would allow us to look at the register as a set of new non-interacting harmonic oscillators. It is immediate to recognize that this is possible through the use of single-oscillator squeezing $`\widehat{S}_j(s_j)=\mathrm{exp}[\frac{i}{2}s_j(\widehat{O}_j^N\widehat{P}_j^N+\widehat{P}_j^N\widehat{O}_j^N)]`$, which realizes $`\widehat{O}_j^Ne^{s_j}\widehat{O}_j^N,\widehat{P}_j^Ne^{s_j}\widehat{P}_j^N`$. We can thus formally write $`\widehat{S}_{coll}^{}\widehat{B}_{coll}^{}\widehat{U}(t)\widehat{B}_{coll}\widehat{S}_{coll}=\widehat{R}_{coll}(t)_{j=2}^N\widehat{R}_j(\varphi _j(t))`$, where $`\widehat{R}_j(\varphi _j)=e^{i\varphi _j[(\widehat{O}_j^N)^2+(\widehat{P}_j^N)^2]}`$ is the phase-space representation of a rotation operator of its angle $`\varphi _j`$ and $`\widehat{S}_{coll}=_{j=2}^N\widehat{S}_j(s_j)`$. By inverting the above relation, we find the formal expression
$$\widehat{U}(t)=\widehat{B}_{coll}\widehat{S}_{coll}\widehat{R}_{coll}(t)\widehat{S}_{coll}^{}\widehat{B}_{coll}^{}.$$
(7)
Analogously to the case of squeezing, the formal collective rotation involves all the oscillators but the one associated to the bare eigen-frequency $`\omega /2`$, which again is a specific feature of the odd $`N`$ case. Moreover, the rotations $`\widehat{R}_j(\varphi _j(t))`$’s contain the entire time dependence of $`\widehat{U}(t)`$. The balancing induced by the squeezing operations imposes, in general, time-independent conditions as it relates the squeezing parameters $`s_j`$’s to the elements of the set $`\{_j^N\}`$. On the contrary, the formal identification of $`\widehat{S}_{coll}^{}\widehat{B}_{coll}^{}\widehat{U}(t)\widehat{B}_{coll}\widehat{S}_{coll}`$ with $`\widehat{R}_{coll}`$ imposes that the rotation angles $`\varphi _j`$’s carry an explicit time dependence. As an example, we consider again $`N=3`$, where we have that
$$s_{2,3}=\frac{1}{4}\mathrm{ln}(2_{2,1}^3/\omega ),\varphi _{2,3}=\frac{t}{2}\sqrt{2_{2,1}^3\omega }.$$
(8)
The decomposition Eq. (7) is a central result of our study. It allows us to provide a clear physical picture of the dynamics occurring within the linear chain, without explicitly solving the dynamical equations of motion of the oscillators plenio . Indeed, once the explicit form of $`\widehat{B}_{coll}`$ is found, one can straighforwardly infer the evolution of the oscillators configuration simply by considering proper squeezing and rotations. This is equivalent to designing formal interferometric setups which could be used for proof-of-principle experiments where, at least for a few elements, the effects of CR terms could be simulated and observed. Motivated by these arguments, in Fig. 3 we show the equivalent interferometer for $`N=3`$.
By inspection, we see that this equivalent configuration results in concatenated Mach-Zehnder interferometers where the oscillators involved are subject to different squeezing and rotation operations. In going from $`N=3`$ to $`N=5`$, the overall concatenated structure of the setup is preserved, with just more oscillators being involved. This is not the case for $`N=4`$, whose equivalent all-optical setup turns out to be more complicated than the one corresponding to $`N=5`$ for instance, with squeezing and rotations involving the entire register, as already stressed. The intrinsic difference between the even and odd $`N`$ cases should now be more evident. The most relevant discrepancy is caused by the absence of the bare frequency $`\omega /2`$ from the spectrum of $`\widehat{H}_{chain}^{q,2l}`$.
The second issue which has to be discussed here in relation to Eq. (7) is the role of the CR terms. It is easy to be convinced that the version of $`\widehat{H}_{chain}`$ where the CR terms are excluded would have a contribution having the form $`\kappa _{j=1}^{N1}\widehat{p}_j\widehat{p}_{j+1}`$. Now the $`p`$-part of the chain Hamiltonian is also non-diagonal, with the same tridiagonal structure in the quadrature operators basis discussed in Eq. (3). Therefore, the orthogonal transformation which diagonalizes $`\widehat{H}_{chain}^q`$ (and the corresponding pattern of BS’s) can be used in order to reduce the $`p`$-part as well, getting the same set of eigen-frequencies. The corresponding eigen-operators are superpositions of just the $`\widehat{p}_j`$ quadrature operators with the same numerical coefficients $`\alpha _{jk}^N`$’s appearing in $`\widehat{O}_j^N`$’s. This implies that, for the odd $`N`$ case, after the application of $`\widehat{B}_{coll}`$, we end up with $`\widehat{B}_{coll}^{}\widehat{H}_{chain}\widehat{B}_{coll}=(\omega /2)[(\widehat{O}_1^N)^2+(\widehat{P}_1^N)^2]+_{j=2}^N_j[(\widehat{O}_j^N)^2+(\widehat{P}_j^N)^2]`$, where $`\widehat{P}_j^N=_{k=1}^N\alpha _{jk}^N\widehat{p}_k`$. Evidently, no squeezing is required in this case as the $`q`$ and $`p`$-parts of the Hamiltonian are already balanced by the diagonalization procedure. Thus, the corresponding time-evolution operation could be immediately reinterpreted as the tensorial product of formal rotation operators, one for each oscillator, showing that in this case the interferometric configurations sketched above are still valid: we only need to remove the squeezing operations.
It should be clear, up to this stage, that the exchange of any information encoded in the elements of the bosonic register occurs entirely by means of the effective collective operation $`\widehat{B}_{coll}`$. The remainder of the decomposition we have found, indeed, involves single-element operations which do not mutually mix the oscillators. Thus, by considering co-rotating terms only, we can see that the structure of $`\widehat{B}_{coll}`$, for a given $`N`$, is unchanged. This observation paves the way to the following consideration: the removal of CR terms from $`\widehat{H}_{chain}`$ prevents the spontaneous creation of excitations in the system. In terms of the equivalent all-optical setups, this means that by preparing the register in a classical initial state, no inter-oscillator entanglement has to be expected, in this case. Indeed, in ref. myungBS it is shown that non-classicality at the inputs of a BS is a fundamental pre-requisite for the entanglement of its outputs. In presence of CR terms, the single-oscillator squeezing provides the necessary non-classicality for inter-oscillator entanglement. We will come back to this point later, when the entanglement generation is quantitatively addressed.
It is worth comparing the decomposition in Eq. (7) with what has been found for a star-shaped bosonic configuration iohelenmyung . In an open linear chain, $`\widehat{B}_{coll}`$ induces multi-body interactions between the element of the chain. In particular, from Fig. 2 we see that an exchange of information is always required between the first and last element in a chain, a feature which holds regardless of $`N`$ commento2 . In the star-shaped configuration, on the other hand, any exchange of information occurs via a preferential way passing through the central component iohelenmyung .
## III Entanglement in an open chain: symmetry of the variance matrix
In order to investigate the dynamics of the entanglement generated among the oscillators in an open chain, we concentrate on Gaussian states and rely on the powerful tools provided by the variance matrix formalism. Indeed, the statistical properties of a Gaussian state, i.e. a state whose characteristic function is Gaussian, are entirely specified by the knowledge of its variance matrix. The variance matrix $`𝐕`$ is defined as $`V_{\alpha \beta }=\{\widehat{x}_\alpha ,\widehat{x}_\beta \}(\alpha ,\beta =1,..,2N)`$, where $`\widehat{x}=\widehat{q},\widehat{p}`$ and, for convenience, we have adopted the ordering of the quadrature operators $`\widehat{𝐱}=(\widehat{q}_1,\widehat{p}_1,..,\widehat{q}_N,\widehat{p}_N)^T`$. Throughout the paper, we assume that the Gaussian peak of each oscillator is at the origin of the respective phase-space. $`𝐕`$ is in one-to-one correspondence with the characteristic function of a Gaussian CV state which, in turns, gives information about the state of the system myungmunro .
When applied to an $`N`$-oscillator input Gaussian state, the operations involved in Eq. (7) give an output state which is also Gaussian. They can be formally described by means of the transformations $`𝒯_{𝐑_j}(\varphi _j)=\mathrm{cos}\varphi _j1𝐥+i\mathrm{sin}\varphi _j𝝈_y`$ for single-oscillator rotation and $`𝒯_{𝐒_j}(s_j)=e^{s_j𝝈_z}`$ for single-oscillator squeezing, where $`𝝈_\alpha `$ ($`\alpha =y,z`$) is the $`\alpha `$-Pauli matrix. For two-oscillator BS we have
$$𝒯_{𝐁_{jk}}(r_{jk},t_{jk})=\left(\begin{array}{cc}t_{jk}1𝐥& r_{jk}1𝐥\\ r_{jk}1𝐥& t_{jk}1𝐥\end{array}\right),$$
(9)
where $`t_{jk},r_{jk}`$ stand for the transmittivity and reflectivity of the BS acting on elements $`j`$ and $`k`$ (with $`t_{jk}^2+r_{jk}^2=1`$myungmunro . Explicitly, these transformations change an input variance matrix $`𝐕`$ to $`𝐕_{\alpha _j}=𝒯_{\alpha _j}^T𝐕𝒯_{\alpha _j}`$ ($`\alpha =𝐒,𝐑`$) for a single-oscillator $`2\times 2`$ variance matrix and $`𝐕_{𝐁_{jk}}=𝒯_{𝐁_{jk}}^T𝐕𝒯_{𝐁_{jk}}`$ for a two-oscillator $`4\times 4`$ variance matrix. From now on, we indicate with $`𝐕_f^N`$ the final variance matrix resulting from the application of all the transformations involved in $`\widehat{U}(t)`$ for a given $`N`$. In this Section we focus the attention onto the case in which all the oscillators are prepared in vacuum state, so that the initial variance matrix of the joint state of the chain is $`𝐕=_{j=1}^N1𝐥_j`$. By using Eq. (7) it can be shown that, for $`N=3`$, the final variance matrix reads
$$𝐕_f^3=\left(\begin{array}{ccc}𝐋_1& 𝐂_{12}& 𝐂_{13}\\ 𝐂_{12}^T& 𝐋_2& 𝐂_{12}\\ 𝐂_{13}^T& 𝐂_{12}^T& 𝐋_1\end{array}\right),$$
(10)
where $`𝐋_1=1𝐥+𝐂_{13}`$ and $`𝐋_2=1𝐥+2𝐂_{13}`$ account for the local properties of the oscillators while $`𝐂_{12}=\frac{1}{\sqrt{2}}_{j=2}^3(1)^{j+1}𝐜_j`$ and $`𝐂_{13}=\frac{1}{2}_{j=2}^3𝐜_j`$ describe the inter-oscillator correlations. We have introduced the elementary correlation matrices (which depend on the effective squeezing and rotations of oscillators $`j=2,3`$)
$$𝐜_j=\left(\begin{array}{cc}e^{2s_j}\mathrm{sin}^2(\varphi _j)\mathrm{sinh}(2s_j)& \frac{1}{2}\mathrm{sin}(2\varphi _j)\mathrm{sinh}(2s_j)\\ \frac{1}{2}\mathrm{sin}(2\varphi _j)\mathrm{sinh}(2s_j)& e^{2s_j}\mathrm{sin}^2(\varphi _j)\mathrm{sinh}(2s_j)\end{array}\right).$$
(11)
It is remarkable in Eq. (10) that the oscillators $`1`$ and $`3`$ have the same local properties, which are different from those of the mediator oscillator $`2`$. Moreover, the correlations between oscillators $`1`$ and $`2`$ appear to be the same as those between $`2`$ and $`3`$, which witnesses an evident degree of symmetry in the bosonic system ruled by Eq. (1). The proportionality of the correlation matrix $`𝐂_{12}`$ to the difference $`𝐜_3𝐜_2`$ is important, in this analysis, and is in striking contrast with the inherent structure of the correlations between the end points of the chain. These observations will be crucial in the upcoming discussion relative to the improvement of the end-to-end entanglement. The structure of Eq. (10) is found to hold for larger registers. Indeed, as still manageable examples, we mention that for $`N=4`$ and $`5`$ the decomposition of $`\widehat{U}(t)`$ is such that
$$𝐕_f^4=\left(\begin{array}{cccc}𝐋_1& 𝐂_{12}& 𝐂_{13}& 𝐂_{14}\\ 𝐂_{12}^T& 𝐋_2& 𝐂_{23}& 𝐂_{13}\\ 𝐂_{13}^T& 𝐂_{23}^T& 𝐋_2& 𝐂_{12}\\ 𝐂_{14}^T& 𝐂_{13}^T& 𝐂_{12}^T& 𝐋_1\end{array}\right),$$
(12)
which extends the symmetry already manifested in $`𝐕_f^3`$. In fact, the symmetry is a general property of $`𝐕_f^N`$: it is straightforward to see that $`𝐕_f^5`$ exhibits symmetry with respect to the central element of the chain, whose local properties are unique in the system. The expressions of $`𝐂_{jk}`$’s in terms of elementary correlation matrices analogous to $`𝐜_j`$ in Eq. (11) are, in general, quite cumbersome.
We address the generation of quantum correlations among the elements of an $`N`$-oscillator open chain as well as a simple strategy suitable for the improvement of the performances of this bosonic system as a long-haul entanglement distributor. The Gaussian preserving nature of the transformations $`𝒯_{𝐑_j}(\varphi _j)`$, $`𝒯_{𝐒_j}(s_j)`$ and $`𝒯_{𝐁_{jk}}(r_{jk},t_{jk})`$ allows us to exploit the well-known necessary and sufficient conditions for the entanglement of two-body CV Gaussian states simon ; myungmunro . The explicit object of our investigation will be the evaluation of the bipartite entanglement between the first and the $`j`$-th oscillator in a chain of $`N`$ oscillators. Therefore, we consider the reduced variance matrices $`v_{1j}`$ of the pairs $`(1,j)`$ which are found from $`𝐕_f`$ by extracting the $`4\times 4`$ submatrices ($`j=2,..,N`$)
$$v_{1j}=\left(\begin{array}{cccc}V_{1,1}& V_{1,2}& V_{1,2j1}& V_{1,2j}\\ V_{2,1}& V_{2,2}& V_{2,2j1}& V_{2,2j}\\ V_{2j1,1}& V_{2j1,2}& V_{2j1,2j1}& V_{2j1,2j}\\ V_{2j,1}& V_{2j,2}& V_{2j,2j1}& V_{2j,2j}\end{array}\right).$$
(13)
As a measure of entanglement we use the logarithmic negativity which provides an upper bound to the entanglement of distillation vidalwerner and is strictly related to the extent to which a given state violates the Peres-Horodecki criterion for separability npt . For bipartite Gaussian states, this entanglement measure can be easily calculated starting from the symplectic spectrum of the partial transposition of the variance matrix $`v_{ab}`$. In the phase-space, the partial transposition with respect to oscillator $`b`$ corresponds to the time-reversal operation which flips the sign of the momenutm quadrature operator of $`b`$. This can be represented by the action of the matrix $`P=1𝐥𝝈_z`$ onto $`v_{ab}`$. We introduce the matrix $`\mathrm{\Sigma }_{ab}=\mathrm{\Sigma }_a\mathrm{\Sigma }_b`$, where $`\mathrm{\Sigma }_m=i𝝈_{y,m}`$ ($`m=a,b`$) is the symplectic matrix of oscillator $`m`$ simon . The symplectic eigen-values of $`v_{ab}^{}=Pv_{ab}P`$ are the eigen-values of $`\left|i\mathrm{\Sigma }_{ab}v_{ab}^{}\right|`$, which are always equal in pairs. By calling $`\gamma _n^{}`$ ($`n=1,2`$) the representative of each pair, the inequality $`\mathrm{min}_n(\gamma _n^{})1`$ is a necessary and sufficient condition for the separability of $`v_{ab}`$. The logarithmic negativity $`\mathrm{\Lambda }_{ng}^{ab}`$ is then evaluated as $`\mathrm{\Lambda }_{ng}^{ab}=_n\mathrm{max}(0,\mathrm{log}_2\gamma _n^{})`$ vidalwerner .
The calculation of $`\mathrm{\Lambda }_{ng}^{1j}`$ for $`N=3`$, $`j[2,3]`$ and $`\kappa /\omega =0.1`$ leads to the plots shown in Fig. 4, where the bipartite entanglement between the three oscillators in the open chain is plotted against the rescaled interaction time $`\tau =\omega t`$. The choice for the ratio $`\kappa /\omega `$ is dictated by the fact that, experimentally, a weak coupling regime of $`\kappa \omega `$ is the only realistic situation plenio ; iohelenmyung .
The periodic behavior of the functions plotted is the signature of the time-dependence of the collective rotation $`\widehat{R}_{coll}(t)`$. As seen in the symmetry of $`𝐕_f^3`$, we have $`\mathrm{\Lambda }_{ng}^{12}=\mathrm{\Lambda }_{ng}^{23}`$ (Fig. 4, dotted curve). The peak of $`\mathrm{\Lambda }_{ng}^{12}`$ occurring at $`\tau =\tau ^{}44.2`$ corresponds to $`\left|\mathrm{sin}\varphi _j\right|>0.998`$. At this instant of time all the off-diagonal elements of $`𝐜_j`$’s do not exceed $`5\times 10^3`$. On the other hand, it is evident that $`\mathrm{\Lambda }_{ng}^{13}`$ (Fig. 4, solid line) is smaller than $`\mathrm{\Lambda }_{ng}^{12}`$ practically for any value of $`\tau `$, entirely disappearing at $`\tau ^{}`$. Thus, despite the CR terms are responsible for the for free generation of entanglement, a passive approach in which the bosonic register evolves freely without external intervention is evidently unsuitable for the creation of a reliable end-to-end entangled channel. On the contrary, almost all the quantum correlations within the system are localized among the nearest-neighbor oscillators (subsystems $`1+2`$ and $`2+3`$). The trend is common to any other case we have checked: $`\mathrm{\Lambda }_{ng}^{12}`$ can be almost an order of magnitude larger than any other $`\mathrm{\Lambda }_{ng}^{1j}`$ (see Fig. 5, for example). Moreover, it is apparent that the behavior of each entanglement function persists by enlarging the register. Only small modifications are observed in $`\mathrm{\Lambda }_{ng}^{1j}`$ when $`NN+1`$, the most evident of which is that $`\mathrm{\Lambda }_{ng}^{1N}`$ becomes non-zero after an increasing time delay. This can be understood by considering the effective all-optical setup (as the one sketched in Fig. 3): when all the oscillators are initially prepared in their vacuum state, the first set of BS’s (on the left hand side of the figure) are ineffective as they superimpose $`|0`$ states. As soon as the squeezing of oscillators $`2,..,N`$ is performed, the second set of BS’s, on the right hand side of the figures, is responsible for the generation and propagation of quantum correlations. Obviously, the number of operations which precede the coupling between $`1`$ and $`N`$, this latter carrying all the necessary non-classicality, increases with the dimension of the register, thus retarding the settlement of their entanglement. Again, the decomposition Eq. (7) shades new light onto the important features of the entanglement dynamics throughout the system, complementing the results highlighted by previous analyses plenio .
## IV Tagging by squeezing
For the purpose of creating an entangled state of the extremal oscillators in the channel, it is certainly desirable to look for strategies which quantitatively improve the entanglement settled between $`1`$ and $`N`$. Moreover, we would like to find out a way to make $`\mathrm{\Lambda }_{ng}^{1N}`$ dominant with respect to any other $`\mathrm{\Lambda }_{ng}^{1j}`$. The approach we are going to follow does not rely on local control over the elements of the chain between the first and the last oscillators. We assume that $`1`$ and $`N`$ are held by two spatially separated parties who can perfectly control the preparation of the respective oscillator and, if required, can measure their state. On the other hand, the interactions between the oscillators in the chain are set in a global way by a potential which collectively addresses all the elements at the same time. This approach is entirely within the rules of the global addressing strategies exemplified by quantum state transfer and phase-covariant cloning in quantum spin chains qst ; cloning and by always-on computational schemes sougato . In this perspective, the chain is seen as a two-terminal device whose intermediate stage is embodied by the $`N2`$ oscillators between the ending terminals $`1`$ and $`N`$. This central section is a black box whose dynamics are out of the grasp.
Intuitively, one would like to magnify the inherent distinction of the pair of oscillators $`1`$ and $`N`$, shown by Eqs. (10) and (12), from the rest of the register. Therefore, any local action performed onto the ending terminals of the device in Fig. 6, has to be designed so that the local properties of oscillators $`1`$ and $`N`$ still remain mutually equal. By considering the analysis performed in Section II and the role that non-classical states have in the entanglement by means of coupler operators myungBS , we look for an initial preparation of the register which can result in a quantitative increase of the end-to-end degree of entanglement. After a close inspection of the decomposition, we conjecture that single-oscillator squeezing operations onto $`1`$ and $`N`$ should improve the degree of entanglement between them.
In order to demonstrate our conejcture, we address the case of $`N=3`$ and we consider the preparation of an initial state whose variance matrix reads $`𝐕=𝐕_{sq,1}1𝐥_2𝐕_{sq,3}`$, where $`𝐕_{sq,j}=e^{2r𝝈_{z,j}}`$ is the variance matrix of a squeezed state with its squeezing parameter $`r`$. The calculation of the logarithmic negativity for the subsystems $`1+2,\mathrm{\hspace{0.17em}1}+3`$ and $`2+3`$ can proceed according to the recipe given in Section III. The corresponding degree of entanglement, from now on, will be indicated as $`\mathrm{\Lambda }_{ng}^{jk,tag}`$ ($`j,k[1,N]`$). The results, for $`r=0.2`$, are shown in Fig. 7 (a).
Due to the symmetry of the particular initial preparation, it is easy to check that we get a final variance matrix having the same general structure as Eq. (10), with suitably modified elementary correlation matrices. In particular, Fig. 7 $`(𝐛)`$ shows the time behavior of the matrix elements of $`𝐂_{12}=𝐂_{23}`$. Differently from what happens for an initially prepared vacuum state, when the ending elements are initially squeezed and for $`\tau 46.8`$, all the elements $`(C_{12})_{jk}`$ simultaneously become very close to zero ($`\left|(C_{12})_{1,1}\right|6\times 10^3,\left|(C_{12})_{2,2}\right|4\times 10^3`$ with $`\left|(C_{12})_{1,2}\right|=\left|(C_{12})_{2,1}\right|=0`$), whereas $`𝐂_{13}`$ (at that value of $`\tau `$) becomes diagonal with matrix elements in the range of $`0.1`$. This accounts for the improvement of the entanglement settled between $`1`$ and $`3`$ with, correspondingly, $`\mathrm{\Lambda }_{ng}^{12,tag}<2\times 10^3`$. The subsystem $`1+3`$ has been tagged by the single-element pre-squeezing to be the preferential pair of oscillators for the entanglement generation within the chain. It is worth stressing that this tagging procedure is possible in virtue of the symmetry existing between the ending elements of the open chain. An analogous conclusion has been drawn in ref. salerno , where a totally symmetric $`N`$-body CV system has been considered in order to point out the possibility of a unitary localization of the entanglement. In our case, however, the CV chain exhibits a degree of symmetry which is inferior to the one treated in ref. salerno . Different pairs of oscillators are characterized by different local and correlation properties, which makes the problem approached in this paper intrinsically different from the one in salerno . Nevertheless, we have shown that entanglement localization is possible with a lower degree of symmetry, which is per se an interesting point.
We can generalize the choice for the initial variance matrix in the tagging procedure to the case of $`N`$ oscillators by considering $`𝐕=𝐕_{sq,1}\left(_{j=2}^{N1}1𝐥_j\right)𝐕_{sq,N}`$. Again, an explicit calculation for the logarithmic negativity can be performed, which leads to the plots shown in Figs. 8 $`(𝐚)`$ and $`(𝐛)`$, for the cases $`N=4`$ and $`5`$. In Fig. 8 (b), the time-range has been restricted to the interesting region where $`\mathrm{\Lambda }_{ng}^{15,tag}\mathrm{\Lambda }_{ng}^{1j,tag},j[2,4]`$ in order to make the plot more transparent.
It is evident that there is always at least one value of $`\tau `$ at which the end-to-end entanglement dominates, making the tagging procedure effective. The amount of pre-required single-oscillator squeezing slightly depends on the dimension of the register and the plots in this paper show those values of $`r`$ at which we have found a good trade-off between the degree of entanglement and the effectiveness of the tagging strategy. As we have stressed before, the even and odd cases are inherently different, as also witnessed by the fact that the value of $`\tau `$ corresponding to an optimized tagging operation is larger for $`N=4`$ than for $`N=5`$. At the same time, $`\mathrm{\Lambda }_{ng}^{14,tag}`$, for $`N=4`$, is roughly proportional to $`r`$, while $`\mathrm{\Lambda }_{ng}^{1N,tag}>r`$ for all the odd $`N`$ cases we have checked.
## V Remarks
We have addressed the problem of long-haul entanglement creation in a register of bosons interacting via a global potential. The dynamics of entanglement can be clearly tracked via the effective decomposition of the time propagator in terms of simple linear optics elements as rotators, single-oscillator squeezers and couplers. This approach has allowed us to spot out a series of interesting features, characterizing the evolution of the quantum correlations settled among the elements of the register. As a result, we have been able to relate the conceptual role played by the CR terms in the entanglement generation process to effective squeezing operations on the elements of the register. The usefulness of this analysis is also witnessed by the design of a tagging protocol for the improvement of the end-to-end entanglement and the simultaneous reduction of any other $`1j`$ ($`j=2,..,N1`$) quantum correlation in an chain of $`N`$ elements. No local control over the central section is required: a proper preparation of the extremal oscillators and a collective interaction are sufficient to achieve the task. We believe this formal approach could be used in order to clarify other aspects related to the role played by the CR terms in entanglement creation, an issue which is certainly relevant especially in many problems of solid-state physics.
###### Acknowledgements.
We acknowledge discussions with Dr. J. Fiurášek. This work has been supported by the UK EPSRC and the Korea Research Foundation (2003-070-C00024).
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# Diffusion in pores and its dependence on boundary conditions
## 1 Introduction
The no-slip boundary condition for a fluid near a solid surface is still under debate . At the macroscopic scale, the no slip boundary condition is a consequence of the microscopic roughness . On the nanometer scale however partial slip is possible, and has indeed been measured experimentally . This issue, which is important both fundamentally and for the conception of microfluidic devices, has motivated a number of theoretical and numerical studies . These studies have highlighted the influence of the fluid-wall interaction and pressure on the slippage . While chemical heterogeneities and surface roughness are expected to decrease slippage , surfaces with special geometries can exhibit a ”super-hydrophobic” state with a strongly increased slippage at the surface that makes fluid dynamics at solids surfaces very sensitive to surface imperfections. Such effects have been evidenced using micro-engineered surfaces in reference .
Slippage is usually accounted for in terms of an extrapolation length, the so-called slip length, here denoted as $`\delta `$ . This is defined as the distance inside the solid wall where the extrapolated flow profile vanishes. More specifically this partial slip boundary condition is written, for the tangential component, $`v_t`$, of the velocity as
$$v_t=\delta \frac{v_t}{n}$$
(1)
with $`n`$ the coordinate in the direction normal to the solid surface. The precise value of this slip length and its dependence on the physical and chemical characteristics of the surface have been investigated in a number of recent experimental studies. In particular, very different values for slip lengths -from a few nanometers to microns - have been reported using different techniques (see e.g. for a review). Many of these techniques are indirect (pressure drop measurements , particle image velocimetry , fluorescence recovery ), or very delicate (surface force apparatus and Atomic force microscopy ). Hence the development of complementary, robust and non-intrusive techniques to investigate the dynamical properties of the solid-liquid interface would provide valuable counterparts of the previous results.
In this manuscript, we discuss how the diffusion of tagged particles between walls is affected by confinement, and how such measurements could be used as a signature of the nature of the boundary conditions .
We develop a theoretical and numerical approach to estimate the roles of confinement and slip on diffusion constrained in a planar or cylindrical pore. We make use of a classical hydrodynamic description, which is expected to be appropriate for colloidal particles, and was previously shown to be also well adapted for molecular diffusion . A numerical approach is used for the general case, and analytical expressions are derived in the high and low confinement regimes.
## 2 Hydrodynamic estimate of the diffusion constant
The quantity of interest is the mean diffusion coefficient of colloidal tracers averaged over the measurement volume. The latter is limited here by the presence of the confining walls. From a theoretical point of view, a moving particle $`𝒫`$ is subjected to a friction force proportional to its velocity. When the motion takes place in a confined volume, the mobility $`\mu `$ depends on the boundary conditions at the confining walls. For a velocity $`\underset{¯}{U}`$ parallel to the boundary, the diffusion coefficient $`D_{}`$ is given by Einstein’s relation
$$D_{}=\mu k_BT$$
(2)
For a particle moving between two flat walls separated by a distance $`H`$, $`D_{}`$ is a function of the particle radius $`a`$, the height $`H`$ and the position $`z`$ of the particle respective to the walls (in the following $`z`$ will be measured by taking the origin at the midplane). The average diffusion coefficient in the direction parallel to the walls is
$`D_{}={\displaystyle \frac{1}{H2a}}{\displaystyle _a^{Ha}}D_{}(z)𝑑z`$ (3)
The next step is to use the so called Stokes Einstein approach, i.e. to estimate the friction force from hydrodynamics. At low Reynolds number, the flow around the particle is governed by the Stokes equations :
$`\eta \mathrm{\Delta }\underset{¯}{V}`$ $`=`$ $`\underset{¯}{}P`$ (4)
$`\underset{¯}{}.\underset{¯}{V}`$ $`=`$ $`0`$ (5)
where $`\underset{¯}{V}`$ is the velocity field, $`P`$ the pressure field and $`\eta `$ the viscosity of the fluid. The boundary conditions are
* fluid at rest at infinity in the unconfined directions :
$$\underset{¯}{V}|_{\mathrm{}}=\underset{¯}{0}$$
(6)
* no slip on the particle surface :
$$\underset{¯}{V}|_𝒫=\underset{¯}{U}$$
(7)
* partial slip on solid walls, expressed by parallel $`\underset{¯}{V}_{}|_𝒮`$ and perpendicular $`\underset{¯}{V}_{}|_𝒮`$ velocities and slip length $`\delta `$ :
$`\delta _{}\underset{¯}{V}_{}|_𝒮\underset{¯}{V}_{}|_𝒮`$ $`=`$ $`\underset{¯}{0}`$ (8)
$`\underset{¯}{V}_{}|_𝒮`$ $`=`$ $`0`$ (9)
This condition is written $`𝒞_𝒮(\underset{¯}{V})=\underset{¯}{0}`$.
The friction force experienced by the particle is then:
$$\underset{¯}{F}=_𝒫\underset{¯}{\underset{¯}{\sigma }}.\underset{¯}{dS}$$
(10)
where $`\underset{¯}{\underset{¯}{\sigma }}=P\underset{¯}{\underset{¯}{I}}+\eta \left(\underset{¯}{}\underset{¯}{U}+^t\underset{¯}{}\underset{¯}{U}\right)`$ is the stress tensor in the fluid.
In the next sections we provide various solutions to this boundary problem : we first start with a numerical “exact” solution of these equations; the latter will be used subsequently as a reference solution for the approximate analytical solutions obtained under various assumptions.
## 3 Numerical Estimates
In this section we first start with a numerical solution of the previous equations, Eqs. (4) to (9). This set of equations was solved numerically with the *FEMLAB*<sup>©</sup> software. A finite domain of size $`2L\times 2L\times H`$ around the particle was considered. The size $`L`$ was chosen large enough compared to $`H`$ to avoid finite size problems : typical values are $`L/a=20`$ for small H, $`L=3H`$ otherwise. Space and time symmetries were taken into account to reduce the meshed domain for faster computations.
From a technical point of view, the *FEMLAB* fluid dynamics module solves the bulk equations as
$`\underset{¯}{}.\underset{¯}{\underset{¯}{\sigma }}=\underset{¯}{0}`$ (11)
$`\underset{¯}{}.\underset{¯}{V}=0`$ (12)
with the stress tensor $`\underset{¯}{\underset{¯}{\sigma }}`$ given above. Boundary conditions are imposed according to Eqs. (7) and (8)-(9).
Once the flow field has been obtained in this geometry, the force is computed according to equation (10). Note that this way of computing the force requires a fine mesh since a differentiation of the velocity field is performed. A better approach would consist in using a weak constraint formulation so that velocity and force are simultaneously computed on the surface. However, such an approach is time and memory consuming and was not used in this work to keep computational time within reasonable bounds.
Typical results are shown in figure 2 where the profile of the local parallel diffusion coefficient is plotted as a function of the altitude in the confining slab. As a case study, we consider the situation in which one of the two walls has a non zero slip length, while the no-slip boundary condition is applied at the second wall. Near the no slip wall, diffusion decreases from its bulk value as a result of the viscous friction and high velocity gradient in the fluid between the particle and the wall. This well known phenomenon is easily explained in term of an image particle (see next section). For a no-slip wall, the image particle moves in the opposite direction thus increasing the viscous force acting on the particle. Near a partially slipping wall diffusion increases from the no slip case and can even be higher than the bulk value. In the limit $`\delta \mathrm{}`$, diffusion reaches a high value that can be estimated using the image particle approach, with the image moving in the same direction as the particle .
We now turn to analytical approximate solutions of the Stokes equation in the previous geometry.
## 4 Analytical expressions in the low confinement (large gap) limit
When the particle is small compared to confinement height, an iterative reflection method can be developed, leading to an analytical expression for the friction force.
In the present work we use this approach in the presence of a single, slipping, wall. Then, summing over the forces due to each wall yields an approximate result for the average diffusion coefficient. A summary of the method is given here, and details are discussed in Appendix Appendix A: First reflected field $`\underset{¯}{V}^1`$..
### 4.1 Reflection method with a single, slipping, wall
The reflection method is an iterative approach , in which the velocity field $`\underset{¯}{V}`$ is expanded in the form
$`\underset{¯}{V}`$ $`=`$ $`\underset{¯}{V}^0+\underset{¯}{V}^1+\underset{¯}{V}^2+\underset{¯}{V}^3+\mathrm{}`$ (13)
with each $`\underset{¯}{V}^n`$ field satisfying the bulk equations (4)-(5). The zero order field, $`\underset{¯}{V}_0`$, is chosen as the flow field around a sphere moving in the bulk :
$$\underset{¯}{V}^0(r,z)=\frac{3}{4}aU\left(\frac{2}{r}\underset{¯}{}(z)\underset{¯}{}(\frac{z}{r})+\frac{a^2}{3}\underset{¯}{}(\frac{z}{r^3})\right)$$
(14)
with $`a`$ the radius of the sphere. This velocity field satisfies the boundary equations on the particle (equation (7)) and at infinity (equation (6)). The method consists in determining $`\underset{¯}{V}^1`$ field such that $`\underset{¯}{V}^0+\underset{¯}{V}^1`$ satisfies the boundary conditions at infinity and on the solid walls (8)-(9) :
$$\{\begin{array}{ccc}\hfill 𝒞_𝒮(\underset{¯}{V}^1)& =& 𝒞_𝒮(\underset{¯}{V}^0)\hfill \\ \hfill \underset{¯}{V}^1|_{\mathrm{}}& =& \underset{¯}{0}\hfill \end{array}$$
Now, at this level of approximation, the boundary condition on the particle $`𝒫`$ is no longer satisfied by $`\underset{¯}{V}^0+\underset{¯}{V}^1`$ and the next order $`\underset{¯}{V}^2`$ is defined from the reflection of $`\underset{¯}{V}^1`$ on the particle as :
$$\{\begin{array}{ccc}\hfill \underset{¯}{V}^2|_𝒫& =& \underset{¯}{V}^1|_𝒫\hfill \\ \hfill \underset{¯}{V}^2|_{\mathrm{}}& =& \underset{¯}{0}\hfill \end{array}$$
The higher moments of the velocity field, $`\underset{¯}{V}^n`$, are built by applying iteratively the boundary condition on the particle and on the flat walls.
### 4.2 Viscous force acting on the particle : a single wall
The friction force experienced by the particle is the sum of individual contributions $`\underset{¯}{F}^n`$ of each reflection :
$$\underset{¯}{F}^n=_𝒫\underset{¯}{\underset{¯}{\sigma }}^n.\underset{¯}{dS}=_𝒫\underset{¯}{}.\underset{¯}{\underset{¯}{\sigma }}^ndV$$
(15)
where $`\underset{¯}{\underset{¯}{\sigma }}^n`$ is the stress tensor in the fluid. For odd reflections, the velocity is regular in the volume of the particle. The momentum equation gives $`\underset{¯}{}.\underset{¯}{\underset{¯}{\sigma }}^n=0`$ in the domain occupied by $`𝒫`$ and the integral vanishes. For even reflections, the Lorentz reciprocal theorem gives algebraically, in the limit of small particles, $`F^{n+2}=\frac{V_𝒪^{n+1}}{V_𝒪^{n1}}F^n`$, where $`V_𝒪^n`$ is defined as the value of the velocity field at the center of the particle. One thus obtains
$$\underset{¯}{F}=\underset{¯}{F}^0\underset{k=0}{\overset{\mathrm{}}{}}\left[\frac{V_𝒪^1}{U}\right]^k=\frac{\underset{¯}{F}^0}{1+\frac{V_𝒪^1}{U}}$$
(16)
with $`\underset{¯}{F}^0=6\pi \eta a\underset{¯}{U}`$. As a consequence, only the velocity of the first reflected field at the center of the particle $`V_𝒪^1`$ is needed to determine mobility and diffusion coefficient. The calculation of this field is described in appendix A.
Equations (16) and (A.7) give the force acting on a particle moving along a single planar wall as a function of the radius of the particle $`a`$, the distance from the wall $`l`$ and the slip length $`\delta `$ :
$`\underset{¯}{F}_{1wall}`$ $`=`$ $`{\displaystyle \frac{6\pi \eta a}{1\frac{a}{z}C\left[\frac{l}{\delta }\right]}}\underset{¯}{U}`$ (17)
where the function $`C`$ is defined as
$$C\left[y\right]=\frac{3}{32}y^2\frac{9}{32}y\frac{3}{8}+\left(\frac{3}{32}y^3+\frac{3}{8}y^2+\frac{3}{8}y\right)E(y)+\frac{3}{2}yE(2y)$$
(18)
with $`E(y)=e^yE_1(y)`$ and $`E_1(y)`$ is the exponential integral function, defined as $`E_1(z)=_z^{\mathrm{}}𝑑te^t/t`$ .
When $`\delta 0`$ (no slip condition), one recovers the well known value
$`C\left[{\displaystyle \frac{l}{\delta }}\right]{\displaystyle \frac{9}{16}}`$ (19)
derived from the method of the image particle : as mentioned above, diffusion decreases near a no slip wall. In the limit $`\delta \mathrm{}`$ (full slip condition),
$`C\left[{\displaystyle \frac{l}{\delta }}\right]{\displaystyle \frac{3}{8}}`$ (20)
and the presence of the wall reduces the friction force, i.e. diffusion increases, as measured experimentally .
Comparisons with numerical simulations using FEMLAB<sup>©</sup> are shown in figure 3. Results are in good agreement down to very small distances $`l/a=1.5`$. At large distance, one recovers the bulk diffusion value as expected.
Moreover, simple and practical approximations can be obtained for the mobility in the limit where the distance to the wall, $`l`$, is large compared to the slip length $`\delta `$. Indeed an asymptotic expansion of $`C\left[y\right]`$ allows to obtain
$$C\left[y\right]=\frac{9}{16}\frac{1}{1+\frac{1}{y}+𝒪\left[\frac{1}{y^2}\right]}$$
(21)
This gives the approximate following form for the friction coefficient
$`\underset{¯}{F}_{1wall}`$ $``$ $`{\displaystyle \frac{6\pi \eta a}{1\frac{9}{16}\frac{a}{l+\delta }}}\underset{¯}{U}`$ (22)
This approximation amounts to replace the distance to the wall $`l`$ by $`l+\delta `$, where the physical meaning of the slip length in terms of an extrapolation length appears quite clearly in this limit. In practice, note that the expression in Eq. (22) leads to values which are within $`5\%`$ to the explicit result in Eq. (17) as soon as $`l/\delta >0.5`$ !
After completing this work, we became aware of a similar calculation by Lauga and Squires who computed the viscous force on a spherical particle close to a wall, using the same reflection method. The use of the ”small particle” approximation corresponds to computing the flow in response to a force applied to a point-like particle, and can be shown to involve errors of order $`(a/h)^3`$ .
#### 4.2.1 Local diffusivity in a confined geometry
In order to compute the friction coefficient for a particle confined between two planar walls, we make the further assumption that each wall contributes independently to the shift in the friction force from its bulk value :
$`\underset{¯}{F}_{2walls}=\underset{¯}{F}_{1wall}(z,\delta )+\underset{¯}{F}_{1wall}(Hz,\delta )\underset{¯}{F}_{bulk}`$ (23)
where $`H`$ is the distance between the two walls and here $`z`$ denotes the distance to the bottom wall. The Einstein equation then yields for the parallel diffusion coefficient at a height $`l`$ :
$`D_{}={\displaystyle \frac{k_BT}{6\pi \eta aU}}{\displaystyle \frac{1}{\frac{1}{1\frac{a}{z}C\left[\frac{z}{\delta }\right]}+\frac{1}{1\frac{a}{Hz}C\left[\frac{Hz}{\delta }\right]}1}}`$ (24)
This expression for the friction coefficient is checked against the “exact” numerical results obtained using the FEMLAB software in Figures 4 and 5. Over the various slip lengths $`\delta `$ and confinement gap $`H`$, the agreement is found to be quite good, within $`6\%`$ as long as the confinement is no too strong ($`h/2a>4`$). It can be observed that Eq. (24) slightly underestimates the diffusion.
Note that a different approximation could be made for the contribution of the two walls, by assuming that the mobility (rather than its inverse) is affected independently by the two walls . This approximation, however, turns out to be less accurate than the previous approximation, in Eq. (24).
When the particle is confined to a cylindrical pore, a similar method can be used and provides an estimate of the viscous force acting on the particle in the low confinement limit (see appendix A). However, as opposed to the planar case, only a numerical estimation of the reflected velocity at the center of the sphere can be reached. An interesting difference between the planar and the cylindrical case, is that in the latter case, except close to the wall where the behavior is similar to a particle moving near a planar wall, the force acting on a tracer particle is never smaller than its bulk value even in the large slip length limit (i.e. the diffusion is reduced). This is due to the necessary recirculation of the fluid around the particle. More precisely, boundary conditions at infinity (no flow) imposes the overall flow rate on a section of the cylinder at zero. In the section centered on the particle, a negative fluid flow rate has to balance the positive flow rate of the particle $`\pi a^2U`$. Hence the viscous force increases from the bulk value even when $`\delta \mathrm{}`$. In the planar geometry, recirculation takes place at infinity in the unconfined directions and this phenomenon does not take place.
## 5 Lubrication theory in the strong confinement limit
When the confinement approaches the particle size ($`H2a`$), the main part of the viscous force is expected to arise from the high velocity gradient in the thin fluid films between each wall and the particle. In these regions, the fluid flow is quasi-parallel to the wall and lubrication theory is expected to provide a good description of the velocity field. An approximation of the force acting on the particle can then be derived.
We assume here that the fluid is confined between a fixed sphere and a solid wall moving at velocity $`\underset{¯}{U}=(U\mathrm{,0,0})`$ (see figure 2). One approximates furthermore the sphere by a paraboloid $`h(r)h_0+\frac{r^2}{2a}`$, with $`r`$ the distance to the axis of symmetry of the paraboloid.
Under the lubrication assumptions , the Stokes equation (4) for the velocity field, $`\underset{¯}{W}=(W_x,W_y,W_z)`$, reduces to :
$`\eta {\displaystyle \frac{^2W_{}(x,y,z)}{z^2}}`$ $`=`$ $`\underset{¯}{}_{}P(x,y)`$ (25)
The boundary conditions are written as
$`{\displaystyle \frac{W_x}{z}}={\displaystyle \frac{W_xU}{\delta }}`$ (26)
$`{\displaystyle \frac{W_y}{z}}={\displaystyle \frac{W_y}{\delta }}`$ (27)
on the wall and
$`\underset{¯}{W}`$ $`=`$ $`\underset{¯}{0}`$ (28)
on the particle. These equations are easily integrated and using the conservation equation $`\underset{¯}{}_{}.\underset{¯}{Q}_{}=0`$ for the flow rate, defined as $`\underset{¯}{Q}_{}=_0^{h(x,y)}\underset{¯}{W}_{}𝑑z`$, one gets the following equation for the pressure, $`P`$ :
$`{\displaystyle \frac{1}{12\eta }}\underset{¯}{}_{}\left({\displaystyle \frac{h^3(h+4\delta )}{h+\delta }}\underset{¯}{}_{}P\right)+\underset{¯}{U}.\underset{¯}{}_{}\left({\displaystyle \frac{h^2}{2(h+\delta )}}\right)=0`$ (29)
A general solution for the pressure can be written in the form $`P=P_{\mathrm{}}+\mathrm{\Pi }(r)\mathrm{cos}(\theta )`$, with {$`r`$,$`\theta `$} the angular coordinates on the planar wall. We could however not find an analytical solution for the previous differential equation. However a “heuristic solution” could be found after some manipulation of the differential equation in the form $`\mathrm{\Pi }(r)=\eta Urb[h(r)]`$, with
$`b[h]={\displaystyle \frac{6}{5h\delta }}{\displaystyle \frac{9}{10}}{\displaystyle \frac{\mathrm{ln}(h)}{\delta ^2}}+{\displaystyle \frac{4}{5}}{\displaystyle \frac{\mathrm{ln}(h+\delta )}{\delta ^2}}+{\displaystyle \frac{1}{10}}{\displaystyle \frac{\mathrm{ln}(h+4\delta )}{\delta ^2}}`$ (30)
We refer to appendix Appendix B: Lubrication approximation for details of the calculations leading to this result. The validity of this approximate expression for the pressure was checked by computing numerically the solution of the full differential equation (29) using a simple ODE solver (Mathematica<sup>©</sup>). One finds that the previous solution for the pressure differs from the “exact” numerical one, from less than a few percents for $`\delta [0,R]`$, $`h_0[\frac{R}{20},R]`$, and over the full range of distance $`r[0,\mathrm{}[`$ (see figure 11 in the appendix). Moreover, in the vanishing slip length limit, $`\delta 0`$, the previous solution for $`b[h]`$ in Eq. (30) reduces to the corresponding exact solution of the differential equation, which can be easily obtained as $`b_0(h)=\frac{6}{5h^2}`$.
Using this previous heuristic solution as a good approximation for the pressure, one may then write the force balance along the x-direction applied on the volume of fluid inside the cylinder $`r<R_c`$ (see figure 7) as
$`F_𝒫=(F_{wall}+{\displaystyle _{r=R_c}}P\underset{¯}{n}.\underset{¯}{x}dS)`$ (31)
At large scales $`R_c\mathrm{}`$, one may verify that the slip effect disappears : $`P(\delta ,r=R_C)P(\delta =0,r=R_C)`$ and $`_{r=R_c}P𝑑S`$ is independent of $`\delta `$ and the dependence of the friction force acting on the particle $`P`$ come from the $`R_C\mathrm{}`$ limit of $`F_wall`$.
A second difficulty however arises with the lubrication calculation : whatever the slip length $`\delta `$, the friction force on the wall, $`F_{wall}=_{wall}\eta _{}W_xdS`$, is found to be logarithmically divergent when $`R_c\mathrm{}`$ . This can be easily verified by inserting in the previous friction force expression the expression $`W_x`$ deduced from the pressure field with Eq. (30) (see also Eq. (B.26) in appendix B). On the other hand, the difference of friction forces, $`\mathrm{\Delta }F_{wall}=F_{wall}[\delta ]F_{wall}[\delta =0]`$, between the finite slip length case and the no-slip case is found to take a finite value, given in eq. (B.28). Note that $`\mathrm{\Delta }F_𝒫=\mathrm{\Delta }F_{wall}\mathrm{\Delta }F`$ (since the second term in Eq. (31) is independent of $`\delta `$ in the $`R_C\mathrm{}`$ limit).
One may however argue that this difference is a physically relevant quantity since the slip effects mainly affect the flow in the region with strongest confinement. One may indeed verify that at the lubrication level, the flows with partial slip reduces to the flow with the no-slip boundary condition in the region far from closest contact.
We have plotted in Figure 8 (left) the result for $`\mathrm{\Delta }F=F_{wall}[\delta ]F_{wall}[\delta =0]`$ (normalized by the bulk value of the force on the particle $`F_{\mathrm{}}=6\pi \eta aU`$) as a function of the minimum gap $`h_0`$ between the sphere and the wall. This result is compared to the FEMLAB calculations in the same geometry. As expected an agreement is found in the small gap limit, where the lubrication approximation is expected to be valid.
Pursuing this calculation, a diffusion coefficient can be obtained. First the friction coefficient on the particle situated at a distance $`l`$ from the wall (with slip length $`\delta `$) can be estimated at this level of approximation by adding to the previous $`\mathrm{\Delta }F`$ the value of the friction force $`F_{wall}[l,\delta =0]`$ computed in the same configuration for a no-slip wall. Then the friction coefficient for the particle confined between two partially slipping wall is estimated by adding the effects of the walls on the friction coefficient, according to Eq. (23). The mobility is finally evaluated by the inverse of the friction coefficient. This procedure is applied in Figure 8 (right), where the numerical (FEMLAB) result has been used for the no-slip friction force $`F_{wall}[l,\delta =0]`$. Here the diffusion coefficient is computed only for the situation where the sphere is at the center of the slab ($`z=H/2`$). This result is compared to “exact” results obtained using a full FEMLAB calculation for the particle confined between two partially slipping walls. Again, the lubrication approximation only yields a correct agreement in the small gap region, and works better for small slip lengths. When the slip length increases, lubrication theory overestimates the mobility : In this case, important contributions to the viscous force are coming from areas far from the confined zones, which are not properly described within the lubrication approach.
The lubrication approach has therefore a quite limited range of application (in the very confined region) but the solution obtained is complementary to the low confinement results which works in the large gap limits.
## 6 Averaged diffusivity and conclusions
We are now in a position to compute the averaged diffusion coefficient over the confined slab, $`D_{}`$ defined in Eq. (3). We consider a geometry where one wall is characterized by a no-slip boundary condition, while a partial slip boundary condition, with a slip length $`\delta `$ applies on the other. This configuration is chosen as to mimic the experimental geometry .
Results are shown in figures 9 and 10 for various values of the confinement $`H`$ and of the slip length $`\delta `$.
Analytical results obtained in the low confinement approximation, combined with the assumption of independent wall contribution, reproduce quite well the trends of the numerical computations. Figure 9 shows that the analytical estimate slightly underestimates diffusion in the low confinement limit, and tends to overestimate it at strong confinements.
In order to observe a significant dependence of diffusion on the slip length, two conditions are required. The particle size should not be much larger than the slip length, and a sufficiently strong confinement is required. With typical values of $`H4a`$, variations of the average diffusion constant of typically 5% to 10% would be expected if the slip length is changed between $`0.1a`$ and $`a`$.
These results therefore suggest that diffusion measurements are quite sensitive to boundary conditions on the solid substrate. This opens new routes to measure slip length on the basis of the thermal motion of colloidal tracers .
Acknowledgments L.B. thanks Yannick Almeras, with whom this work was initiated.
## Appendix A: First reflected field $`\underset{¯}{V}^1`$.
### A.1 Planar geometry
For a particle moving at a distance $`l`$ from a planar surface, the general form of $`\underset{¯}{V}^1`$ satisfying equations (4)-(5) and (6) is given by :
$`V_x^1={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{2\pi }}\left[\left(2\mathrm{cos}^2(u)(k|z|+1)\right)\mathrm{\Theta }^1+ik\mathrm{cos}(u)\mathrm{{\rm Y}}^1+zk^2\mathrm{cos}^2(u)\mathrm{\Xi }^1\right]`$
$`\times e^{ik(x\mathrm{cos}(u)+y\mathrm{sin}(u))k|z|}dkdu`$ (A.1)
$`V_y^1={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{2\pi }}i\mathrm{sin}(u)\left[i\mathrm{cos}(u)(k|z|+1)\mathrm{\Theta }^1+k\mathrm{{\rm Y}}^1izk^2\mathrm{cos}(u)\mathrm{\Xi }^1\right]`$
$`\times e^{ik(x\mathrm{cos}(u)+y\mathrm{sin}(u))k|z|}dkdu`$ (A.2)
$`V_z^1={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{2\pi }}\left[izk\mathrm{cos}(u)\mathrm{\Theta }^1k\mathrm{{\rm Y}}^1+ik\mathrm{cos}(u)(k|z|+1)\mathrm{\Xi }^1\right]`$
$`\times e^{ik(x\mathrm{cos}(u)+y\mathrm{sin}(u))k|z|}dkdu`$ (A.3)
where $`\mathrm{\Theta }^1`$, $`\mathrm{{\rm Y}}^1`$ et $`\mathrm{\Xi }^1`$ are functions of $`k`$ and $`u`$ and $`(x,y,z)`$ are cartesian coordinates centered on the particle with $`x`$ along the particle velocity and $`z`$ normal to the wall.
The initial field $`\underset{¯}{V}^0`$ is written in this form with ($`\mathrm{\Theta }^0(k,u)=\frac{3}{4}aU`$, $`\mathrm{{\rm Y}}^0(k,u)=\frac{1}{4}a^3Uik\mathrm{cos}(u)`$, $`\mathrm{\Xi }^0(k,u)=0`$). In the limit of a small particle, $`\mathrm{{\rm Y}}^0(k,u)<<\mathrm{\Theta }^0(k,u)`$.
Functions ($`\mathrm{\Theta }^1`$, $`\mathrm{{\rm Y}}^1`$, $`\mathrm{\Xi }^1`$) are determined as the unique solution of $`𝒞_𝒮(\underset{¯}{V}^0+\underset{¯}{V}^1)=\underset{¯}{0}`$ :
$`\mathrm{\Theta }^1`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{k\delta 1}{k\delta +1}}aUe^{2kl}`$ (A.4)
$`\mathrm{{\rm Y}}^1`$ $`=`$ $`{\displaystyle \frac{3i}{2}}{\displaystyle \frac{\left(2lk^2\delta ^2+k^2l^2\delta +k\delta l\delta +kl^2l\right)}{2k^2\delta ^2+1+3k\delta }}Ua\mathrm{cos}(u)e^{2kl}`$ (A.5)
$`\mathrm{\Xi }^1`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{\left(k\delta l+\delta +l\right)}{2k^2\delta ^2+1+3k\delta }}aUe^{2kl}`$ (A.6)
Integration of $`V_x^1(\mathrm{0,0,0})`$ gives
$`V_𝒪^1={\displaystyle \frac{a}{l}}C\left[{\displaystyle \frac{l}{\delta }}\right]U`$ (A.7)
with $`C\left[y\right]=\frac{3}{32}y^2\frac{9}{32}y\frac{3}{8}+\left(\frac{3}{32}y^3+\frac{3}{8}y^2+\frac{3}{8}y\right)E(y)+\frac{3}{2}yE(2y)`$. $`E(y)=e^yEi(1,y)`$ and $`Ei(1,y)`$ is the exponential integral function.
### A.2 Cylindrical geometry
For a particle moving in a cylinder, the general solution for the reflected field $`\underset{¯}{V}^1`$ in $`(r,\varphi ,z)`$ cylindrical coordinates is
$`\underset{¯}{V}^1(r,\varphi ,z)`$ $`=`$ $`{\displaystyle \frac{3aU}{2\pi }}{\displaystyle \underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle _0^+\mathrm{}}𝑑\lambda \left(\begin{array}{c}a_k(\lambda ,r)\mathrm{cos}(k\varphi )\mathrm{sin}(\lambda z)\\ b_k(\lambda ,r)\mathrm{sin}(k\varphi )\mathrm{sin}(\lambda z)\\ c_k(\lambda ,r)\mathrm{cos}(k\varphi )\mathrm{cos}(\lambda z)\end{array}\right)_{(\underset{¯}{e}_r,\underset{¯}{e}_\varphi ,\underset{¯}{e}_z)}`$ (A.11)
$`a_k(\lambda ,r)`$ $`=`$ $`{\displaystyle \frac{k}{\lambda r}}\mathrm{\Omega }_k^1(\lambda )I_k(\lambda r)+\mathrm{\Psi }_k^1(\lambda )I_k^{}(\lambda r)+\lambda r\mathrm{\Pi }_k^1(\lambda )I_k^{\prime \prime }(\lambda r)`$
$`b_k(\lambda ,r)`$ $`=`$ $`\mathrm{\Omega }_k^1(\lambda )I_k^{}(\lambda r){\displaystyle \frac{k}{\lambda r}}\mathrm{\Psi }_k^1(\lambda )I_k(\lambda r)k\mathrm{\Pi }_k^1(\lambda )I_k^{}(\lambda r)+{\displaystyle \frac{k}{\lambda r}}\mathrm{\Pi }_k^1(\lambda )I_k(\lambda r)`$
$`c_k(\lambda ,r)`$ $`=`$ $`\mathrm{\Psi }_k^1(\lambda )I_k(\lambda r)+\lambda r\mathrm{\Pi }_k^1(\lambda )I_k^{}(\lambda r)+\mathrm{\Pi }_k^1(\lambda )I_k(\lambda r)`$
where $`I_k`$, $`I_k^{}`$ et $`I_k^{\prime \prime }`$ are the first order modified Bessel functions and their derivatives.
Bulk field $`\underset{¯}{V}^0`$ for a particle moving along the cylinder axis, at a distance $`b`$ from it, is expressed in such a form as :
$`\underset{¯}{V}^0(r,\varphi ,z)`$ $`=`$ $`{\displaystyle \frac{3aU}{2\pi }}{\displaystyle \underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle _0^+\mathrm{}}𝑑\lambda \left(\begin{array}{c}\alpha _k(\lambda ,r)\mathrm{cos}(k\varphi )\mathrm{sin}(\lambda z)\\ \beta _k(\lambda ,r)\mathrm{sin}(k\varphi )\mathrm{sin}(\lambda z)\\ \gamma _k(\lambda ,r)\mathrm{cos}(k\varphi )\mathrm{cos}(\lambda z)\end{array}\right)_{(\underset{¯}{e}_r,\underset{¯}{e}_\varphi ,\underset{¯}{e}_z)}`$ (A.16)
$`\alpha _k(\lambda ,r)`$ $`=`$ $`\left(\lambda r+{\displaystyle \frac{k^2}{\lambda r}}\right)K_k(\lambda r)I_k(\lambda b)+\lambda bK_k^{}(\lambda r)I_k^{}(\lambda b)`$
$`\beta _k(\lambda ,r)`$ $`=`$ $`k\left(K_k^{}(\lambda r)I_k(\lambda b)+{\displaystyle \frac{b}{r}}K_k(\lambda r)I_k^{}(\lambda b)\right)`$
$`\gamma _k(\lambda ,r)`$ $`=`$ $`2K_k(\lambda r)I_k(\lambda b)+\lambda rK_k^{}(\lambda r)I_k(\lambda b)+\lambda bK_k(\lambda r)I_k^{}(\lambda b)`$
where $`K_k`$, $`K_k^{}`$ et $`K_k^{\prime \prime }`$ are the second order modified Bessel functions and their derivatives.
$`\mathrm{\Omega }_k^1(\lambda )`$, $`\mathrm{\Psi }_k^1(\lambda )`$ and $`\mathrm{\Pi }_k^1(\lambda )`$ are the unique solution of
$$𝒞_𝒮(\underset{¯}{V}^0+\underset{¯}{V}^1)=\underset{¯}{0}$$
(A.18)
on the cylindrical wall $`r=R`$. The analytical expression is then used to compute numerically $`V_𝒫^1`$ as :
$$V_𝒪^1=\frac{3aU}{2\pi }\underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}_0^+\mathrm{}\left[\left(\mathrm{\Psi }_k(\lambda )+\mathrm{\Pi }_k(\lambda )\right)I_k(\lambda b)+\mu b\mathrm{\Pi }_k(\lambda )I_k^{}(\lambda b)\right]𝑑\lambda $$
(A.19)
## Appendix B: Lubrication approximation
Using reduced variables $`r=\sqrt{2h_0a}\stackrel{~}{r}`$, $`h=h_0\stackrel{~}{h}`$, $`\delta =h_0\stackrel{~}{\delta }`$ and $`P=\eta U\frac{\sqrt{2h_0a}}{h_0^2}\stackrel{~}{p}`$, mass conservation is (for compactness w remove the $`\stackrel{~}{}`$ for the reduced variables):
$`\underset{¯}{}_{}\left({\displaystyle \frac{h^3(h+4\delta )}{12(h+\delta )}}\underset{¯}{}_{}p\right)+\underset{¯}{e}_x.\underset{¯}{}_{}\left({\displaystyle \frac{h^2}{2(h+\delta )}}\right)=0`$ (B.20)
with $`h(r)=1+r^2`$.
Assuming $`p(r)=p_{\mathrm{}}+rb(r)\mathrm{cos}(\theta )`$, equation (B.20) becomes
$`{\displaystyle \frac{}{r}}\left[r\alpha (h(r)){\displaystyle \frac{}{r}}\left(rb(r)\right)\right]+\alpha (h(r))b(r)+2r\beta (h(r))=0`$ (B.21)
with $`\alpha (h)=\frac{h^3(h+4\delta )}{12(h+\delta )}`$ and $`\beta (h)=\frac{h^2+2\delta h}{2(h+\delta )^2}`$.
We could not find an exact solution for this equation. However in order to proceed further, we have tried to construct in a heuristic way a good approximation to the solution to avoid purely numerical solutions. We have proceeded as follows. First $`b`$ is assumed to depend functionaly on $`h(r)`$, as $`b[h(r)]`$. Expressing $`p=P_{\mathrm{}}+xb[h]`$ ($`x=r\mathrm{cos}\theta `$) in Eq. (B.20), this equation rewrites
$`\underset{¯}{}_{}\left(\alpha (h)\left(b(h)\underset{¯}{e}_x+x{\displaystyle \frac{b}{h}}\underset{¯}{}_{}h\right)\right)+2\beta (h)x=0`$ (B.22)
with $`\underset{¯}{}_{}h=\{2x\mathrm{,2}y\}`$ in cartesian coordinates. An heuristic solution is found by assuming $`\alpha (h),\beta (h)`$ and $`\frac{b}{h}`$ as constant the previous equation, which amounts to replace the previous equation by
$`6x{\displaystyle \frac{b}{h}}+2𝒜\beta (h)x=0`$ (B.23)
The constant $`𝒜`$ is adjusted so that the exact no-slip solution of Eq. (B.20), $`b_0(h)=\frac{6}{5h^2}`$, is recovered. The solution of the Eq. (B.23) with $`𝒜=6/5`$ is
$`b(h)={\displaystyle \frac{6}{5h\delta }}{\displaystyle \frac{9}{10}}{\displaystyle \frac{\mathrm{ln}(h)}{\delta ^2}}+{\displaystyle \frac{4}{5}}{\displaystyle \frac{\mathrm{ln}(h+\delta )}{\delta ^2}}+{\displaystyle \frac{1}{10}}{\displaystyle \frac{\mathrm{ln}(h+4\delta )}{\delta ^2}}`$ (B.24)
which indeed reduces to the no-slip solution $`b_0(h)=\frac{6}{5h^2}`$ when $`\delta 0`$. Note also that in the limit $`h>>\delta `$, one also recovers $`b(h)b_0(h)`$ : the pressure is independent of $`\delta `$ far from the particle.
The reduced viscous force acting on the wall is
$`F_{wall}={\displaystyle \frac{W_x}{z}r𝑑r𝑑\theta }`$ (B.25)
$`W_x`$ is determined from Stokes equation $`\frac{^2W_x}{z^2}=\frac{p}{x}`$ along with the boundary conditions and yields
$`F_{wall}=3\pi {\displaystyle _0^{\mathrm{}}}\left[{\displaystyle \frac{1}{6}}(b(r)+{\displaystyle \frac{}{r}}(rb(r))){\displaystyle \frac{h^2}{h+\delta }}+{\displaystyle \frac{2}{3}}{\displaystyle \frac{1}{h+\delta }}\right]r𝑑r+Cste`$ (B.26)
Deviation of the viscous force from $`\delta =0`$ case is, back with dimensionalized variables :
$`\mathrm{\Delta }F=6\pi \eta aU{\displaystyle _0^{\mathrm{}}}\left[{\displaystyle \frac{h1}{6}}\left(b_0(h)b(h){\displaystyle \frac{h(h+2\delta )}{(h+\delta )^2}}\right){\displaystyle \frac{2}{6}}{\displaystyle \frac{\delta }{h(h+\delta )}}\right]𝑑h`$ (B.27)
This expression can be exactly computed for the approximated $`b(h)`$ given above :
$`\mathrm{\Delta }F`$ $`=`$ $`{\displaystyle \frac{6\pi \eta aU}{360\delta ^2}}(36\delta 12\delta ^2+10\pi ^2\delta ^2+54\delta \mathrm{ln}\left({\displaystyle \frac{1}{\delta }}\right)54\delta ^2\mathrm{ln}\left({\displaystyle \frac{1}{\delta }}\right)`$ (B.28)
$`+27\delta ^2\mathrm{ln}\left({\displaystyle \frac{1}{\delta }}\right)^2+3\delta ^2\mathrm{ln}\left({\displaystyle \frac{1}{3\delta }}\right)^224\mathrm{ln}(1+\delta )+134\delta \mathrm{ln}(1+\delta )`$
$`190\delta ^2\mathrm{ln}(1+\delta )24\delta ^2\mathrm{ln}(1+\delta )^254\delta \mathrm{ln}\left({\displaystyle \frac{1+\delta }{\delta }}\right)+54\delta ^2\mathrm{ln}\left({\displaystyle \frac{1+\delta }{\delta }}\right)`$
$`3\mathrm{ln}\left(1+4\delta \right)26\delta \mathrm{ln}\left(1+4\delta \right)56\delta ^2\mathrm{ln}\left(1+4\delta \right)`$
$`6\delta ^2\mathrm{ln}\left(1+\delta \right)\mathrm{ln}\left(1+4\delta \right)+6\delta ^2\mathrm{ln}\left(1+\delta \right)\mathrm{ln}\left({\displaystyle \frac{1+4\delta }{3\delta }}\right)`$
$`+54\delta ^2\mathrm{Li}_2\left[{\displaystyle \frac{1}{\delta }}\right]+6\delta ^2\mathrm{Li}_2\left[{\displaystyle \frac{1\delta }{3\delta }}\right])`$
with $`\mathrm{Li}_2(x)`$ the dilogarithm function defined as $`\mathrm{Li}_2(z)=_{k=1,\mathrm{}}z^k/k^2`$ .
FIGURE CAPTIONS
Figure 1 : Geometry of the present calculations. A tracer particle with radius $`a`$ diffuses in a slab with thickness $`H`$.
Figure 2 : Numerical estimates of the reduced diffusion coefficient of a particle moving between a partially slipping wall ($`\delta =1`$ : full line, $`\delta =100`$ : dotted line) at $`z=0`$ and a no-slip wall at $`z=H`$, as a function of the position of the particle. From left to right, $`H/a=\mathrm{3,5,8,12,17}and22`$.
Figure 3 : Diffusion coefficient near a single planar wall as a function of the distance $`l`$, for various slip length. Numerical results (solid lines) are compared with the analytical solution (dashed line) in the low confinement limit $`l>>a`$. The slip length $`\delta `$ increases from bottom to top.
Figure 4 : Local diffusivity computed using the approximate analytical results Eq. (24) for $`\delta /a=10^1`$ (dahsed line), compared to the numerical results (solid line). See figure 2 for details and notations.
Figure 5 : Same as in Fig. 4 but for $`\delta /a=10^1`$. See figure 2 for details and notations.
Figure 6 : Flow description in the lubrication limit in the thin confined film
Figure 7 : Sketch of the force balance in the volume $`r<R_c`$.
Figure 8 : Numerical test of the lubrication calculations : (left) plot of the friction force difference $`\mathrm{\Delta }F=F_{wall}(\delta )F_{wall}(\delta =0)`$ normalized by the bulk value $`F_{\mathrm{}}=6\pi \eta aU`$, for a single wall, as a function of the distance $`l`$ to the wall. The solid line is the FEMLAB calculation, while the dashed line is the lubrication estimate; (right) Diffusion coefficient for a particle in the middle plane of the confined geometry between 2 identical partially slipping walls, in the high confinement limit $`H/2a1`$. A good agreement between the numerical and lubrication calculations is found when $`\delta 0`$ and $`H/2a1`$.
Figure 9 : Mean diffusion coefficient between a no-slip wall and a partially slipping wall ($`\delta `$) as a function of the gap $`H`$ for various slip lengths $`\delta `$. Full line : numerical results, dashed line : low confinement approximation (average of Eq. (24)).
Figure 10 : Same as figure 9, but plotted as a function of slip length $`\delta `$ for fixed confinement $`H`$. Full line : numerical results, dashed line : low confinement approximation (average of Eq. (24)).
Figure 11 : Negative pressure $`\mathrm{\Pi }(r)=rb(r)`$ rescaled by $`P_0=\eta U/a`$, as a function of the radial distance $`r`$. The minimum gap $`h_0`$ between the sphere and the solid surface is $`h_0=0.1a`$ and the slip length is $`\delta =a`$. The solid line is the numerical solution of the equation for the pressure using a ODE solver (Mathematica <sup>©</sup>). The dashed line is the approximate solution, Eq. (B.24), see text.
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# Quantum effective potential for 𝑈(1) fields on 𝐒²_𝐿×𝐒²_𝐿
## 1 Introduction
Fuzzy approximations of spacetime (like lattice regularizations) are designed for the study of gauge theories in the nonperturbative regime using Monte-Carlo simulations. They consist in replacing continuous manifolds by matrix algebras. The resulting field theory will thus only have a finite number of degrees of freedom and hence it is regularized. The claim is that this method has the advantage -in contrast with lattice- of preserving all continuous symmetries of the original action at the classical level .
Field theory on the fuzzy sphere is the most studied example in the literature. In perturbation theory it is shown that scalar field theories on $`𝐒_L^2`$ suffer from the UV-IR mixing problem . Moreover it is shown that there exists new nonperturbative phenomena which are missing in the commutative theory. For example a novel phase has been discovered in scalar field theories on $`𝐒_L^2`$ (the so-called non-uniform phase or matrix phase) which has no commutative analogue . This new phase was also observed in three dimensions . Generalization to $`4`$dimensional fuzzy spaces and their scalar field theories were undertaken in .
The quantum properties of the gauge field on the fuzzy sphere have been studied in . In the effective action was computed to one loop for $`U(1)`$ gauge fields. It was shown that the model contains a gauge invariant UV-IR mixing in the limit $`L\mathrm{}`$, i.e the effective action does not go over to the commutative action in the continuum limit. Furthermore a first order phase transition was observed at one-loop from the fuzzy sphere phase to a matrix phase where the sphere collapses. This transition was previously detected in Monte Carlo simulation of the model reported in . In some sense the one-loop result for the $`U(1)`$ model is exact.
It was also shown in that the UV-IR mixing and the matrix phase both disappear in the limit where we send the mass of the normal scalar component of the gauge field on $`𝐒_L^2`$ to infinity. This means in particular that the nonperturbative $`𝐒_L^2`$-to-matrix phase transition is a reflection of the UV-IR mixing seen in perturbation theory and that this latter finds its origin in the coupling of the normal scalar field to the two dimensional gauge field. The differential calculus on the fuzzy sphere is intrinsically $`3`$dimensional and as a consequence there is no a gauge-covariant splitting of the $`3`$dimensional fuzzy gauge field into its normal and tangent components on $`𝐒_L^2`$; hence the action will necessarily involve the interaction of the two fields. This result (among many others) was confirmed recently in our Monte Carlo simulation of the model where we have also found a novel third-order one-plaquette-like phase transition which the model undergoes and which we can also trace to the coupling of the normal scalar field. The full phase diagram of the model will be reported elsewhere .
The main goal of this article is to study the phase structure of $`U(1)`$ gauge theories on fuzzy $`𝐒_L^2\times 𝐒_L^2`$. The advantage of considering $`𝐒_L^2\times 𝐒_L^2`$ is that one can use all the machinery of the well known $`SU(2)`$ polarization tensors. Other studies of noncommutative gauge theories on $`4`$dimensional fuzzy spaces have already appeared .
This article is organized as follows. In section $`2`$ we give a brief description of the geometry of fuzzy $`𝐒_L^2\times 𝐒_L^2`$. In section $`3`$ we introduce fuzzy gauge fields and we write down the action we will study in this article. In section $`4`$ we compute the effective potential. In section $`5`$ we show the existence of a first order $`𝐒_L^2\times 𝐒_L^2`$-to-matrix phase transition in exact analogy with the two-dimensional case and we derive the critical line. Section $`6`$ contains the conclusion.
## 2 Fuzzy $`𝐒_L^2\times 𝐒_L^2`$
Fuzzy $`𝐒_L^2\times 𝐒_L^2`$ is the simpliest $`4`$ dimensional fuzzy space. It is a finite dimensional matrix approximation of the cartesian product of two continuous spheres. This fuzzy space is defined by a sequence of Connes triples
$$𝐒_L^2\times 𝐒_L^2=\{Mat_{\left(L+1\right)^2},H_L,\mathrm{\Delta }_L\}.$$
(1)
$`Mat_{\left(L+1\right)^2}`$ is the matrix algebra of dimension $`\left(L+1\right)^2`$ and $`\mathrm{\Delta }_L`$ is a suitable Laplacian acting on matrices which encodes the geometry of the space. It is defined by
$$\mathrm{\Delta }_L[L_{AB},[L_{AB},]]=_{AB}^2$$
(2)
where $`L_{AB}`$, with $`A,B=\overline{1,4}`$, are the generators of the irreducible representation $`(\frac{L}{2},\frac{L}{2})`$ of $`SO\left(4\right)`$. The generators $`L_{AB}`$ (with $`L_{AB}=L_{BA}`$) satisfy the commutation relations
$`[L_{AB},L_{CD}]`$ $`=`$ $`f_{ABCDEF}L_{EF}`$ (3)
$``$ $`\delta _{BC}L_{AD}\delta _{BD}L_{AC}+\delta _{AD}L_{BC}\delta _{AC}L_{BD}.`$
$`H_L`$ in (1) is the Hilbert space (with inner product $`<M,N>=\frac{1}{\left(L+1\right)^2}Tr\left(M^{}N\right)`$) which is associated with the irreducible representation $`(\frac{L}{2},\frac{L}{2})`$ of $`SO(4)`$.
Since $`SO(4)=\left[SU(2)\times SU(2)\right]/Z_2`$ we can introduce $`SU(2)`$ (mutually commuting) generators $`L_a^{(1)}`$ and $`L_a^{(2)}`$ by $`2L_a^{\left(1\right)}=\frac{1}{2}ϵ_{abc}L_{bc}+L_{a4}`$ and $`2L_a^{\left(2\right)}=\frac{1}{2}ϵ_{abc}L_{bc}L_{a4}`$ with $`a=1,2,3`$ and $`ϵ_{abc}`$ is the three dimensional Levi-Civita tensor. Then it can be easily shown that the two $`SO(4)`$ quadratic Casimir can be rewritten in the form (where $`ϵ_{ABCD}`$ is the four dimensional Levi-Civita tensor)
$`L_{AB}^2`$ $`=`$ $`4\left[(L_a^{(1)})^2+(L_a^{(2)})^2\right]=2L(L+2)8c_2`$
$`ϵ_{ABCD}L_{AB}L_{CD}`$ $`=`$ $`8\left[(L_a^{(1)})^2(L_a^{(2)})^2\right]0.`$ (4)
Similarly the Laplacian $`_{AB}^2`$ reads in terms of the three dimensional indices as follows
$$_{AB}^2=4\left[\left(_a^{\left(1\right)}\right)^2+\left(_a^{\left(2\right)}\right)^2\right],$$
(5)
where $`_a^{\left(1\right)}[L_a^{\left(1\right)},]`$ and $`_a^{\left(2\right)}[L_a^{\left(2\right)},]`$. For $`𝐒_L^2\times 𝐒_L^2`$ the algebra $`Mat_{\left(L+1\right)^2}`$ is generated by the coordinate operators
$$x_a^{\left(1\right)}=R_1\frac{L_a^{\left(1\right)}}{\sqrt{c_2}},x_a^{\left(2\right)}=R_2\frac{L_a^{\left(2\right)}}{\sqrt{c_2}}$$
(6)
which satisfy
$$\underset{a=1}{\overset{3}{}}\left(x_a^{\left(i\right)}\right)^2=R_i^2\mathrm{𝟏},[x_a^{\left(i\right)},x_b^{\left(j\right)}]=\frac{iR_i}{\sqrt{c_2}}\delta _{ij}ϵ_{abc}x_c^{\left(i\right)},i=1,2.$$
(7)
In the limit $`L\mathrm{}`$ keeping $`R_1`$ and $`R_2`$ fixed we recover the commutative algebra of functions on $`𝐒^2\times 𝐒^2`$. If we also choose to scale the radii $`R_1`$ and $`R_2`$ such as for example $`\theta _1^2=R_1^2/L_1`$ and $`\theta _2^2=R_2^2/L_2`$ are kept fixed we obtain the non-commutative Moyal-Weyl space $`_{\theta _1}^2\times _{\theta _2}^2`$ .
The algebra of matrices $`Mat_{(L+1)^2}`$ can be decomposed under the action of the two $`SU(2)`$ of $`SO(4)`$ as $`Mat_{\left(L+1\right)}Mat_{\left(L+1\right)}`$. As a consequence a general function on $`𝐒_L^2\times 𝐒_L^2`$ can be expanded in terms of polarization tensors as follows
$`\varphi ={\displaystyle \underset{k_1=0}{\overset{L}{}}}{\displaystyle \underset{m_1=k_1}{\overset{k_1}{}}}{\displaystyle \underset{k_2=0}{\overset{L}{}}}{\displaystyle \underset{m_2=k_2}{\overset{k_2}{}}}\varphi _{k_1m_1k_2m_2}\widehat{Y}_{k_1m_1}\widehat{Y}_{k_2m_2}.`$ (8)
## 3 Fuzzy gauge fields
$`U(n)`$ gauge field on $`𝐒_L^2\times 𝐒_L^2`$ can be associated with a set of six hermitian matrices $`D_{AB}Mat_{n(L+1)^2}`$ ($`D_{AB}=D_{BA}`$) which transform homogeneously under the action of the group, i.e
$$D_{AB}UD_{AB}U^1,UU\left(n(L+1)^2\right).$$
(9)
In this paper we will be mainly interested in $`U(1)`$ theory on $`S_L^2\times S_L^2`$. The action is given by (with $`Tr_L=\frac{1}{\left(L+1\right)^2}Tr`$, $`g`$ is the gauge coupling constant and $`m`$ is the mass of the normal components of the gauge field )
$`S`$ $`=`$ $`{\displaystyle \frac{1}{16g^2}}\left\{{\displaystyle \frac{1}{4}}Tr_L[D_{AB},D_{CD}]^2+{\displaystyle \frac{i}{3}}f_{ABCDEF}Tr_L[D_{AB},D_{CD}]D_{EF}\right\}`$ (10)
$`+`$ $`{\displaystyle \frac{m^2}{8g^2L_{AB}^2}}Tr_L(D_{AB}^2L_{AB}^2)^2+{\displaystyle \frac{m^2}{32g^2L_{AB}^2}}Tr_L(ϵ_{ABCD}D_{AB}D_{CD})^2.`$
The equations of motion are given by
$`i[D_{CD},F_{AB,CD}]+{\displaystyle \frac{4m^2}{\sqrt{c_2}}}\{D_{AB},\mathrm{\Phi }_1+\mathrm{\Phi }_2\}+{\displaystyle \frac{m^2}{\sqrt{c_2}}}\{ϵ_{ABCD}D_{CD},\mathrm{\Phi }_1\mathrm{\Phi }_2\}=0.`$ (11)
As we will see shortly $`F_{AB,CD}=i[D_{AB},D_{CD}]+f_{ABCDEF}D_{EF}`$ can be interpreted as the curvature of the gauge field on fuzzy $`𝐒_L^2\times 𝐒_L^2`$ whereas $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ (defined by $`D_{AB}^2L_{AB}^2=8\sqrt{c_2}(\mathrm{\Phi }_1+\mathrm{\Phi }_2)`$ and $`ϵ_{ABCD}D_{AB}D_{CD}=16\sqrt{c_2}(\mathrm{\Phi }_1\mathrm{\Phi }_2)`$) can be interpreted as the normal components of the gauge field on $`𝐒_L^2\times 𝐒_L^2`$.
The most obvious non-trivial solution of the equations of motion (11) must satisfy $`F_{AB,CD}=0`$, $`D_{AB}^2=L_{AB}^2`$ and $`ϵ_{ABCD}D_{AB}D_{CD}=0`$ (or equivalently $`F_{AB}=0`$, $`\mathrm{\Phi }_i=0`$). This solution is clearly given by the generators $`L_{AB}`$ of the irreducible representation $`(\frac{L}{2},\frac{L}{2})`$ of $`SO(4)`$, viz
$`D_{AB}=L_{AB}.`$ (12)
As it turns out this is also the absolute minimum of the model. By expanding $`D_{AB}`$ around this vacuum as $`D_{AB}=L_{AB}+A_{AB}`$ and substituting back into the action (10) we obtain a $`U(1)`$ gauge field $`A_{AB}`$ on $`𝐒_L^2\times 𝐒_L^2`$ with the correct transformation law under the action of the group, namely $`A_{AB}UA_{AB}U^1+U_{AB}U^1`$. The matrices $`D_{AB}`$ are thus the covariant derivatives on $`𝐒_L^2\times 𝐒_L^2`$. The curvature $`F_{AB,CD}`$ in terms of $`A_{AB}`$ takes the usual form $`F_{AB,CD}=i_{AB}A_{CD}i_{CD}A_{AB}+f_{ABCDEF}A_{EF}+i[A_{AB},A_{CD}]`$. The normal scalar fields in terms of $`A_{AB}`$ are on the other hand given by $`8\sqrt{c_2}(\mathrm{\Phi }_1+\mathrm{\Phi }_2)=L_{AB}A_{AB}+A_{AB}L_{AB}+A_{AB}^2`$ and $`16\sqrt{c_2}(\mathrm{\Phi }_1\mathrm{\Phi }_2)=ϵ_{ABCD}(L_{AB}A_{CD}+A_{AB}L_{CD}+A_{AB}A_{CD})`$.
We can verify this conclusion explicitly by introducing the matrices $`D_a^{(1)}=L_a^{(1)}+A_a^{(1)}`$ and $`D_a^{(2)}=L_a^{(2)}+A_a^{(2)}`$ defined by
$`D_a^{(1)}{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{2}}ϵ_{abc}D_{bc}+D_{a4}\right],D_a^{(2)}{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{2}}ϵ_{abc}D_{bc}D_{a4}\right].`$ (13)
Clearly $`D_a^{(1)}`$ ($`A_a^{(1)}`$) and $`D_a^{(2)}`$ ($`A_a^{(2)}`$) are the components of $`D_{AB}`$ ($`A_{AB}`$) on the two spheres respectively. The curvature becomes $`F_{ab}^{(i,j)}=i_a^{(i)}A_b^{(j)}i_b^{(j)}A_a^{(i)}+\delta _{ij}ϵ_{abc}A_c^{(i)}+i[A_a^{(i)},A_b^{(j)}]`$ whereas the normal scalar fields become $`2\sqrt{c_2}\mathrm{\Phi }_i=(D_a^{(i)})^2c_2=L_a^{(i)}A_a^{(i)}+A_a^{(i)}L_a^{(i)}+(A_a^{(i)})^2`$. In terms of this three dimensional notation the action (10) reads
$`S=S^{(1)}+S^{(2)}+{\displaystyle \frac{1}{2g^2}}Tr_L\left(F_{ab}^{(1,2)}\right)^2.`$ (14)
$`S^{(1)}`$ and $`S^{(2)}`$ are the actions for the $`U(1)`$ gauge fields $`A_a^{(1)}`$ and $`A_a^{(2)}`$ on a single fuzzy sphere $`𝐒_L^2`$. They are given by
$`S^{(i)}`$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}Tr_L\left(F_{ab}^{(i,i)}\right)^2{\displaystyle \frac{1}{2g^2}}ϵ_{abc}Tr_L\left[{\displaystyle \frac{1}{2}}F_{ab}^{(i,i)}A_c^{(i)}{\displaystyle \frac{i}{6}}[A_a^{(i)},A_b^{(i)}]A_c^{(i)}\right]+{\displaystyle \frac{2m^2}{g^2}}Tr_L\mathrm{\Phi }_i^2.`$
It is immediately clear that in the continuum limit $`L\mathrm{}`$ the action (14) describes the interaction of a genuine $`4`$d gauge field with the normal scalar fields $`\mathrm{\Phi }_i=n_a^{(i)}A_a^{(i)}`$ where $`n_a^{(i)}`$ is the unit normal vector to the $`i`$-th sphere. The parameter $`m`$ is precisely the mass of these scalar fields. Let us also remark that in this limit the $`3`$dimensional fields $`A_a^{(i)}`$ decompose as $`A_a^{(i)}=(A_a^{(i)})^T+n_a^{(i)}\mathrm{\Phi }_i`$ where $`(A_a^{(i)})^T`$ are the tangent $`2`$dimensional gauge fields. Since the differential calculus on $`𝐒_L^2\times 𝐒_L^2`$ is intrinsically $`6`$dimensional we can not decompose the fuzzy gauge field in a similar (gauge-covariant) fashion and as a consequence we can not write an action on the fuzzy $`𝐒_L^2\times 𝐒_L^2`$ which will only involve the desired $`4`$dimensional gauge field.
## 4 Quantum effective potential
The partition function of the theory depends on $`3`$ parameters, the Yang-Mills coupling constant $`g`$, the mass $`m`$ of the normal scalar fields, and the size $`L`$ of the matrices, viz
$`Z_L[J,g,m]={\displaystyle \underset{A<B=1}{\overset{4}{}}\left[dX_{AB}\right]e^{S[X]+Tr_L\left[J_{AB}X_{AB}\right]}}.`$ (16)
In the background field method the field is decomposed as $`X_{AB}=D_{AB}+Q_{AB}`$ where $`D_{AB}`$ is the background we are interested in studying and $`Q_{AB}`$ stands for the fluctuation field. We add the usual gauge fixing and Faddeev-Popov terms given by
$`S_{g.f}+S_{gh}={\displaystyle \frac{1}{32g^2}}Tr_L{\displaystyle \frac{[D_{AB},Q_{AB}]^2}{\xi }}+{\displaystyle \frac{1}{16g^2}}Tr_Lc[D_{AB},[D_{AB},b]].`$ (17)
Performing the Gaussian path integral we obtain the one-loop effective action
$$\mathrm{\Gamma }\left[D_{AB}\right]=S\left[D_{AB}\right]+\frac{1}{2}Tr_6TR\mathrm{log}\mathrm{\Omega }_{ABCD}TR\mathrm{log}𝒟_{AB}^2.$$
(18)
$`\mathrm{\Omega }_{ABCD}`$ is defined by
$`\mathrm{\Omega }_{ABCD}`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒟_{EF}^2\delta _{AB,CD}\left(1{\displaystyle \frac{1}{\xi }}\right)𝒟_{AB}𝒟_{CD}2i_{ABCD}+{\displaystyle \frac{4m^2}{L_{AB}^2}}\mathrm{\Omega }_{ABCD}^{\left(1\right)},`$ (19)
where $`\delta _{AB,CD}=\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC}`$, and
$`\mathrm{\Omega }_{ABCD}^{\left(1\right)}`$ $`=`$ $`(D_{EF}^2L_{EF}^2)\delta _{AB,CD}+{\displaystyle \frac{1}{2}}(ϵ_{EFGH}D_{EF}D_{GH})ϵ_{ABCD}`$ (20)
$`𝒟_{AB}𝒟_{CD}\stackrel{~}{𝒟}_{AB}\stackrel{~}{𝒟}_{CD}+4D_{AB}D_{CD}+4\stackrel{~}{D}_{AB}\stackrel{~}{D}_{CD}.`$
The notation $`𝒟_{AB}`$ and $`_{ABCD}`$ means that the covariant derivative $`D_{AB}`$ and the curvature $`F_{ABCD}`$ act by commutators, i.e $`𝒟_{AB}(M)=[D_{AB},M]`$, $`_{ABCD}(M)=[F_{ABCD},M]`$ where $`M`$ is an element of $`Mat_{(L+1)^2}`$. Wwe have also introduced the notation $`\stackrel{~}{D}_{AB}\frac{1}{2}ϵ_{ABCD}D_{CD}`$. $`TR`$ is the trace over the $`4`$ indices corresponding to the left and right actions of operators on matrices. $`Tr_6`$ is the trace associated with the action of $`SU(2)\times SU(2)`$.
The main goal of this article is to check the stability of the solution (12), in other words to check whether or not the fuzzy space $`𝐒_L^2\times 𝐒_L^2`$ is stable under quantum fluctuations. Towards this end it is sufficient to consider only the background field $`D_{AB}=\varphi L_{AB}`$ where the order parameter $`\varphi `$ plays the role of the radius of the two spheres of $`𝐒_L^2\times 𝐒_L^2`$. Therefore the computation of the effective action reduces to the computation of the effective potential $`V_{\mathrm{eff}}(\varphi )\mathrm{\Gamma }[\varphi L_{AB}]`$. The classical potential is given by
$$VS[\varphi L_{AB}]=\frac{L(L+2)}{g^2}\left(\frac{1}{4}\varphi ^4\frac{1}{3}\varphi ^3+\frac{1}{4}m^2\left(\varphi ^21\right)^2\right).$$
(21)
The effective potential (in the gauge $`\xi =1`$) is given by
$`V_{\mathrm{eff}}`$ $`=`$ $`V+{\displaystyle \frac{1}{2}}Tr_6TR\mathrm{log}\varphi ^2TR\mathrm{log}\varphi ^2+{\displaystyle \frac{1}{2}}Tr_6TR\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ABCD}`$ (22)
$`=`$ $`V+4(L+1)^4\mathrm{log}\varphi +{\displaystyle \frac{1}{2}}Tr_6TR\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ABCD}.`$
We are only interested in the $`\varphi `$dependence of the operator $`\stackrel{~}{\mathrm{\Omega }}`$ which is defined by
$`\stackrel{~}{\mathrm{\Omega }}_{ABCD}={\displaystyle \frac{1}{2}}_{EF}^2\delta _{AB,CD}+2i\left(1{\displaystyle \frac{1}{\varphi }}\right)f_{ABCDEF}_{EF}+{\displaystyle \frac{4m^2}{L_{AB}^2}}\stackrel{~}{\mathrm{\Omega }}_{ABCD}^{(1)},`$ (23)
where
$`\stackrel{~}{\mathrm{\Omega }}_{ABCD}^{\left(1\right)}`$ $`=`$ $`\left(1{\displaystyle \frac{1}{\varphi ^2}}\right)L_{EF}^2\delta _{AB,CD}_{AB}_{CD}\stackrel{~}{}_{AB}\stackrel{~}{}_{CD}+4L_{AB}L_{CD}+4\stackrel{~}{L}_{AB}\stackrel{~}{L}_{CD}.`$
We will need to use the following identities
$`X_{AB}Y_{AB}`$ $`=`$ $`4\left(X_a^{\left(1\right)}Y_a^{\left(1\right)}+X_a^{\left(2\right)}Y_a^{\left(2\right)}\right),`$
$`f_{ABCDEF}Tr\left[X_{AB}Y_{CD}Z_{EF}\right]`$ $`=`$ $`16ϵ_{abc}Tr\left[X_a^{\left(1\right)}Y_b^{\left(1\right)}Z_c^{\left(1\right)}+X_a^{\left(2\right)}Y_b^{\left(2\right)}Z_c^{\left(2\right)}\right].`$ (25)
The matrices $`X_a^{\left(i\right)}`$ ($`Y_a^{\left(i\right)}`$) are related to the matrices $`X_{AB}`$ ($`Y_{AB}`$) by equations of the form (13). Using these identities we can express the last term in (22) in the following way
$`{\displaystyle \frac{1}{2}}Tr_6TR\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ABCD}`$ $`=`$ $`{\displaystyle 𝑑X_{AB}e^{TrX_{AB}\stackrel{~}{\mathrm{\Omega }}_{ABCD}X_{CD}}}`$ (26)
$`=`$ $`\left[{\displaystyle 𝑑X_a^{(1)}e^{2TrX_a^{(1)}\stackrel{~}{\mathrm{\Omega }}_{ab}X_b^{(1)}}}\right]^2`$
$`=`$ $`Tr_3TR\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ab}.`$
The contributions coming from the two spheres are equal and hence the factor of $`1`$ (instead of $`\frac{1}{2}`$) in front of the last logarithm. $`Tr_3`$ is the trace associated with the action of $`SU(2)`$ on the two dimensional sphere. The Laplacian $`\stackrel{~}{\mathrm{\Omega }}_{ab}`$ is defined by
$`\stackrel{~}{\mathrm{\Omega }}_{ab}`$ $`=`$ $`2_{AB}^2\delta _{ab}+16\left(1{\displaystyle \frac{1}{\varphi }}\right)iϵ_{abc}_c^{(1)}+8m^2\stackrel{~}{\mathrm{\Omega }}_{ab}^{(1)},`$
$`\stackrel{~}{\mathrm{\Omega }}_{ab}^{(1)}`$ $`=`$ $`4P_{ab}^{(1)}{\displaystyle \frac{1}{c_2}}_a^{(1)}_b^{(1)}+2\left(1{\displaystyle \frac{1}{\varphi ^2}}\right)\delta _{ab}.`$ (27)
$`P_{ab}^{(1)}`$ is the normal projector on the fuzzy sphere defined by $`P_{ab}^{(1)}=x_a^{\left(1\right)}x_b^{\left(1\right)}`$ where $`x_a^{\left(1\right)}`$ are the coordinate operators defined in (6) with $`R_1=R_2=1`$. The presence of this projector means in particular that we can not diagonalize in the polarization tensors basis. However, in order to have an idea of the phase structure of the model, we can expand around $`m=0`$. This approximation was more than sufficient in the two-dimensional case as discussed in great detail in . Therefore it is convenient to separate the logarithm term as
$$\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ab}=\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ab}^{(0)}+\mathrm{log}\left(1+8m^2\left(\frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}\right)_{ac}\stackrel{~}{\mathrm{\Omega }}_{cb}^{(1)}\right).$$
(28)
$`\stackrel{~}{\mathrm{\Omega }}_{ab}^{(0)}`$ is clearly equal to $`\stackrel{~}{\mathrm{\Omega }}_{ab}`$ when $`m^2=0`$. This operator can be trivially diagonalized in the vector polarization tensors basis $`(\widehat{Y}_{l_1}^{j_1M_1})_a`$ on the first sphere tensor product the scalar polarization tensors basis $`\widehat{Y}_{l_2m_2}`$ on the second sphere. Indeed by introducing the total angular momentum on the two-dimensional sphere $`𝒥_a^{(1)}=_a^{(1)}+\theta _a^{(1)}`$ where $`\theta _a^{(1)}`$ are the generators of $`SU(2)`$ in the spin $`1`$ irreducible representation we can rewrite $`\stackrel{~}{\mathrm{\Omega }}_{ab}^{(0)}`$ in the following form
$$\frac{1}{8}\stackrel{~}{\mathrm{\Omega }}_{ab}^{(0)}=(_c^{(1)})^2\delta _{ab}+(_c^{(2)})^2\delta _{ab}\left(1\frac{1}{\varphi }\right)\left[(𝒥_c^{(1)})_{ab}^2(_c^{(1)})^2\delta _{ab}2\delta _{ab}\right].$$
(29)
Hence it is convenient to use the following expansion for the matrices $`X_a^{\left(1\right)}`$ in (26)
$$X_a^{\left(1\right)}=\underset{j_1M_1\mathrm{}_1}{}\underset{\mathrm{}_2m_2}{}q_{\mathrm{}_2m_2}^{j_1M_1\mathrm{}_1}\left(\widehat{Y}_\mathrm{}_1^{j_1M_1}\right)_a\widehat{Y}_{\mathrm{}_2m_2}.$$
(30)
Thus
$`Tr_3TR\mathrm{log}\stackrel{~}{\mathrm{\Omega }}_{ab}^{(0)}={\displaystyle \underset{\mathrm{}_1j_1\mathrm{}_2}{}}\left(2j_1+1\right)\left(2\mathrm{}_2+1\right)\mathrm{log}\left[12\left(1{\displaystyle \frac{1}{\varphi }}\right){\displaystyle \frac{j_1\left(j_1+1\right)\mathrm{}_1\left(\mathrm{}_1+1\right)2}{\mathrm{}_1\left(\mathrm{}_1+1\right)+\mathrm{}_2\left(\mathrm{}_2+1\right)}}\right]`$
(31)
In the limit $`L\mathrm{}`$ it is easily verifiable (for example by making an expansion in $`1\frac{1}{\varphi }`$) that this term is subleading compared to $`L^4`$. The second contribution in the limit $`m0`$ is given by
$`Tr_3TR\mathrm{log}\left(1+8m^2\left({\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ac}\stackrel{~}{\mathrm{\Omega }}_{cb}^{(1)}\right)32m^2Tr_3TR\left({\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ac}x_c^{(1)}x_b^{(1)}`$
$`{\displaystyle \frac{8m^2}{c_2}}Tr_3TR\left({\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ac}_c^{(1)}_b^{(1)}+16m^2\left(1{\displaystyle \frac{1}{\varphi ^2}}\right)Tr_3TR\left({\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ab}.`$
In the large $`L`$ limit it is possible to show (see the appendix) that all terms in (LABEL:contri2) are subleading compared to the $`L^4`$ behaviour seen in the second term in (22) and hence the full one-loop quantum contribution to the effective potential is given by the logarithmic potential in (22). Thus as long as we are in the region of the phase space near $`m0`$ the effective potential behaves in the large $`L`$ limit as follows
$$\frac{V_{\mathrm{eff}}}{4L^4}=\frac{1}{4g^2L^2}\left(\frac{1}{4}\varphi ^4\frac{1}{3}\varphi ^3+\frac{1}{4}m^2\left(\varphi ^21\right)^2\right)+\mathrm{log}\varphi .$$
(33)
This result is to be compared with the quantum effective potential for $`U(1)`$ gauge fields on a single fuzzy sphere $`𝐒_L^2`$ computed in which is given explicitly by
$`{\displaystyle \frac{V_{\mathrm{eff}}}{L^2}}={\displaystyle \frac{1}{2g^2}}\left[{\displaystyle \frac{1}{4}}\varphi ^4{\displaystyle \frac{1}{3}}\varphi ^3+{\displaystyle \frac{1}{4}}m^2(\varphi ^21)^2\right]+\mathrm{log}\varphi .`$ (34)
$`V_{\mathrm{eff}}(\varphi )=2c_2N^2\alpha ^4\left[{\displaystyle \frac{1}{4}}\varphi ^4{\displaystyle \frac{1}{3}}\varphi ^3\right]+4c_2\mathrm{log}\varphi +\mathrm{subleading}\mathrm{terms}.`$ (35)
## 5 The $`𝐒_L^2\times 𝐒_L^2`$-to-matrix phase transition
The second term in the potential (33) is not convex. This implies that there is a competition between the classical potential and the logarithmic term which depends on the values of $`m`$ and $`g`$. The equation of motion $`\frac{V_{\mathrm{eff}}}{\varphi }=0`$ will admit in general two real solutions where the one with the least energy can be identified with the fuzzy $`𝐒_L^2\times 𝐒_L^2`$ solution (12). This equation of motion reads
$`(1+m^2)\varphi ^4\varphi ^3m^2\varphi ^2+4g^2L^2=0.`$ (36)
The quantum solution is found to be very close to $`1`$, viz
$`\varphi =1{\displaystyle \frac{4g^2L^2}{1+2m^2}}+O((g^2L^2)^2).`$ (37)
However this is only true up to an upper value of the gauge coupling constant $`g`$ (for every fixed value of $`m`$) beyond which the equation of motion ceases to have any real solutions. At this value the fuzzy $`𝐒_L^2\times 𝐒_L^2`$ collapses under the effect of quantum fluctuations and we cross to a pure matrix phase. In other words we can not define a gauge theory everywhere in the phase space. As we will see below when the mass $`m`$ is sent to infinity it is more difficult to reach the matrix phase and hence the presence of the mass makes the fuzzy $`𝐒_L^2\times 𝐒_L^2`$ solution (12) more stable.
The critical value can be computed by requiring that both the first and the second derivatives of the potential $`V_{\mathrm{eff}}`$ with respect to $`\varphi `$ vanish. In other words, for every fixed value of $`m`$ the critical point is defined at the point $`(g_{},m)`$ of the phase space where we go from a bounded potential to an unbounded potential. Solving for the critical value we get the results
$`\varphi _{}={\displaystyle \frac{3}{8(1+m^2)}}\left[1+\sqrt{1+{\displaystyle \frac{32m^2(1+m^2)}{9}}}\right],`$ (38)
and
$`2g_{}^2L^2={\displaystyle \frac{1}{2}}(1+m^2)\varphi _{}^4+{\displaystyle \frac{1}{2}}\varphi _{}^3+{\displaystyle \frac{m^2}{2}}\varphi _{}^2.`$ (39)
In the particular case of $`m^2=0`$ the critical value is
$$g_{}^2L^2=\frac{1}{2}\left(\frac{3}{8}\right)^3.$$
(40)
Extrapolating to large values of the mass ($`m\mathrm{}`$) we obtain the scaling behaviour
$`g_{}^2L^2={\displaystyle \frac{m^2+\sqrt{2}1}{16}}.`$ (41)
In figure $`1`$ we plot the phase diagram defined by this equation<sup>1</sup><sup>1</sup>1Notice that if we allow $`m^2`$ to take negative values, the gauge coupling constant $`g_{}^2`$ will be a more complicated function of $`m^2`$. However we are only interested in positive values of $`m^2`$ for which the behaviour of $`g_{}^2`$ as a function of $`m^2`$ is the straight line (41) which can be deduced from the large $`m^2`$ behaviour of (38) and (39). . As we increase the value of the coupling constant $`g`$ (for a fixed value of $`m^2`$) there exists a critical point $`g_{}`$ where the fuzzy $`𝐒_L^2\times 𝐒_L^2`$ solution becomes unstable and thus the minimum (12) disappears. Similarly as the value of the mass squared $`m^2`$ increases (for a fixed value of the coupling constant $`g`$) there is a critical point $`m_{}^2`$ where $`𝐒_L^2\times 𝐒_L^2`$ collapses. Clearly the value of $`m_{}^2`$ is found by inverting equation (41) , viz
$`m_{}^2=16g^2L^2+1\sqrt{2}.`$ (42)
Finally we remark that as the value of $`m^2`$ increases it is more difficult to reach the transition point, in fact when $`m^2\mathrm{}`$ the critical value $`g_{}^2`$ approaches infinity.
## 6 Conclusion
We have described the qualitative behaviour of a first order phase transition which occurs in a $`U(1)`$ gauge theory on $`𝐒_L^2\times 𝐒_L^2`$. Using the one-loop effective potential (33) of this theory we found that there exists values of the gauge coupling constant $`g`$ and the mass $`m`$ for which the fuzzy $`𝐒_L^2\times 𝐒_L^2`$ solution (12) is not stable. Thus for these values a $`U(1)`$ gauge theory on $`𝐒_L^2\times 𝐒_L^2`$ is not well defined. This means in particular that the model (10) can be used to approximate $`U(1)`$ gauge field theories on $`𝐒^2\times 𝐒^2`$ only deep inside the fuzzy sphere phase. However it is obvious from the critical line (41) that when the mass $`m`$ of the two normal scalar fields on $`𝐒_L^2\times 𝐒_L^2`$ goes to infinity it is more difficult to reach the transition line. Therefore we can say that our main goal of defining a nonperturbative regularization of a $`U(1)`$ gauge theory on $`𝐒^2\times 𝐒^2`$ is achieved. Generalization to $`U(n)`$ with and without fermions should be straightforward as long as we are only interested in the effective potential.
### Acknowledgements
The authors P. Castro-Villarreal and R. Delgadillo-Blando would like to thank Denjoe O’Connor for his supervision through the course of this study. B. Ydri would like to thank Denjoe O’Connor for his extensive discussions and critical comments while this research was in progress. The work of P. C. V. was supported by CONACYT México Grant No.43891-Y. The work of R. D. B. is supported by CONACYT México.
## Appendix A Evaluation of $`I_1`$, $`I_2`$ and $`I_3`$.
In this appendix we show that the $`3`$ terms in (4) are subleading compared to $`L^4`$. Let us define
$`I_1={\displaystyle \frac{1}{L^4}}Tr_3TR\left({\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ab},I_2={\displaystyle \frac{1}{L^6}}Tr_3TR\left({\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ab}_b^{\left(1\right)}_c^{\left(1\right)},`$
$`I_3={\displaystyle \frac{1}{L^4}}Tr_3TR\left({\displaystyle \frac{4}{\stackrel{~}{\mathrm{\Omega }}^{(0)}}}\right)_{ab}x_b^{\left(1\right)}x_c^{\left(1\right)}.`$ (43)
We evaluate these traces by using the base of polarization tensors $`\left(\widehat{Y}_\mathrm{}_1^{j_1M_1}\right)_a\widehat{Y}_{\mathrm{}_2m_2}`$. Using the identity $`_a_b=^2\delta _{ab}(\theta )_{ab}(\theta )_{ab}^2`$ the eigenvalues of the operator $`_a_b`$ are given by
$`\eta _{\mathrm{}_1j_1}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(j_1\left(j_1+1\right)\mathrm{}_1\left(\mathrm{}_1+1\right)\right)^2{\displaystyle \frac{1}{2}}\left(j_1\left(j_1+1\right)+\mathrm{}_1\left(\mathrm{}_1+1\right)\right),`$ (44)
whereas the eigenvalues of $`\stackrel{~}{\mathrm{\Omega }}_{ab}^{(0)}`$ given by (29) are
$`\lambda _\mathrm{}_2^{\mathrm{}_1j_1}`$ $`=`$ $`8\left(\mathrm{}_1\left(\mathrm{}_1+1\right)+\mathrm{}_2\left(\mathrm{}_2+1\right)\right)+8{\displaystyle \frac{1\varphi }{\varphi }}\left(j_1\left(j_1+1\right)\mathrm{}_1\left(\mathrm{}_1+1\right)2\right).`$ (45)
Using these facts the two quantities $`I_1`$ and $`I_2`$ can be shown to be given by
$$I_1=\frac{1}{L^4}\underset{\mathrm{}_1j_1\mathrm{}_2}{}\frac{\left(2l_2+1\right)\left(2j_1+1\right)}{\lambda _\mathrm{}_2^{\mathrm{}_1j_1}},I_2=\frac{1}{L^6}\underset{\mathrm{}_1j_1\mathrm{}_2}{}\frac{\left(2\mathrm{}_2+1\right)\left(2j_1+1\right)\eta _{\mathrm{}_1j_1}}{\lambda _\mathrm{}_2^{\mathrm{}_1j_1}}.$$
(46)
In order to evaluate $`I_3`$ we notice the fact that $`x_a^{(i)}`$ is proportional to $`\left(\widehat{𝐘}_1^{00}\right)_a`$ thus by using the algebra of vectorial polarization tensors we get the identity
$`\text{Tr}\left\{\left(𝐘_{\mathrm{}}^{jM}𝐘_1^{00}\right)\left(𝐘_1^{00}𝐘_{\mathrm{}}^{+jM}\right)\right\}=(L+1)(2\mathrm{}+1)\left\{\begin{array}{ccc}1& \mathrm{}& j\\ \frac{L}{2}& \frac{L}{2}& \frac{L}{2}\end{array}\right\}^2.`$ (49)
The final result for $`I_3`$ is
$$I_3=\frac{2}{L^4}\underset{\mathrm{}_1j_1\mathrm{}_2}{}\frac{2l_2+1}{\lambda _\mathrm{}_2^{\mathrm{}_1j_1}}\left[\left(L+1\right)\left(2\mathrm{}_1+1\right)\left\{\begin{array}{ccc}1& \mathrm{}_1& j\\ \frac{L}{2}& \frac{L}{2}& \frac{L}{2}\end{array}\right\}^2\right].$$
(50)
The large $`L`$ behaviour of $`I_1`$, $`I_2`$ and $`I_3`$ can be studied with the help of the different identities of . The first sum in (46) diverges at most as $`L^2`$ in the continuum $`L\mathrm{}`$ limit and hence $`I_1`$ converges to zero as $`1/L^2`$. On the other hand the sum in $`I_2`$ behaves at most as $`L^4`$ thus the whole expression goes to zero as $`1/L^2`$. For $`I_3`$ we can check that the sum goes as $`L`$, i.e $`I_3`$ appraoches $`0`$ as $`1/L^3`$.
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# Large Non-perturbative Effects of Small Δ𝑚²₂₁/Δ𝑚²₃₁ and sin{𝜃₁₃} on Neutrino Oscillation and CP Violation in Matter
## 1 Introduction
Last year, the direct evidence for neutrino oscillation was found in three kinds of experiments, namely atmospheric neutrino experiment , reactor neutrino experiment and K2K experiment . In these experiments, the dip of neutrino oscillation and the energy dependence of the probability, were observed. The possibilities of neutrino decay and neutrino decoherence are excluded by these results and it is found that the only solution for the solar and the atmospheric neutrino problem is neutrino oscillation. The observation of the dip also means that the neutrino experiments herald in a new era of precise measurements, because the effect, which disappears by averaging out the time-varying part on the neutrino energy, has been observed in these experiments for the first time . Solar neutrino parameters have been also accurately determined by recent neutrino experiments such as SNO and KamLAND .
From the results of the past experiments, it was found that the solar neutrino deficit is explained by the Large Mixing Angle (LMA) MSW solution ,
$`\mathrm{\Delta }m_{21}^28.0\times 10^5\mathrm{eV}^2,\mathrm{sin}^22\theta _{12}0.8,`$ (1.1)
where the mass squared difference is defined by $`\mathrm{\Delta }m_{ij}^2=m_i^2m_j^2`$. It was obtained that
$`|\mathrm{\Delta }m_{31}^2|2.0\times 10^3\mathrm{eV}^2,\mathrm{sin}^22\theta _{23}1.0`$ (1.2)
from the atmospheric neutrino experiment . Furthermore, the upper bound of the 1-3 mixing angle, $`\mathrm{sin}\theta _{13}`$ is given by
$`\mathrm{sin}^22\theta _{13}0.2`$ (1.3)
from the reactor experiment . The next step for neutrino physics is the determination of $`\mathrm{sin}\theta _{13}`$, the sign of $`\mathrm{\Delta }m_{31}^2`$ and CP phase $`\delta `$. In particular, the measurement of the leptonic CP phase is one of the most important themes from the viewpoint of the origin of the universe. CP violation has been investigated also in quark sector for the first time and the Kobayashi-Maskawa theory has been established . However, it has been found that the CP violation in quark sector is too small to generate the sufficient baryon number in the universe , because the electroweak symmetry breaking is not the first phase transition as the Higgs particle is too heavy. This means that the origin of baryon asymmetry of the universe is not a CP violation from the KM phase and an extra source of CP violation is needed. One of the alternatives is the generation of a baryon number due to the leptonic CP violation . The possibility of this scenario has been investigated by many researchers .
In order to attain the next step, the long baseline experiments like superbeam experiments and neutrino factory experiments are planned. In these experiments, the earth matter effects disturb the observation of the CP violation because the matter in the earth is not CP invariant and generate the effects of fake CP violation. Therefore, it is very important to understand the earth matter effects in neutrino oscillation experiments.
Here, summarizing the results of the atmospheric, solar and reactor neutrino experiments, there are two small parameters
$`\alpha ={\displaystyle \frac{\mathrm{\Delta }m_{21}^2}{\mathrm{\Delta }m_{31}^2}}`$ $``$ $`0.04,`$ (1.4)
$`s_{13}=\mathrm{sin}\theta _{13}`$ $``$ $`0.23.`$ (1.5)
The magnitude of these small parameters is most important for measuring the CP violation, because it cannot be observed, if one of these parameters vanishes. As the LMA MSW solution was chosen to explain the results of the solar neutrino experiments, $`\alpha `$ reduced to the largest value compared to other solutions. This means that the LMA MSW solution opens the door for measuring the leptonic CP violation. If $`s_{13}`$ is too small, it will be impossible to observe the CP violation. Therefore the magnitude of $`s_{13}`$ controls whether the leptonic CP violation can be observed or not.
Let us briefly review the approximate formulas using the small parameter $`\alpha `$ or $`s_{13}`$ and the related papers. At first using the perturbation of oscillation probability in $`\alpha `$, the magnitude of the fake CP violation by the matter effects has been investigated in . Furthermore, by expanding the matter potential to the Fourier mode, it has been shown in that the mode with large wavelength mainly contributes to the oscillation probability. Higher order perturbative calculations have been performed by . The perturbation in $`s_{13}`$ has been investigated in . The perturbation in both $`\alpha `$ and $`s_{13}`$ has also been studied in and this method has been extended to all channels in .
Next let us review the remarkable features related to the leptonic CP violation. In the case of constant matter density, the notable identity $`\stackrel{~}{J}\stackrel{~}{\mathrm{\Delta }}_{12}\stackrel{~}{\mathrm{\Delta }}_{23}\stackrel{~}{\mathrm{\Delta }}_{31}=J\mathrm{\Delta }_{12}\mathrm{\Delta }_{23}\mathrm{\Delta }_{31}`$ has been found in , where $`J`$ is the Jarlskog factor related to the leptonic CP violation, $`\mathrm{\Delta }_{ij}`$ means $`\mathrm{\Delta }m_{ij}^2/(2E)`$ and tilde stands for the quantities in matter. In addition, it has been pointed out that the oscillation probability in matter almost coincide with that in vacuum under the certain condition, which is called vacuum mimicking phenomena, and the method to solve the problem on the fake CP violation by using the phenomena is discussed in detail . Furthermore, it can be applied to the future long baseline experiments by using the statistical method explained by .
In a previous series of papers we have considered the three generation neutrino oscillation in matter and have shown that the CP dependence of the oscillation probabilities are derived exactly . In the case that $`\nu _e`$ is included in the initial or final state, the CP dependence is given by
$`P(\nu _e\nu _e)`$ $`=`$ $`C_{ee},`$ (1.6)
$`P(\nu _\alpha \nu _\beta )`$ $`=`$ $`A_{\alpha \beta }\mathrm{cos}\delta +B_{\alpha \beta }\mathrm{sin}\delta +C_{\alpha \beta },`$ (1.7)
and in the case that both the initial and final state are $`\nu _\alpha ,\nu _\beta =\nu _\mu ,\nu _\tau `$, the CP dependence is given by
$`P(\nu _\alpha \nu _\beta )`$ $`=`$ $`A_{\alpha \beta }\mathrm{cos}\delta +B_{\alpha \beta }\mathrm{sin}\delta +C_{\alpha \beta }+D_{\alpha \beta }\mathrm{cos}2\delta +E_{\alpha \beta }\mathrm{sin}2\delta ,`$ (1.8)
where the coefficients $`A_{\alpha \beta }E_{\alpha \beta }`$ are independent of the CP phase. We have also shown that these coefficients can be calculated exactly in constant matter and then the approximate formulas are derived in a simple way . Furthermore, we proposed a new method for approximating these coefficients in the case of non-constant matter density , and then applied it to the earth matter .
In this paper, at first within the framework of two generations, it has been shown that perturbation of the small mixing angle is not effective near the MSW resonance point. This means that the non-perturbative effects by the small mixing angle is important in the MSW resonance region. Next, we consider non-perturbative effects of $`\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{sin}\theta _{13}`$ in the three generation neutrino oscillation. The importance of the non-perturbative effects is shown by comparing the exact numerical calculation with the perturbative expansion of the small parameters. Furthermore, we consider the method for deriving the approximate formulas in which the non-perturbative effects are taken into account. In our previous paper , the approximate formulas for $`P(\nu _e\nu _\mu )`$ have been derived. These formulas are effective for both MSW resonance regions. However, there is a problem because this method cannot be extended to other channels $`P(\nu _\mu \nu _\tau )`$ and so on. In order to solve this problem, we assume the two natural conditions, $`\theta _{23}=45^{}`$ and the symmetric matter potential. Under these conditions, we derive the approximate formulas for all channels, including non-perturbative effects of the two small parameters. These formulas are useful to solve the problem of parameter degeneracy.-
## 2 Non-perturbative Effect by Small Mixing Angle
In this section, we discuss the perturbative expansion of a small mixing angle in two generation neutrino oscillation. Although we discussed the perturbation of small parameters in our previous papers , in order to clarify the physical meaning, we consider the perturbation due to a small mixing angle within the framework of two generations. Then, we show that the perturbation breaks down in the MSW resonance region even if the mixing angle is small.
### 2.1 MSW Resonance of Probability in Two Generations
In this subsection, we consider the two generation neutrino oscillation and we choose the energy region and the baseline length in which the MSW resonance occurs. Let us start from the Hamiltonian in constant matter
$`H`$ $`=`$ $`O\mathrm{diag}(0,\mathrm{\Delta })O^T+\mathrm{diag}(a,0)`$ (2.1)
$`=`$ $`\stackrel{~}{O}\mathrm{diag}(\lambda _1,\lambda _2)\stackrel{~}{O}^T,`$ (2.2)
where the matter potential is defined by $`a=\sqrt{2}G_FN_e`$. $`G_F`$ is the Fermi constant and $`N_e`$ is the electron density in matter. The matrix $`O`$ is mixing matrix as
$`O=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),`$ (2.5)
where $`\mathrm{\Delta }=\mathrm{\Delta }m^2/2E`$ and the quantities with tilde stand for the quantities in matter. Diagonalizing (2.1) to (2.2), the effective masses $`\lambda _i(i=1,2)`$ and effective mixing angle $`\stackrel{~}{\theta }`$ are determined. If we use the notation $`\stackrel{~}{\mathrm{\Delta }}=\lambda _2\lambda _1`$ as the mass squared difference, there is a relation between the mass squared difference and the mixing angles as
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}}{\mathrm{\Delta }}}={\displaystyle \frac{\mathrm{sin}2\theta }{\mathrm{sin}2\stackrel{~}{\theta }}}=\sqrt{\left(\mathrm{cos}2\theta {\displaystyle \frac{a}{\mathrm{\Delta }}}\right)^2+\mathrm{sin}^22\theta }.`$ (2.6)
Using these quantities in matter, the oscillation probability is given by
$`P=\mathrm{sin}^22\stackrel{~}{\theta }\mathrm{sin}^2{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}L}{2}}.`$ (2.7)
The oscillating part with $`L/E`$ of this probability becomes large if the condition
$`\mathrm{sin}{\displaystyle \frac{\stackrel{~}{\mathrm{\Delta }}L}{2}}1`$ (2.8)
is satisfied. On the other hand, the condition for the maximal effective mixing angle is given by
$`\mathrm{sin}2\stackrel{~}{\theta }1.`$ (2.9)
In the case of small mixing angle, this condition is rewritten as $`a=\mathrm{\Delta }\mathrm{cos}2\theta \mathrm{\Delta }`$, and furthermore we define the resonance energy by
$`E{\displaystyle \frac{\mathrm{\Delta }m^2}{2a}}.`$ (2.10)
We also define the resonance length by
$`L{\displaystyle \frac{1}{a\mathrm{sin}\theta }}.`$ (2.11)
For the case of $`\mathrm{sin}\theta =0.16`$, which is the upper bound in the CHOOZ experiment, the resonance length is roughly estimated as 10000 km. This means that in near future it is impossible to realize the long baseline experiments such that the baseline length from beam source to the detector is nearly equal to the resonance length. However, it has been shown that matter effects exist even if the baseline length is shorter than the resonance length. Therefore, we use $`L=6000`$ km as the baseline length in the later sections.
### 2.2 Perturbation due to Small Mixing Angle
Next, let us consider the expansion of the effective mass $`\stackrel{~}{\mathrm{\Delta }}`$ and the effective mixing angle $`\mathrm{sin}2\stackrel{~}{\theta }`$ by a small mixing angle $`\mathrm{sin}\theta `$. We show that although the effective mass and the effective mixing angle diverge in the MSW resonance energy region, the oscillation probability, which is a function of these two quantities, converges.
At first, the effective mass is expanded as
$`\stackrel{~}{\mathrm{\Delta }}=|\mathrm{\Delta }a|+{\displaystyle \frac{2a\mathrm{\Delta }}{|\mathrm{\Delta }a|}}\mathrm{sin}^2\theta +{\displaystyle \frac{a^2\mathrm{\Delta }^2}{2|\mathrm{\Delta }a|^3}}\mathrm{sin}^4\theta +\mathrm{}.`$ (2.12)
One can see from this result that other terms than the first term diverge. The higher order term have larger divergence near the MSW resonance. The effective mixing angle is expanded as
$`\mathrm{sin}2\stackrel{~}{\theta }={\displaystyle \frac{\mathrm{\Delta }\mathrm{sin}2\theta }{|\mathrm{\Delta }a|}}\left(1{\displaystyle \frac{2a\mathrm{\Delta }}{(\mathrm{\Delta }a)^2}}\mathrm{sin}^2\theta +{\displaystyle \frac{3a^2\mathrm{\Delta }^2}{2(\mathrm{\Delta }a)^4}}\mathrm{sin}^4\theta +\mathrm{}\right).`$ (2.13)
The condition
$`\mathrm{sin}\theta <{\displaystyle \frac{|\mathrm{\Delta }a|}{2\sqrt{a\mathrm{\Delta }}}}`$ (2.14)
is needed for $`\mathrm{sin}2\stackrel{~}{\theta }`$ to converge the finite value. However, this condition cannot be satisfied in the MSW resonance region defined by $`\mathrm{\Delta }a`$, even if $`\mathrm{sin}\theta `$ is small. This means that the above perturbation series diverges. In the expansion for the effective mass and the effective mixing angle, the coefficients become large, even if these quantities are expanded by the small mixing angle.
Next, let us consider the oscillation probability and let us demonstrate that the oscillation probability reaches a finite value, where the divergences due to the effective mass and the effective mixing angle are canceled out by each other. Substituting (2.12) and (2.13) into (2.7), we obtain
$`P`$ $``$ $`{\displaystyle \frac{\mathrm{\Delta }^2\mathrm{sin}^22\theta }{(\mathrm{\Delta }a)^2}}\mathrm{sin}^2{\displaystyle \frac{(\mathrm{\Delta }a)L}{2}}`$ (2.15)
$`+`$ $`{\displaystyle \frac{\mathrm{\Delta }^2\mathrm{sin}^22\theta }{(\mathrm{\Delta }a)^2}}\left[{\displaystyle \frac{4a\mathrm{\Delta }\mathrm{sin}^2\theta }{(\mathrm{\Delta }a)^2}}\mathrm{sin}^2{\displaystyle \frac{(\mathrm{\Delta }a)L}{2}}+{\displaystyle \frac{a\mathrm{\Delta }L\mathrm{sin}^2\theta }{\mathrm{\Delta }a}}\mathrm{sin}(\mathrm{\Delta }a)L\right]+\mathrm{}.`$
In the limit, $`\mathrm{\Delta }a`$, it is found that the oscillation probability becomes finite as
$`P`$ $``$ $`\mathrm{cos}^2\theta \left(\mathrm{sin}^2\theta a^2L^2{\displaystyle \frac{1}{3}}\mathrm{sin}^4\theta a^4L^4+\mathrm{}\right).`$ (2.16)
From this equation, the oscillation probability becomes finite and the perturbation is a good approximation if
$`L<{\displaystyle \frac{1}{a\mathrm{sin}\theta }}.`$ (2.17)
As you see from (2.11), this is the condition that the baseline length is shorter than the resonance length.
Next, let us investigate the magnitude of non-perturbative effects numerically. We use the following parameters, $`\mathrm{\Delta }m^2=2.0\times 10^3\mathrm{eV}^2`$ and $`\mathrm{sin}\theta =0.16`$. We set the baseline length, $`L=6000\mathrm{km}`$ and the energy region, $`1\mathrm{GeV}E50\mathrm{GeV}`$, to include the MSW resonance energy. Furthermore we choose a density of $`\rho =4\mathrm{g}/\mathrm{cm}^3`$.
At first, in figure 1(a) the level crossing of two eigenvalues is plotted. It is shown that the crossing energy is about 6-7 GeV, which corresponds to the MSW resonance energy. Next, in figure 1(b) we compare the oscillation probability calculated by perturbation with the one by numerical calculation. These figures show that the perturbation breaks down around the MSW resonance energy. The results of this subsection are summarized as
1. The perturbative expansion in the small mixing angle breaks down around the MSW resonance because the perturbation because the perturbation series diverges. The coefficients of this expansion become larger around the MSW resonance. The divergence included in the effective mass cancels with that in the effective mixing angle, and as a result, the value of the oscillation probability reaches a finite value. Term of eq. (2.12) and (2.13) cancel with each other.
2. Although the divergences of the effective mass and the effective mixing angle in the perturbative expansion cancel in the oscillation probability, the finite value of the probability differs from that by numerical calculation. The perturbation around the MSW resonance energy becomes a good approximation, if the baseline length is shorter than the resonance length as seen from (2.17). However, we need to take higher order terms of the perturbation into account, when the baseline length is longer, namely when non-perturbative effects become important.
## 3 Extension of method to Approximate Oscillation Probabilities
In this section, we consider the matter effects in three generation neutrino oscillation. At first, we review that the 2-3 mixing angle $`\theta _{23}`$ and the CP phase $`\delta `$ can be separated from matter effects in the oscillation probability . This means that the matter effects appear through the remained four parameters. Furthermore, these four parameters can be separated to two set of parameters and each set is related to the phenomena in low and high energy as
$`(\theta _{12},\mathrm{\Delta }m_{21}^2):\mathrm{Low}\mathrm{Energy}\mathrm{Phenomenon}`$ (3.1)
$`(\theta _{13},\mathrm{\Delta }m_{31}^2):\mathrm{High}\mathrm{Energy}\mathrm{Phenomenon}.`$ (3.2)
This separation means that the parameters for the solar neutrino and those for the atmospheric neutrino are almost independent to each other. We propose the method deriving the approximate formulas simply by using this feature.
### 3.1 Definition of Low and High Energy Regions
In this subsection, we define the low energy and the high energy Hamiltonians in the small quantity limit when $`s_{13}`$ or $`\alpha `$ approximate zero. Although these Hamiltonian have been already introduced in our earlier papers , we review them here, as they are used in later section.
It is noted that $`H(t)`$ satisfies the relation
$$H(t)=O_{23}\mathrm{\Gamma }H^{}(t)\mathrm{\Gamma }^{}O_{23}^T$$
(3.3)
where $`H^{}`$ is given by
$$H^{}=O_{13}O_{12}\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})O_{12}^TO_{13}^T+\mathrm{diag}(a(t),0,0).$$
(3.4)
This means that the 1-2 and 1-3 mixing angles are separated from the 2-3 mixing and the CP phase, as explained in detail in Appendix A. In this Appendix A, we derive the same result as that derived from this section from another point of view. Taking the limit $`s_{13}0`$, the Hamiltonian reduces to the two generation Hamiltonian as
$`H^{\mathrm{}}`$ $`=`$ $`\underset{s_{13}0}{lim}H^{}`$ (3.5)
$`=`$ $`O_{12}\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})O_{12}^T+\mathrm{diag}(a(t),0,0)`$ (3.6)
$`=`$ $`\left(\begin{array}{ccc}\mathrm{\Delta }_{21}s_{12}^2+a(t)& \mathrm{\Delta }_{21}s_{12}c_{12}& 0\\ \mathrm{\Delta }_{21}s_{12}c_{12}& \mathrm{\Delta }_{21}c_{12}^2& 0\\ 0& 0& \mathrm{\Delta }_{31}\end{array}\right)`$ (3.10)
This means that the third generation is now separated from the first and the second generation. As seen from this Hamiltonian (3.10), the components except for $`H_{\tau \tau }^{\mathrm{}}`$, depend only on $`(\theta _{12},\mathrm{\Delta }_{21})`$. We call $`H^{\mathrm{}}`$ the low energy Hamiltonian. On the other hand, taking the limit $`\alpha 0`$, the Hamiltonian reduces to the two generation Hamiltonian as
$`H^h`$ $`=`$ $`\underset{\alpha 0}{lim}H^{}`$ (3.11)
$`=`$ $`O_{13}\mathrm{diag}(0,0,\mathrm{\Delta }_{31})O_{13}^T+\mathrm{diag}(a(t),0,0)`$ (3.12)
$`=`$ $`\left(\begin{array}{ccc}\mathrm{\Delta }_{31}s_{13}^2+a(t)& 0& \mathrm{\Delta }_{31}s_{13}c_{13}\\ 0& 0& 0\\ \mathrm{\Delta }_{31}s_{13}c_{13}& 0& \mathrm{\Delta }_{31}c_{13}^2\end{array}\right).`$ (3.16)
This means that the second generation is also separated from the two others. This Hamiltonian (3.16) is expressed by only the parameters $`(\theta _{13},\mathrm{\Delta }_{31})`$. We call $`H^h`$ high energy Hamiltonian. Next, let us define the high and low energy regions described by $`H^h`$ and $`H^{\mathrm{}}`$. We first calculate the MSW resonance energy because the MSW effect is the most important in matter effects. In the case of $`L=6000\mathrm{km}`$, which we use later, the average matter potential is calculated as $`\rho =3.91\mathrm{g}/\mathrm{cm}^3`$. By using this value, we obtain the high energy MSW resonance as $`E^h=\mathrm{\Delta }m_{31}^2/a5\mathrm{GeV}`$ and the low energy MSW resonance as $`E^{\mathrm{}}=\mathrm{\Delta }m_{21}^2/a0.2\mathrm{GeV}`$. From these results, we regard $`E1\mathrm{G}\mathrm{e}\mathrm{V}`$ as the boundary energy of low and high energy regions. Therefore, we define the high as $`E1\mathrm{GeV}`$ and the low energy regions as $`E1\mathrm{GeV}`$.
### 3.2 Order Counting of Amplitude on $`\alpha `$ and $`s_{13}`$
In this subsection, we investigate how the amplitude $`S^{}`$, which is defined by the primed Hamiltonian (3.4), depends on the two small parameters $`\alpha `$ and $`s_{13}`$. Before, we have already clarified some general features of $`S^{}`$ related to the order of $`\alpha `$ and $`s_{13}`$, and the dependences on $`s_{13}`$ and $`\alpha `$ for particular amplitudes $`S_{\mu e}^{}`$ and $`S_{\tau e}^{}`$ have been given in our previous papers . We investigate now the dependences on $`s_{13}`$ and $`\alpha `$ for all amplitudes.
At first, we represent the explicit form of the Hamiltonian, when the 2-3 mixing angle and the CP phase are factored out as
$`H^{}(t)`$ $`=`$ $`O_{13}O_{12}\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})O_{12}^TO_{13}^T+\mathrm{diag}(a(t),0,0)`$ (3.17)
$`=`$ $`\left(\begin{array}{ccc}\mathrm{\Delta }_{21}c_{13}^2s_{12}^2+\mathrm{\Delta }_{31}s_{13}^2+a(t)& \mathrm{\Delta }_{21}c_{13}s_{12}c_{12}& \mathrm{\Delta }_{21}c_{13}s_{13}s_{12}^2+\mathrm{\Delta }_{31}s_{13}c_{13}\\ \mathrm{\Delta }_{21}c_{13}s_{12}c_{12}& \mathrm{\Delta }_{21}c_{12}^2& \mathrm{\Delta }_{21}s_{13}s_{12}c_{12}\\ \mathrm{\Delta }_{21}c_{13}s_{13}s_{12}^2+\mathrm{\Delta }_{31}s_{13}c_{13}& \mathrm{\Delta }_{21}s_{13}s_{12}c_{12}& \mathrm{\Delta }_{21}s_{13}^2s_{12}^2+\mathrm{\Delta }_{31}c_{13}^2\end{array}\right).`$ (3.21)
The components of this Hamiltonian depend on $`\alpha `$ and $`s_{13}`$ as
$`H^{}(t)`$ $`=`$ $`\left(\begin{array}{ccc}O(1)& O(\alpha )& O(s_{13})\\ O(\alpha )& O(\alpha )& O(\alpha s_{13})\\ O(s_{13})& O(\alpha s_{13})& O(1)\end{array}\right).`$ (3.25)
From this result, we can see that non-diagonal components are small compared to the diagonal components. We also understand that $`H_{\mu \tau }^{}`$ is the smallest component and $`H_{e\mu }^{},H_{e\tau }^{}`$ are the next smaller components. We should note the salient feature that the result of this order counting holds in $`H^n`$ for arbitrary $`n`$. Namely, we obtain
$`(H^{}(t))^n`$ $`=`$ $`\left(\begin{array}{ccc}O(1)& O(\alpha )& O(s_{13})\\ O(\alpha )& O(\alpha ^2)& O(\alpha s_{13})\\ O(s_{13})& O(\alpha s_{13})& O(1)\end{array}\right)\mathrm{for}n=1,2,\mathrm{}.`$ (3.29)
According to this result, the order of the amplitude $`S^{}(t)`$ for two small parameters $`\alpha `$ and $`s_{13}`$ is given by
$`S^{}(t)=\mathrm{T}\mathrm{exp}\left\{i{\displaystyle H^{}(t)𝑑t}\right\}=\left(\begin{array}{ccc}O(1)& O(\alpha )& O(s_{13})\\ O(\alpha )& O(1)& O(\alpha s_{13})\\ O(s_{13})& O(\alpha s_{13})& O(1)\end{array}\right).`$ (3.33)
This result is almost the same as that of the original Hamiltonian. Furthermore, we consider the general features derived from the original Hamiltonian. The $`\theta _{13}`$ dependence of this Hamiltonian is described as
$$H^{}=\left(\begin{array}{ccc}\mathrm{even}& \mathrm{even}& \mathrm{odd}\\ \mathrm{even}& \mathrm{even}& \mathrm{odd}\\ \mathrm{odd}& \mathrm{odd}& \mathrm{even}\end{array}\right)$$
(3.34)
and this dependence does not change for $`(H^{})^n`$, because
$$(H^{})^n=\left(\begin{array}{ccc}\mathrm{even}& \mathrm{even}& \mathrm{odd}\\ \mathrm{even}& \mathrm{even}& \mathrm{odd}\\ \mathrm{odd}& \mathrm{odd}& \mathrm{even}\end{array}\right)\mathrm{for}n=1,2,\mathrm{}.$$
(3.35)
Due to this result, the amplitude $`S^{}(t)`$ has the same structure,
$`S^{}=\mathrm{T}\mathrm{exp}\left\{i{\displaystyle H^{}(t)𝑑t}\right\}=\left(\begin{array}{ccc}\mathrm{even}& \mathrm{even}& \mathrm{odd}\\ \mathrm{even}& \mathrm{even}& \mathrm{odd}\\ \mathrm{odd}& \mathrm{odd}& \mathrm{even}\end{array}\right).`$ (3.39)
This is a general feature, which holds in arbitrary matter profile.
### 3.3 Proposal of Simple Method
In the previous subsection, we have shown the general features (3.33) and (3.39) for the amplitude $`S^{}(t)`$ related to the almost vanishing parameters $`s_{13}`$ and $`\alpha `$. However, we cannot calculate $`S^{}(t)`$ by using only this features. In this subsection, we propose a generalized method for the calculation. Let us consider if there is an approximation available for both region, low and high energy. After expanding the amplitude $`S^{}`$ on the two small parameters $`\alpha `$ and $`s_{13}`$, we can arrange this as
$`S^{}`$ $`=`$ $`O(1)+O(\alpha )+O(s_{13})+O(\alpha ^2)+O(\alpha s_{13})+O(s_{13}^2)+\mathrm{}`$
$`=`$ $`\left(O(1)+O(\alpha )+O(\alpha ^2)+\mathrm{}\right)+\left(O(1)+O(s_{13})+O(s_{13}^2)+\mathrm{}\right)`$
$``$ $`O(1)+O(\alpha s_{13})+O(\alpha ^2s_{13})+O(\alpha s_{13}^2)+\mathrm{}`$ (3.41)
$`=`$ $`\underset{s_{13}0}{lim}S^{}+\underset{\alpha 0}{lim}S^{}\underset{\alpha ,s_{13}0}{lim}S^{}+O(\alpha s_{13})+O(\alpha ^2s_{13})+\mathrm{}`$ (3.42)
$`=`$ $`S^{\mathrm{}}+S^hS^d+O(\alpha s_{13})+O(\alpha ^2s_{13})+O(\alpha s_{13}^2)+\mathrm{},`$ (3.43)
where $`S^{\mathrm{}},S^h`$ and $`S^d`$ are defined by
$`S^{\mathrm{}}`$ $`=`$ $`\underset{s_{13}0}{lim}S^{}=\mathrm{T}\mathrm{exp}\left\{i{\displaystyle H^{\mathrm{}}𝑑t}\right\}`$ (3.44)
$`S^h`$ $`=`$ $`\underset{\alpha 0}{lim}S^{}=\mathrm{T}\mathrm{exp}\left\{i{\displaystyle H^h𝑑t}\right\}`$ (3.45)
$`S^d`$ $`=`$ $`\underset{\alpha ,s_{13}0}{lim}S^{}=\mathrm{diag}(\mathrm{exp}\left\{i{\displaystyle a(t)𝑑t}\right\},1,e^{i\mathrm{\Delta }_{31}L}),`$ (3.46)
respectively. $`S^{\mathrm{}}`$ ($`S^h`$) corresponds to the amplitudes, which gives the main contribution in low (high) energy. The term $`S^d`$ counts twice, because contributions to this term comes from both, low energy and high energy terms. Therefore, we subtract this contribution and approximate the amplitude as
$`S^{}S^{\mathrm{}}+S^hS^d`$ (3.47)
ignoring higher order terms. Let us discuss this approximation, which is used to derive our main result here.
In (3.3)-(3.43), the higher order terms in $`\alpha `$ and $`s_{13}`$ are included in $`S^{\mathrm{}}`$ and $`S^h`$. The reason for including the higher order terms is to take into account non-perturbative effects, which become important in the low and high energy MSW resonance region as discussed in section 2. On the other hand, we ignore those higher order terms, which are proportional to both $`\alpha `$ and $`s_{13}`$. For example, in the case of second order of the small parameters, $`\alpha `$ and $`s_{13}`$, we ignore only the mixed $`O(\alpha s_{13})`$ term among the three terms with second order $`O(\alpha ^2),O(s_{13}^2)`$ and $`O(\alpha s_{13})`$. This procedure is more appropriate than the usual perturbation, because both non-perturbative effects on a small $`\alpha `$ in the low energy region and on a small $`s_{13}`$ in the high energy region can be included in our approximation. However, as the derivation of the approximation (3.47) is not exact, we need to check this later numerically. In the previous subsection, the parity of the matrix elements related to $`s_{13}`$ has been derived. The equations (3.33), (3.39) and (3.47) lead to the magnitude of the correction for the amplitudes as
$`S_{\mu e}^{}`$ $`=`$ $`S_{\mu e}^{\mathrm{}}+O(\alpha s_{13}^2)`$ (3.48)
$`S_{\tau e}^{}`$ $`=`$ $`S_{\tau e}^h+O(\alpha s_{13})`$ (3.49)
$`S_{\tau \mu }^{}`$ $`=`$ $`O(\alpha s_{13})`$ (3.50)
$`S_{ee}^{}`$ $`=`$ $`S_{ee}^{\mathrm{}}+S_{ee}^hS_{ee}^d+O(\alpha s_{13}^2)`$ (3.51)
$`S_{\mu \mu }^{}`$ $`=`$ $`S_{\mu \mu }^{\mathrm{}}+O(\alpha s_{13}^2)`$ (3.52)
$`S_{\tau \tau }^{}`$ $`=`$ $`S_{\tau \tau }^h+O(\alpha s_{13}^2).`$ (3.53)
If we ignore the higher order terms which are proportional to both, $`\alpha `$ and $`s_{13}`$, in these equations, we obtain approximate formulas by using the two generation amplitudes. The main term for $`S_{\mu e}^{},S_{\mu \mu }^{}`$ is approximated by the low-energy amplitude as seen from (3.48) and (3.52). On the other hand, the main terms for $`S_{\tau e}^{}`$ and $`S_{\tau \tau }^{}`$ are approximated by the high-energy amplitude as derived from (3.49) and (3.53), and these features come from eq. (3.33). As seen from (3.48)-(3.53), these are expressed by only two generation amplitudes and have the advantage of simplicity. The precision of the approximation depends on the values of $`s_{13}`$ and $`\alpha `$. If the value of $`s_{13}`$ is smaller than the upper bound derived by the CHOOZ experiment, the precision of approximation becomes better. It should be mentioned that the method described in this subsection does not need the assumption of constant matter density.
Next, we show that the results using the approximate formulas (3.48)-(3.53) are in excellent agreement with the numerical calculations. We choose the Preliminary Reference Earth Model (PREM) as an earth matter density model and compare the amplitudes in all channels calculated from our approximate formulas with the numerical calculation. Here, $`\mathrm{\Delta }m_{21}^2=8.3\times 10^5\mathrm{eV}^2,\mathrm{\Delta }m_{31}^2=2.0\times 10^3\mathrm{eV}^2,\mathrm{sin}^22\theta _{12}=0.8`$ and $`\mathrm{sin}\theta _{13}=0.23`$ are chosen. Furthermore, we set the baseline length as $`L=6000\mathrm{km}`$, a length, for which the MSW effect becomes significant, and the energy region as $`1\mathrm{GeV}E20\mathrm{GeV}`$, for which the MSW resonance energy appears.
We compare our formulas with the numerical calculation in Figure 2. One can see in the following that some remarkable features occur. At first, the four amplitudes $`|S_{\mu e}^{}|,|S_{ee}^{}|,|S_{\mu \mu }^{}|`$ and $`|S_{\tau \tau }^{}|`$ coincide with the numerical calculation with a good precision. This happens, because there is no first order correction of $`s_{13}`$ from (3.48) and (3.51)-(3.53). Next, the low-energy part of $`|S_{\tau e}^{}|`$ differs from the numerical calculation only a little, which can be understood from the eq. (3.49). Furthermore, our approximation for $`|S_{\tau \mu }^{}|`$ is not at all in agreement with the numerical calculation. Although the value of this amplitude is exactly zero in our approximation as seen from (3.50), the actual magnitude of this amplitude attains $`0.02`$ in the low energy region from Figure 1. It is noted that this value is almost the same as the value expected from the order counting $`O(\alpha s_{13})0.01`$. Next, we would like to derive the approximate formulas of the oscillation probabilities from the amplitudes obtained here, however, there is a problem. As seen from eqs. (A.32)-(A.49) in Appendix A, we cannot obtain the approximate formulas for the CP dependence of the probabilities $`P(\nu _\mu \nu _\mu ),P(\nu _\mu \nu _\tau )`$ and $`P(\nu _\tau \nu _\tau )`$. The reason is that the CP dependence in these channels is directly proportional to $`S_{\mu \tau }^{}`$. However there is a method to calculate these indirectly by using the unitarity, even if we cannot directly obtain the amplitude $`S_{\mu \tau }^{}`$, as we will show in section 4.
### 3.4 Discussion
In this subsection, let us reconsider the method proposed in the previous subsection in more detail. In (3.3)-(3.43), we ignored the terms of the order $`O(\alpha s_{13})`$ for the amplitude $`S^{}`$. The reader probably wonder, why we ignore the terms of order $`O(\alpha s_{13})`$ for the amplitude $`S^{}`$, but not for other quantities, like for example $`H^{}`$ and $`P`$. Let us demonstrate the case of using the physical quantity $`Q`$. Expanding $`Q`$ on $`\alpha `$ and $`s_{13}`$, we obtain
$`Q=Q^{\mathrm{}}+Q^hQ^d+O(\alpha s_{13})+O(\alpha ^2s_{13})+\mathrm{}`$ (3.54)
by the same procedure as (3.3)-(3.43). If we neglect the higher order terms like $`O(\alpha s_{13})`$, we can approximate $`Q`$ as
$`QQ^{\mathrm{}}+Q^hQ^d.`$ (3.55)
As in the case of the approximated amplitude defined in the previous subsection, $`Q^{\mathrm{}}=Q^{\mathrm{}}(\theta _{12},\mathrm{\Delta }_{21})`$ is the main term in low-energy and $`Q^h=Q^h(\theta _{13},\mathrm{\Delta }_{31})`$ is the main term in high-energy. $`Q^d`$ is the double counting term. It is a method to be able to take into account non-perturbative effects in both of the two MSW resonance regions. In principle, this method is effective whatever we choose for the quantity $`Q`$, there is just a difference in simplicity and the magnitude of error, as discussed in the following.
We consider the Hamiltonian $`H^{}`$ as $`Q`$. Namely, $`H^{}`$ can be approximated as
$`H^{}H^{\mathrm{}}+H^hH^d,`$ (3.56)
where the double counting term is given by
$`H^d=\mathrm{diag}(a(t),0,\mathrm{\Delta }_{31}).`$ (3.57)
There is a problem, because approximation became too simple: The form of the solution for the amplitude is given by
$`S^{}\mathrm{T}\mathrm{exp}\left\{i{\displaystyle (H^{\mathrm{}}+H^hH^d)𝑑t}\right\},`$ (3.58)
and we cannot simplify this amplitude without calculation of the commutator of $`H^{\mathrm{}}`$ and $`H^h`$. Thus, the direct application of our method for the Hamiltonian needs other approximations to estimate the amplitude and this is not effective from the point of the simplicity. Especially, the amplitudes cannot be calculated within the framework of the two generation approximation although the precision of this approximation was good.
Next, let us consider the probability $`P`$ as the quantity $`Q`$. In this case, we can approximate as
$`PP^{\mathrm{}}+P^hP^d,`$ (3.59)
where $`P^{\mathrm{}}`$ and $`P^h`$ are given by
$`P^{\mathrm{}(h)}=\left|\mathrm{T}\mathrm{exp}\left\{i{\displaystyle H^{\mathrm{}(h)}𝑑t}\right\}\right|^2,`$ (3.60)
and $`P^d`$ is the identity matrix. As an example, we consider $`P(\nu _e\nu _\mu )`$. The CP phase $`\delta `$ dependence is given by
$`P(\nu _e\nu _\mu )=A_{e\mu }\mathrm{cos}\delta +B_{e\mu }\mathrm{sin}\delta +C_{e\mu },`$ (3.61)
where the coefficients $`A_{e\mu }`$ and $`B_{e\mu }`$ determine the magnitude of the CP violation. On the other hand, the CP violation becomes zero in the limit, $`\alpha 0`$ or $`s_{13}0`$, as seen from
$`A_{e\mu }=O(\alpha s_{13}),B_{e\mu }=O(\alpha s_{13}).`$ (3.62)
Namely, we obtain
$`A_{e\mu }^{\mathrm{}}=A_{e\mu }^h=A_{e\mu }^d=0,B_{e\mu }^{\mathrm{}}=B_{e\mu }^h=B_{e\mu }^d=0`$ (3.63)
and therefore we cannot calculate quantities like the CP violation, because it is the effects of three generations in this approximation. This result holds for all channels.
To summarize this subsection, if we choose the Hamiltonian $`H^{}`$ as $`Q`$, the precision of approximation is good, but the calculation is not so simple compared to the exact calculation. If we choose the probability $`P`$ as $`Q`$, we cannot calculate three generation effects like CP violation. On the other hand, if we choose the amplitude $`S^{}`$ as $`Q`$, we can calculate the three generation effects like CP violation within the framework of two generation approximation.
## 4 Approximate Formulas for Oscillation Probabilities
In this section, we calculate the CP dependent terms from $`\nu _\mu `$ to $`\nu _\mu `$ and so on, not determined by the method in the previous section, by using the unitarity. After that, we derive the approximate formulas of the oscillation probabilities $`P(\nu _\alpha \nu _\beta )`$ in arbitrary matter profile without using $`S_{\mu \tau }^{}`$ directly. Namely, we derive the approximate formulas in all channels by our new method.
### 4.1 Unitarity Relation
We cannot calculate the amplitude $`S_{\mu \tau }^{}`$ in the method introduced in the previous section. The reason is that the amplitude $`S_{\mu \tau }^{}`$ is a very small quantity, which has an order of $`O(\alpha s_{13})`$. As seen from (A.32)-(A.49) in Appendix A, it seems that the approximate formulas, including CP violation, of three channels, $`P(\nu _\mu \nu _\mu ),P(\nu _\tau \nu _\tau )`$ and $`P(\nu _\mu \nu _\tau )`$ cannot be derived without directly calculating the amplitude $`S_{\mu \tau }^{}`$. However in this subsection we show, that we can derive these probabilities without directly calculating this amplitude, if we assume the two natural conditions,
$`s_{23}`$ $``$ $`c_{23},`$ (4.1)
$`S_{\alpha \beta }^{}`$ $``$ $`S_{\beta \alpha }^{}.`$ (4.2)
The first condition is supported by the best fit value of atmospheric neutrino experiments and the second condition holds in one-dimensional models of the earth matter density like PREM or ak-135f. Accordingly, the error due to the difference between these conditions and the real situations is considered to be relatively small. We perform the analysis under these two conditions in the following.
At first, we obtain
$`B_{\mu \mu }=2\mathrm{I}\mathrm{m}[(S_{\mu \mu }^{}c_{23}^2+S_{\tau \tau }^{}s_{23}^2)^{}(S_{\tau \mu }^{}S_{\mu \tau }^{})]c_{23}s_{23}=0`$ (4.3)
from (A.34) and (4.2) in the case of the symmetric matter density. In the same way, we obtain
$`B_{\tau \tau }=2\mathrm{I}\mathrm{m}[(S_{\mu \mu }^{}s_{23}^2+S_{\tau \tau }^{}c_{23}^2)^{}(S_{\tau \mu }^{}S_{\mu \tau }^{})]c_{23}s_{23}=0`$ (4.4)
from (A.40) and (4.2). Furthermore, in the case of the symmetric matter density and the maximal 2-3 mixing angle 45, we also obtain
$`A_{\mu \tau }=2\mathrm{R}\mathrm{e}[(S_{\mu \mu }^{}S_{\tau \tau }^{})^{}(S_{\tau \mu }^{}c_{23}^2S_{\mu \tau }^{}s_{23}^2)]c_{23}s_{23}=0`$ (4.5)
from (A.45) and (4.2). Let us here consider now, how the oscillation probabilities are derived, which are related to the amplitude $`S_{\mu \tau }^{}`$ but have not been determined in the previous section,. At first, in the probability,
$`P(\nu _\mu \nu _\mu )=A_{\mu \mu }\mathrm{cos}\delta +B_{\mu \mu }\mathrm{sin}\delta +C_{\mu \mu }+D_{\mu \mu }\mathrm{cos}2\delta +E_{\mu \mu }\mathrm{sin}2\delta ,`$ (4.6)
the coefficient proportional to $`\mathrm{cos}\delta `$ can be calculated as
$`A_{\mu \mu }=A_{\mu e}A_{\mu \tau }A_{e\mu }2\mathrm{R}\mathrm{e}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]s_{23}c_{23}`$ (4.7)
from (4.5) and the unitarity relation. Next, let us turn to the probability $`P(\nu _\tau \nu _\tau )`$. In the probability,
$`P(\nu _\tau \nu _\tau )=A_{\tau \tau }\mathrm{cos}\delta +B_{\tau \tau }\mathrm{sin}\delta +C_{\tau \tau }+D_{\tau \tau }\mathrm{cos}2\delta +E_{\tau \tau }\mathrm{sin}2\delta ,`$ (4.8)
the coefficient of $`\mathrm{cos}\delta `$ can be calculated as
$`A_{\tau \tau }=A_{\tau e}A_{\tau \mu }A_{e\tau }2\mathrm{R}\mathrm{e}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]s_{23}c_{23}`$ (4.9)
from (4.5) and the unitarity relation.
Finally, let us calculate the probability $`P(\nu _\mu \nu _\tau )`$. In the probability,
$`P(\nu _\mu \nu _\tau )=A_{\mu \tau }\mathrm{cos}\delta +B_{\mu \tau }\mathrm{sin}\delta +C_{\mu \tau }+D_{\mu \tau }\mathrm{cos}2\delta +E_{\mu \tau }\mathrm{sin}2\delta ,`$ (4.10)
the coefficient of $`\mathrm{sin}\delta `$ becomes
$`B_{\mu \tau }=B_{\mu e}B_{\mu \mu }B_{e\mu }2\mathrm{I}\mathrm{m}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]s_{23}c_{23}`$ (4.11)
from (4.3) and the unitarity relation. We can derive the probability up to the second order of two small parameters by using the unitarity relation although we cannot directly calculate $`S_{\mu \tau }^{}`$ in the previous method. In addition, the coefficients of $`\mathrm{sin}2\delta `$ and $`\mathrm{cos}2\delta `$, $`D`$ and $`E`$, have an order of
$`D=O(\alpha ^2s_{13}^2),E=O(\alpha ^2s_{13}^2)`$ (4.12)
in these three channels as derived from (A.36), (A.37), (A.42), (A.43), (A.48) and (A.49) and are expected to be small. Actually, these coefficients have the second order of $`S_{\mu \tau }^{}`$, and the values are about $`(0.02)^20.0004`$ from Figure 1 in the high energy region related with long baseline experiments. So we ignore them in the following section.
### 4.2 Approximate Formulas in All Channels
In this subsection, we present the approximate formulas which are useful in arbitrary matter density profile. Ignoring the higher order terms of $`\alpha `$ and $`s_{13}`$ than the second order, we can present the oscillation probabilities for all channels with the amplitudes calculated in two generations.
At first, let us derive the approximate formulas for $`P(\nu _e\nu _\mu )`$ and $`P(\nu _e\nu _\tau )`$. The approximate formula for $`P(\nu _e\nu _\mu )`$ has already been derived in our previous paper . We only have to replace the amplitudes $`S_{\mu e}^{}`$ and $`S_{\tau e}^{}`$ in three generations into $`S_{\mu e}^{\mathrm{}}`$ and $`S_{\tau e}^h`$ in two generations. From (A.24)-(A.31) and (3.48)-(3.49), we obtain
$`P(\nu _e\nu _\mu )`$ $`=`$ $`A_{e\mu }\mathrm{cos}\delta +B_{e\mu }\mathrm{sin}\delta +C_{e\mu }`$ (4.13)
$`A_{e\mu }`$ $``$ $`2\mathrm{R}\mathrm{e}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.14)
$`B_{e\mu }`$ $``$ $`2\mathrm{I}\mathrm{m}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.15)
$`C_{e\mu }`$ $``$ $`|S_{\mu e}^{\mathrm{}}|^2c_{23}^2+|S_{\tau e}^h|^2s_{23}^2,`$ (4.16)
$`P(\nu _e\nu _\tau )`$ $`=`$ $`A_{e\tau }\mathrm{cos}\delta +B_{e\tau }\mathrm{sin}\delta +C_{e\tau }`$ (4.17)
$`A_{e\tau }`$ $``$ $`2\mathrm{R}\mathrm{e}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.18)
$`B_{e\tau }`$ $``$ $`2\mathrm{I}\mathrm{m}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.19)
$`C_{e\tau }`$ $``$ $`|S_{\mu e}^{\mathrm{}}|^2s_{23}^2+|S_{\tau e}^h|^2c_{23}^2,.`$ (4.20)
Eqs. (4.13)-(4.16) are the same as those derived in our previous paper . Next, let us derive the approximate formulas for $`P(\nu _e\nu _e)`$. Using (3.51) directly, we obtain
$`P(\nu _e\nu _e)`$ $`=`$ $`C_{ee}=|S_{ee}^{}|^2`$ (4.21)
$``$ $`|S_{ee}^{\mathrm{}}+S_{ee}^hS_{ee}^d|^2.`$ (4.22)
On the other hand, we obtain
$`P(\nu _e\nu _e)`$ $`=`$ $`C_{ee}=1C_{e\mu }C_{e\tau }`$ (4.23)
$``$ $`1|S_{\mu e}^{\mathrm{}}|^2|S_{\tau e}^h|^2,`$ (4.24)
by using the unitarity relation. This is a different approximate formula than (4.22). Thus, there are two kinds of expressions (4.22) and (4.24) for $`P(\nu _e\nu _e)`$. We checked numerically that the expression (4.24) has a better precision than the expression (4.22). Furthermore, the expression (4.24) easily reproduces the approximate formula derived with double expansion up to the second order of two small parameters in ref. (second order formula). So we use the expression (4.24) in the following.
Next, let us derive the approximate formula for $`P(\nu _\mu \nu _\tau )`$. At first we calculate the terms independent of the CP phase in this calculation. We can approximate
$`C_{\mu \tau }=|S_{\mu \tau }^{}|^2s_{23}^4+|S_{\tau \mu }^{}|^2c_{23}^4+|S_{\mu \mu }^{}S_{\tau \tau }^{}|^2c_{23}^2s_{23}^2|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2c_{23}^2s_{23}^2`$ (4.25)
from (A.47) and (3.52)-(3.53), where we ignore the terms proportional to $`|S_{\mu \tau }^{}|^2=O(\alpha ^2s_{13}^2)`$. This leads to the approximated probability as
$`P(\nu _\mu \nu _\tau )`$ $`=`$ $`B_{\mu \tau }\mathrm{sin}\delta +C_{\mu \tau }`$ (4.26)
$`B_{\mu \tau }`$ $``$ $`2\mathrm{I}\mathrm{m}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.27)
$`C_{\mu \tau }`$ $``$ $`|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2c_{23}^2s_{23}^2.`$ (4.28)
Next, let us derive the approximate formulas for $`P(\nu _\mu \nu _\mu )`$ and $`P(\nu _\tau \nu _\tau )`$. From (A.35) and (3.52)-(3.53), we obtain
$`C_{\mu \mu }`$ $`=`$ $`|S_{\mu \mu }^{}c_{23}^2+S_{\tau \tau }^{}s_{23}^2|^2+(|S_{\mu \tau }^{}|^2+|S_{\tau \mu }^{}|^2)c_{23}^2s_{23}^2`$ (4.29)
$``$ $`|S_{\mu \mu }^{\mathrm{}}c_{23}^2+S_{\tau \tau }^hs_{23}^2|^2,`$ (4.30)
where we neglect the terms proportional to $`|S_{\mu \tau }^{}|^2=O(\alpha ^2s_{13}^2)`$. On the other hand, we obtain another expression by using the unitarity relation as
$`C_{\mu \mu }`$ $`=`$ $`1C_{\mu e}C_{\mu \tau }`$ (4.31)
$``$ $`1|S_{\mu e}^{\mathrm{}}|^2c_{23}^2|S_{\tau e}^h|^2s_{23}^2|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2c_{23}^2s_{23}^2`$ (4.32)
This seems to be different from (4.30) at a glance, but we confirmed that (4.30) and (4.32) are the same expression by using the unitarity relation. In the following, we use the expression (4.32) for the reason that this easily reproduces the second order formula and we can check the unitarity. In the same way, $`C_{\tau \tau }`$ is given by
$`C_{\tau \tau }=1C_{e\tau }C_{\mu \tau }1|S_{\mu e}^{\mathrm{}}|^2s_{23}^2|S_{\tau e}^h|^2c_{23}^2|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2c_{23}^2s_{23}^2`$ (4.33)
from the unitarity relation. From the result obtained in subsection 4.1, the approximate formulas for $`P(\nu _\mu \nu _\mu )`$ and $`P(\nu _\tau \nu _\tau )`$ are given by
$`P(\nu _\mu \nu _\mu )`$ $`=`$ $`A_{\mu \mu }\mathrm{cos}\delta +C_{\mu \mu }`$ (4.34)
$`A_{\mu \mu }`$ $``$ $`2\mathrm{R}\mathrm{e}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.35)
$`C_{\mu \mu }`$ $``$ $`1|S_{\mu e}^{\mathrm{}}|^2c_{23}^2|S_{\tau e}^h|^2s_{23}^2|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2c_{23}^2s_{23}^2`$ (4.36)
$`P(\nu _\tau \nu _\tau )`$ $`=`$ $`A_{\tau \tau }\mathrm{cos}\delta +C_{\tau \tau }`$ (4.37)
$`A_{\tau \tau }`$ $``$ $`2\mathrm{R}\mathrm{e}[S_{\mu e}^{\mathrm{}}S_{\tau e}^h]c_{23}s_{23},`$ (4.38)
$`C_{\tau \tau }`$ $``$ $`1|S_{\mu e}^{\mathrm{}}|^2s_{23}^2|S_{\tau e}^h|^2c_{23}^2|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2c_{23}^2s_{23}^2.`$ (4.39)
These results are one of the main results of this paper. In all channels, we can present the probabilities including the CP violation by using the amplitudes calculated in two generations. It is noted that the CP violating terms due to the existence of three generations can be calculated from the two generation amplitudes.
Next, let us compare the approximate formulas (4.13)-(4.20), (4.24), (4.26)-(4.28) and (4.34)-(4.39) with the numerical calculations. We take the PREM (Preliminary Reference Earth Model) as the earth matter density profile and compare the approximated values of all probabilities with those calculated numerically. We use the same parameters as those used in fig. 1 and $`\mathrm{sin}2\theta _{23}=1`$, $`\delta =90^{}`$. We set the baseline length, $`L=6000\mathrm{km}`$ and the energy region, $`1\mathrm{GeV}E20\mathrm{GeV}`$, to include the high energy MSW resonance.
We compare our approximate formulas with the numerical calculation in figure 3. One can see some features from this figure. The approximated value of probabilities for $`P(\nu _e\nu _\mu ),P(\nu _e\nu _\tau )`$ and $`P(\nu _e\nu _e)`$ coincide to the numerical values very well, on the other hand, the remaining three channels of probabilities $`P(\nu _\mu \nu _\tau ),P(\nu _\mu \nu _\mu )`$ and $`P(\nu _\tau \nu _\tau )`$ show a small difference between the approximate and the numerical value. However, the difference is not so large as in figure 2 and as a first step the result is sufficiently accurate.
## 5 Comparison of Our Results with Second Order Formulas
In this section, we concretely calculate the amplitudes by using the approximate formulas derived in the previous section for the case of constant matter and show that simple approximate formulas can be obtained. Finally, we demonstrate that the approximate formula derived with double expansion up to the second order of the two small parameters (second order formulas) are largely different from the exact values in the MSW resonance region under the condition that the baseline length is longer than the oscillation length.
### 5.1 Approximate Formulas for Amplitudes
In the previous section, we have given a method for approximation of the probabilities in three generations by amplitudes in two generations. In this subsection, we use the constant matter density profile in order to compare our method with other method and investigate the non-perturbative effects. As seen from (4.13)-(4.20), (4.24), (4.26)-(4.28) and (4.34)-(4.39), we only have to calculate four kinds of amplitudes, namely $`S_{\mu e}^{\mathrm{}},S_{\mu \mu }^{\mathrm{}},S_{\tau e}^h`$ and $`S_{\tau \tau }^h`$.
The low-energy approximate formulas are obtained by taking the limit $`s_{13}=0`$ and from
$`H^{\mathrm{}}`$ $`=`$ $`O_{12}\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})O_{12}^T+\mathrm{diag}(a,0,0)`$ (5.1)
$`=`$ $`O_{12}^{\mathrm{}}\mathrm{diag}(\lambda _1^{\mathrm{}},\lambda _2^{\mathrm{}},\mathrm{\Delta }_{31})(O_{12}^{\mathrm{}})^T.`$ (5.2)
The effective masses $`\lambda _i^{\mathrm{}}(i=1,2)`$ and the effective mixing angle $`\theta _{12}^{\mathrm{}}`$ are determined by the diagonalization of (5.1) to (5.2). If we define the mass squared difference in matter as $`\mathrm{\Delta }_{21}^{\mathrm{}}=\lambda _2^{\mathrm{}}\lambda _1^{\mathrm{}}`$, we obtain the relation
$`{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}}{\mathrm{\Delta }_{21}}}={\displaystyle \frac{\mathrm{sin}2\theta _{12}}{\mathrm{sin}2\theta _{12}^{\mathrm{}}}}=\sqrt{\left(\mathrm{cos}2\theta _{12}{\displaystyle \frac{a}{\mathrm{\Delta }_{21}}}\right)^2+\mathrm{sin}^22\theta _{12}}.`$ (5.3)
Therefore, the amplitude is calculated as
$`S_{\mu e}^{\mathrm{}}`$ $`=`$ $`i\mathrm{sin}2\theta _{12}^{\mathrm{}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{21}+a}{2}}L\right)`$ (5.4)
$`S_{\mu \mu }^{\mathrm{}}`$ $`=`$ $`\left(\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}i\mathrm{cos}2\theta _{12}^{\mathrm{}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\right)\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{21}+a}{2}}L\right)`$ (5.5)
by substituting (5.2) into (3.44). On the other hand, the approximate formulas in high energy are obtained by taking the limit $`\alpha =0`$ and we get
$`H^h`$ $`=`$ $`O_{13}\mathrm{diag}(0,0,\mathrm{\Delta }_{31})O_{13}^T+\mathrm{diag}(a,0,0)`$ (5.6)
$`=`$ $`O_{13}^h\mathrm{diag}(\lambda _1^h,0,\lambda _3^h)(O_{13}^h)^T.`$ (5.7)
The effective masses $`\lambda _i^h(i=1,3)`$ and the effective mixing angle $`\theta _{13}^h`$ are determined by the diagonalization of (5.6) to (5.7). If we define the mass squared difference in matter as $`\mathrm{\Delta }_{31}^h=\lambda _3^h\lambda _1^h`$, we obtain the relation
$`{\displaystyle \frac{\mathrm{\Delta }_{31}^h}{\mathrm{\Delta }_{31}}}={\displaystyle \frac{\mathrm{sin}2\theta _{13}}{\mathrm{sin}2\theta _{13}^h}}=\sqrt{\left(\mathrm{cos}2\theta _{13}{\displaystyle \frac{a}{\mathrm{\Delta }_{31}}}\right)^2+\mathrm{sin}^22\theta _{13}}.`$ (5.8)
Accordingly, the amplitude can be calculated by substituting (5.7) into (3.45) as
$`S_{\tau e}^h`$ $`=`$ $`i\mathrm{sin}2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{31}+a}{2}}L\right)`$ (5.9)
$`S_{\tau \tau }^h`$ $`=`$ $`\left(\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}i\mathrm{cos}2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\right)\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{31}+a}{2}}L\right).`$ (5.10)
As seen from (5.4) and (5.9), $`S_{\mu e}^{\mathrm{}}`$ and $`S_{\tau e}^h`$ have simple forms, but the expressions of $`S_{\mu \mu }^{\mathrm{}}`$ and $`S_{\tau \tau }^h`$ are more complex than (5.5) and (5.10).
### 5.2 Approximate Formulas for Probabilities
In this subsection, we derive the approximate formulas of the oscillation probabilities in constant matter by using the result of the previous section.
At first, let us consider the case of including electron neutrino in the initial or final state. In this case, the probability for any channel can be calculated almost in the same way. The probability $`P(\nu _e\nu _e)`$ is obtained by substituting (5.4) and (5.9) into (4.24) as
$`P(\nu _e\nu _e)=1\mathrm{sin}^22\theta _{12}^{\mathrm{}}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\mathrm{sin}^22\theta _{13}^h\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}.`$ (5.11)
The probability $`P(\nu _e\nu _\mu )`$ is obtained by substituting (5.4) and (5.9) into (4.14)-(4.16) as
$`P(\nu _e\nu _\mu )`$ $`=`$ $`A_{e\mu }\mathrm{cos}\delta +B_{e\mu }\mathrm{sin}\delta +C_{e\mu }`$ (5.12)
$`A_{e\mu }`$ $``$ $`\mathrm{sin}2\theta _{12}^{\mathrm{}}\mathrm{sin}2\theta _{23}\mathrm{sin}2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}}`$ (5.13)
$`B_{e\mu }`$ $``$ $`\mathrm{sin}2\theta _{12}^{\mathrm{}}\mathrm{sin}2\theta _{23}\mathrm{sin}2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}}`$ (5.14)
$`C_{e\mu }`$ $``$ $`c_{23}^2\mathrm{sin}^22\theta _{12}^{\mathrm{}}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}+s_{23}^2\mathrm{sin}^22\theta _{13}^h\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}.`$ (5.15)
The remaining probability $`P(\nu _e\nu _\tau )`$ can be calculated from the unitarity relation. Next, let us calculate the probabilities for the case, that not all electron neutrinos in the initial and final state are included. Also in this case, the probability for any channel can be calculated almost in the same way. Accordingly, as an example, we calculate the probability for muon neutrino to tau neutrino,
$`P(\nu _\mu \nu _\tau )`$ $`=`$ $`B_{\mu \tau }\mathrm{sin}\delta +C_{\mu \tau }`$ (5.16)
$`B_{\mu \tau }`$ $``$ $`\mathrm{sin}2\theta _{12}^{\mathrm{}}\mathrm{sin}2\theta _{23}\mathrm{sin}2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}}`$ (5.17)
$`C_{\mu \tau }`$ $``$ $`|S_{\mu \mu }^{\mathrm{}}S_{\tau \tau }^h|^2s_{23}^2c_{23}^2.`$ (5.18)
At first, we use the relations, $`\mathrm{cos}2\theta _{12}^{\mathrm{}}=2\mathrm{cos}^2\theta _{12}^{\mathrm{}}1`$ and $`\mathrm{cos}2\theta _{13}^h=12\mathrm{sin}^2\theta _{13}^h`$, and we rewrite $`S_{\mu \mu }^{\mathrm{}}`$ and $`S_{\tau \tau }^h`$ as
$`S_{\mu \mu }^{\mathrm{}}`$ $``$ $`\left[\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}}{2}}L\right)2i\mathrm{cos}^2\theta _{12}^{\mathrm{}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\right]\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{21}+a}{2}}L\right)`$ (5.19)
$`S_{\tau \tau }^h`$ $``$ $`\left[\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\right)+2i\mathrm{sin}^2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\right]\mathrm{exp}\left(i{\displaystyle \frac{\mathrm{\Delta }_{31}+a}{2}}L\right).`$ (5.20)
Then, arranging $`C_{\mu \tau }`$ in the order of the effective mixing angles $`\mathrm{cos}\theta _{12}^{\mathrm{}}`$ and $`\mathrm{sin}\theta _{13}^h`$, we obtain
$`C_{\mu \tau }^1`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{(\mathrm{\Delta }_{21}^{\mathrm{}}+\mathrm{\Delta }_{31}^h+\mathrm{\Delta }_{32})L}{4}}`$ (5.21)
$`C_{\mu \tau }^{2a}`$ $`=`$ $`2\mathrm{sin}^22\theta _{23}\mathrm{cos}^2\theta _{12}^{\mathrm{}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{21}^{\mathrm{}}+\mathrm{\Delta }_{31}^h+\mathrm{\Delta }_{32})L}{4}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{21}^{\mathrm{}}\mathrm{\Delta }_{31}^h\mathrm{\Delta }_{32})L}{4}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}`$ (5.22)
$`C_{\mu \tau }^{2b}`$ $`=`$ $`2\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{21}^{\mathrm{}}+\mathrm{\Delta }_{31}^h+\mathrm{\Delta }_{32})L}{4}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{21}^{\mathrm{}}\mathrm{\Delta }_{31}^h+\mathrm{\Delta }_{32})L}{4}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}`$ (5.23)
$`C_{\mu \tau }^3`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{cos}^4\theta _{12}^{\mathrm{}}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}+\mathrm{sin}^22\theta _{23}\mathrm{sin}^4\theta _{13}^h\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}`$ (5.24)
$`+`$ $`2\mathrm{sin}^22\theta _{23}\mathrm{cos}^2\theta _{12}^{\mathrm{}}\mathrm{sin}^2\theta _{13}^h\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}^{\mathrm{}}L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}^hL}{2}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}}.`$
As we show in the following section, the reason of arranging the terms like this is, because the second order formulas can be easily derived. In order to derive the second order formulas, it is sufficient to use $`C_{\mu \tau }^1`$, $`C_{\mu \tau }^{2a}`$ and $`C_{\mu \tau }^{2b}`$. We can also calculate the other channels $`P(\nu _\mu \nu _\mu )`$ and $`P(\nu _\tau \nu _\tau )`$ in the same way. In a recent study, it was found that the channels $`P(\nu _\mu \nu _\mu )`$ and $`P(\nu _\tau \nu _\tau )`$ are largely affected by the earth matter in the long baseline .
We can see from these expressions that the approximate formulas are rather complex for the case not including electron neutrino in the initial or final state. We also understand from these formulas how matter affects the probabilities. Thus, the formulas are expected to be useful for studying matter effects.
### 5.3 Large Non-perturbative Effects of small $`\alpha `$ and $`s_{13}`$
In this subsection, we compare the approximate formulas obtained in the previous subsection with the second order formulas numerically and it is shown that the latter have a large difference from the numerical value in the MSW resonance region.
The second order formulas are approximated by the main terms of the expansion and are widely used by many authors for their simplicity. In refs. , the formula for $`P(\nu _e\nu _\mu )`$ has been derived and later on all probabilities were presented in ref. . As examples, the probabilities, $`P(\nu _e\nu _\mu )`$ and $`P(\nu _\mu \nu _\tau )`$, are taken. For the other channels of probabilities, similar expressions have been obtained. In all channels similar results were obtained from comparison with numerical calculations. The second order formula for $`P(\nu _e\nu _\mu )`$ is given by
$`P(\nu _e\nu _\mu )`$ $`=`$ $`A_{e\mu }\mathrm{cos}\delta +B_{e\mu }\mathrm{sin}\delta +C_{e\mu },`$ (5.25)
$`A_{e\mu }`$ $``$ $`\alpha s_{13}\mathrm{sin}2\theta _{12}\mathrm{sin}2\theta _{23}{\displaystyle \frac{2\mathrm{\Delta }_{31}^2}{a(\mathrm{\Delta }_{31}a)}}\mathrm{sin}{\displaystyle \frac{aL}{2}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}}`$ (5.26)
$`B_{e\mu }`$ $``$ $`\alpha s_{13}\mathrm{sin}2\theta _{12}\mathrm{sin}2\theta _{23}{\displaystyle \frac{2\mathrm{\Delta }_{31}^2}{a(\mathrm{\Delta }_{31}a)}}\mathrm{sin}{\displaystyle \frac{aL}{2}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}}`$ (5.27)
$`C_{e\mu }`$ $``$ $`\alpha ^2c_{23}^2\mathrm{sin}^22\theta _{12}{\displaystyle \frac{\mathrm{\Delta }_{31}^2}{a^2}}\mathrm{sin}^2{\displaystyle \frac{aL}{2}}+s_{13}^2s_{23}^2{\displaystyle \frac{4\mathrm{\Delta }_{31}^2}{(\mathrm{\Delta }_{31}a)^2}}\mathrm{sin}^2{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}.`$ (5.28)
Comparing our approximate formulas (5.13)-(5.15) with the second order formulas (5.26)-(5.28), each term corresponds one by one. Actually, the second order formulas (5.26)-(5.28) are derived by expanding our approximate formulas (5.13)-(5.15) up to the second order in $`\alpha `$ and $`s_{13}`$ . Next, the second order formula for $`P(\nu _\mu \nu _\tau )`$ which has been already derived in ref. is
$`P(\nu _\mu \nu _\tau )`$ $`=`$ $`A_{\mu \tau }\mathrm{cos}\delta +B_{\mu \tau }\mathrm{sin}\delta +C_{\mu \tau }`$ (5.29)
$`A_{\mu \tau }`$ $``$ $`\alpha s_{13}\mathrm{sin}^22\theta _{23}\mathrm{sin}2\theta _{12}\mathrm{cos}2\theta _{23}{\displaystyle \frac{2\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}`$ (5.30)
$`\times `$ $`\left[{\displaystyle \frac{a}{\mathrm{\Delta }_{31}}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}{\displaystyle \frac{\mathrm{\Delta }_{31}}{a}}\mathrm{sin}{\displaystyle \frac{aL}{2}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\right]`$
$`B_{\mu \tau }`$ $``$ $`\alpha s_{13}\mathrm{sin}2\theta _{12}\mathrm{sin}2\theta _{23}{\displaystyle \frac{2\mathrm{\Delta }_{31}^2}{a(\mathrm{\Delta }_{31}a)}}\mathrm{sin}{\displaystyle \frac{aL}{2}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{32}L}{2}},`$ (5.31)
and $`C_{\mu \tau }`$ is given by
$`C_{\mu \tau }`$ $``$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}`$ (5.32)
$``$ $`\alpha \mathrm{sin}^22\theta _{23}\mathrm{cos}^2\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)\mathrm{sin}\mathrm{\Delta }_{31}L+\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{cos}^4\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)^2\mathrm{cos}\mathrm{\Delta }_{31}L`$
$``$ $`\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{2a}}\right)\left[\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\mathrm{sin}{\displaystyle \frac{aL}{2}}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{a}}\right){\displaystyle \frac{\mathrm{\Delta }_{31}L}{4}}\mathrm{sin}(\mathrm{\Delta }_{31}L)\right]`$
$``$ $`s_{13}^2\mathrm{sin}^22\theta _{23}{\displaystyle \frac{2\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}\left[\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{aL}{2}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}\right){\displaystyle \frac{aL}{4}}\mathrm{sin}(\mathrm{\Delta }_{31}L)\right].`$
In the next section we show that this formula (5.32) can be also derived from our formulas (5.21)-(5.24). It is noted that the second order formula (5.32) for $`C_{\mu \tau }`$ is rather complex. Furthermore comparing our approximate formula (5.17)-(5.24) with the second order formula (5.30)-(5.32), we see that $`A_{\mu \tau }`$ is not included in our formula. The reason is, that $`A_{\mu \tau }=0`$ in the case of maximal mixing angle $`\mathrm{sin}2\theta _{23}=1`$ and there is no way of calculating this by the method described in this paper. If we consider $`\mathrm{cos}2\theta _{23}`$ as a small parameter like $`\alpha `$ and $`s_{13}`$, this $`A_{\mu \tau }`$ has the magnitude of $`O(\alpha s_{13}\mathrm{cos}2\theta _{23})`$. Therefore, $`A_{\mu \tau }`$ is proportional to the third order of small parameters and is expected to be neglectable. This means that our formula is not largely affected by the error due to this term, which cannot be derived from our method. However, as this error affects the precision measurement of $`\mathrm{sin}\theta _{23}`$ by the atmospheric neutrino experiments in future, the improvement of this point is important future work. The formulas for the other channels are given in ref. . The second order formulas are effective under the following two conditions.
The first one is for the neutrino energy and is given by
$`E0.45\mathrm{GeV}\left({\displaystyle \frac{\mathrm{\Delta }m_{21}^2}{10^4\mathrm{eV}^2}}\right)\left({\displaystyle \frac{3\mathrm{g}/\mathrm{cm}^3}{\rho }}\right).`$ (5.33)
The second one is for the baseline length and is given by
$`L8000\mathrm{km}\left({\displaystyle \frac{E}{\mathrm{GeV}}}\right)\left({\displaystyle \frac{10^4\mathrm{eV}^2}{\mathrm{\Delta }m_{21}^2}}\right).`$ (5.34)
These conditions come from the utilization of perturbative expansion on the two small parameters. The detailed discussion are given in . These approximate formulas are used for the purpose of understanding of the results obtained by numerical calculations . However, as shown in the next figure, these formulas have large difference from the true value in the MSW resonance region, which is considered to be the most important region.
Next, let us compare our formulas (5.12)-(5.23) with the second order formulas (5.25)-(5.32) in all channels by numerical calculation. In order to see the magnitude of the error, we also compare two kinds of approximate formulas with the exact values. We set the baseline length as $`L=6000\mathrm{km}`$, where the MSW effect becomes large, and the energy region as $`1\mathrm{GeV}E20\mathrm{GeV}`$, where the MSW resonance energy is included. Furthermore, the second order formulas are derived only in the case of constant matter, so we choose the average density $`\rho =3.91\mathrm{g}/\mathrm{cm}^3`$ of the earth calculated by the PREM. Note that two conditions (5.33) and (5.34) are satisfied in these region.
We compare the probabilities calculated by our approximate formulas in all channels with those by the second order formulas in addition to numerical calculation in figure 4. One can see the following points from this figure. The second order formulas show large differences from the numerical values around 5 GeV, where the high energy MSW resonance occurs. In other energy regions they are in good coincidence. The value of $`P(\nu _e\nu _e)`$ has the largest difference, the probabilities $`P(\nu _e\nu _\mu )`$ and $`P(\nu _e\nu _\tau )`$ have the next largest difference, the values of $`P(\nu _\mu \nu _\mu )`$ and $`P(\nu _\tau \nu _\tau )`$ have also significantly large difference, but only the probability $`P(\nu _\mu \nu _\tau )`$ has a small difference. In addition, these figures show that the difference between the second order formulas and the numerical calculation exists even in the two applicable regions (5.33) and (5.34). Although we do not show a figure, the difference between our approximate formulas (5.12)-(5.23) and the second order formulas (5.25)-(5.32) become more clear out of the two applicable regions (5.33) and (5.34). The reason is that our approximate formulas (5.12)-(5.23) are applicable even for the case that the condition (5.33) or (5.34) for energy and baseline length does not hold, as confirmed from the comparison with the exact numerical calculation. However, the second order formulas are good approximations, when the neutrino energy is not near the resonance energy, even if the baseline length is long.
## 6 Non-perturbative Effects of Small Parameters $`\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{sin}\theta _{13}`$
In this section, we investigate the reason for the difference between the second order formula,which contains the approximation with double expansion up to the second order of two small parameters, and the numerical calculation around the MSW resonance region as explained in the previous subsection.
We discuss the non-perturbative effects of small mixing angle more detailed than in section 2.
### 6.1 Derivation of the Second Order Formulas
In this subsection, we investigate how the second order formulas are approximated expanding on $`\alpha =\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{31}^2`$ In the previous paper, we have discussed the probability $`P(\nu _e\nu _\mu )`$, so we calculate the second order formula for $`P(\nu _\mu \nu _\tau )`$ here. The method of calculation is basically the same but the calculation itself becomes slightly complex, because we need to calculate the effective mass and the effective mixing angle up to the second order of $`\alpha `$ and $`s_{13}`$ in the case of $`P(\nu _\mu \nu _\tau )`$. In this point, the calculation is not straightforward compared with that of $`P(\nu _e\nu _\mu )`$ but the method of approximation is the same. Note that $`C_{\mu \tau }^1`$ in (5.21) does not include the effective mixing angle. For this reason, it is sufficient to expand the effective mixing angle up to the zeroth order, but we need to expand the effective mass up to the second order to calculate the probability up to the second order in $`\alpha `$ and $`s_{13}`$, The effective mixing angle is expanded up to the zeroth order as
$`\mathrm{cos}\theta _{12}^{\mathrm{}}`$ $``$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}2\theta _{12}^{\mathrm{}}\alpha \mathrm{sin}2\theta _{12}{\displaystyle \frac{\mathrm{\Delta }_{31}}{2a}}`$ (6.1)
$`\mathrm{sin}\theta _{13}^h`$ $``$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}2\theta _{13}^hs_{13}{\displaystyle \frac{\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}`$ (6.2)
and the effective mass is expanded up to the second order as
$`\mathrm{\Delta }_{21}^{\mathrm{}}`$ $``$ $`a\alpha \mathrm{cos}2\theta _{12}\mathrm{\Delta }_{31}+\alpha ^2\mathrm{sin}^22\theta _{12}{\displaystyle \frac{\mathrm{\Delta }_{31}^2}{2a}}`$ (6.3)
$`\mathrm{\Delta }_{31}^h`$ $``$ $`\mathrm{\Delta }_{31}a+s_{13}^2{\displaystyle \frac{2\mathrm{\Delta }_{31}a}{\mathrm{\Delta }_{31}a}}.`$ (6.4)
Here, we should emphasize the following points. Eqs. (6.1) and (6.3) obtained by the expansion in $`\alpha `$ diverge in the vacuum limit $`a0`$ and eqs. (6.2) and (6.4) obtained by the expansion in $`s_{13}`$ diverge in the high energy MSW resonance limit $`a\mathrm{\Delta }_{31}`$. As shown in the following, these divergences cancel and the probability has a finite value. Expanding $`C_{\mu \tau }^1`$ in (5.21) up to the second order, we obtain
$`C_{\mu \tau }^1`$ $``$ $`C_{\mu \tau }^{1a}+C_{\mu \tau }^{1b}+C_{\mu \tau }^{1c}+C_{\mu \tau }^{1d}+C_{\mu \tau }^{1e}`$ (6.5)
$`C_{\mu \tau }^{1a}`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}`$ (6.6)
$`C_{\mu \tau }^{1b}`$ $`=`$ $`\alpha \mathrm{sin}^22\theta _{23}\mathrm{cos}^2\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)\mathrm{sin}\mathrm{\Delta }_{31}L`$ (6.7)
$`C_{\mu \tau }^{1c}`$ $`=`$ $`\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{cos}^4\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)^2\mathrm{cos}\mathrm{\Delta }_{31}L`$ (6.8)
$`C_{\mu \tau }^{1d}`$ $`=`$ $`\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}^2L}{8a}}\right)\mathrm{sin}\mathrm{\Delta }_{31}L`$ (6.9)
$`C_{\mu \tau }^{1e}`$ $`=`$ $`s_{13}^2\mathrm{sin}^22\theta _{23}\left({\displaystyle \frac{a\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}{\displaystyle \frac{L}{2}}\right)\mathrm{sin}\mathrm{\Delta }_{31}L.`$ (6.10)
We also expand $`C_{\mu \tau }^{2a}`$ in (5.22) and $`C_{\mu \tau }^{2b}`$ in (5.23) as
$`C_{\mu \tau }^{2a}`$ $``$ $`2\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{sin}^22\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{2a}}\right)^2\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\mathrm{sin}{\displaystyle \frac{aL}{2}}`$ (6.11)
$`C_{\mu \tau }^{2b}`$ $``$ $`2s_{13}^2\mathrm{sin}^22\theta _{23}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}\right)^2\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{aL}{2}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}.`$ (6.12)
Finally, we obtain (5.32) arranging these result order by order.
Here, let us consider the applicable region of the second order formulas. $`C_{\mu \tau }^{1d}`$ diverges in the limit $`a0`$ and $`C_{\mu \tau }^{1e}`$ diverges in the limit $`a\mathrm{\Delta }_{31}`$. $`C_{\mu \tau }^{2a}`$ also diverges in the limit $`a0`$ and $`C_{\mu \tau }^{2b}`$ diverges in the limit $`a\mathrm{\Delta }_{31}`$. The divergences in $`a0`$ and in $`a\mathrm{\Delta }_{31}`$ come from the expansion of the effective masses (6.3) and (6.4) respectively. It seems that the second order formulas do not reduce to those in vacuum due to the divergence in $`a0`$ and furthermore do not reduce to those in the high energy MSW resonance point due to the divergence in $`a\mathrm{\Delta }_{31}`$. However, when we consider the pair
$`C_{\mu \tau }^{1d}+C_{\mu \tau }^{2a}=\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{sin}^22\theta _{12}{\displaystyle \frac{\mathrm{\Delta }_{31}}{2a}}\left[\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\mathrm{sin}{\displaystyle \frac{aL}{2}}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{a}}\right){\displaystyle \frac{\mathrm{\Delta }_{31}L}{4}}\mathrm{sin}(\mathrm{\Delta }_{31}L)\right],`$ (6.13)
the divergence in $`a0`$ cancel and the value converges. The obtained finite value is given by
$`\underset{a0}{lim}(C_{\mu \tau }^{1d}+C_{\mu \tau }^{2a})=\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{sin}^22\theta _{12}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)^2\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}.`$ (6.14)
Before we expand, $`C_{\mu \tau }^1`$ and $`C_{\mu \tau }^2`$ have finite values in the limit $`a0`$ and $`a\mathrm{\Delta }_{31}`$. However, the divergence appears in expansion of $`\alpha `$ and $`s_{13}`$. The cancellation of these divergences occurs between $`C_{\mu \tau }^{1d}`$ and $`C_{\mu \tau }^{2a}`$. This means that the cancellation occurs between the different terms and result in finite values, respectively, at first, which is an interesting result. Considering the pair as
$`C_{\mu \tau }^{1e}+C_{\mu \tau }^{2b}=s_{13}^2\mathrm{sin}^22\theta _{23}{\displaystyle \frac{2\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}\left[\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{aL}{2}}\mathrm{sin}{\displaystyle \frac{(\mathrm{\Delta }_{31}a)L}{2}}\left({\displaystyle \frac{\mathrm{\Delta }_{31}}{\mathrm{\Delta }_{31}a}}\right){\displaystyle \frac{aL}{4}}\mathrm{sin}(\mathrm{\Delta }_{31}L)\right],`$ (6.15)
the divergence in the limit $`a\mathrm{\Delta }_{31}`$ cancels and the value converges. The finite value is given by
$`\underset{a\mathrm{\Delta }_{31}}{lim}(C_{\mu \tau }^{1e}+C_{\mu \tau }^{2b})=s_{13}^2\mathrm{sin}^22\theta _{23}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)\left[(\mathrm{\Delta }_{31}L)\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{sin}(\mathrm{\Delta }_{31}L)\right].`$ (6.16)
The cancellation of these divergences occurs between the different terms $`C_{\mu \tau }^{1e}`$ and $`C_{\mu \tau }^{2b}`$, which is also a remarkable result.
We have shown that the second order formulas have finite values in the limit $`a0`$ and $`a\mathrm{\Delta }_{31}`$, but it is not always the same as that in the numerical calculation. Actually, the difference in fig. 4 in the limit $`a\mathrm{\Delta }_{31}`$, shows that the second order formulas have finite values but they are not in accordance with those in the numerical calculation. In order to study this, we compare the three quantities, the numerical calculation, our formulas and the second order formulas. We can learn the differences mainly in the vacuum limit $`a0`$ and the high energy MSW resonance limit $`a\mathrm{\Delta }_{31}`$ from the comparison.
At first, let us consider the vacuum limit $`a0`$. Furthermore, to simplify the discussion, we consider the case of $`s_{13}0`$. The second order formulas in the limits $`a0`$ and $`s_{13}0`$ are given by
$`\underset{a,s_{13}0}{lim}C_{\mu \tau }^{(\mathrm{double})}`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\alpha \mathrm{sin}^22\theta _{23}\mathrm{cos}^2\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)\mathrm{sin}\mathrm{\Delta }_{31}L`$ (6.17)
$`+`$ $`\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{cos}^4\theta _{12}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)^2\mathrm{cos}\mathrm{\Delta }_{31}L`$
$``$ $`\alpha ^2\mathrm{sin}^22\theta _{23}\mathrm{sin}^22\theta _{12}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)^2\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}.`$
Next, taking the limit $`a0`$ and $`s_{13}0`$ in our formulas, we obtain
$`\underset{a,s_{13}0}{lim}C_{\mu \tau }^{(\mathrm{exact})}`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}`$ (6.18)
$``$ $`2\mathrm{sin}^22\theta _{23}\mathrm{cos}^2\theta _{12}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{(\mathrm{\Delta }_{21}\mathrm{\Delta }_{31})L}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Delta }_{21}L}{2}}.`$
Expanding the oscillating part of (6.18) in our formula, it leads to (6.17) obtained from the second order formula. The condition for the expansion on the oscillating part for sufficiently good approximation is
$`L<{\displaystyle \frac{2}{\mathrm{\Delta }_{21}}}.`$ (6.19)
Next, let us consider the high energy MSW resonance limit $`a\mathrm{\Delta }_{31}`$. In order to simplify the discussion, we take the high energy MSW resonance limit $`a\mathrm{\Delta }_{31}`$ under the condition $`\alpha 0`$. In the high energy MSW resonance limit of the second order formula, we obtain
$`\underset{a\mathrm{\Delta }_{31},\alpha 0}{lim}C_{\mu \tau }^{(\mathrm{double})}`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}`$ (6.20)
$`+`$ $`s_{13}^2\mathrm{sin}^22\theta _{23}\left({\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\right)\left[(\mathrm{\Delta }_{31}L)\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }_{31}L}{2}}\mathrm{sin}(\mathrm{\Delta }_{31}L)\right].`$
Next, taking the limit $`a\mathrm{\Delta }_{31}`$ and $`\alpha 0`$ in our formulas, we obtain
$`\underset{a\mathrm{\Delta }_{31},\alpha 0}{lim}C_{\mu \tau }^{(\mathrm{exact})}`$ $`=`$ $`\mathrm{sin}^22\theta _{23}\mathrm{sin}^2{\displaystyle \frac{(1+s_{13})\mathrm{\Delta }_{31}L}{2}}`$ (6.21)
$``$ $`\mathrm{sin}^22\theta _{23}(1s_{13}^2)\mathrm{sin}{\displaystyle \frac{(1+s_{13})\mathrm{\Delta }_{31}L}{2}}\mathrm{cos}{\displaystyle \frac{(1s_{13})\mathrm{\Delta }_{31}L}{4}}\mathrm{sin}(s_{13}\mathrm{\Delta }_{31}L).`$
By expanding the oscillating part obtained from our formula (6.21), it is shown that this coincides with that from the second order formula (6.20). The condition for the expansion of the oscillating part for a sufficient approximation is given by
$`L<{\displaystyle \frac{2}{s_{13}\mathrm{\Delta }_{31}}}={\displaystyle \frac{2}{s_{13}a}}.`$ (6.22)
If the baseline length is shorter than that obtained from above condition, the second order formula becomes a good approximation. We obtain the following results about the perturbative expansion on the small parameters $`\alpha `$ and $`s_{13}`$.
1. The perturbative expansion in $`\alpha `$ actually corresponds to the expansion in $`\mathrm{\Delta }_{21}/a`$. This constrains the applicable energy for the approximate formulas. If we expand in the parameter $`\mathrm{\Delta }_{21}/a`$, the effective mass $`\mathrm{\Delta }_{21}^{\mathrm{}}`$ and the effective mixing angle $`\mathrm{sin}2\theta _{12}^{\mathrm{}}`$ diverge in the vacuum limit $`a0`$. However, these divergences cancel out each other in the calculation of the oscillation probability. Thus, the probability has a finite value, but the value largely differs from the numerical calculation in low-energy. The magnitude of this difference becomes large and serious in the case of small mixing angles and in low-energy long baseline experiments.
2. If we expand in the small mixing angle $`s_{13}`$, the effective mass $`\mathrm{\Delta }_{31}^h`$ and the effective mixing angle $`\mathrm{sin}2\theta _{13}^h`$ diverge in the MSW resonance energy limit $`a\mathrm{\Delta }_{31}`$. However, these divergences also cancel each other out in the calculation of the oscillation probability. Thus, the probability has a finite value, but the value largely differs from the numerical calculation in the high-energy MSW resonance region. This means that the second order formulas cannot be used in the high energy MSW resonance region.
In two generations, we can calculate the oscillation probabilities exactly by solving the second order equation. So, we do not need the perturbative expansion. On the other hand, the construction of the approximate formulas applicable to arbitrary matter density profile is very difficult in three generations. Therefore, we need to expand on the small parameters $`\alpha `$ and $`s_{13}`$.
### 6.2 Discussion
We have shown that the double expansion formulas up to the second order in the two small parameters $`\alpha `$ and $`s_{13}`$ does not give a good approximation in the MSW resonance region. This is because the coefficients of the small parameters have large values in the MSW resonance region. In this subsection, let us discuss some methods proposed up to present to solve this problem. The Hamiltonian $`H^{}`$ is written by four parameters. The two parameters $`(\mathrm{\Delta }m_{21}^2,\theta _{12})`$ control the physics mainly in the low-energy region and the other two parameters $`(\mathrm{\Delta }m_{31}^2,\theta _{13})`$ control the physics mainly in the high-energy region. In other words, the magnitude of $`\alpha `$ determines low-energy phenomena and the magnitude of $`s_{13}`$ determines high-energy phenomena. Both of these parameters are very small but the energy region, where the expansion converges, is different. This means that we need to treat the applicable energy region carefully when we expand on these two parameters. There are several methods in order to take into account the higher order terms of $`\alpha `$ and $`s_{13}`$ for example
1. exact formulas in constant matter density profile
2. reduction formulas taking into account the two generation part exactly
In the first method, there does not exist any error generated from the perturbative expansion, because of the exact treatment of both $`\alpha `$ and $`s_{13}`$ . Furthermore, non-perturbative effects can be easily investigated by using these exact formulas. The second method was introduced in our previous paper . In this method, we try to include the higher order terms of $`\alpha `$ and $`s_{13}`$ partially, except for the terms including the product of two small parameters. This method includes the higher order terms of $`\alpha `$ and $`s_{13}`$ and is simply and applicable even in the case of arbitrary matter density profile . Although this method uses only the second order approximation of the amplitude, it has the notable feature that the third order (three generation) effects such as CP violation can be calculated.
## 7 Summary
In this paper, we consider the method how to approximates the neutrino oscillation probabilities in matter under three generations and the obtained results are summarized as following.
1. In the framework of two generation neutrino oscillation, we discuss the applicable region of the perturbative expansion on the small mixing angle in matter. The result of the perturbation differs largely from the exact numerical calculation in the MSW resonance point. This means that non-perturbative effects are important even for the neutrino oscillation in two generations.
2. We extend the method to calculate the approximate formulas, in which non-perturbative effects of the small parameters $`\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{sin}\theta _{13}`$, to all channels. Under the conditions, $`\theta _{23}=45^{}`$ and the symmetric matter density profile, we derive simple approximate formulas of the probabilities in all channels by using the unitary relation. Although all these approximate formulas are expressed by the amplitudes calculated within the framework of two generations, it has a notable feature that the three generation effects such as CP violation can also be calculated.
3. In the three generation neutrino oscillation with matter, we investigate non-perturbative effects of the two small parameters $`\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{sin}\theta _{13}`$. We compare our approximate formulas with those from the double expansion, which include the terms up to the second order in the low and high energy MSW resonance regions. The obtained result is that the second order formulas show large differences from the exact numerical calculation, which means that non-perturbative effects of the small $`\mathrm{\Delta }m_{21}^2/\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{sin}\theta _{13}`$ become important in the MSW resonance region.
Finally, we describe two problems that we could not fully address in this paper, and which are tasks for future research.
1. The approximate formulas in this paper are derived by using the condition $`\theta _{23}=45^{}`$, which is the center value obtained from the atmospheric neutrino experiments. However, but differences from this value may exist within 90% confidence level.
2. The condition for the symmetric matter density is satisfied in the 1-dimensional models, like the PREM and the ak135f, but the actual matter density, for example, that from J-PARC to Beijing is not symmetric . Therefore, our aim for future work is, to derive more sophisticated approximate formulas that hold not only in symmetric matter but in arbitrary matter as well.
To solve the above two problems are the future works This is now included in the upper sentence.
## Acknowledgement
We are grateful to H. Yokomakura, and T. Yoshikawa for useful discussions and careful reading of our manuscript. We would like to thank Prof. Wilfried Wunderlich (Nagoya Inst. Technology) for helpful comments and advice on English expressions.
## Appendix A General Feature of CP Dependence
In this appendix we calculate the coefficients of the probabilities in detail. We show that the 2-3 mixing angle and the CP phase are not affected by matter, from a different point of view as described in our previous paper . This result means that we only have to consider the matter effects on four parameters $`(\mathrm{\Delta }m_{21}^2,\theta _{12})`$ and $`(\mathrm{\Delta }m_{31}^2,\theta _{13})`$. By using this result, we can understand the matter effects in three generations, which become complex compared with that in two generations.
### A.1 Remarkable Features of Effective Masses
In this subsection, we show that $`(\theta _{23},\delta )`$ do not affect the effective mass in three generation Hamiltonian. If we express the effective Hamiltonian in matter as
$`H=U\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})U^{}+\mathrm{diag}(a,0,0),`$ (A.1)
the equation of eigenvalue is given by
$`\mathrm{det}(tH)`$ $`=`$ $`t^3(\mathrm{\Delta }_{21}+\mathrm{\Delta }_{31}+a)t^2`$ (A.2)
$`+`$ $`(\mathrm{\Delta }_{21}\mathrm{\Delta }_{31}+a(\mathrm{\Delta }_{21}(1|U_{e2}|^2)+\mathrm{\Delta }_{31}(1|U_{e3}|^2)))ta\mathrm{\Delta }_{21}\mathrm{\Delta }_{31}|U_{e1}|^2=0,`$
and by solving this equation, we obtain the effective masses as
$`\lambda _1`$ $`=`$ $`{\displaystyle \frac{A}{3}}{\displaystyle \frac{1}{3}}\sqrt{A^23B}S{\displaystyle \frac{\sqrt{3}}{3}}\sqrt{A^23B}\sqrt{1S^2}`$ (A.3)
$`\lambda _2`$ $`=`$ $`{\displaystyle \frac{A}{3}}{\displaystyle \frac{1}{3}}\sqrt{A^23B}S+{\displaystyle \frac{\sqrt{3}}{3}}\sqrt{A^23B}\sqrt{1S^2}`$ (A.4)
$`\lambda _3`$ $`=`$ $`{\displaystyle \frac{A}{3}}+{\displaystyle \frac{2}{3}}\sqrt{A^23B}S`$ (A.5)
, where $`A,B,C`$ and $`S`$ are defined by
$`A`$ $`=`$ $`\mathrm{\Delta }_{21}+\mathrm{\Delta }_{31}+a`$ (A.6)
$`B`$ $`=`$ $`\mathrm{\Delta }_{21}\mathrm{\Delta }_{31}+a[\mathrm{\Delta }_{21}(1|U_{e2}|^2)+\mathrm{\Delta }_{31}(1|U_{e3}|^2)]`$ (A.7)
$`C`$ $`=`$ $`a\mathrm{\Delta }_{21}\mathrm{\Delta }_{31}|U_{e1}|^2`$ (A.8)
$`S`$ $`=`$ $`\mathrm{cos}\left[{\displaystyle \frac{1}{3}}\mathrm{arccos}\left({\displaystyle \frac{2A^39AB+27C}{2\sqrt{(A^23B)^3}}}\right)\right].`$ (A.9)
These effective masses depend only on the following three vacuum mixing angles
$`|U_{e1}|=c_{12}c_{13},|U_{e2}|=s_{12}c_{13},|U_{e3}|=s_{13}.`$ (A.10)
One can see from these equalities that the effective masses are independent of the 2-3 mixing angle $`\theta _{23}`$ and the CP phase $`\delta `$. Next, let us consider this result from a different point of view.
### A.2 Decomposition of 2-3 mixing and CP Phase from Hamiltonian
In this section, we separate $`\theta _{23}`$ and $`\delta `$ from the Hamiltonian and we study the dependence of the amplitudes on the two small parameters $`\alpha `$ and $`s_{13}`$. The Standard Parametrization is defined by
$`U=O_{23}\mathrm{\Gamma }O_{13}\mathrm{\Gamma }^{}O_{12},`$ (A.11)
where the CP phase matrix $`\mathrm{\Gamma }`$ is given by
$$\mathrm{\Gamma }=\mathrm{diag}(1,1,e^{i\delta }).$$
(A.12)
The CP phase matrix $`\mathrm{\Gamma }`$ and the 1-2 mixing matrix $`O_{12}`$ are commutable as
$$[\mathrm{\Gamma },O_{12}]=[\mathrm{\Gamma },\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})]=0.$$
(A.13)
Therefore, the Hamiltonian can be separated as
$`H(t)=U\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})U^{}+\mathrm{diag}(a(t),0,0)=O_{23}\mathrm{\Gamma }H^{}(t)\mathrm{\Gamma }^{}O_{23}^T,`$ (A.14)
where $`H^{}(t)`$ is defined by
$$H^{}(t)=O_{13}O_{12}\mathrm{diag}(0,\mathrm{\Delta }_{21},\mathrm{\Delta }_{31})O_{12}^TO_{13}^T+\mathrm{diag}(a(t),0,0).$$
(A.15)
This means that the 2-3 mixing and the CP phase can be separated from the part which includes the matter effects $`a(t)`$.
In the case of constant matter density profile, we obtain
$`\mathrm{det}(\lambda H)=\mathrm{det}(\lambda H^{}),`$ (A.16)
the 2-3 mixing angle and the CP phase do not affect the eigenvalue equation. Accordingly, the effective masses are independent of the 2-3 mixing angle and the CP phase, which coincide with the result obtained in the previous subsection.
### A.3 Exact CP and 2-3 mixing Dependence of Oscillation Probabilities
Here, let us consider the case in which we apply the above discussion used in the Hamiltonian to the amplitude. Solving the Schrodinger eq. for the amplitude in matter, we obtain
$$S(t)=\mathrm{T}\mathrm{exp}\left\{iH(t)𝑑t\right\}.$$
(A.17)
By using this, we obtain
$`S(t)=\mathrm{T}\mathrm{exp}\left\{i{\displaystyle O_{23}\mathrm{\Gamma }H^{}(t)\mathrm{\Gamma }^{}O_{23}^T𝑑t}\right\}=O_{23}\mathrm{\Gamma }\mathrm{T}\mathrm{exp}\left\{i{\displaystyle H^{}(t)𝑑t}\right\}\mathrm{\Gamma }^{}O_{23}^T=O_{23}\mathrm{\Gamma }S^{}(t)\mathrm{\Gamma }^{}O_{23}^T`$ (A.18)
from (A.14). Therefore, $`S(t)`$ satisfies
$$S(t)=O_{23}\mathrm{\Gamma }S^{}(t)\mathrm{\Gamma }^{}O_{23}^T.$$
(A.19)
From this equation, we obtain
$`P(\nu _e\nu _e)`$ $`=`$ $`C_{ee},`$ (A.20)
$`P(\nu _\alpha \nu _\beta )`$ $`=`$ $`A_{\alpha \beta }\mathrm{cos}\delta +B_{\alpha \beta }\mathrm{sin}\delta +C_{\alpha \beta },`$ (A.21)
when the initial or final state is $`\nu _e`$, and
$`P(\nu _\alpha \nu _\beta )`$ $`=`$ $`A_{\alpha \beta }\mathrm{cos}\delta +B_{\alpha \beta }\mathrm{sin}\delta +C_{\alpha \beta }+D_{\alpha \beta }\mathrm{cos}2\delta +E_{\alpha \beta }\mathrm{sin}2\delta ,`$ (A.22)
in the case of $`\nu _\alpha ,\nu _\beta =\nu _\mu ,\nu _\tau `$ . The final result is given by
$`P(\nu _e\nu _e)`$ $`=`$ $`C_{ee}=|S_{ee}^{}|^2,`$ (A.23)
$`P(\nu _e\nu _\mu )`$ $`=`$ $`A_{e\mu }\mathrm{cos}\delta +B_{e\mu }\mathrm{sin}\delta +C_{e\mu },`$ (A.24)
$`A_{e\mu }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[S_{\mu e}^{{}_{}{}^{}}S_{\tau e}^{}]c_{23}s_{23},`$ (A.25)
$`B_{e\mu }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[S_{\mu e}^{{}_{}{}^{}}S_{\tau e}^{}]c_{23}s_{23},`$ (A.26)
$`C_{e\mu }`$ $`=`$ $`|S_{\mu e}^{}|^2c_{23}^2+|S_{\tau e}^{}|^2s_{23}^2,`$ (A.27)
$`P(\nu _e\nu _\tau )`$ $`=`$ $`A_{e\tau }\mathrm{cos}\delta +B_{e\tau }\mathrm{sin}\delta +C_{e\tau },`$ (A.28)
$`A_{e\tau }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[S_{\mu e}^{{}_{}{}^{}}S_{\tau e}^{}]c_{23}s_{23},`$ (A.29)
$`B_{e\tau }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[S_{\mu e}^{{}_{}{}^{}}S_{\tau e}^{}]c_{23}s_{23},`$ (A.30)
$`C_{e\tau }`$ $`=`$ $`|S_{\mu e}^{}|^2s_{23}^2+|S_{\tau e}^{}|^2c_{23}^2,`$ (A.31)
$`P(\nu _\mu \nu _\mu )`$ $`=`$ $`A_{\mu \mu }\mathrm{cos}\delta +B_{\mu \mu }\mathrm{sin}\delta +C_{\mu \mu }+D_{\mu \mu }\mathrm{cos}2\delta +E_{\mu \mu }\mathrm{sin}2\delta ,`$ (A.32)
$`A_{\mu \mu }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[(S_{\mu \mu }^{}c_{23}^2+S_{\tau \tau }^{}s_{23}^2)^{}(S_{\tau \mu }^{}+S_{\mu \tau }^{})]c_{23}s_{23},`$ (A.33)
$`B_{\mu \mu }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[(S_{\mu \mu }^{}c_{23}^2+S_{\tau \tau }^{}s_{23}^2)^{}(S_{\tau \mu }^{}S_{\mu \tau }^{})]c_{23}s_{23},`$ (A.34)
$`C_{\mu \mu }`$ $`=`$ $`|S_{\mu \mu }^{}|^2c_{23}^4+|S_{\tau \tau }^{}|^2s_{23}^4+(|S_{\mu \tau }^{}|^2+|S_{\tau \mu }^{}|^2+2\mathrm{R}\mathrm{e}[S_{\mu \mu }^{{}_{}{}^{}}S_{\tau \tau }^{}])c_{23}^2s_{23}^2,`$ (A.35)
$`D_{\mu \mu }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[S_{\tau \mu }^{{}_{}{}^{}}S_{\mu \tau }^{}]c_{23}^2s_{23}^2,`$ (A.36)
$`E_{\mu \mu }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[S_{\tau \mu }^{{}_{}{}^{}}S_{\mu \tau }^{}]c_{23}^2s_{23}^2,`$ (A.37)
$`P(\nu _\tau \nu _\tau )`$ $`=`$ $`A_{\tau \tau }\mathrm{cos}\delta +B_{\tau \tau }\mathrm{sin}\delta +C_{\tau \tau }+D_{\tau \tau }\mathrm{cos}2\delta +E_{\tau \tau }\mathrm{sin}2\delta ,`$ (A.38)
$`A_{\tau \tau }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[(S_{\mu \mu }^{}s_{23}^2+S_{\tau \tau }^{}c_{23}^2)^{}(S_{\tau \mu }^{}+S_{\mu \tau }^{})]c_{23}s_{23},`$ (A.39)
$`B_{\tau \tau }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[(S_{\mu \mu }^{}s_{23}^2+S_{\tau \tau }^{}c_{23}^2)^{}(S_{\tau \mu }^{}S_{\mu \tau }^{})]c_{23}s_{23},`$ (A.40)
$`C_{\tau \tau }`$ $`=`$ $`|S_{\mu \mu }^{}|^2s_{23}^4+|S_{\tau \tau }^{}|^2c_{23}^4+(|S_{\mu \tau }^{}|^2+|S_{\tau \mu }^{}|^2+2\mathrm{R}\mathrm{e}[S_{\mu \mu }^{{}_{}{}^{}}S_{\tau \tau }^{}])c_{23}^2s_{23}^2,`$ (A.41)
$`D_{\tau \tau }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[S_{\tau \mu }^{{}_{}{}^{}}S_{\mu \tau }^{}]c_{23}^2s_{23}^2,`$ (A.42)
$`E_{\tau \tau }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[S_{\tau \mu }^{{}_{}{}^{}}S_{\mu \tau }^{}]c_{23}^2s_{23}^2,`$ (A.43)
$`P(\nu _\mu \nu _\tau )`$ $`=`$ $`A_{\mu \tau }\mathrm{cos}\delta +B_{\mu \tau }\mathrm{sin}\delta +C_{\mu \tau }+D_{\mu \tau }\mathrm{cos}2\delta +E_{\mu \tau }\mathrm{sin}2\delta ,`$ (A.44)
$`A_{\mu \tau }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[(S_{\mu \mu }^{}S_{\tau \tau }^{})^{}(S_{\tau \mu }^{}c_{23}^2S_{\mu \tau }^{}s_{23}^2)]c_{23}s_{23},`$ (A.45)
$`B_{\mu \tau }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[(S_{\mu \mu }^{}S_{\tau \tau }^{})^{}(S_{\tau \mu }^{}c_{23}^2+S_{\mu \tau }^{}s_{23}^2)]c_{23}s_{23},`$ (A.46)
$`C_{\mu \tau }`$ $`=`$ $`|S_{\mu \tau }^{}|^2s_{23}^4+|S_{\tau \mu }^{}|^2c_{23}^4+(|S_{\mu \mu }^{}|^2+|S_{\tau \tau }^{}|^22\mathrm{R}\mathrm{e}[S_{\mu \mu }^{{}_{}{}^{}}S_{\tau \tau }^{}])c_{23}^2s_{23}^2,`$ (A.47)
$`D_{\mu \tau }`$ $`=`$ $`2\mathrm{R}\mathrm{e}[S_{\tau \mu }^{{}_{}{}^{}}S_{\mu \tau }^{}]c_{23}^2s_{23}^2,`$ (A.48)
$`E_{\mu \tau }`$ $`=`$ $`2\mathrm{I}\mathrm{m}[S_{\tau \mu }^{{}_{}{}^{}}S_{\mu \tau }^{}]c_{23}^2s_{23}^2.`$ (A.49)
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# Diagrammatic approach in the variational coupled-cluster method
## I introduction
A microscopic quantum many-body theory is mainly to study correlations between the constituent particles of a quantum system. One of the most successful quantum many-body theories is the method of correlated basis functionals (CBF) which includes real-space many-body correlation functions in the ground state and employs similar techniques as in classical statistical mechanics to calculate the corresponding distribution functions. The CBF has proved to be one of very few many-body theories capable of dealing with strongly correlated boson systems in liquid phase. Another successful quantum many-body theory is the coupled-cluster method (CCM) in which excitation operators with respect to an uncorrelated model state are employed to construct the many-body correlations in an exponentiated operator factor in the ket ground-state. The distribution functions are not needed in the CCM as the Hamiltonian expectation value is straightforwardly calculated as a finite order polynomial in terms of the correlation coefficients. This is because the bra ground-state in the CCM is not the hermitian conjugate of the ket ground-state but is in a simple, linear form . The CCM has proved to be one of most powerful techniques in calculating ground-state energy for many non-liquid fermion systems such as atoms, molecules and electron gas. However, a general many-body theory capable of dealing with strongly correlated fermion systems in liquid phase is still needed.
In our earlier papers , by an application to bipartite quantum antiferromagnetic lattice systems, we have extended the CCM to a variational formalism where, in contrast to the traditional CCM, the bra and ket ground-states are hermitian conjugate to one-another. Two sets of the distribution functions were introduced to evaluate Hamiltonian expectation. We developed an algebraic scheme to calculate these distribution functions through self-consistent sets of equations. In this paper, we present an alternative scheme based on diagrams to calculate these distribution functions. As in the CBF, a generating functional is introduced and calculated by a linked-cluster expansion in terms of diagrams. These diagrams are constructed according to a few simple rules and using only three basic elements: (a) dots representing the ket-state correlation coefficients, (b) exchange lines representing the bra-state correlation coefficients, and (c) Pauli line representing Pauli exclusion principle manifested by the spin-1/2 operator property $`(s^\pm )^2=0`$ in our spin model example. In this fashion, infinite diagrams can be resummed in a straightforward manner. As a simple application in our spin model, the spin-wave theory (SWT) is reproduced by including all so-called ring diagrams without any Pauli line. Approximations beyond SWT by resummations including Pauli lines are also given. One such approximation, which includes all so-called super-ring diagrams by a resummation of infinite Pauli lines in addition to resummations of ring diagrams, produces a convergent, precise result for the order-parameter of the one-dimensional isotropic model, contrast to the well-known divergence of SWT. Furthermore, the diagrammatic analysis discussed here forms a basis for a possible combination of the variational CCM and CBF. Such a unified theory may prove to be capable of dealing with strongly-correlated fermion system in liquid phase.
## II The variational coupled-cluster method
We take the spin-1/2 antiferromagnetic XXZ model on a bipartite lattice as our example. The Hamiltonian is given by
$$H=\frac{1}{2}\underset{l,\rho }{}H_{l,l+n}=\frac{1}{2}\underset{l,n}{}\left(As_l^zs_{l+n}^z+\frac{1}{2}s_l^+s_{l+n}^{}+\frac{1}{2}s_l^{}s_{l+n}^+\right),$$
(1)
where $`A>0`$ is the anisotropy constant, the index $`l`$ runs over all lattice sites, $`n`$ runs over all $`z`$ nearest-neighbour sites, and $`s^\pm `$ are the usual spin raising $`(+)`$ and lowering $`()`$ operators. As in our earlier work using the traditional CCM , we take the model state as the classical Néel state with alternating spin-up and spin-down sublattices. As before, we shall exclusively use index $`i`$ for the spin-up sublattice and the index $`j`$ for the spin-down sublattice. The many-spin correlations in its ground state of Eq. (1) can then be included by considering the excited states with respect to the uncorrelated model state. These excited states are constructed by applying the so-called configuration creation operators $`C_I^{}`$ to the Néel model state with the nominal index $`I`$ labelling these operators. In our spin model, the operators $`C_I^{}`$ are given by any combination of the spin-flip operators to the Néel state $`s_i^{}`$ and $`s_j^+`$; the index $`I`$ in this case corresponds to the collection of the lattice indices ($`i`$’s and $`j`$’s). The hermitian conjugate operators of $`C_I^{}`$ are the configuration destruction operator $`C_I`$, given by any combination of $`s_i^+`$ and $`s_j^{}`$. For example, the two-spin flip creation operator is given by $`C_{ij}^{}=s_i^{}s_j^+/2s`$, and their destruction counterpart, $`C_{ij}=s_i^+s_j^{}/2s`$, where $`s`$ is the spin quantum number. Although we are mainly interested in $`s=1/2`$ in this article, we keep the factor of $`1/2s`$ for the purpose of comparison with the large-$`s`$ expansion.
As discussed in details in our earlier paper , we use Coester representation for both the ket and bra ground-states and write
$$|\mathrm{\Psi }=e^S|\mathrm{\Phi },S=\underset{I}{}F_IC_I^{};\stackrel{~}{\mathrm{\Psi }}|=\mathrm{\Phi }|e^{\stackrel{~}{S}},\stackrel{~}{S}=\underset{I}{}\stackrel{~}{F}_IC_I,$$
(2)
with
$$\underset{I}{}F_IC_I^{}=\underset{n=1}{\overset{N/2}{}}\underset{i_1\mathrm{},j_1\mathrm{}}{}f_{i_1\mathrm{},j_1\mathrm{}}\frac{s_{i_1}^{}\mathrm{}s_{i_n}^{}s_{j_1}^+\mathrm{}s_{j_n}^+}{(2s)^n},$$
(3)
for the ket state and the corresponding hermitian conjugate of Eq. (3) for the bra state, using notation $`\stackrel{~}{F}_I=\stackrel{~}{f}_{i_1\mathrm{},j_1\mathrm{}}`$ for the bra-state coefficients. The coefficients $`\{F_I,\stackrel{~}{F}_I\}`$ are then determined by the usual variational equations as
$$\frac{\delta H}{\delta \stackrel{~}{F}_I}=\frac{\delta H}{\delta F_I}=0,H\frac{\stackrel{~}{\mathrm{\Psi }}|H|\mathrm{\Psi }}{\stackrel{~}{\mathrm{\Psi }}|\mathrm{\Psi }}.$$
(4)
We define the so-called bare distribution functions as
$$g_I=C_I,\stackrel{~}{g}_I=C_I^{},$$
(5)
where we have exchanged the definition of $`g_I`$ with that of $`\stackrel{~}{g}_I`$ as compared with those in Ref. 7 for purely notational reason. The Hamiltonian expectation $`H`$ is shown, in general, to be a function containing up to linear terms in $`g_I`$ and $`\stackrel{~}{g}_I`$ and finite order polynomial in $`F_I`$ (or in $`\stackrel{~}{F}_I`$) in Eq. (21) of Ref. 7:
$$H=(g_I,\stackrel{~}{g}_I,F_I)=(\stackrel{~}{g}_I,g_I,\stackrel{~}{F}_I).$$
(6)
Two systematic schemes have been developed for calculating the distribution functions of Eqs. (5): one is algeraic and the other is diagrammatic. In the algebraic approach, by taking the advantage of the properties of the operators, it is straightforward to derive the following self-consistent sets of equations for the distribution functions
$$g_I=G(\stackrel{~}{g}_J,F_J),\stackrel{~}{g}_I=G(g_J,\stackrel{~}{F}_J),$$
(7)
where $`G`$ is a function containing up to linear terms in $`\stackrel{~}{g}_J`$ (or $`g_J`$) and finite order polynomial in $`F_J`$ (or $`\stackrel{~}{F}_J`$). Eqs. (7) are solved for $`g_I`$ and $`\stackrel{~}{g}_I`$ as a function of $`F_I`$ and $`\stackrel{~}{F}_I`$. The variational Eqs. (4) are then carried to determined the optimum $`F_I`$ and $`\stackrel{~}{F}_I`$. In this algebraic calculation, direct comparison with the traditional CCM can be made. It is shown that the CCM is a linear approximation to one set of distributions as simply $`\stackrel{~}{g}_I\stackrel{~}{F}_I`$, which is a poor approximation for the spin-spin correlation function and low-lying excitations. More detailed comparison was given in Ref. 7. Here we present the diagrammatic scheme similar to that in CBF to calculate these distribution functions. We like to point out that Eqs. (2) and (4)-(7) are the main general equations of the variational CCM.
## III Diagrammatic representation of generating functional
In this section, we calculate the bare distribution functions $`g_I`$ and $`\stackrel{~}{g}_I`$ of Eq. (5) by employing a diagrammatic scheme. As a demonstration, we consider a simple truncation approximation in which the correlation operators $`S`$ and $`\stackrel{~}{S}`$ of Eqs. (2) retain only the two-spin flip operators as (the so-called SUB2 approximation as defined in Ref. 10),
$$S\underset{ij}{}f_{ij}C_{ij}^{}=\underset{ij}{}f_{ij}\frac{s_i^{}s_j^+}{2s},\stackrel{~}{S}\underset{ij}{}\stackrel{~}{f}_{ij}C_{ij}=\underset{ij}{}\stackrel{~}{f}_{ij}\frac{s_i^+s_j^{}}{2s}.$$
(8)
Using the usual angular momentum commutations $`[s_l^z,s_l^{}^\pm ]=\pm s_l^\pm \delta _{ll^{}}`$, $`[s_l^+,s_l^{}^{}]=2s_l^z\delta _{ll^{}}`$, and the Néel state eigenequations, $`s_i^z|\mathrm{\Phi }=s|\mathrm{\Phi },s_j^z|\mathrm{\Phi }=s|\mathrm{\Phi }`$, it is a straightforward calculation to derive expectation value of any physical operators in terms distribution functions of Eqs. (5). In this approximation, for example, the expectation value of Eq. (1) is given by
$$H_{ij}=As_i^zs_j^z+\frac{1}{2}(g_{ij}+\stackrel{~}{g}_{ij}),$$
(9)
where $`s_i^zs_j^z`$ is calculated as
$$s_i^zs_j^z=s^2+s\left(\underset{i^{}}{}\rho _{i^{}j}+\underset{j^{}}{}\rho _{ij^{}}\right)\left(\underset{i^{}j^{}}{}\rho _{ij^{},i^{}j}+\rho _{ij}\right),$$
(10)
$`\rho _{ij}`$ is the usual full one-body distribution function defined as
$$\rho _{ij}f_{ij}\stackrel{~}{g}_{ij}=f_{ij}\frac{s_i^{}s_j^+}{2s},$$
(11)
and where $`\rho _{ij,i^{}j^{}}`$ is the full two-body distribution function define as
$$\rho _{ij,i^{}j^{}}f_{ij}f_{i^{}j^{}}\stackrel{~}{g}_{ij,i^{}j^{}}=f_{ij}f_{i^{}j^{}}\frac{s_i^{}s_j^+s_i^{}^{}s_j^{}^+}{(2s)^2}.$$
(12)
The order parameter is given by
$$s_i^z=s\rho ,$$
(13)
where $`\rho =_j\rho _{ij}`$, taking the advantage of translational invariance.
We define a generating functional $`W`$ in the usual fashion as,
$$W\mathrm{ln}\stackrel{~}{\mathrm{\Psi }}|\mathrm{\Psi },$$
(14)
so that the bare and full distribution functions can be simply expressed as functional derivatives of $`W`$. For example, the one-body and two-body bare functions are given by
$$\stackrel{~}{g}_1=C_1^{}=\frac{\delta W}{\delta f_1},\stackrel{~}{g}_{12}=C_1^{}C_2^{}=\frac{\delta ^2W}{\delta f_1\delta f_2}+\stackrel{~}{g}_1\stackrel{~}{g}_2,$$
(15)
where, for simplicity, we have employed notation $`1(i_1,j_1)`$ so that $`f_1=f_{i_1j_1}`$ etc.; and the structure function $`S_{12}`$ has the usual relation as in the CBF as
$$S_{12}f_1\frac{\delta \rho _2}{\delta f_1}=\rho _1\delta _{12}+\rho _{12}\rho _1\rho _2,$$
(16)
where $`\rho _1=\rho _{i_1j_1}`$, etc.
We now write $`W`$ in terms of a linked-cluster expansion
$$W=\mathrm{sum}\mathrm{of}\mathrm{all}\mathrm{linked}\mathrm{cluster}\mathrm{contributions}.$$
(17)
The main task of this section is to find a diagrammatic scheme to categorize this expansion. We first expand the ket-state operator in the simplified notation, $`e^S=1+S+\frac{1}{2!}S^2+\mathrm{}=1+f_1C_1^{}+\frac{1}{2!}f_1f_2C_1^{}C_2^{}+\mathrm{}`$, where in the last equation, the summation over all indices is understood. The normalization integral,
$$\stackrel{~}{\mathrm{\Psi }}|\mathrm{\Psi }=1+\stackrel{~}{f}_1^{}f_1C_1^{}C_1^{}+\frac{1}{(2!)^2}\stackrel{~}{f}_2^{}\stackrel{~}{f}_1^{}f_1f_2C_2^{}C_1^{}C_1^{}C_2^{}+\mathrm{},$$
(18)
can be evaluated straightforwardly for the first few terms. In the above series, the primed indices are used for bra state expansion. We notice that each term of Eq. (18) contains equal number of creation and destruction operators (otherwise, the expectation is zero).
The first-order expectation is easily calculated as $`C_1^{}C_1^{}=\frac{1}{(2s)^2}\mathrm{\Phi }|s_{j_1^{}}^{}s_{i_1^{}}^+s_{i_1}^{}s_{j_1}^+|\mathrm{\Phi }=\delta _{i_1^{}i_1}\delta _{j_1^{}j_1}`$. Hence we have, writing out the summation explicitly,
$$1\mathrm{s}\mathrm{t}\mathrm{order}=\underset{1}{}f_1\stackrel{~}{f}_1.$$
(19)
The calculation of the second-order expectation $`C_2^{}C_1^{}C_1^{}C_2^{}`$ is slightly more complicated. We first consider the case of $`12`$ (i.e., $`i_1i_2`$ and $`j_1j_2`$). There are four nonzero terms
$`\left(\delta _{i_1^{}i_1}\delta _{i_2^{}i_2}+\delta _{i_1^{}i_2}\delta _{i_2^{}i_1}\right)\left(\delta _{j_1^{}j_1}\delta _{j_2^{}j_2}+\delta _{j_1^{}j_2}\delta _{j_2^{}j_1}\right)`$
$`=`$ $`\left(\delta _{i_1^{}i_1}\delta _{i_2^{}i_2}+\delta _{i_1^{}i_2}\delta _{i_2^{}i_1}\right)\left(\delta _{j_1^{}j_1}\delta _{j_2^{}j_2}+\delta _{j_1^{}j_2}\delta _{j_2^{}j_1}\right).`$
The cases when $`i_1=i_2`$ and/or $`j_1=j_2`$ can be easily accounted for by introducing a factor involving the usual delta functions as,
$`\left(1{\displaystyle \frac{1}{2s}}\delta _{i_1i_2}\right)\left(1{\displaystyle \frac{1}{2s}}\delta _{j_1j_2}\right)=1+\mathrm{\Delta }_{12},`$
with a definition,
$$\mathrm{\Delta }_{12}\frac{1}{\left(2s\right)}\left(\delta _{i_1i_2}+\delta _{j_1j_2}\right)+\frac{1}{\left(2s\right)^2}\delta _{i_1i_2}\delta _{j_1j_2}.$$
(20)
This is because $`\left(s_i^{}\right)^2=\left(s_j^+\right)^2=0`$ for $`s=1/2`$, a manifestation of Pauli exclusion principle. The second-order contribution is hence derived as
$`{\displaystyle \frac{1}{\left(2!\right)^2}}\stackrel{~}{f}_2^{}\stackrel{~}{f}_1^{}f_1f_2\left(\delta _{i_1^{}i_1}\delta _{i_2^{}i_2}+\delta _{i_1^{}i_2}\delta _{i_2^{}i_1}\right)\left(\delta _{j_1^{}j_1}\delta _{j_2^{}j_2}+\delta _{j_1^{}j_2}\delta _{j_2^{}j_1}\right)\left(1+\mathrm{\Delta }_{12}\right)`$
$`=`$ $`{\displaystyle \frac{1}{2!}}\left[\left(f_1\stackrel{~}{f}_1\right)\left(f_2\stackrel{~}{f}_2\right)+f_1f_2\stackrel{~}{f}_{i_1j_2}\stackrel{~}{f}_{i_2j_1}\right]\left(1+\mathrm{\Delta }_{12}\right),`$
where the second term inside the square brackets clearly represents the so-called exchange contributions. We notice that we did not consider explicitly the Pauli exclusion principle for the bra-state operators in the above derivations as the delta functions for the ket state operators also take this principle into account due to the fact that each of the bra-state operators always need to match one of the ket-state operators in order to give nonzero contribution. This is also true for higher-order terms. We also notice the expression of Eq. (20) is in fact also correct for spin quantum number $`s>1/2`$ because, for a general $`s`$,
$`{\displaystyle \frac{1}{2(2s)^2}}\mathrm{\Phi }|(s_i^+)^2(s_i^{})^2|\mathrm{\Phi }=\left(1{\displaystyle \frac{1}{2s}}\right),`$
etc. For higher-order terms in the expansion of Eq. (18), the extension of Pauli exclusion principle can be simply written as a product of two-body factors as
$$\underset{n>m}{}\left(1+\mathrm{\Delta }_{nm}\right).$$
(21)
We notice that the above product in general is not exact any more but an approximation for $`s>1/2`$ as the three-body effects (e.g., from $`(s_{i_1}^{})^3`$ when $`i_1=i_2=i_3`$) have been ignored.
In order to extend to higher-order calculations including the exchange contributions, we need a systematic graph representation. For this purpose, as shown in Fig. 1, we use a solid dot to represent the ket state coefficient $`f_1`$ with $`1=(i_1j_1)`$ as defined earlier; a (directed) exchange line drawing from $`i_1`$ to $`j_2`$ to represent the bra state coefficient $`\stackrel{~}{f}_{i_1j_2}`$; and a Pauli (dashed) line drawing between any two dots to represent delta function $`\mathrm{\Delta }_{12}`$ of Eq. (20). With these graphic notations, a linked contribution is represented by a connected diagram. After the detailed calculations up to 5th order, we have established the following simple and complete rules for construction of these diagrams in the normalization integral of Eq. (18):
* The $`k`$th-order contribution consists of all possible diagrams involving $`k`$ dots;
* In each diagram the number of dots equal to number of exchange lines;
* A dot is always connected by exchange lines (leaving and coming) hence exchange lines always form loops;
* Between any pair of dots one can draw at most one Pauli line;
* The contribution of each diagram is divided by its symmetry factor;
* Summations over all indices involved.
We first consider the case without any Pauli line, $`\mathrm{\Delta }_{nm}=0`$. (This is equivalent to turning spin operators to boson operators as will be shown later.) For example, the first-order contribution of Eq. (19) is simply a dot with an exchange line leaving and coming as shown as diagram $`a`$ in Fig. 2, where the direction of exchange line is clockwise as in most other diagrams (we therefore do not show arrows of exchange lines explicitly). The second-order contribution with $`\mathrm{\Delta }_{12}=0`$ is given by two diagrams $`b`$ and $`c`$ in Fig. 2, namely $`\frac{1}{2!}(b+c)`$. The 3rd-order contribution is calculated as $`\frac{1}{3!}(a+3b+2c)`$ and is shown in Fig. 3, where factor $`3`$ for diagram $`b`$ is due to the three equivalent diagrams by rotation and the factor $`2`$ for diagram $`c`$ comes from the two equivalent diagram with opposite directions, one clockwise the other counter-clockwise (this is referred to as parity symmetry). In similar fashion one can write down the 4th-order contribution as shown in Fig. 4 for the corresponding diagrams as
$$4\mathrm{t}\mathrm{h}\mathrm{order}=\frac{1}{4!}\left(a+6b+8c+3d+6e\right),$$
(22)
where the coefficient numbers are the symmetry factors of the corresponding diagrams. For example, the factor $`6`$ for diagram $`e`$ is due to the fact that there are three equivalent diagrams each with parity symmetry factor of $`2`$. The 5th-order contributions include 7 independent diagrams, as shown in Fig. 5, namely
$$5\mathrm{t}\mathrm{h}\mathrm{order}=\frac{1}{5!}\left(a+10b+20c+15d+30e+20f+24g\right).$$
(23)
We notice that, in all these results, the last term represents a ring diagram with $`k`$ dots in the $`k`$th order contribution. We use $`R_k`$ to represent this ring diagram with the symmetry factor $`(k1)!/k!=1/k`$. For example, the 4th-order ring contribution is, writing out the summations explicitly
$$R_4=\frac{1}{4}\underset{1,2,3,4}{}f_1\stackrel{~}{f}_{i_1j_2}f_2\stackrel{~}{f}_{i_2j_3}f_3\stackrel{~}{f}_{i_3j_4}f_4\stackrel{~}{f}_{i_4j_1}.$$
(24)
Furthermore, the other terms in these $`k`$-order contributions are simply a product of smaller ring contributions. This property can be extended to higher order. (For this purpose one needs to apply symmetric group $`𝒮_n`$ to count the number of diagrams. See, for example, Ref. 11). We are now in position to write all contributions without any Pauli line in terms of these ring diagrams. The normalization integral of Eq. (18) is then written as
$`\stackrel{~}{\mathrm{\Psi }}|\mathrm{\Psi }_{\mathrm{\Delta }_{nm}=0}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}{\displaystyle \frac{1}{\nu _1!}}(R_1)^{\nu _1}{\displaystyle \frac{1}{\nu _2!}}(R_2)^{\nu _2}\mathrm{}{\displaystyle \frac{1}{\nu _k!}}(R_k)^{\nu _k}`$
$`=`$ $`\mathrm{exp}(R_1+R_2+R_3+\mathrm{}).`$
The corresponding generating functional $`W^{}`$ without any Pauli line is simply
$$W^{}W|_{\mathrm{\Delta }_{nm}=0}=\underset{k=1}{\overset{\mathrm{}}{}}R_k.$$
(25)
To include Pauli lines (i.e. $`\mathrm{\Delta }_{nm}0`$), we use notation $`L_k`$ to represents the contribution of all linked $`k`$-clusters and write
$$W=\mathrm{ln}\stackrel{~}{\mathrm{\Psi }}|\mathrm{\Psi }=L_1+L_2+L_3+\mathrm{}.$$
(26)
Using the simple rules discussed earlier, without much difficulty, we can list all $`k`$-cluster contributions of $`L_k`$ in terms of a ring diagram $`R_k`$ plus all possible ways of drawing Pauli lines between any pair of $`k`$ dots of rings, including those pairs of dots between rings and those pairs of dots inside rings. In Fig. 6 we list all 3rd-order contributions in $`L_3`$ except $`R_1,R_2`$ and $`R_3`$.
## IV Diagram resummations, spin-wave theory and beyond
We first consider all diagrams without any Pauli line, namely all the ring diagram contributions $`R_k`$ with $`k=1,2,\mathrm{}`$, and show that the spin-wave theory is thus reproduced. As can be seen from Eq. (20), these ring diagrams represent the first order approximation in the large-$`s`$ limit. In fact, in this limit, operators $`s_i^{}`$ and $`s_j^+`$ behave like bosons as $`s_i^{}\sqrt{2s}a_i^{}`$, $`s_j^+\sqrt{2s}b_j^{}`$ . The corresponding wavefunction by Eq. (8) becomes the spin-wave function as
$$|\mathrm{\Psi }|\mathrm{\Psi }_{sw}=\mathrm{exp}\left(\underset{ij}{}f_{ij}a_i^{}b_j^{}\right)|\mathrm{\Phi }=\underset{q}{}\mathrm{exp}\left(f_qa_q^{}b_q^{}\right)|\mathrm{\Phi },$$
(27)
where the Néel state $`|\mathrm{\Phi }`$ should be considered as the vacuum state for the two sets of bosons $`a_i^{}`$ and $`b_j^{}`$ and where, in the last equation, we have made Fourier transformations using the translational symmetry as,
$`a_i^{}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{N}}}{\displaystyle \underset{𝐪}{}}e^{i𝐪𝐫_i}a_q^{},b_j^{}=\sqrt{{\displaystyle \frac{2}{N}}}{\displaystyle \underset{𝐪}{}}e^{i𝐪𝐫_j}b_q^{}`$
$`f_{ij}`$ $`=`$ $`{\displaystyle \frac{2}{N}}{\displaystyle \underset{𝐪}{}}e^{i𝐪(𝐫_j𝐫_i)}f_q,`$
with summation over $`𝐪`$ restricted to the magnetic zone. The normalization integral of Eq. (27) can be easily calculated as the wavefunction is uncoupled in $`q`$-space. Using expansion $`\mathrm{exp}(f_qa_q^{}b_q^{})=_n(f_qa_q^{}b_q^{})^n/n!`$ and a simple algebra $`\mathrm{\Phi }|a_q^n(a_q^{})^n|\mathrm{\Phi }=n!`$, we have the following well-known result (see, for example, Chapter 2 of Ref. 12),
$$\mathrm{\Psi }_{sw}|\mathrm{\Psi }_{sw}=\underset{q}{}\frac{1}{1\stackrel{~}{f}_qf_q}.$$
(28)
The corresponding generating functional is hence given by,
$$W_{sw}=\mathrm{ln}\mathrm{\Psi }_{sw}|\mathrm{\Psi }_{sw}=\underset{q}{}\mathrm{ln}\left(1\stackrel{~}{f}_qf_q\right)=\underset{q}{}\left[\stackrel{~}{f}_qf_q+\frac{1}{2}(\stackrel{~}{f}_qf_q)^2+\mathrm{}\right],$$
(29)
which is precisely the result of Eq. (25) after Fourier transformation, namely
$$W^{}=W_{sw}.$$
(30)
Distribution functions without any Pauli line can be easily calculated using functional derivatives of Eq. (15)-(16) with diagrammatic representation. For example, the one-body bare distribution function, $`\stackrel{~}{g}_1^{}=\delta W^{}/\delta f_1`$, is simply represented by Fig. 7, where the action of partial derivative is equivalent to unfolding the ring. Writing out the summations explicitly, we have the expansion of Fig. 7 as,
$$\stackrel{~}{g}_1^{}=\stackrel{~}{f}_1+\underset{2}{}\stackrel{~}{f}_{i_1j_2}f_2\stackrel{~}{f}_{i_2j_1}+\underset{2,3}{}\stackrel{~}{f}_{i_1j_2}f_2\stackrel{~}{f}_{i_2j_3}f_3\stackrel{~}{f}_{i_3j_1}+\mathrm{},$$
(31)
and similar expansion for $`g_1^{}`$. A close inspection of $`\stackrel{~}{g}_1^{}`$ and $`g_1^{}`$ expansions yields self-consistency equations as
$$\stackrel{~}{g}_1^{}=\stackrel{~}{f}_1+\underset{2}{}\stackrel{~}{f}_{i_1j_2}g_2^{}\stackrel{~}{f}_{i_2j_1},g_1^{}=f_1+\underset{2}{}f_{i_1j_2}\stackrel{~}{g}_2^{}f_{i_2j_1},$$
(32)
agreed exactly with Eq. (31) of Ref. 7 in this SWT approximation. The two-body functions in this approximation can also be easily obtained in this fashion as given in Ref. 7. The spontaneous magnetization of Eq. (13) is given by $`s_i^z=s\rho ^{}`$, with $`\rho ^{}`$ given by
$$\rho ^{}=\underset{j}{}\rho _{ij}^{}=\underset{q}{}\frac{\stackrel{~}{f}_qf_q}{1\stackrel{~}{f}_qf_q}=\frac{1}{2}\underset{q}{}\left(\frac{1}{\sqrt{1\gamma _q^2/A^2}}1\right),$$
(33)
where we have used the reproduced SWT results of Ref. 7,
$$\stackrel{~}{f}_q=f_q=\frac{A}{\gamma _q}(\sqrt{1\gamma _q^2/A^2}1),\gamma _q=\frac{1}{z}\underset{n}{}e^{i𝐪𝐫_𝐧},$$
(34)
where $`z`$ is the coordination number and $`n`$ is the nearest-neighbour index of the bipartite lattice. For one-dimensional (1D) model at isotropic point $`A=1`$, the integral of Eq. (33) diverges, contrast to the well-known exact result of $`\rho =1/2`$ for $`s=1/2`$ by Bethe ansatz (see Ref. 10 for references).
To go beyond SWT, we need to include Pauli lines. Using the similar resummation technique as discussed above, we express the expansion of bare one-body distribution function $`\stackrel{~}{g}_1`$ in terms of diagrams as shown in Fig. 8, similar to the expansion in Chap. 9 of Ref. 12 and in Ref. 13, after multiplying $`f_1`$ on both sides of the equation,
$$\rho _1=f_1\stackrel{~}{g}_1=f_1\frac{\delta W}{\delta f_1}=\mathrm{Fig}\mathrm{.\hspace{0.33em}8},$$
(35)
where we have done all resummations of ring diagrams as in Eq. (31) and hence all exchange lines in the diagrams of Fig. 8 are now function $`\stackrel{~}{g}_{ij}^{}`$, not the original exchange line function $`\stackrel{~}{f}_{ij}`$. For a simple approximation, we consider the first two diagrams of Fig. 8 as
$$\rho _1\rho _1^{}+\rho _1^{}\underset{2}{}\mathrm{\Delta }_{12}\rho _2^{}.$$
(36)
We notice that, after ignoring the higher-order $`\frac{1}{(2s)^2}`$ term, Eq. (36) (without the common $`f_1`$ factor) agrees with the expression $`s_i^+s_j^{}=s_i^{}s_j^+`$ of SWT in Ref. 9 and with Eq. (31) of our earlier paper Ref. 7 (after changing the sign of both 5th and 6th terms in the equation as they were typos). After summing over index $`j_1`$ with $`\rho =_{j_1}\rho _{i_1j_1}`$, we have, using Eq. (20) for $`\mathrm{\Delta }_{12}`$,
$$\rho =\rho ^{}\frac{2}{2s}(\rho ^{})^2+\frac{1}{(2s)^2}\underset{j}{}\left(\rho _{ij}^{}\right)^2.$$
(37)
For $`s=1/2`$ and isotropic point $`A=1`$, we obtain $`\rho 0.127`$ for the square lattice and $`0.067`$ for the cubic lattice. They are smaller than $`\rho ^{}=0.197`$ and $`0.078`$ of SWT respectively. This is not surprising because SWT is known to have over estimated the quantum fluctuations. The best numerical values for the square lattice vary from $`\rho =0.16`$ to $`0.19`$, including results from extrapolation of high-order localized CCM calculations . For the 1D isotropic model, however, $`\rho `$ of Eq. (37) diverges as $`\rho ^{}`$ diverges as mentioned earlier.
We next consider an approximation involving higher-order Pauli lines by including all higher-order diagrams similar to that of Eq. (36), as shown in Fig. 9. This infinite series can again be resummed in a closed form as a self-consistency equation, equivalent to replacing $`\rho _2^{}`$ in Eq. (36) by $`\rho _2`$ itself as
$$\rho _1=\rho _1^{}+\rho _1^{}\underset{2}{}\mathrm{\Delta }_{12}\rho _2.$$
(38)
The resummation in Eq. (38) is similar to the resummation of rings in Eqs. (31)-(32), we therefore refer it as super-ring resummation. The numerical results for $`\rho `$ thus obtained at the isotropic point for high dimensions improve slightly, as $`\rho =0.145`$ for the square lattice and $`0.068`$ for the cubic lattice. However, for the 1D isotropic model, Eq. (38) produces a convergent, precise number $`\rho =1/2`$. This is interesting indeed, as the divergence of SWT has troubled theorists for many years. It is worth mentioning that the traditional CCM SUB2 approximation also produced a convergent result for the 1D model but at $`A=0.373`$, not at the isotropic point $`A=1`$. We leave more discussion to the following section, and leave detailed calculations including other higher-order terms and resummations in two-body function $`\rho _{12}`$ and structure function $`S_{12}`$ of Eq. (16) somewhere else.
## V Discussion
In this article, we present a diagrammatic scheme for the calculations of distribution functions of the variational CCM, as an alternative to the algebraic scheme published in our earlier papers . The results of SWT are reproduced by an approximation which resums all ring diagrams without any Pauli line. Approximations beyond SWT can also easily be made by including diagrams with Pauli lines. One such approximation, which includes all super-ring diagrams by a resummation of infinite Pauli lines in addition to resummations of all ring diagrams, produces a convergent, precise number for the order parameter of the 1D isotropic model, contrast to the divergence of SWT. This cure of SWT divergence is also interesting to 2D models (including square and triangle lattices) as naive higher-order calculations within the framework of SWT are also likely to produce divergent results, despite the fact that the first order results are reasonable. We believe that similar resummations of super-ring diagrams as Fig. 9 and Eq. (38) may provide a solution for such divergent problems. We leave more detailed calculations to somewhere else.
It is also possible to include in the ground state higher-order many-body correlations such as 4-spin-flip operators, in additional to the 2-spin-flip operators of Eqs. (8). Furthermore, as demonstrated here by the diagrammatic approach, a direct link between our variational CCM and the powerful CBF has now been established, as both rely on determination of distribution functions through functional derivatives of a generating functional. In particular, as given by Eq. (13), particle density $`\rho `$ in CBF is equivalent to the order parameter of our spin models as $`s_i^z=s\rho `$. Its diagrammatic expansions in two theories are similar (see Chap. 9 of Ref. 12 and Ref. 13 for more CBF details). For 2D and 3D lattice models, the values of density $`\rho `$ are small compared with $`s`$. Such spin systems can therefore be described as dilute gases (dilute gases of quasiparticle magnons of spin waves). For the isotropic 1D model, density $`\rho `$ is saturated, corresponding to the order parameter equal to zero, a critical value. Our approximation including a resummation of super-ring diagrams is capable of reproducing precisely such number. It is also interesting to know that our diagrammatic analysis of the variational CCM is for the translational invariance lattice system while similar analysis in CBF is for inhomogeneous systems .
Clearly, for more accurate results in general, we need to include correlations between those quasiparticles in our ground state, and the CBF is well known to be one of most effective theories for dealing with such particle correlations (even when they are very strong as in a Helium-4 quantum liquid ) by systematic calculations of the important two-body distribution functions. We therefore propose a unified trial wavefunction $`|\mathrm{\Psi }_U`$ as, including a generalized Jastrow correlation operator $`S^0`$ involving quasiparticle density operator $`s^z`$,
$$|\mathrm{\Psi }_U=e^{S^0/2}|\mathrm{\Psi },S^0=\underset{ij}{}f_{ij}^0s_i^zs_j^z,$$
(39)
where $`\{f_{ij}^0\}`$ are the new additional variational parameters and $`|\mathrm{\Psi }`$ is our variational CCM state of Eq. (2). The diagrammatic scheme as discussed in this article is useful for calculating the expansion of the new generating functional of Eq. (39). We have made progress in such calculations and wish to report results soon. We also believe such a unified many-body theory may prove to be capable of dealing with strongly correlated fermion systems in general.
###### Acknowledgements.
I am grateful to R.F. Bishop for introducing the CCM to me. Useful discussions with J. Arponen, R.F. Bishop, F. Coester, and H. Kümmel are also acknowledged.
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# Amenability, Tubularity, and Embeddings into ℛ^𝜔
## 1. Introduction
One version of Voiculescu’s free entropy involves matricial microstates. Given an $`n`$-tuple of selfadjoint elements $`X=\{x_1,\mathrm{},x_n\}`$ in a von Neumann algebra $`M`$ with trace $`\phi `$, denote by $`\mathrm{\Gamma }(X;m,k,\gamma )`$ the set of all $`n`$-tuples of selfadjoint $`k\times k`$ matrices $`(a_1,\mathrm{},a_n)`$ such that for any $`1pm`$ and $`1i_1,\mathrm{},i_pn,`$
$$|tr_k(a_{i_1}\mathrm{}a_{i_p})\phi (x_{i_1}\mathrm{}x_{i_p})|<\gamma $$
where $`tr_k`$ denotes the normalized trace on the $`k\times k`$ matrices. Such an $`n`$-tuple $`(a_1,\mathrm{},a_n)`$ is called a microstate for $`X`$ and the sets $`\mathrm{\Gamma }(X;m,k,\gamma )`$ are called the microstate spaces of $`X`$.
One can use microstates and random matrices to obtain nonisomorphism results (for an overview see ). At some point all of these kinds of arguments rely on the fact that amenability (=hyperfiniteness=injectivity by ; see also - and ) forces an extremely rigid condition on the microstate spaces. When $`X^{\prime \prime }`$ is amenable then the following is true: given $`ϵ>0`$ there exist $`m`$ and $`\gamma >0`$ such that for any two elements $`\xi =(\xi _1,\mathrm{},\xi _n),\eta =(\eta _1,\mathrm{},\eta _n)\mathrm{\Gamma }(X;m,k,\gamma )`$ there exists a $`k\times k`$ unitary $`u`$ such that for each $`1in,`$ $`|u\xi _iu^{}\eta _i|_2<ϵ.`$ Roughly speaking then, a microstate space of $`X`$ is the neighborhood of the unitary orbit of a single microstate of the space.
In this note we observe the converse. To be more exact, suppose $`X^{\prime \prime }`$ embeds into $`^\omega `$ and satisfies the aforementioned geometric property: for any $`ϵ>0`$ there exist $`m`$ and $`\gamma >0`$ such that if $`\xi =(\xi _1,\mathrm{},\xi _n),\eta =(\eta _1,\mathrm{},\eta _n)\mathrm{\Gamma }(X;m,k,\gamma ),`$ then there exists a $`k\times k`$ unitary $`u`$ satisfying $`|u\xi _iu^{}\eta _i|_2<ϵ,1in.`$ We show that $`X^{\prime \prime }`$ is semidiscrete in the sense of and thus, by amenable.
We also notice that this characterization is equivalent to the condition that any two embeddings $`\sigma ,\pi `$ of $`X^{\prime \prime }`$ into $`^\omega `$ are conjugate by a unitary $`u^\omega `$. In the course of the proof we will introduce seemingly weaker notions of tubularity: finite tubularity, quasitubularity, and finite quasitubularity. All of these notions will be shown to coincide with amenability as well.
## 2. Some other characterizations of Amenability
Throughout $`M`$ will be a von Neumann algebra with tracial, faithful, normal state $`\phi `$ and $`X=\{x_1,\mathrm{},x_n\}`$ is a finite generating set of selfadjoints of $`M`$. Fix a nontrivial ultrafilter $`\omega `$ of $``$; there exists an obvious net $`i:\mathrm{\Lambda }`$ such that for any bounded sequence of complex numbers, $`c_n_{n=1}^{\mathrm{}}`$, $`\omega (c_n_{n=1}^{\mathrm{}})=lim_{\lambda \mathrm{\Lambda }}c_{i(\lambda )}.`$ $``$ will be a fixed copy of the hyperfinite $`\mathrm{II}_1`$-factor, and $`^\omega `$ will be the associated ultraproduct construction with respect to $`\omega `$ and $``$. $`Q:\mathrm{}^{\mathrm{}}()^\omega `$ is the obvious quotient map. By ”embedding” we mean a normal, injective $``$-homorphism which preserves units and the traces.
We will be constantly working with finite tuples of operators and for this reason introduce the following, somewhat abusive notation. For an ordered $`n`$-tuple $`\xi =\{\xi _1,\mathrm{},\xi _n\}`$ in a von Neumann algebra $`N`$ and a unitary $`uN`$ $`u\xi u^{}=\{u\xi _1u^{},\mathrm{},u\xi _nu^{}\}`$. If $`F`$ is a map from $`N`$ into some set $`S`$, then we carelessly write $`F(\xi )`$ for $`\{F(\xi _1),\mathrm{},F(\xi _n)\}`$. If $`\eta =\{\eta _1,\mathrm{},\eta _n\}N`$, then $`\xi \eta =\{\xi _1\eta _1,\mathrm{},\xi _n\eta _n\}`$ and $`|\xi |_2=(_{i=1}^n\tau (\xi _i^2))^{\frac{1}{2}}`$ where $`\tau `$ is a tracial state on $`N`$. We will assume throughout that $`M`$ embeds into $`^\omega `$. For $`k,`$ $`M_k()`$ denotes the $`k\times k`$ matrices and $`U_k`$ is the unitary group of $`M_k().`$
###### Definition 2.1.
$`X`$ is $`N`$-tubular ($`N`$) if for any $`ϵ>0`$ there exist $`m`$ and $`\gamma >0`$, such that if $`\xi _1,\mathrm{},\xi _{N+1}\mathrm{\Gamma }(X;m,k,\gamma )`$, then there exists a $`uU_k`$ satisfying $`|u\xi _iu^{}\xi _j|_2<ϵ`$ for some $`1i<jN+1`$. $`X`$ is finitely tubular if $`X`$ is $`N`$-tubular for some $`N`$. $`X`$ is simply tubular if $`X`$ is $`1`$-tubular.
Thus, $`X`$ is finitely tubular if it can be encapsulated in the unitary orbits of no more than $`N`$ of its microstates. Clearly, tubularity coincides with the definition given in the introduction. We have a similar notion for quasitubularity:
###### Definition 2.2.
$`X`$ is $`N`$-quasitubular ($`N`$) if for any $`ϵ>0`$ there exist $`m`$ and $`\gamma >0`$, such that for any $`\xi _1,\mathrm{},\xi _{N+1}\mathrm{\Gamma }(X;m,k,\gamma )`$ there exists a $`p`$ (dependent on $`\xi _1,\mathrm{},\xi _{N+1}`$) and a unitary $`u`$ in $`M_k()M_p()`$ satisfying $`|u(\xi _iI_p)u^{}\xi _jI_p|_2<ϵ`$ for some $`1i<jN+1`$. $`X`$ is finitely quasitubular if $`X`$ is $`N`$-quasitubular for some $`N`$. $`X`$ is simply quasitubular if $`X`$ is $`1`$-quasitubular.
###### Remark 2.3.
Obviously if $`X`$ is $`N`$-tubular, then $`X`$ is $`N`$-quasitubular.
###### Lemma 2.4.
If $`X`$ is tubular, then any two embeddings $`\sigma ,\pi `$ of $`X^{\prime \prime }`$ into $`^\omega `$ are conjugate by a unitary in $`^\omega `$.
###### Proof.
Suppose $`X`$ is tubular and $`\sigma ,\pi :X^{\prime \prime }^\omega `$ are two embeddings. Find algebras $`A_k`$ such that for each $`k`$, $`A_kM_{r(k)}()`$ for some $`r(k)`$ and such that for each $`1jn`$ there exist sequences $`y_j=y_{jk}_{k=1}^{\mathrm{}}`$, $`z_j=z_{jk}_{k=1}^{\mathrm{}}\mathrm{}^{\mathrm{}}()`$ satisfying $`y_{jk},z_{jk}A_k`$ for each $`k`$, $`\pi (x_j)=Q(y_j),`$ and $`\sigma (x_j)=Q(z_j).`$
For each $`p`$ there exists by tubularity a corresponding $`m(p)`$ such that for any $`k`$, if $`\xi ,\eta \mathrm{\Gamma }(F;m(p),k,m(p)^1),`$ then there exists a $`uU_k`$ satisfying $`|u\xi u^{}\eta |_2<p^1.`$ For each $`k`$ pick a unitary $`w_kA_kM_{r(k)}()`$ satisfying
$$\underset{1jn}{\mathrm{max}}|w_ky_{jk}w_k^{}z_{jk}|_2=\underset{wU(A_k)}{inf}\left(\underset{1jn}{\mathrm{max}}|wy_{jk}w^{}z_{jk}|_2\right)$$
where $`U(A_k)`$ ($`U_k`$) is the unitary group of $`A_k`$. Define $`w=w_k_{k=1}^{\mathrm{}}`$; $`w\mathrm{}^{\mathrm{}}(M).`$ Given an $`N`$, for $`\lambda `$ sufficiently large $`(y_{1i(\lambda )},\mathrm{},y_{ni(\lambda )}),(z_{1i(\lambda )},\mathrm{},z_{ni(\lambda )})\mathrm{\Gamma }(F;m(N),i(\lambda ),m(N)^1)`$ which in turn implies that for such $`\lambda `$ and all $`1jn`$
$$|w_{i(\lambda )}y_{i(\lambda )j}w_{i(\lambda )}^{}z_{i(\lambda )j}|_2<N^1.$$
Set $`u=Q(w)^\omega `$. We have that for all $`1jn,`$ $`\sigma (x_j)=u\pi (x_j)u^{}`$ which completes the proof. ∎
###### Lemma 2.5.
If any two embeddings $`\sigma `$ and $`\pi `$ of $`M`$ into $`^\omega `$ are conjugate by a unitary in $`^\omega `$, then $`X`$ is quasitubular.
###### Proof.
Suppose by contradiction that $`X`$ is not quasitubular. For some $`ϵ_0>0`$ and any $`m`$ and $`\gamma >0`$ there exists an $`N`$ and $`\xi ,\eta \mathrm{\Gamma }(X;m,N,\gamma )`$ such that for all $`p,`$
$$\underset{uU_{Np}}{inf}|u(\xi I_p)u^{}\eta I_p|_2>ϵ_0.$$
Thus, for each $`m`$ we can find a corresponding $`N_m`$ and $`\xi _m,\eta _m\mathrm{\Gamma }(X;m,N_m,m^1)`$ such that for any $`k`$ $`inf_{uU_{kN_m}}|u(\xi _mI_k)u^{}\eta _mI_k|_2>ϵ_0.`$ Without loss of generality we may assume that the operator norms of any of the elements in $`\xi _m`$ or $`\eta _m`$ are strictly less than $`C=\mathrm{max}_{xX}x+1`$. For each $`m`$, $`=M_{N_m}()_m`$ where $`_m`$; define $`x_m=\xi _mI`$ and $`y_m=\eta _mI`$ with respect to this decomposition of $``$ and set $`x=x_m_{m=1}^{\mathrm{}}`$ and $`y=y_m_{m=1}^{\mathrm{}}`$. It is not too hard to see that we can find two embeddings $`\pi ,\sigma :M^\omega `$ satisfying $`\pi (X)=Q(x)`$ and $`\sigma (X)=Q(y)`$.
By hypothesis there exists a unitary $`u^\omega `$ satisfying $`\sigma (x)=u\pi (x)u^{}`$ for all $`xM`$. We can find some $`w=w_m_{m=1}^{\mathrm{}}\mathrm{}^{\mathrm{}}()`$ such that $`Q(w)=u`$. Because $`u`$ is a unitary we can assume that for each $`m`$, $`w_m`$ is a unitary. The condition $`\sigma (x)=u\pi (x)u^{}`$ implies that there exists a $`\lambda \mathrm{\Lambda }`$ such that
$$|w_{i(\lambda )}(\xi _{i(\lambda )}I)w_{i(\lambda )}^{}\eta _{i(\lambda )}I|_2=|w_{i(\lambda )}x_{i(\lambda )}w_{i(\lambda )}^{}y_{i(\lambda )}|_2<ϵ_0/3C.$$
Now, $`w_{i(\lambda )}=M_{N_{i(\lambda )}}()R_{i(\lambda )}`$ and by standard approximations we can find some $`p_0`$, a unital $``$-algebra $`A_\lambda M_{p_0}()`$ with $`A_\lambda _{i(\lambda )}`$, and a unitary $`uM_{N_{i(\lambda )}}A_\lambda `$ satisfying $`|uw_{i(\lambda )}|_2<ϵ_0/3C.`$ It is then clear that
$$\underset{vU_{i(\lambda )p_0}}{inf}|v(\xi _{i(\lambda )}I_{p_0})v^{}\eta _{i(\lambda )}I_{p_0}|_2|u(\xi _{i(\lambda )}I_{p_0})u^{}\eta _{i(\lambda )}I_{p_0}|_2<ϵ_0.$$
which contradicts the initial hypothesis. ∎
We now present a lemma which is undoubtedly known but which we will prove for completeness. Recall that a von Neumann algebra $`N`$ is semidiscrete if there exist nets $`\varphi _j_{j\mathrm{\Omega }}`$, and $`\psi _j_{j\mathrm{\Omega }}`$ of unital completely positive maps $`\varphi _j:NM_{n_j}()`$, $`\psi _j:M_{n_j}()N`$, $`n_j`$ such that for any $`xN`$, $`(\varphi _j\psi _j)(x)x`$ $`\sigma `$-weakly. This is not the original definition of semidiscreteness found in (which demands that the maps only have finite rank), but it is equivalent to that definition by and . Semidiscreteness, which was introduced in , is yet another characterization of amenability ( again).
###### Lemma 2.6.
Suppose $`A`$ is a tracial von Neumann algebra. Assume that for some finite set of generators $`F`$ of $`A`$ and any $`ϵ>0`$ there exist an embedding $`\pi _ϵ`$ of $`A`$ into a tracial von Neumann algebra $`A_ϵ`$ and a unital, finite dimensional algebra $`B_ϵA_ϵ`$ with the property that every element of $`\pi _ϵ(F)`$ is contained in the $`ϵ`$-neighborhood of $`B_ϵ`$ with respect to the $`||_2`$-norm of $`A_ϵ`$. Then $`A`$ is semidiscrete, and thus amenable.
###### Proof.
The hypothesis implies that for any finite set $`S`$ of $`A`$ and any $`ϵ>0`$ there exists an embedding $`\pi _ϵ`$ of $`A`$ into a tracial von Neumann algebra $`A_ϵ`$ and a finite dimensional unital subalgebra $`B_ϵ`$ of $`A_ϵ`$ such that every element of $`\pi _ϵ(S)`$ is contained in the $`ϵ`$-neighborhood of $`B_ϵ`$ with respect to $`||_2`$. This is because the elements of $`\pi _ϵ(F)`$ can be approximated in $`||_2`$-norm by elements in $`B_ϵ`$ with operator norms no bigger than the maximum of the operator norms of elements in $`F`$ (one uses conditional expectations as we do below) and because multiplication is $`||_2`$-continuous on operator norm bounded sets. This gives the implication for $`S`$ consisting of polynomials of elements from $`F`$ and the general case follows immediately.
To complete the proof it suffices to show that for any finite $`SA`$ we can construct sequences of u.c.p. maps $`\varphi _n:AM_{k_n}()`$ and $`\psi _n:M_{k_n}()A`$, $`k_n`$, such that for any normal linear functional $`f`$ on $`N,`$ $`f((\varphi _n\psi _n)(x))f(x)`$ as $`n\mathrm{}`$. Thus, let the finite subset $`S`$ of $`A`$ be given. By the first paragraph for each $`n`$ we can find a tracial von Neumann algebra $`A_n`$, an embedding $`\pi _n:AA_n`$ and a finite dimensional, unital subalgebra $`B_n`$ of $`A_n`$ such that every element of $`S`$ is contained in the $`n^1`$-neighborhood of $`B_n`$ with respect to the $`||_2`$-norm of $`A_n`$. Define $`\varphi _n:AB_n`$ to be the composition of $`\pi _n`$ with the conditional expectation $`E_n`$ of $`A_n`$ onto $`B_n`$. Define $`F_n`$ to be the conditional expection of $`A_n`$ onto $`\pi _n(A)`$ restricted to $`B_n`$; so $`F_n:B_n\pi _n(A)`$. Define $`\psi _n:B_nA`$ to be the composition of $`F_n`$ with $`\pi _n^1`$. Obviously for any $`xS`$
$$|(\psi _n\varphi _n)(x)x|_2=|F_n(E_n(\pi _n(x)))\pi _n(x)|_2|E_n(\pi _n(x))\pi _n(x)|_2<n^1.$$
Now $`(\psi _n\varphi _n)(x)_{n=1}^{\mathrm{}}`$ is a sequence in $`A`$ uniformly bounded in the operator norm and $`(\psi _n\varphi _n)(x)x`$ in the $`||_2`$-norm. This implies $`(\psi _n\varphi _n)(x)x`$ $`\sigma `$-weakly for every $`xS.`$
###### Remark 2.7.
Suppose $`N`$ is a von Neumann algebra, and $`AN`$ is a von Neumann subalgebra. By the above lemma if for some finite set of generators $`F`$ of $`A`$ and any $`ϵ>0`$ there exists a finite dimensional algebra $`BN`$ such that every element of $`F`$ is contained in the $`ϵ`$-neighborhood of $`B`$ with respect to $`||_2`$, then $`A`$ is amenable.
###### Lemma 2.8.
If $`X`$ is quasitubular, then $`M`$ is amenable.
###### Proof.
For each $`m`$ there exist a $`k(m)`$ and a $`\xi _m=(\xi _{1m},\mathrm{},\xi _{nm})\mathrm{\Gamma }(X;m,k(m),m^1)`$. $`=_1_2_3`$ where $`_i`$, $`i=1,2,3.`$ Define for $`1jn`$, $`y_j=I\xi _{jm}I_{m=1}^{\mathrm{}}\mathrm{}^{\mathrm{}}()`$ where $`I\xi _{jm}I_1_m_3_1_2_3=`$ and $`M_{k(m)}()B_mR_2.`$ It is clear that there exists a (trace preserving) embedding $`\pi :X^{\prime \prime }R^\omega `$ satisfying $`\pi (x_j)=Q(y_j)`$. Given $`ϵ>0`$ we can find by quasitubularity an $`m_0`$ such that for any $`k`$ and $`\xi ,\eta \mathrm{\Gamma }(X;m_0,k,m_0^1)`$ there exists a corresponding $`p(k)`$ and unitary $`u`$ of $`M_k()M_{p(k)}()`$ satisfying $`|u(\xi I_{p(k)})u^{}\eta I_{p(k)}|_2<ϵ`$. Fix once and for all, an $`N`$ such that there exists a $`\zeta \mathrm{\Gamma }(X;m_0,N,m_0^1).`$ We can regard $`\zeta `$ as a subset of $`_1`$ by finding a copy $`A`$ of the $`N\times N`$ matrices in $`_1`$ and for each $`m`$ consider $`\zeta II,I\xi I(A_m)_3`$. For each $`m`$ we can find an algebra $`D_m`$ isomorphic to a full matrix algebra, and a unitary $`u_m`$ of $`A_mD_m`$ such that
$$|u_m(\zeta II)u_m^{}I\xi I|_2<ϵ.$$
Define $`\rho :A^\omega `$ by $`\rho (x)=Qu_m(xII)u_m^{}_{m=1}^{\mathrm{}}`$. We have now shown that for $`ϵ>0`$, every element of $`\pi (X)`$ is within the $`||_2`$ $`ϵ`$-ball of the finite dimensional full matrix algebra $`\rho (A)`$.
By the remark this implies that the von Neumann algebra generated by $`\pi (X)`$ is semidiscrete, and hence, amenable. Since $`\pi `$ is an isomorphism, $`X^{\prime \prime }=M`$ is amenable. ∎
At this point we have already demonstrated the equivalence claimed in the introduction. We will now prove $`ϵ`$ more. Notice that the content of the observation below is the implication that finite quasitubularity implies amenability. This will be very much like the proof above modulo some technicalities. We could have gone straight to the following more technical argument without proving Lemma 2.8, but the proof of Lemma 2.8 has the advantage of being more lucid.
###### Lemma 2.9.
The following are equivalent:
* (1) $`M`$ is amenable.
* (2) $`X`$ is tubular.
* (3) $`X`$ is finitely tubular.
* (4) $`X`$ is quasitubular.
* (5) $`X`$ is finitely quasitubular.
* (6) Any two embeddings of $`M`$ into $`^\omega `$ are conjugate by a unitary in $`^\omega `$.
###### Proof.
By $`M`$ is amenable iff $`M`$ is hyperfinite and thus by , $`X`$ must be tubular. Hence, (1) $``$ (2). By Lemma 2.4, (2) $``$ (6) and by Lemma 2.5 (6) $``$ (4). Clearly (4) $``$ (5).
I will now show (5) $``$ (1). Suppose $`X`$ is finitely quasitubular. Find the smallest $`N`$ for which $`X`$ is $`N`$-quasitubular. By Lemma 2.6 we can assume $`N>1`$. $`X`$ is not $`(N1)`$-quasitubular, which implies that we can find some $`ϵ_1>0`$, such that for any $`m`$ there exists some $`k`$, and $`\xi _1,\mathrm{},\xi _N\mathrm{\Gamma }(X;m,k,m^1)`$ such that for any $`p`$ and unitary $`u`$ of $`M_k()M_p()`$,
$$\underset{1i<jN}{\mathrm{min}}|u(\xi _iI_p)u^{}\xi _jI_p|_2>ϵ_1.$$
Let $`ϵ_1>ϵ>0`$ be given. There exists an $`m_0`$ such that if $`\xi _1,\mathrm{},\xi _{N+1}\mathrm{\Gamma }(X;m_0,k,m_0^1)`$ then there exists some $`p`$ and unitary $`u`$ in $`M_{k_0}()M_p()`$ ($`u`$ and $`p`$ dependent on the $`\xi _i`$) satisfying $`|u(\xi _iI_p)u^{}\eta |_2<ϵ`$ for some $`1i<jN+1`$. On the other hand, by our initial remark, there is some $`k(0)`$, and $`\eta _1,\mathrm{},\eta _N\mathrm{\Gamma }(X;m_0,k(0),m_0^1)`$ such that for any $`p`$ and unitary $`u`$ of $`M_{k(0)}()M_p()`$,
(1) $`\underset{1i<jN}{\mathrm{min}}|u(\xi _iI_p)u^{}\xi _jI_p|_2>ϵ_1>ϵ.`$
Now there exists a sequence $`\zeta _m_{m=1}^{\mathrm{}}`$ such that for each $`m`$ there exists a $`k(m)`$ with $`\zeta _m\mathrm{\Gamma }(X;m,k(m),m^1)`$. Identifying $`M_{k(0)k(m)}()`$ with $`M_{k(0)}()M_{k(m)}()`$ for all $`mm_0`$ we have that $`\xi _1I_{k(m)},\mathrm{},\xi _NI_{k(m)},I_{k(0)}\zeta _m\mathrm{\Gamma }(X;m_0,k(0)k(m),m_0^1)`$. (1) and the $`N`$-quasitubularity of $`X`$ implies that there must exist an $`1i_mN`$, a $`p_m`$, and a unitary $`u_m`$ of $`M_{k(0)k(m)}()M_{p_m}()`$ satisfying
(2) $`|u_m(\xi _{i_m}I_{p_m})u_m^{}\zeta _mI_{p_m}|_2<ϵ.`$
$`i_m_{m=1}^{\mathrm{}}`$ is a sequence taking integral values between $`1`$ and $`N`$ and thus $`i_m=i`$ for some $`1iN`$ and infinitely many $`m`$. Without loss of generality assume $`i_{m_q}=1`$ for some increasing sequence $`m_q_{q=1}^{\mathrm{}}`$ of $``$.
For each $`q`$ set $`\theta _q=\zeta _{m_q}\mathrm{\Gamma }(X;m_q,k(m_q),m_q^1)`$. $`=_1_2_3`$ where $`_i`$, $`i=1,2,3.`$ Define $`Y=I\theta _qI_{q=1}^{\mathrm{}}(\mathrm{}^{\mathrm{}}())^n`$ where $`I\theta _qI_1_q_3_1_2_3=`$ and $`M_{k(m_q)}()B_qR_2.`$ Because $`lim_q\mathrm{}m_q=\mathrm{}`$, it is clear that there exists a (trace preserving) embedding $`\pi :X^{\prime \prime }R^\omega `$ satisfying $`\pi (X)=Q(Y)`$. We can regard $`\xi _1`$ as a subset of $`_1`$ by fixing a copy $`A`$ of the $`k(0)\times k(0)`$ matrices in $`_1`$ and for each $`q`$ consider $`\xi _1II,I\theta _qI(A_m)_3`$. For each $`q`$ (2) provides an algebra $`D_q`$ isomorphic to a full matrix algebra, and a unitary $`v_q`$ of $`A_qD_q`$ such that
$$|v_q(\xi _1II)v_q^{}I\theta _qI|_2<ϵ.$$
Define $`\rho :A^\omega `$ by $`\rho (x)=Qv_q(xII)v_q^{}_{q=1}^{\mathrm{}}`$. We have shown that every element of $`\pi (X)`$ is within the $`||_2`$ $`ϵ`$-ball of the finite dimensional full matrix algebra $`\rho (A)`$. Thus, for every $`ϵ>0`$ there exists an isomorphic copy of $`M`$ in $`^\omega `$ such that $`X`$ identified in $`^\omega `$ with respect to this embedding is within the $`||_2`$-$`ϵ`$ ball of a finite dimensional subalgebra of $`^\omega `$. It follows that $`M`$ is amenable by Lemma 2.6. (1) follows.
We now have the equivalence of conditions (1), (2), (4), (5) and (6). To conclude, (3) $``$ (5) = (2) $``$ (3) so condition (3) is equivalent to all the other conditions as well. ∎
###### Remark 2.10.
It should be fairly obvious to the reader by now that conditions (1)-(6) of Lemma 2.9 are equivalent to the condition that there exist finitely many embeddings $`\pi _1,\mathrm{},\pi _n`$ of $`M`$ into $`^\omega `$ such that for any other embedding $`\sigma `$ of $`M`$ into $`^\omega `$, there exists a $`1in`$ and a unitary $`u`$ in $`^\omega `$ such that $`\sigma `$ is conjugate to $`\pi _i`$ via $`u`$.
Acknowledgements. I would like to thank Nate Brown, Ed Effros, Sorin Popa, and Dimitri Shlyakhtenko for useful conversations. I am especially grateful to Nate for simplifying my original argument and to Ed for the discussions on injectivity.
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# Bessel beam propagation: Energy localization and velocity
## Abstract
The propagation of a Bessel beam (or Bessel-X wave) is analyzed on the basis of a vectorial treatment. The electric and magnetic fields are obtained by considering a realistic situation able to generate that kind of scalar field. Specifically, we analyze the field due to a ring-shaped aperture over a metallic screen on which a linearly polarized plane wave impinges. On this basis, and in the far field approximation, we can obtain information about the propagation of energy flux and the velocity of the energy.
The motion of a Bessel beam is of great interest in physics both for its characteristic as a non-diffracting beam dur87 ; dur87-1 ; zio ; spr ; dur91 ; tan ; rei , and for its implications with regard to the topic of superluminality saa ; mug00 ; ale ; bes ; zam ; saa1 . Extended studies have been devoted to these subjects from both an experimental and a theoretical point of view. However, in spite of the many efforts devoted to this topic, no definite answer has been found about the amount of the energy transfer and its velocity.
As far as Bessel beam propagation is concerned, the problem is mainly related to the difficulty in finding the vectorial field that describes this system, while, on the contrary, the field in the scalar approximation is well-known.
The purpose of the present work is to investigate the propagation of a Bessel beam on the basis of a vectorial treatment. This kind of approach allows us to obtain information regarding the propagation of the energy flux and the energy mean velocity.
A Bessel beam consists of a set of plane waves with directions of propagation s = $`\alpha _1`$i\+ $`\beta _1`$j \+ $`\gamma _1`$k which makes the same angle $`\theta _0(0\theta _0<\pi /2`$) with the $`z`$-axis (hence $`\gamma _1=\mathrm{cos}\theta _0`$ for all the plane waves). In spherical coordinates ($`\rho ,\theta ,\phi `$), the direction of propagation is specified by
$`\alpha _1=\mathrm{sin}\theta _0\mathrm{cos}\phi ,\beta _1=\mathrm{sin}\theta _0\mathrm{sin}\phi \gamma _1=\mathrm{cos}\theta _0.`$ (1)
Thus, for propagation in vacuum or air, each one of these waves, at the point $`x,y,z,`$ can be written as
$$u(P)=u_0d\phi \mathrm{exp}[ik_0(\alpha _1x+\beta _1y+\gamma _1z)]\mathrm{exp}(i\omega t),$$
(2)
where $`u_0d\phi `$ is the amplitude of the elementary wave, $`\omega `$ is the angular frequency, $`k_0=\omega /c`$ is the wavenumber, and $`x,y`$ and $`z`$ denotes the Cartesian coordinates of $`P`$. In cylindrical coordinates $`\rho ,\psi ,z`$ around the $`z`$-axis
$`x=\rho \mathrm{cos}\psi ,y=\rho \mathrm{sin}\psi ,zz,`$
and the total field $`U`$, given by the superposition of all the waves (2), can be obtained by integrating over $`d\phi `$, that is,
$`U=`$ $`u_0`$ $`{\displaystyle _0^{2\pi }}\mathrm{exp}[ik_0(\alpha _1x+\beta _1y+\gamma _1z)]\mathrm{exp}(i\omega t)𝑑\phi `$
$`=`$ $`2\pi u_0J_0(k_0\rho \mathrm{sin}\theta _0)\mathrm{exp}\left(i\omega {\displaystyle \frac{z}{c}}\mathrm{cos}\theta _0\right)\mathrm{exp}(i\omega t)`$ (3)
where $`J_0`$ denotes the zero-order Bessel function of first kindwat .
The scalar field of Eq. (3) is known as a Bessel beam (or Bessel-X wave), the unusual features of which are
\- that it does not change its shape during propagation, since its amplitude is independent of $`z`$note ;
\- that it propagates in the $`z`$ direction with phase and group velocities $`v=c/\mathrm{cos}\theta _0`$ larger than $`c`$saa ; mug00 ; ale ; note1 .
Both the above mentioned characteristics can be analyzed in detail by means of a vectorial treatment, since the scalar field (3) represents an approximation of an electromagnetic field, and only a knowledge of the vectorial field (and of the Poynting vector in particular) can provide detailed information about the energy propagation.
Vectorial fields with amplitude proportional to the zero-order Bessel function can be found in different ways. However, in order to derive just the vectorial field describing a system which has Eq. (3) as scalar approximation, we have to consider a realistic situation that is able to generate a field of that kind. For this purpose, let us consider the system of Fig. 1, which consists of a ring-shaped aperture, of radius $`r`$, over a metallic screen on which a linearly polarized plane wave impinges (from the left). The ring is placed on the focal plane of the converging system $`C`$ with focal length $`f\lambda `$, $`\lambda `$ being the wavelength. Let us consider the impinging electric filed to be polarized in the i direction, and the thickness $`d`$ of the ring to be very small with respect to $`\lambda `$. We may assume the element $`dA`$ of the ring, at $`S(r,\phi ,f)`$, to behave as an elementary dipole, parallel to i, with amplitude proportional to $`d\phi `$ (this requires a suitable choice of the transparency of the ring at $`S`$). Thus, the associated field at the optical center $`O`$ of $`C`$ has the well-known characteristics of the far field radiated by an elementary dipole, that is to be a spherical wave centered at $`S`$, with the electric field e in the meridional plane of the dipole through $`O`$, and perpendicular to the direction R from $`S`$ to $`O`$. We can write
$`𝐑`$ $`=r\mathrm{cos}\phi 𝐢r\mathrm{sin}\phi 𝐣+f𝐤,`$ (4)
$`𝐞`$ $`=e_x𝐢+e_y𝐣+e_z𝐤.`$
The two characteristics mentioned above can be written as
$`𝐞𝐑=r\mathrm{cos}\phi e_xr\mathrm{sin}\phi e_y+fe_z=0`$
$`(𝐑\times 𝐞)𝐢=fe_y+r\mathrm{sin}\phi e_z=0,`$ (5)
where we disregarded a phase factor $`\mathrm{exp}(ikR)`$. By solving Eqs. (5) we obtain
$`e_y=r\mathrm{cos}\phi e_x\left({\displaystyle \frac{r\mathrm{sin}\phi }{r^2\mathrm{sin}^2\phi +f^2}}\right)`$ (6)
$`e_z=r\mathrm{cos}\phi e_x\left({\displaystyle \frac{f}{r^2\mathrm{sin}^2\phi +f^2}}\right).`$ (7)
At this point, we can assume that the field emerging from the optical converging system is a plane wave propagating in the direction of R, with amplitude proportional to the amplitude at $`C`$ of the incident field, namely note2
$`𝐞`$ $`=`$ $`\left(e_x𝐢+e_y𝐣+e_z𝐤\right)\mathrm{exp}[ik(\alpha x+\beta y+\gamma z)]`$ (8)
$`𝐡`$ $`=`$ $`{\displaystyle \frac{1}{Z}}\left[{\displaystyle \frac{}{}}(\alpha 𝐢+\beta 𝐣+\gamma 𝐤)\times \left(e_x𝐢+e_y𝐣+e_z𝐤\right)\right]\mathrm{exp}[ik(\alpha x+\beta y+\gamma z)]`$ (9)
where $`\alpha =\mathrm{cos}\phi \mathrm{sin}\theta _0,\beta =\mathrm{sin}\phi \mathrm{sin}\theta _0,\gamma =\mathrm{cos}\theta _0`$ are the director cosines, $`Z`$ is the free-space impedance, and the temporal factor $`\mathrm{exp}(i\omega t)`$ is omitted. The total electric field $`𝐄`$ will be given by the superposition of all e$`d\phi `$ contribution arising from the dipoles, and results in
$`𝐄`$ $`=`$ $`{\displaystyle _0^{2\pi }}𝐞𝑑\phi ={\displaystyle _0^{2\pi }}(e_x𝐢+e_y𝐣+e_z𝐤)\mathrm{exp}\left\{ik\left[\rho \left({\displaystyle \frac{r}{R}}\right)\mathrm{cos}(\phi \psi )+\mathrm{cos}\theta _0z\right]\right\}𝑑\phi `$ (10)
$`=`$ $`\mathrm{exp}\left(ik\mathrm{cos}\theta _0z\right){\displaystyle _0^{2\pi }}(e_x𝐢+e_y𝐣+e_z𝐤)\mathrm{exp}\left[ik\rho \mathrm{sin}\theta _0\mathrm{cos}(\phi \psi )\right]𝑑\phi .`$
With reference to Eqs. (6) and (7), it is expedient to choose
$$e_x=\frac{e_0}{f^2}(r^2\mathrm{sin}^2\phi +f^2):$$
(11)
this condition can be experimentally obtained by a suitable choice of the transparency of the ring as a function of $`\phi `$. Thus, by substituting Eqs. (6), (7) and (11) into Eq. (10), and by recalling that $`r=f\mathrm{tan}\theta _0`$, we finally obtain (calculations are rather cumbersome but of no difficulty)
$`E_x`$ $`=`$ $`2\pi e_0e^{i\xi z}\left\{J_0(\eta \rho )+\mathrm{tan}^2\theta _0\left[\left(J_0(\eta \rho ){\displaystyle \frac{J_1(\eta \rho )}{\eta \rho }}\right)\mathrm{cos}^2\psi \left(J_0(\eta \rho ){\displaystyle \frac{2J_1(\eta \rho )}{\eta \rho }}\right)\right]\right\}`$ (12)
$`E_y`$ $`=`$ $`2\pi e_0e^{i\xi z}\left[{\displaystyle \frac{\mathrm{sin}2\psi }{2}}\mathrm{tan}^2\theta _0\left(J_0(\eta \rho ){\displaystyle \frac{2J_1(\eta \rho )}{\eta \rho }}\right)\right]`$ (13)
$`E_z`$ $`=`$ $`2\pi ie_0e^{i\xi z}\left[\mathrm{tan}\theta _0\mathrm{cos}\psi J_1(\eta \rho )\right],`$ (14)
where $`\xi =k\mathrm{cos}\theta _0,\eta =k\mathrm{sin}\theta _0`$, and $`J_1`$ denotes the first-order Bessel function of first kind.
Equation (12) (the main contribution of the electric field) describes a field different from the scalar field of Eq. (3) because of the presence of the term depending on $`\mathrm{tan}\theta _0`$. However, for $`\theta _0\pi /2(rf)`$, as in the present case, this term is negligible. We note that also the dependence on $`\psi `$ (which is absent in the scalar approximation) is negligible, and may be due to the approximation indicated in note2 . In Fig. 2 we report the normalized value of $`E_x,E_y`$ and $`E_z`$ vs $`\rho `$, for $`\theta _0=10^{}`$, together with the scalar field of Eq. (3). The scalar field is practically coincident with $`E_x`$. Therefore, we can conclude that the vectorial field derived above has Eq. (3) as its scalar approximation, at least for $`rf`$.
The magnetic field can be derived by Eq. (9), and results in
$`H_x`$ $`=`$ $`0`$ (15)
$`H_y`$ $`=`$ $`{\displaystyle \frac{2\pi }{Z}}e_0e^{i\xi z}{\displaystyle \frac{1}{\mathrm{cos}\theta _0}}J_0(\eta \rho )`$ (16)
$`H_z`$ $`=`$ $`i{\displaystyle \frac{2\pi }{Z}}e_0e^{i\xi z}\mathrm{sin}\psi {\displaystyle \frac{\mathrm{sin}\theta _0}{\mathrm{cos}^2\theta _0}}J_1(\eta \rho )`$ (17)
From a knowledge of the electric and magnetic fields, we are now in a position to evaluate the mean density of the energy flux which is defined as one half of the real part of the complex Poynting vectorjac ; str
$$𝐒=\frac{1}{2}\mathrm{Re}\left(𝐄\times 𝐇^{}\right).$$
(18)
For the fields (12)-(14) and (15)-(17), it turns out that $`S`$ has only the $`k`$-component
$$S_z=\frac{1}{2}\left(E_xH_y^{}\right),$$
(19)
that is the propagation of the energy flux occurs only in the $`z`$-direction, in accordance with the information given by the scalar field (3). Moreover, since the flux is independent of $`z`$, the energy propagates with no deformation.
In Fig. 3, the behavior of the energy flux (19) is shown as a function of $`\rho `$ for a few values of $`\theta _0`$. We note that, for $`\theta _0`$ very small (nearly plane wave) the flux is nearly independent of $`\rho `$, while when the beam originates the flux increases by increasing $`\theta _0`$, and tends to concentrate near $`\rho =0`$, that is, along the $`z`$-axis. Thus, for small values of $`\rho `$ (that is, in the proximity of the $`z`$-axis), the power supplied by a Bessel beam is always greater than the one due to a plane wave.
As for the velocity of the energy, from Eq. (3) it follows that in the scalar approximation the dependence of the field on $`t`$ and $`z`$ occurs only through the quantity $`(z/c)\mathrm{cos}\theta _0t`$ and, therefore, the field propagates with velocity $`v=c/\mathrm{cos}\theta _0`$. On the basis of these arguments, it could be concluded that also the energy propagates with a velocity $`v_e`$ greater than $`c`$. In the vectorial treatment we can evaluate the energy velocity $`v_e`$ as jac ; str
$$v_e=\frac{S_z}{\frac{1}{4}\left(\epsilon EE^{}+\mu HH^{}\right)},$$
(20)
where the quantity $`(1/4)(\epsilon EE^{}+\mu HH^{})`$ is the total mean density of energy which can be evaluated with the help of Eqs. (12)-(14) and (15)-(17).
In Fig. 4, we report the normalized velocity of energy as a function of the radial coordinate $`\rho `$, for $`\theta _0=10^{}`$. The velocity is found to be equal to $`c`$ from $`\rho =0`$ up to near to the first zero of the Bessel function: that is, the beam moves like an almost rigid system, in spite of its dependence on $`\rho `$ and $`\psi `$. In the proximity of the first zero of the Bessel function, the velocity decreases and tends to zero. Naturally, the zero in the velocity does not represent a stop of the motion but, more simply, the absence of energy flux. In this situation, the concept of velocity has no physical meaning.
Some remarks must be made on this surprising result. In fact, we recall that for propagation in vacuum “if an energy density is associated with the magnitude of the wave $`\mathrm{}\mathrm{}`$ the transport of energy occurs with the group velocity, since that is the rate of which the pulse travel along”jac7-8 . If the definition of the energy velocity as given by Eq. (20) is applicable also to a Bessel beam (or, more generally, to localized waves), it is not clear what kind of physical mechanism makes the energy velocity different from the phase and group ones.
We wish to recall that the present analysis was performed for an ideal system in the far field approximation. For a real system, we have to take into account the finite dimension of the converging system, which limits the field depth and introduces diffractive effects. The role of diffraction, together with the analysis of the near field (as in real experimental situations), make the problem much more complicated, and is beyond the purpose of this paper.
Acknowledgments
Special thanks are due to Laura Ronchi Abbozzo for useful suggestions and discussions.
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# Hurst exponent estimation of locally self-similar Gaussian processes using sample quantiles
## 1 Introduction
Many naturally occuring phenomena can be effectively modelled using self-similar processes. Among the simplest models, one can consider the fractional Brownian motion introduced in the statistics community by Mandelbrot and Van Ness (1968). Fractional Brownian motion can be defined as the only centered Gaussian process, denoted by $`(X(t))_t`$, with stationary increments and with variance function $`v()`$, given by $`v(t)=\sigma ^2|t|^{2H}`$, for all $`t`$. The fractional Brownian motion is an $`H`$-self-similar process, that is for all $`c>0`$, $`\left(X(ct)\right)_t\stackrel{d}{=}c^H\left(X(t)\right)_t`$ (where $`\stackrel{d}{=}`$ means equal in finite-dimensional distributions) with autocovariance function behaving like $`𝒪(|k|^{2H2})`$ as $`|k|+\mathrm{}`$. So the discretized increments of the fractional Brownian motion (called the fractional Gaussian noise) constitute a short-range dependent process, when $`H<1/2`$, and a long-range dependent process, when $`H>1/2`$. The index $`H`$ also characterizes the path regularity since the fractal dimension of the fractional Brownian motion is equal to $`D=2H`$. According to the context (long-range dependent processes, self-similar processes,…), a very large variety of estimators of the parameter $`H`$ has been investigated. The reader is referred to Beran (1994), Coeurjolly (2000) or Bardet et al. (2003) for an overview of this problem. Among the most often used estimators we have: methods based on the variogram, on the log$``$periodogram e.g. Geweke and Porter-Hudak (1983) in the context of long-range dependent processes, maximum likelihood estimator (and Whittle estimator) when the model is parametric e.g. fractional Gaussian noise, methods based on the wavelet decomposition e.g. Flandrin (1992) or Stoev et al. (2006) and the references therein, and on discrete filtering studied by Kent and Wood (1997), Istas and Lang (1997) and Coeurjolly (2001). We are mainly interested in the last one, which has several similarities with the wavelet decomposition method. Following Constantine and Hall (1994), Kent and Wood (1997), Istas and Lang (1997), in the case when the process is observed at times $`i/n`$ for $`i=1,\mathrm{},n`$, this method is adapted to a larger class than the fractional Brownian motion, namely the class of centered Gaussian processes with stationary increments that are locally self-similar (at zero). A process $`(X(t))_t`$ is said to be locally self-similar (at zero) if its variance function, denoted by $`v()`$, satisfies
$$v(t)=𝐄(X(t)^2)=\sigma ^2|t|^{2H}\left(1+r(t)\right),\text{ with }r(t)=o\left(1\right)\text{ as }|t|0,$$
(1)
for some $`0<H<1`$. An estimator of $`H`$ is derived by using the stationarity of the increments and the local behavior of the variance function. When observing the process at regular subdivisions, the stationarity of the increments is crucial since the method based on discrete filtering (and the one based on the wavelet decomposition) essentially uses the fact that the variance of the increments can be estimated by the sample moment of order 2. We do not believe that this framework could be valid for the estimation of the Hurst exponent of Riemann-Liouville’s process, e.g. Alòs et al. (1999) which is an $`H`$-self-similar centered Gaussian process but with increments satisfying only some kind of local stationarity, see Remark 2 for more details.
Let us be more specific on the construction of the wavelet decomposition method, see e.g. Flandrin (1992): the authors noticed that the variance of the wavelet coefficient at a scale say $`j`$ behaves like $`2^{j(2H1)}`$. An estimator of $`H`$ is then derived by regressing the logarithm of sample moment of order 2 at each scale against $`\mathrm{log}(j)`$ for various scales. This procedure exhibits good properties since it is also proved that the more vanishing moments the wavelet has the observations are more decorrelated. And so asymptotic results are quite easy to obtain. However, Stoev et al. (2006) illustrate the fact that this kind of estimator is very sensitive to additive outliers and to non-stationary artefacts. Therefore, they mainly propose to replace at each scale, the sample moment of order 2, by the sample median of the squared coefficients. This procedure, for which the authors assert that no theoretical result is available, is clearly more robust.
The main objective of this paper is to extend the procedure proposed by Stoev et al. (2006) by deriving semi-parametric estimators of the parameter $`H`$, using discrete filtering methods, for the class of processes defined by (1). The procedure is extended in the sense that we consider either convex combinations of sample quantiles or trimmed-means. Moreover, we provide convergence results. The key-ingredient is a Bahadur representation of sample quantiles obtained in a certain dependence framework. Let $`𝒀=(Y(1),\mathrm{},Y(n))`$ be a vector of $`n`$ i.i.d. random variables with cumulative distribution function $`F`$, as well denote by $`\xi (p)`$ and $`\widehat{\xi }\left(p\right)`$ the quantile respectively the sample quantile of order $`p`$. By assuming that $`F^{}(\xi (p))>0`$ and $`F^{\prime \prime }(\xi (p))`$ exists, Bahadur proved that as $`n+\mathrm{}`$,
$$\widehat{\xi }\left(p\right)\xi (p)=\frac{p\widehat{F}\left(p\right)}{f(\xi (p)}+r_n,$$
with $`r_n=𝒪_{a.s.}\left(n^{3/4}\mathrm{log}(n)^{3/4}\right)`$. Using a law of iterated logarithm’s type result, Kiefer obtained the exact rate $`n^{3/4}\mathrm{log}\mathrm{log}(n)^{3/4}`$. Extensions of the above results to dependent random variables have been pursued in Sen and Ghosh (1972) for $`\varphi `$mixing variables, in Yoshihara (1995) for strongly mixing variables, and recently in Wu (2005) for short-range and long-range dependent linear processes, following works of Hesse (1990) and Ho and Hsing (1996). Our contribution is to provide a Bahadur representation for sample quantiles in another context that is for non-linear functions of Gaussian processes with correlation function decreasing as $`k^\alpha L(k)`$ for some $`\alpha >0`$ and some slowly varying function $`L()`$. The bounds for $`r_n`$ are obtained under the same assumption as those used by Bahadur (1966).
The paper is organized as follows. In Section 2, we give some basic notations and some background on discrete filtering. In Section 3, we derive semi-parametric estimators of the parameter $`H`$, when a single sample path of a process defined by (1) is observed over a discrete grid of the interval $`[0,1]`$. Section 4 presents the main results: Bahadur representations and asymptotic results for our estimators. In Section 5 are presented some numerical computations to compare the theoretical asymptotic variance of our estimators and a simulation study is also given. In particular, we illustrate the relative efficiency with respect to Whittle estimator and the fact that such estimators are more robust than classical ones. Finally, proofs of differents results are presented in Section 6.
## 2 Some notations and some background on discrete filtering
Given some random variable $`Y`$, $`F_Y()`$ denotes the cumulative distribution function of $`Y`$ and $`\xi _Y(p)`$ the quantile of order $`p`$, $`0<p<1`$. If $`F_Y()`$ is absolutely continuous with respect to Lebesgue measure, the probability density function is denoted by $`f_Y()`$. The cumulative distribution (resp. probability density) function of a standard Gaussian variable is denoted by $`\mathrm{\Phi }()`$ (resp. $`\varphi ()`$). Based on the observation of a vector $`𝒀=(Y(1),\mathrm{},Y(n))`$ of $`n`$ random variables distributed as $`Y`$, the sample cumulative distribution function and the sample quantile of order $`p`$ are respectively denoted by $`\widehat{F}_Y(;𝒀)`$ and $`\widehat{\xi }_Y(p;𝒀)`$ or simply by $`\widehat{F}(;𝒀)`$ and $`\widehat{\xi }(p;𝒀)`$. Finally, for some measurable function $`g()`$, we denote by $`𝒈\mathbf{(}𝒀\mathbf{)}`$ the vector of length $`n`$ with real components $`g(Y(i))`$, for $`i=1,\mathrm{},n`$.
A sequence of real numbers $`u_n`$ is said to be $`𝒪\left(v_n\right)`$ (resp. $`o\left(v_n\right)`$) for an other sequence of real numbers $`v_n`$, if $`u_n/v_n`$ is bounded (resp. converges to 0 as $`n+\mathrm{}`$). A sequence of random variables $`U_n`$ is said to be $`𝒪_{a.s.}\left(v_n\right)`$ (resp. $`o_{a.s.}\left(v_n\right)`$) if $`U_n/v_n`$ is almost surely bounded (resp. if $`U_n/v_n`$ converges towards 0 with probability 1).
The statistical model corresponds to a discretized version $`𝑿=\left(X(i/n)\right)_{i=1,\mathrm{},n}`$ of a locally self-similar Gaussian process defined by (1).
One of the ideas of our method is to construct some estimators by using some properties of the variance of the increments of $`𝑿`$ or the variance of the increments of order 2 of $`𝑿`$. While considering the increments of $`𝑿`$ is conventional since the associated sequence is stationary, considering the increments of order 2 (or of a higher order) could be stranger. However, the main interest relies upon the fact that the observations of the latter resulting sequences are less correlated than those of the simple increments’ sequence. All these vectors can actually be seen as special discrete filtering of the vector $`𝑿`$. Let us now specify some general background on discrete filtering and its consequence on the correlation structure. The vector $`𝒂`$ is a filter of length $`\mathrm{}+1`$ and of order $`\nu 1`$ with real components if
$$\underset{q=0}{\overset{\mathrm{}}{}}q^ja_q=0,\text{ for }j=0,\mathrm{},\nu 1\text{ and }\underset{q=0}{\overset{\mathrm{}}{}}q^\nu a_q0.$$
For example, $`𝒂=(1,1)`$ (resp. $`𝒂=(1,2,1)`$) is a filter with order 1 (resp. 2). Let $`𝑿^𝒂`$ be the series obtained by filtering $`𝑿`$ with $`𝒂`$, then:
$$X^𝒂\left(\frac{i}{n}\right)=\underset{q=0}{\overset{\mathrm{}}{}}a_qX\left(\frac{iq}{n}\right)\text{ for }i\mathrm{}+1.$$
Applying in turn the filter $`𝒂=(1,1)`$ and $`𝒂=(1,2,1)`$ leads to the increments of $`𝑿`$, respectively the increments of $`𝑿`$ of order 2. One may also consider other filters such as Daubechies wavelet filters, e.g. Daubechies (1992).
The following assumption is needed by different results presented hereafter:
$`𝑨𝒔𝒔𝒖𝒎𝒑𝒕𝒊𝒐𝒏𝑨_\mathrm{𝟏}\mathbf{(}𝒌\mathbf{)}:`$ for $`i=1,\mathrm{},k`$
$$v^{(i)}(t)=\sigma ^2\beta (i)|t|^{2Hi}+o\left(|t|^{2Hi}\right)$$
with $`\beta (i)=2H(2H1)\mathrm{}(2Hi+1)`$ (where $`k1`$ is an integer).
This assumption assures that the variance function $`v()`$ is sufficiently smooth around 0. It allows us to assert that the correlation structure of a locally self-similar discretized and filtered Gaussian process can be compared to the one of the fractional Brownian motion. This is announced more precisely in the following Lemma.
###### Lemma 1 (e.g. Kent and Wood (1997))
Let $`𝐚`$ and $`𝐚^{\mathbf{}}`$ be two filters of length $`\mathrm{}+1`$ and $`\mathrm{}^{}+1`$, of order $`\nu `$ and $`\nu ^{}1`$. Then we have:
$`𝐄\left(X^𝒂\left({\displaystyle \frac{i}{n}}\right)X^𝒂^{\mathbf{}}\left({\displaystyle \frac{i+j}{n}}\right)\right)`$ $`=`$ $`{\displaystyle \frac{\sigma ^2}{2}}{\displaystyle \underset{q,q^{}=0}{\overset{\mathrm{}}{}}}a_qa_q^{}^{}v\left({\displaystyle \frac{qq^{}+j}{n}}\right)`$ (2)
$`=`$ $`\gamma _n^{𝒂,𝒂^{\mathbf{}}}(j)\left(1+\delta _n^{𝒂,𝒂^{\mathbf{}}}(j)\right),`$
with
$$\gamma _n^{𝒂,𝒂^{\mathbf{}}}(j)=\frac{\sigma ^2}{n^{2H}}\gamma ^{𝒂,𝒂^{\mathbf{}}}(j),\gamma ^{𝒂,𝒂^{\mathbf{}}}(j)=\frac{1}{2}\underset{q,q^{}=0}{\overset{\mathrm{}}{}}a_qa_q^{}^{}|qq^{}+j|^{2H}$$
(3)
and
$$\delta _n^{𝒂,𝒂^{\mathbf{}}}(j)=\frac{_{q,q^{}}a_qa_q^{}|qq^{}+j|^{2H}\times r\left(\frac{qq^{}+j}{n}\right)}{\gamma ^{𝒂,𝒂^{\mathbf{}}}(j)}.$$
(4)
Moreover, as $`|j|+\mathrm{}`$
$$\gamma ^{𝒂,𝒂^{\mathbf{}}}(j)=𝒪\left(\frac{1}{|j|^{2H\nu \nu ^{}}}\right).$$
(5)
Finally, under Assumption $`𝐀_\mathrm{𝟏}\mathbf{(}𝛎\mathbf{+}𝛎^{\mathbf{}}\mathbf{)}`$, as $`n+\mathrm{}`$
$$\delta _n^{𝒂,𝒂^{\mathbf{}}}(j)=o\left(1\right).$$
(6)
###### Remark 1
In the case of the fractional Brownian motion the sequence $`\delta _n`$ is equal to 0, whereas it converges towards 0 for more general locally self-similar Gaussian processes, such as the Gaussian processes with stationary increments and with variance function $`v(t)=1\mathrm{exp}(|t|^{2H})`$ or $`v(t)=\mathrm{log}(1+|t|^{2H})`$ for which Assumption $`𝐀_\mathrm{𝟏}\mathbf{(}𝐤\mathbf{)}`$ is satisfied (for every $`k1`$).
###### Remark 2
The stationarity of the increments and the local self-similarity required on the process $`X()`$ are important, if the process is observed at times $`i/n`$ for $`i=1,\mathrm{},n`$. The crucial result of Lemma 1 is that the variance function of the filtered series behaves asymptotically as $`\gamma _n^𝐚(0)`$. It seems to be difficult to relax the constraint of stationarity. Consider for example the Riemann-Liouville’s process, e.g. Alòs et al. (1999). This process is a Gaussian process which is $`H`$-self similar Gaussian but with increments satisfying only some kind of local stationarity. Following the computations of Lim (2001), the variance of the increments’ series of the Riemann-Liouville’s process is equal to
$$𝐄\left(\left(X\left(\frac{i+1}{n}\right)X\left(\frac{i}{n}\right)\right)^2\right)=\frac{1}{n^{2H}}\frac{1}{\mathrm{\Gamma }(H+1/2)^2}\left\{I+\frac{1}{2H}\right\},$$
with $`I=_0^i\left((1+u)^{H1/2}u^{H1/2}\right)^3𝑑u+_0^{i/n}u^{2H1}𝑑u.`$ This integral cannot be asymptotically independent of time. Note that this could be the case if the process is observed at irregular subdivisions. This question has not been investigated.
Define $`𝒀^𝒂`$ as the normalized vector $`𝑿^𝒂`$ with variance $`1`$. The covariance between $`Y^𝒂(i/n)`$ and $`Y^𝒂^{\mathbf{}}(i+j/n)`$ is denoted by $`\rho _n^{𝒂,𝒂^{\mathbf{}}}(j)`$. Under Assumption $`𝑨_\mathrm{𝟏}\mathbf{(}𝝂\mathbf{+}𝝂^{\mathbf{}}\mathbf{)}`$, the following equivalence holds as $`n+\mathrm{}`$
$$\rho _n^{𝒂,𝒂^{\mathbf{}}}(j)\rho ^{𝒂,𝒂^{\mathbf{}}}(j)=\frac{\gamma ^{𝒂,𝒂^{\mathbf{}}}(j)}{\sqrt{\gamma ^{𝒂,𝒂}(0)\gamma ^{𝒂^{\mathbf{}},𝒂^{\mathbf{}}}(0)}}.$$
(7)
When $`𝒂=𝒂^{\mathbf{}}`$, we set, for the sake of simplicity $`\gamma _n^𝒂()=\gamma _n^{𝒂,𝒂}()`$, $`\delta _n^𝒂()=\delta _n^{𝒂,𝒂}()`$, $`\rho _n^𝒂()=\rho _n^{𝒂,𝒂}()`$, $`\gamma ^{𝒂,𝒂}()=\gamma ^𝒂()`$ and $`\rho ^𝒂()=\rho ^{𝒂,𝒂}()`$.
## 3 New estimators of $`H`$
### 3.1 Estimators based on a convex combination of sample quantiles
Let $`(𝒑,𝒄)=(p_k,c_k)_{k=1,\mathrm{},K}((0,1)\times ^+)^K`$ for an integer $`1K<+\mathrm{}`$. Define the following statistics based on a convex combination of sample quantiles:
$$\widehat{\xi }(𝒑,𝒄;𝑿^𝒂)=\underset{k=1}{\overset{K}{}}c_k\widehat{\xi }(p_k;𝑿^𝒂),$$
(8)
where $`c_k,k=1,\mathrm{},K`$ are positive real numbers such that $`_{k=1}^Kc_k=1`$. For example, this corresponds to the sample median when $`K=1,𝒑=1/2,𝒄=1`$ , to a mean of quartiles when $`K=2,𝒑=(1/4,3/4),𝒄=(1/2,1/2)`$ . Consider the following computation: from Lemma 1, we have, as $`n+\mathrm{}`$
$$\widehat{\xi }(𝒑,𝒄;𝑿^𝒂)\frac{\sigma ^2}{n^{2H}}\gamma ^𝒂(0)\widehat{\xi }(𝒑,𝒄;𝒀^𝒂).$$
###### Remark 3
It may be expected that $`\widehat{\xi }(𝐩,𝐜;𝐘^𝐚)`$ converges towards a constant as $`n+\mathrm{}`$. In itself, this result is not interesting, since two parameters remain unknown: $`\sigma ^2`$ and $`H`$ and thus, it is impossible to derive an estimator of $`H`$.
Remark 3 suggests that we have to use at least two filters. Among all available filters, let us consider the sequence $`(𝒂^m)_{m1}`$ defined by
$$a_i^m=\{\begin{array}{cc}a_j\hfill & \text{if }i=jm\hfill \\ 0\hfill & \text{otherwise }\hfill \end{array}\text{for }i=0,\mathrm{},m\mathrm{},$$
which is none other than the filter $`𝒂`$ dilated $`m`$ times. For example, if the filter $`𝒂=𝒂^1`$ corresponds to the filter $`(1,2,1)`$, then $`𝒂^2=(1,0,2,0,1)`$, $`𝒂^3=(1,0,0,2,0,0,1)`$, …As noted by Kent and Wood (1997) or Istas and Lang (1997), the filter $`𝒂^m`$, of length $`m\mathrm{}+1`$, is of order $`\nu `$ and has the following interesting property :
$$\gamma ^{𝒂^m}(0)=m^{2H}\gamma ^𝒂(0).$$
(9)
From Lemma 1, this simply means that $`𝐄\left(𝑿^{𝒂^𝒎}(i/n)^2\right)=m^{2H}𝐄\left(𝑿^𝒂(i/n)^2\right)`$, exhibiting some kind of self-similarity property of the filtered coefficients. As specified in the introduction, the same property can be pointed out in the context of wavelet decomposition.
Our methods, that exploit the nice property (9), are based on a convex combination of sample quantiles $`\widehat{\xi }(𝒑,𝒄;𝒈\mathbf{(}𝑿^{𝒂^𝒎}\mathbf{)})`$ for two positive functions $`g()`$: $`g()=||^\alpha `$ for $`\alpha >0`$ and $`g()=\mathrm{log}||`$. For such functions $`g()`$ we manage, by using some property established in Lemma 1, to define some very simple estimators of the Hurst exponent through a simple linear regression. Other choices of the function $`g()`$ have not been investigated in this paper. At this stage, let us specify that our methods extend the one proposed by Stoev et al. (2006); indeed they only consider the statistic $`\widehat{\xi }(𝒑,𝒄;𝒈\mathbf{(}𝑿^{𝒂^𝒎}\mathbf{)})`$ for $`𝒑=1/2`$, $`𝒄=1`$, $`g()=()^2`$, that is the sample median of the squared coefficients. From (3) and (9), we have
$`\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}^𝜶)`$ $`=`$ $`𝐄\left((X^{𝒂^𝒎}(1/n))^2\right)^{\alpha /2}\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)`$ (10)
$`=`$ $`m^{\alpha H}{\displaystyle \frac{\sigma ^\alpha }{n^{\alpha H}}}\gamma ^𝒂(0)^{\alpha /2}\left(1+\delta _n^{a^m}(0)\right)^{\alpha /2}\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶),`$
and
$`\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}𝐄(X^{𝒂^𝒎}(1/n))^2+\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})`$ (11)
$`=`$ $`H\mathrm{log}(m)+\mathrm{log}\left({\displaystyle \frac{\sigma ^2}{n^{2H}}}\gamma ^𝒂(0)\right)`$
$`+{\displaystyle \frac{1}{2}}\mathrm{log}\left(1+\delta _n^{a^m}(0)\right)+\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}).`$
Denote by $`\kappa _H=n^{2H}\sigma ^2\gamma ^𝒂(0)`$. Equations (10) and (11) can be rewritten as
$`\mathrm{log}\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}^𝜶)`$ $`=`$ $`\alpha H\mathrm{log}(m)+\mathrm{log}\left(\kappa _H^{\alpha /2}\xi _{|Y|^\alpha }(𝒑,𝒄)\right)+\epsilon _m^\alpha ,`$ (12)
$`\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|})`$ $`=`$ $`H\mathrm{log}(m)+\mathrm{log}(\kappa _H)+\xi _{\mathrm{log}|Y|}(𝒑,𝒄)+\epsilon _m^{\mathrm{log}}`$ (13)
with the random variables $`\epsilon _m^\alpha `$ and $`\epsilon _m^{\mathrm{log}}`$ respectively defined by
$$\epsilon _m^\alpha =\mathrm{log}\left(\frac{\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)}{\xi _{|Y|^\alpha }(𝒑,𝒄)}\right)+\frac{\alpha }{2}\mathrm{log}\left(1+\delta _n^{a^m}(0)\right),$$
(14)
and
$$\epsilon _m^{\mathrm{log}}=\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{\mathrm{log}|Y|}(𝒑,𝒄)+\frac{1}{2}\mathrm{log}\left(1+\delta _n^{a^m}(0)\right)$$
(15)
where, for some random variable $`Z`$, $`\xi _Z(𝒑,𝒄)=_{k=1}^Kc_k\xi _Z(p_k)`$. We decide to rewrite Equations (10) and (11) as (12) and (13), since we expect that $`\epsilon _m^\alpha `$ and $`\epsilon _m^{\mathrm{log}}`$ converge (almost surely) towards 0 as $`n+\mathrm{}`$.
From Remark 3, two estimators of $`H`$ can be defined through a linear regression of $`\left(\mathrm{log}\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}^𝜶)\right)_{m=1,\mathrm{},M}`$ and $`\left(\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|})\right)_{m=1,\mathrm{},M}`$ on $`\left(\mathrm{log}m\right)_{m=1,\mathrm{},M}`$ for some $`M2`$. These estimators are denoted by $`\widehat{H}^\alpha `$ and $`\widehat{H}^{\mathrm{log}}`$. By denoting $`𝑨`$ the vector of length $`M`$ with components $`A_m=\mathrm{log}m\frac{1}{M}_{m=1}^M\mathrm{log}(m)`$, $`m=1,\mathrm{},M`$, we have explicitly from (12) and (13) and the definition of least squares estimates (see e.g. Antoniadis et al. (1992)):
$`\widehat{H}^\alpha `$ $`=`$ $`{\displaystyle \frac{𝑨^T}{\alpha 𝑨^2}}\left(\mathrm{log}\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}^𝜶)\right)_{m=1,\mathrm{},M},`$ (16)
$`\widehat{H}^{\mathrm{log}}`$ $`=`$ $`{\displaystyle \frac{𝑨^T}{𝑨^2}}\left(\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|})\right)_{m=1,\mathrm{},M},`$ (17)
where $`𝒛`$ for some vector $`𝒛`$ of length $`d`$ denotes the norm defined by $`\left(_{i=1}^dz_i^2\right)^{1/2}`$.
We can point out that $`\widehat{H}^\alpha `$ and $`\widehat{H}^{\mathrm{log}}`$ are independent of the scaling coefficient $`\sigma ^2`$.
### 3.2 Estimators based on trimmed means
Let $`0<\beta _1\beta _2<1`$ and $`𝜷=(\beta _1,\beta _2)`$, denote by $`\overline{𝒈\mathbf{(}𝑿^𝒂\mathbf{)}}^{(𝜷)}`$ the $`𝜷`$trimmed mean of the vector $`𝒈\mathbf{(}𝑿^𝒂\mathbf{)}`$ given by
$$\overline{𝒈\mathbf{(}𝑿^𝒂\mathbf{)}}^{(𝜷)}=\frac{1}{n[n\beta _2][n\beta _1]}\underset{[n\beta _1]+1}{\overset{n[n\beta _2]}{}}\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_{(i),n},$$
where $`\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_{(1),n}\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_{(2),n}\mathrm{}\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_{(2),n}`$ are the order statistics of $`\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_1,\mathrm{},\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_n`$. It is well-known that $`\left(𝒈\mathbf{(}𝑿^𝒂\mathbf{)}\right)_{(i),n}=\widehat{\xi }(\frac{i}{n};𝒈\mathbf{(}𝑿^𝒂\mathbf{)})`$. Hence, by following the ideas of the previous section, one may obtain
$`\mathrm{log}\left(\overline{\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}^𝜶}^{(𝜷)}\right)`$ $`=`$ $`\alpha H\mathrm{log}(m)+\mathrm{log}\left(\kappa _H^{\alpha /2}\overline{|Y|^\alpha }^{(𝜷)}\right)+\epsilon _m^{\alpha ,tm},`$ (18)
$`\overline{\mathrm{𝐥𝐨𝐠}\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}}^{(𝜷)}`$ $`=`$ $`H\mathrm{log}(m)+\mathrm{log}(\kappa _H)+\overline{\mathrm{log}|Y|}^{(𝜷)}+\epsilon _m^{\mathrm{log},tm}`$ (19)
with
$$\epsilon _m^{\alpha ,tm}=\overline{\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶}^{(𝜷)}\overline{|Y|^\alpha }^{(𝜷)}+\frac{\alpha }{2}\mathrm{log}\left(1+\delta _n^{a^m}(0)\right),$$
(20)
and
$$\epsilon _m^{\mathrm{log},tm}=\overline{\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}}^{(𝜷)}\overline{\mathrm{log}|Y|}^{(𝜷)}+\frac{1}{2}\mathrm{log}\left(1+\delta _n^{a^m}(0)\right),$$
(21)
where for some random variable $`Z`$, $`\overline{Z}^{(𝜷)}`$ is referring to
$$\overline{Z}^{(𝜷)}=\frac{1}{1\beta _2\beta _1}_{\beta _1}^{1\beta _2}\xi _Z(p)𝑑p.$$
(22)
As in the previous section, two estimators of $`H`$, denoted by $`\widehat{H}^{\alpha ,tm}`$ and $`\widehat{H}^{\mathrm{log},tm}`$, is derived through a log-linear regression
$`\widehat{H}^{\alpha ,tm}`$ $`=`$ $`{\displaystyle \frac{𝑨^T}{\alpha 𝑨^2}}\left(\overline{\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}^𝜶}^{(𝜷)}\right)_{m=1,\mathrm{},M}.`$ (23)
$`\widehat{H}^{\mathrm{log},tm}`$ $`=`$ $`{\displaystyle \frac{𝑨^T}{𝑨^2}}\left(\overline{\mathrm{𝐥𝐨𝐠}\mathbf{|}𝑿^{𝒂^𝒎}\mathbf{|}}^{(𝜷)}\right)_{m=1,\mathrm{},M}.`$ (24)
###### Remark 4
The estimator referred to the “estimator based on the quadratic variations” in the simulation study and studied with the same formalism by Coeurjolly (2001) corresponds to the estimator $`\widehat{H}^{\alpha ,tm}`$ with $`\alpha =2`$, $`\beta _1=\beta _2=0`$.
## 4 Main results
To simplify the presentation of different results, consider the two following assumptions on different parameters involved in the estimation procedures
$`𝑨𝒔𝒔𝒖𝒎𝒑𝒕𝒊𝒐𝒏𝑨_\mathrm{𝟐}\mathbf{(}𝒑\mathbf{,}𝒄\mathbf{)}:`$ $`𝒂`$ is a filter of order $`\nu 1`$, $`\alpha `$ is a positive real number, $`𝒑`$ (resp. $`𝒄`$) is a vector of length $`K`$ (for some $`1K<+\mathrm{}`$) such that $`0<p_k<1`$ (resp. $`c_k>0`$ and $`_{k=1}^Kc_k=1`$), $`M`$ is an integer $`2`$.
$`𝑨𝒔𝒔𝒖𝒎𝒑𝒕𝒊𝒐𝒏𝑨_\mathrm{𝟑}\mathbf{(}𝜷\mathbf{)}:`$ $`𝒂`$ is a filter of order $`\nu 1`$, $`\alpha `$ is a positive real number, $`𝜷=(\beta _1,\beta _2)`$ is such that $`0<\beta _1\beta _2<1`$, $`M`$ is an integer $`2`$.
Since $`A^T\left(\mathrm{log}(m)\right)_{m=1,\mathrm{},M}=A^2`$ and $`A^T\mathrm{𝟏}=0`$ (where $`\mathrm{𝟏}=(1)_{m=1,\mathrm{},M}`$), we have
$$\widehat{H}^\alpha H=\frac{𝑨^T}{\alpha 𝑨^2}𝜺^𝜶\text{ and }\widehat{H}^{\mathrm{log}}H=\frac{𝑨^T}{𝑨^2}𝜺^{\mathrm{𝐥𝐨𝐠}},$$
(25)
and
$$\widehat{H}^{\alpha ,tm}H=\frac{𝑨^T}{\alpha 𝑨^2}𝜺^{𝜶\mathbf{,}𝒕𝒎}\text{ and }\widehat{H}^{\mathrm{log},tm}H=\frac{𝑨^T}{𝑨^2}𝜺^{\mathrm{𝐥𝐨𝐠}\mathbf{,}𝒕𝒎},$$
(26)
where $`𝜺^𝜶=\left(\epsilon _m^\alpha \right)_{m=1,\mathrm{},M}`$, $`𝜺^{\mathrm{𝐥𝐨𝐠}}=\left(\epsilon _m^{\mathrm{log}}\right)_{m=1,\mathrm{},M}`$, $`𝜺^{𝜶\mathbf{,}𝒕𝒎}=\left(\epsilon _m^{\alpha ,tm}\right)_{m=1,\mathrm{},M}`$ and $`𝜺^{\mathrm{𝐥𝐨𝐠}\mathbf{,}𝒕𝒎}=\left(\epsilon _m^{\mathrm{log},tm}\right)_{m=1,\mathrm{},M}`$. Hence, in order to study the convergence of different estimators, it is sufficient to obtain some convergence results of sample quantiles $`\widehat{\xi }(p,𝒈\mathbf{(}𝒀^𝒂\mathbf{)})`$ for some function $`g()`$ and some filter $`𝒂`$. Therefore, we first establish a Bahadur representation of sample quantiles for some non-linear function of Gaussian sequences with correlation function decreasing as $`k^\alpha `$, for some $`\alpha >0`$. In fact, the existing litterature on nonlinear function of Gaussian sequences (e.g. Taqqu (1977)) allows us to slighlty extend this framework by considering correlation function decreasing as $`k^\alpha L(k)`$, for some slowly varying function $`L()`$.
### 4.1 Bahadur representation of sample quantiles
Let us recall some important definitions on Hermite polynomials. The $`j`$-th Hermite polynomial (for $`j0`$) is defined for $`t`$ by
$$H_j(t)=\frac{(1)^j}{\varphi (t)}\frac{d^j\varphi (t)}{dt^j}.$$
(27)
The Hermite polynomials form an orthogonal system for the Gaussian measure. More precisely, we have $`𝐄\left(H_j(Y)H_k(Y)\right)=j!\delta _{j,k}`$. For a measurable function $`g()`$ defined on $``$ for which $`𝐄(g(Y)^2)<+\mathrm{}`$, the following expansion holds
$$g(t)=\underset{j\tau }{}\frac{c_j}{j!}H_j(t)\text{ with }c_j=𝐄\left(g(Y)H_j(Y)\right),$$
where the integer $`\tau `$ defined by $`\tau =inf\left\{j0,c_j0\right\}`$, is called the Hermite rank of the function $`g`$. Note that this integer plays an important role. For example, it is related to the correlation of $`g(Y_1)`$ and $`g(Y_2)`$ (for $`Y_1`$ and $`Y_2`$ two standard gaussian variables with correlation $`\rho `$) since $`𝐄(g(Y_1)g(Y_2))=_{k\tau }\frac{(c_k)^2}{k!}\rho ^k\rho ^\tau g_{L^2(d\varphi )}`$.
In order to obtain a Bahadur representation (see e.g. Serfling (1980)), we have to ensure that $`F_{g(Y)}^{}(\xi (p))>0`$ and $`F_{g(Y)}^{\prime \prime }()`$ exists and is bounded in a neighborhood of $`\xi (p)`$. This is achieved if the function $`g()`$ satisfies the following assumption (see e.g. Dacunha-Castelle and Duflo (1982), p.33).
$`𝑨𝒔𝒔𝒖𝒎𝒑𝒕𝒊𝒐𝒏𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}\mathbf{:}`$ there exist $`U_i`$, $`i=1,\mathrm{},L`$, disjoint open sets such that $`U_i`$ contains a unique solution to the equation $`g(t)=\xi _{g(Y)}(p)`$, such that $`F_{g(Y)}^{}(\xi (p))>0`$ and such that $`g`$ is a $`𝒞^2`$diffeomorphism on $`_{i=1}^LU_i`$.
Note that under this assumption
$$F_{g(Y)}^{}(\xi _{g(Y)}(p))=f_{g(Y)}(\xi _{g(Y)}(p))=\underset{i=1}{\overset{L}{}}\frac{\varphi (g_i^1(\xi (p)))}{g^{}(g_i^1(\xi (p)))},$$
where $`g_i()`$ is the restriction of $`g()`$ on $`U_i`$. Now, define, for some real $`u`$, the function $`h_u()`$ by:
$$h_u(t)=\mathrm{𝟏}_{\{g(t)u\}}(t)F_{g(Y)}(u).$$
(28)
We denote by $`\tau (u)`$ the Hermite rank of $`h_u()`$. For the sake of simplicity, we set $`\tau _p=\tau (\xi _{g(Y)}(p))`$. For some function $`g()`$ satisfying Assumption $`𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}`$, we denote by
$$\overline{\tau }_p=\underset{\gamma _{i=1}^Lg(U_i)}{inf}\tau (\gamma ),$$
(29)
that is the minimal Hermite rank of $`h_u()`$ for $`u`$ in a neighborhood of $`\xi _{g(Y)}(p)`$.
###### Theorem 2
Let $`\left\{Y(i)\right\}_{i=1}^+\mathrm{}`$ be a stationary (centered) gaussian process with variance 1, and correlation function $`\rho ()`$ such that, as $`i+\mathrm{}`$
$$|\rho (i)|L(i)i^\alpha ,$$
(30)
for some $`\alpha >0`$ and some slowly varying function at infinity $`L(s),s0`$. Then, under Assumption $`𝐀_\mathrm{𝟒}\mathbf{(}𝛏\mathbf{(}𝐩\mathbf{)}\mathbf{)}`$, we have almost surely, as $`n+\mathrm{}`$
$$\widehat{\xi }(p;𝒈\mathbf{(}𝒀\mathbf{)})\xi _{g(Y)}(p)=\frac{p\widehat{F}(\xi _{g(Y)}(p);𝒈\mathbf{(}𝒀\mathbf{)})}{f_{g(Y)}(\xi _{g(Y)}(p))}+𝒪_{a.s.}\left(r_n(\alpha ,\overline{\tau }_p)\right),$$
(31)
the sequence $`\left(r_n(\alpha ,\overline{\tau }_p)\right)_{n1}`$ being defined by
$$r_n(\alpha ,\overline{\tau }_p)=\{\begin{array}{ccc}n^{3/4}\mathrm{log}(n)^{3/4}\hfill & \text{if}\hfill & \alpha \overline{\tau }_p>1,\hfill \\ n^{3/4}\mathrm{log}(n)^{3/4}L_{\overline{\tau }_p}(n)^{3/4}\hfill & \text{if}\hfill & \alpha \overline{\tau }_p=1,\hfill \\ n^{1/2\alpha \overline{\tau }_p/4}\mathrm{log}(n)^{\overline{\tau }_p/4+1/2}L(n)^{\overline{\tau }_p/4}\hfill & \text{if}\hfill & 2/3<\alpha \overline{\tau }_p<1,\hfill \\ n^{\alpha \overline{\tau }_p}\mathrm{log}(n)^{\overline{\tau }_p}L(n)^{\overline{\tau }_p}\hfill & \text{if}\hfill & 0<\alpha \overline{\tau }_p2/3,\hfill \end{array}$$
(32)
where for some $`\tau 1`$, $`L_\tau (n)=_{|i|n}|\rho (i)|^\tau `$.
Note that if $`L()`$ is an increasing function, $`L_\tau (n)=𝒪\left(\mathrm{log}(n)L(n)^\tau \right)`$.
###### Remark 5
Without giving any details here, let us say that the behaviour of the sequence $`r_n(,)`$ is related to the characteristic (short-range or long-range dependence) of the process $`\left\{h_u(Y(i))\right\}_{i=1}^+^{\mathrm{}}`$ for $`u`$ in a neighborhood of $`\xi _{g(Y)}(p)`$. In the case $`\alpha \overline{\tau }_p>1`$, corresponding to short-range dependent processes, the result is similar to the one proved by Bahadur, see e.g. Serfling (1980), in the i.i.d. case. For short-range dependent linear processes, using a law of iterated logarithm’s type result Wu (2005) obtained a sharper bound, that is $`n^{3/4}\mathrm{log}\mathrm{log}(n)^{3/4}`$. This bound is obtained under the assumption that $`F^{}()`$ and $`F^{\prime \prime }()`$ exist and are uniformly bounded. For long-range dependent processes ($`\alpha \overline{\tau }_p1`$), we can observe that the rate of convergence is always lower than $`n^{3/4}\mathrm{log}(n)^{3/4}`$ and that the dominant term $`n^{3/4}`$ is obtained when $`\alpha \overline{\tau }_p1`$.
We now propose a uniform Bahadur type representation of sample quantiles. Such a representation has an application in the study of trimmed-mean. For $`0<p_0p_1<1`$ consider the following assumption which extends $`𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}`$
$`𝑨𝒔𝒔𝒖𝒎𝒑𝒕𝒊𝒐𝒏𝑨_\mathrm{𝟓}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{,}𝒑_\mathrm{𝟏}\mathbf{)}\mathbf{:}`$ there exists $`U_i`$, $`i=1,\mathrm{},L`$, disjoint open sets such that $`U_i`$ contains a solution to the equation $`g(t)=\xi _{g(Y)}(p)`$ for all $`p_0pp_1`$, such that $`F_{g(Y)}^{}(\xi (p))>0`$ for all $`p_0pp_1`$ and such that $`g`$ is a $`𝒞^2`$diffeomorphism on $`_{i=1}^LU_i`$.
Under the previous assumption, define
$$\tau _{p_0,p_1}=\underset{\gamma _{i=1}^Lg(U_i)}{inf}\tau (\gamma ).$$
(33)
###### Theorem 3
Under the conditions of Theorem 2 and Assumption $`𝐀_\mathrm{𝟓}\mathbf{(}𝐩_\mathrm{𝟎}\mathbf{,}𝐩_\mathrm{𝟏}\mathbf{)}`$, we have almost surely, as $`n+\mathrm{}`$
$$\underset{p_0pp_1}{sup}\left|\widehat{\xi }(p;𝒈\mathbf{(}𝒀\mathbf{)})\xi _{g(Y)}(p)\frac{p\widehat{F}(\xi _{g(Y)}(p);𝒈\mathbf{(}𝒀\mathbf{)})}{f_{g(Y)}(\xi _{g(Y)}(p))}\right|=𝒪_{a.s.}\left(r_n(\alpha ,\tau _{p_0,p_1})\right).$$
(34)
###### Remark 6
To obtain convergence results of estimators of $`H`$, some results are needed concerning sample quantiles of the form $`\widehat{\xi }(p;𝐠\mathbf{(}𝐘^{𝐚^𝐦}\mathbf{)})`$, with $`g()=||`$. Lemma 14 asserts that the Hermite rank $`\tau _p`$ of the function $`h_{\xi _{g(Y)}(p)}()`$ with $`g()=||`$, is equal to 2 for all $`0<p<1`$. Moreover, for all $`0<p<1`$ and for all $`0<p_0p_1<1`$, Assumptions $`𝐀_\mathrm{𝟒}\mathbf{(}𝛏\mathbf{(}𝐩\mathbf{)}\mathbf{)}`$ and $`𝐀_\mathrm{𝟓}\mathbf{(}𝐩_\mathrm{𝟎}\mathbf{,}𝐩_\mathrm{𝟏}\mathbf{)}`$ are satisfied, and we have $`\overline{\tau }_p=\tau _{p_0,p_1}=2`$. Since from Lemma 1, the correlation function of $`𝐘^{𝐚^m}`$ satisfies (30) with $`\alpha =2\nu 2H`$ and $`L()=1`$, by applying Theorem 2, the sequence $`r_n(,)`$ is then given by
$$r_n(2\nu 2H,2)=n^{3/4}\mathrm{log}(n)^{3/4},\text{ if }\nu 2$$
(35)
and for $`\nu =1`$
$$r_n(22H,2)=\{\begin{array}{cc}n^{3/4}\mathrm{log}(n)^{3/4}\hfill & \text{if }0<H<3/4,\hfill \\ n^{3/4}\mathrm{log}(n)^{3/2}\hfill & \text{if }H=3/4,\hfill \\ n^{1/2(1H)}\mathrm{log}(n)\hfill & \text{if }3/4<H<5/6,\hfill \\ n^{2(22H)}\mathrm{log}(n)^2\hfill & \text{if }5/6H<1.\hfill \end{array}$$
(36)
### 4.2 Convergence results of estimators of $`H`$
In order to specify convergence results, we make the following assumption concerning the remainder term of the variance function $`v()`$.
$`𝑨𝒔𝒔𝒖𝒎𝒑𝒕𝒊𝒐𝒏𝑨_\mathrm{𝟔}\mathbf{(}𝜼\mathbf{)}:`$ there exists $`\eta >0`$ such that $`v(t)=\sigma ^2|t|^{2H}\left(1+𝒪\left(|t|^\eta \right)\right),`$ as $`|t|0.`$
The first result concentrates itself on estimators $`\widehat{H}^\alpha `$ and $`\widehat{H}^{\mathrm{log}}`$ based on a convex combination of sample quantiles.
###### Theorem 4
Under Assumptions $`𝐀_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}𝛎\mathbf{)}`$, $`𝐀_\mathrm{𝟐}\mathbf{(}𝐩\mathbf{,}𝐜\mathbf{)}`$ and $`𝐀_\mathrm{𝟔}\mathbf{(}𝛈\mathbf{)}`$,
$`(i)`$ we have almost surely, as $`n+\mathrm{}`$
$$\widehat{H}^\alpha H=\{\begin{array}{cc}𝒪\left(n^\eta \right)+𝒪_{a.s.}\left(n^{1/2}\mathrm{log}(n)\right)\hfill & \text{if }\nu >H+\frac{1}{4},\hfill \\ 𝒪\left(n^\eta \right)+𝒪_{a.s.}\left(n^{1/2}\mathrm{log}(n)^{3/2}\right)\hfill & \text{if }\nu =1,H=\frac{3}{4},\hfill \\ 𝒪\left(n^\eta \right)+𝒪_{a.s.}\left(n^{2(1H)}\mathrm{log}(n)\right)\hfill & \text{if }\nu =1,\frac{3}{4}<H<1.\hfill \end{array}$$
(37)
A similar result holds for $`\widehat{H}^{\mathrm{log}}`$.
$`(ii)`$ the mean squared errors (MSE) of $`\widehat{H}^\alpha `$ satisfies
$$\text{MSE}\left(\widehat{H}^\alpha H\right)=𝒪\left(v_n(2\nu 2H)\right)+𝒪\left(r_n(2\nu 2H,2)^2\right)+𝒪\left(n^{2\eta }\right).$$
(38)
The sequence $`r_n(2\nu 2H,2)`$ is given by (35) and (36) and the sequence $`v_n()`$ is defined by
$$v_n(2\nu 2H)=\{\begin{array}{cc}n^1\hfill & \text{if }\nu >H+\frac{1}{4},\hfill \\ n^1\mathrm{log}(n)\hfill & \text{if }\nu =1,H=\frac{3}{4},\hfill \\ n^{4(1H)}\hfill & \text{if }\nu =1,\frac{3}{4}<H<1.\hfill \end{array}$$
(39)
Again, the same result holds for $`MSE\left(\widehat{H}^{\mathrm{log}}H\right)`$.
$`(iii)`$ if the filter $`𝐚`$ is such that $`\nu >H+1/4`$, and if $`\eta >1/2`$, then we have the following convergence in distribution, as $`n+\mathrm{}`$
$$\sqrt{n}\left(\widehat{H}^\alpha H\right)𝒩(0,\sigma _\alpha ^2)\text{ and }\sqrt{n}\left(\widehat{H}^{\mathrm{log}}H\right)𝒩(0,\sigma _0^2),$$
(40)
where $`\sigma _\alpha ^2`$ is defined for $`\alpha 0`$ by
$$\sigma _\alpha ^2=\underset{i}{}\underset{j1}{}\frac{1}{(2j)!}\left(\underset{k=1}{\overset{K}{}}\frac{H_{2j1}(q_k)c_k}{q_k}\pi _k^\alpha \right)^2𝑩^T\underset{¯}{𝑹}(i,j)𝑩.$$
(41)
The vector $`𝐁`$ is defined by $`𝐁={\displaystyle \frac{𝐀^T}{𝐀^2}}`$, and the real numbers $`q_k`$ and $`\pi _k^\alpha `$ are defined by
$$q_k=\mathrm{\Phi }^1\left(\frac{1+p_k}{2}\right)\text{ and }\pi _k^\alpha =\frac{(q_k)^\alpha }{_{j=1}^Kc_j(q_j)^\alpha }.$$
(42)
Finally, the matrix $`\underset{¯}{𝐑}(i,j)`$, defined for $`i`$ and $`j1`$, is a $`M\times M`$ matrix whose $`(m_1,m_2)`$ entry is
$$\left(\underset{¯}{𝑹}(i,j)\right)_{m_1,m_2}=\rho ^{𝒂^{𝒎_\mathrm{𝟏}},𝒂^{𝒎_\mathrm{𝟐}}}(i)^{2j},$$
(43)
where $`\rho ^{𝐚^{𝐦_\mathrm{𝟏}},𝐚^{𝐦_\mathrm{𝟐}}}()`$ is the correlation function defined by (7).
###### Remark 7
The expression of the variance $`\sigma _\alpha ^2`$ given by (41) could appear to be very complicated. However, given some vectors $`𝐩`$ and $`𝐜`$ and some integer $`M`$, it does not take unreasonnable effort to compute it for each value of $`H`$ by truncating the two series. This issue is investigated in Section 5 to compare the different parameters.
###### Remark 8
Let us discuss the result (38). The first term, $`𝒪\left(v_n\right)`$, is due to the variance of the sample cumulative distribution function. The second term, $`𝒪\left(r_n^2\right)`$ is due to the departure of $`\widehat{\xi }\left(p\right)\xi (p)`$ from $`\widehat{F}\left(\xi (p)\right)p`$. We leave the reader to check that
$$𝒪\left(r_n(2\nu 2H,2)^2\right)+𝒪\left(v_n(2\nu 2H)\right)=\{\begin{array}{cc}𝒪\left(v_n(2\nu 2H)\right)\hfill & \text{ if }\nu H+\frac{1}{4},\hfill \\ 𝒪\left(r_n(2\nu 2H,2)^2\right)\hfill & \text{ if }\nu <H+\frac{1}{4}.\hfill \end{array}$$
Finally, the third one, $`𝒪\left(n^{2\eta }\right)`$ is a bias term due to the misspecification of the variance function $`v()`$ around 0.
###### Remark 9
If $`K=1`$, we have, for every $`\alpha >0`$,
$$\sigma _\alpha ^2=\sigma _0^2=\underset{i}{}\underset{j1}{}\frac{H_{2j1}(q)^2}{q^2(2j)!}𝑩^T\underset{¯}{𝑹}(i,j)𝑩.$$
Assume $`𝐀_\mathrm{𝟔}\mathbf{(}𝛈\mathbf{)}`$ with $`\eta >1/2`$ which allows to neglict the bias term with respect to the variance one. The result (40) is proved by using some general central limit theorem obtained in this dependence context by Arcones (1994), which is available as soon as $`\rho ^𝐚()^2`$ is summable. Therefore, if only $`𝐀_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}\mathbf{)}`$ is assumed, the filter $`𝐚`$ cannot exceed 1 (and then correspond to $`𝐚=(1,1)`$) and, due to (5), the result (40) is valid only for $`0<H<3/4`$. As a practical point of view, one observes that for such a filter and large values of $`H`$, the estimators have very big variance. Note that if $`𝐀_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}𝛎\mathbf{)}`$ can be assumed for $`\nu >1`$, then the asymptotic normality is valid for all the values of $`H`$.
The next result asserts the link between $`\widehat{H}^{\mathrm{log}}`$ and $`\widehat{H}^\alpha `$.
###### Corollary 5
Let $`(\alpha _n)_{n1}`$ be a sequence such that $`\alpha _n0`$, as $`n+\mathrm{}`$. Then, under conditions of Theorem 4 $`(ii)`$, the following convergence in distribution holds, as $`n+\mathrm{}`$
$$\sqrt{n}\left(\widehat{H}_n^{\alpha _n}H\right)𝒩(0,\sigma _0^2).$$
(44)
The following theorem presents the analog results obtained for the estimators $`\widehat{H}^{\alpha ,tm}`$ and $`\widehat{H}^{\mathrm{log},tm}`$ based on trimmed-means.
###### Theorem 6
Under Assumptions $`𝐀_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}𝛎\mathbf{)}`$, $`𝐀_\mathrm{𝟑}\mathbf{(}𝛃\mathbf{)}`$ and $`𝐀_\mathrm{𝟔}\mathbf{(}𝛈\mathbf{)}`$, properties $`(i)`$ and $`(ii)`$ of Theorem 4 hold for the estimator $`\widehat{H}^{\alpha ,tm}`$ and $`\widehat{H}^{\mathrm{log},tm}`$ with the same rates of convergences.
$`(iii)`$ if the filter $`𝐚`$ is such that $`\nu >H+1/4`$ and if $`\eta >1/2`$, then, under the notations of Theorem 4, we have the following convergence in distribution, as $`n+\mathrm{}`$
$$\sqrt{n}\left(\widehat{H}^{\alpha ,tm}H\right)𝒩(0,\sigma _{\alpha ,tm}^2)\text{ and }\sqrt{n}\left(\widehat{H}^{\mathrm{log},tm}H\right)𝒩(0,\sigma _{0,tm}^2),$$
(45)
where $`\sigma _{\alpha ,tm}^2`$ is defined for $`\alpha 0`$ by
$$\sigma _{\alpha ,tm}^2=\underset{i}{}\underset{j1}{}\frac{1}{(2j)!}\left(\frac{_{\beta _1}^{1\beta _2}H_{2j1}(q)q^{\alpha 1}𝑑p}{_{\beta _1}^{1\beta _2}q^\alpha 𝑑p}\right)^2𝑩^T\underset{¯}{𝑹}(i,j)𝑩,$$
(46)
with $`q=\mathrm{\Phi }^1\left(\frac{1+p}{2}\right)`$.
## 5 Numerical computation and simulations
### 5.1 Asymptotic constants $`\sigma _\alpha ^2`$ and $`\sigma _{\alpha ,tm}^2`$
In order to compare the different estimators, we intend to compute the asymptotic constants $`\sigma _\alpha ^2`$ and $`\sigma _{\alpha ,tm}^2`$ defined by (41) and (46) for various set of parameters ($`𝒂,𝒑,𝒄,𝜷,M`$). For this work, both series defining $`\sigma _\alpha ^2`$ and $`\sigma _{\alpha ,tm}^2`$ are truncated ($`|i|200`$, $`j150`$). Figure 2 illustrates a part of this work. We can propose the following general remarks:
$``$ Among all filters tested, the best one seems to be
$$𝒂^{}=\{\begin{array}{cc}inc1\hfill & \text{ if }0<H<3/4,\hfill \\ db4\hfill & \text{ otherwise.}\hfill \end{array}$$
where $`inc1`$ and $`db4`$ respectively denote the filter $`(1,1)`$ and the Daubechies wavelet filter with two zero moments explicitly given by
$$db4=(0.4829629,0.8365763,0.22414386,0.12940952).$$
$``$ Choice of $`M`$: increasing $`M`$ seems to reduce the asymptotic constant $`\sigma _\alpha ^2`$. Obviously, a too large $`M`$ increases the bias since $`\widehat{\xi }(𝒑,𝒄;𝒈\mathbf{(}𝑿^{𝒂^𝑴}\mathbf{)})`$ or $`\overline{𝒈\mathbf{(}𝑿^{𝒂^𝑴}\mathbf{)}}^{(𝜷)}`$ are estimated with $`NM\mathrm{}`$ observations. We recommend setting it to the value 5.
$``$ We did not manage (theoretically and numerically since series defining (41) and (46) are truncated) to determine the optimal value of $`\alpha `$. However, for examples considered, it should be near the value $`2`$.
$``$ Again, this is quite difficult to know theoretically and numerically which choice of $`𝒑`$ is optimal. What we observed is that, for fixed parameters $`𝒂`$, $`M`$ and $`\alpha `$, the asymptotic constants are very close to each other.
$``$ Choice of $`p`$ in the case of a single quantile (see Figure 2): the optimal $`p`$ seems to be near the value $`90\%`$. However, $`p=1/2`$, corresponding to the estimator based on the median, leads to good results.
$``$ Choice of $`\beta _1=\beta _2=\beta `$ for the estimators based on trimmed-means (see Figure 2): obviously the constant grows with $`\beta `$ but we can point out that estimators based on $`10\%`$trimmed-means are very competitive with the ones obtained by quadratic variations ($`\beta =0`$).
### 5.2 Simulation
A short simulation study is proposed in Table 1 and Figure 1 for $`n=1000`$ and $`H=0.8`$. We consider two locally self-similar Gaussian processes whose variance functions are in turn $`v(t)=|t|^{2H}`$ (fractional Brownian motion) and $`v(t)=1\mathrm{exp}(|t|^{2H})`$. To generate sample paths discretized over a grid $`[0,1]`$, we use the method of circulant matrix (see Wood and Chan (1994)), which is particularly fast, even for large sample sizes. Various versions of estimators are considered and compared with classical ones, that is the one based on quadratic variations, Coeurjolly (2001), and the Whittle estimator, Beran (1994). In order to illustrate the robustness of our estimators, we also applied them to contaminated version of sample path processes. We obtain a new sample path discretized at times $`i/n`$ and denoted by $`X^C(i/n)`$ for $`i=1,\mathrm{},n`$ through the following model
$$X^C(i/n)=X(i/n)+U(i)V(i),$$
(47)
where $`U(i)`$, $`i=1,\mathrm{},n`$ are Bernoulli independent variables $`(0.005)`$, and $`V(i)`$, $`i=1,\mathrm{},n`$ are independent centered Gaussian variables with variance $`\sigma _C^2(i)`$ such that the signal noise ratio at time $`i/n`$ is equal to 20 dB. As a general conclusion of Table 1, one can say that all versions of our estimators are very competitive with classical ones when the processes are observed without contamination and they seem to be particularly robust to additive outliers. Both bias and variance are approximately unchanged. This is clearly not the case for classical estimators. Indeed, concerning quadratic variations’ method, the estimation procedure is based on the estimation of $`𝐄((X^{𝒂^m}(1/n))^2)`$ by sample mean of order 2 of $`(𝑿^{𝒂^m})^2`$, Coeurjolly (2001)), that is particularly sensitive to additive outliers. Bad results of Whittle estimator can be explained by the fact that maximum likelihood methods are also non-robust methods.
## 6 Proofs
We denote by $`||||_{L^2(d\varphi )}`$ (resp. $`||||_\mathrm{}^q`$) the norm defined by $`h_{L^2(d\varphi )}=E(h(Y)^2)^{1/2}`$ for some measurable function $`h()`$ (resp. $`(_i|u_i|^q)^{1/2}`$ for some sequence $`(u_i)_i`$). In order to simplify the presentation of proofs, we use the notations $`F()`$, $`\xi ()`$, $`f()`$, $`\widehat{F}()`$ and $`\widehat{\xi }()`$ instead of $`F_{g(Y)}()`$, $`\xi _{g(Y)}()`$, $`f_{g(Y)}()`$, $`\widehat{F}_{g(Y)}(;𝒈\mathbf{(}𝒀\mathbf{)})`$ and $`\widehat{\xi }_{g(Y)}(;𝒈\mathbf{(}𝒀\mathbf{)})`$ respectively. For some real $`x`$, $`[x]`$ denotes the integer part of $`x`$. Finally, $`\lambda `$ denotes a generic positive constant.
### 6.1 Sketch of the proof of Theorem 2
We give here a brief explanation of the strategy to prove Theorem 2. This proof follows exactly the one proposed by Serfling (1980) in the i.i.d. case. One starts by writing
$$\frac{p\widehat{F}\left(\xi (p)\right)}{f(\xi (p))}\left(\widehat{\xi }\left(p\right)\xi (p)\right)=A(p)+B(p)+C(p),$$
with
$`A(p)`$ $`=`$ $`{\displaystyle \frac{p\widehat{F}\left(\widehat{\xi }\left(p\right)\right)}{f(\xi (p))}}`$ (48)
$`B(p)`$ $`=`$ $`{\displaystyle \frac{\widehat{F}\left(\widehat{\xi }\left(p\right)\right)\widehat{F}\left(\xi (p)\right)\left(F(\widehat{\xi }\left(p\right))F(\xi (p))\right)}{f(\xi (p))}}`$ (49)
$`C(p)`$ $`=`$ $`{\displaystyle \frac{F(\widehat{\xi }\left(p\right))F(\xi (p))}{f(\xi (p))}}\left(\widehat{\xi }\left(p\right)\xi (p)\right).`$ (50)
From the definition of sample quantile, we have almost surely, see e.g. Serfling (1980), $`A(p)=𝒪_{a.s.}\left(n^1\right)`$. Now, in order to control the term $`C(p)`$, Taylor’s Theorem is used and a control of $`\widehat{\xi }\left(p\right)\xi (p)`$ is needed. The latter one is done by Lemma 10 which exhibits the sequence $`\epsilon _n(\alpha ,\tau _p)`$ such that $`\widehat{\xi }\left(p\right)\xi (p)=𝒪_{a.s.}\left(\epsilon _n(\alpha ,\tau _p)\right)`$. Then, in order to control $`B(p)`$ it is sufficient to control the random variable
$$S_n(\xi (p),\epsilon _n(\alpha ,\tau _p))=\underset{|x|\epsilon _n(\alpha ,\tau _p)}{sup}\left|\mathrm{\Delta }(\xi (p)+x)\mathrm{\Delta }(\xi (p))\right|,$$
with $`\mathrm{\Delta }()=\widehat{F}()F()`$. This result is detailed in Lemma 11. In order to specify the rate explicited by Theorem 2, we present and prove Lemmas 10 and 11. Some preliminary results, given by Lemma 7, Corollary 8 and Lemma 9, are needed. Among other things, Lemma 7 and Corollary 8 propose some inequalities for controlling the sample mean of non-linear function of Gaussian sequences with correlation function satisfying (30).
### 6.2 Auxiliary Lemmas for the proof of Theorem 2
###### Lemma 7
Let $`\left\{Y(i)\right\}_{i=1}^+\mathrm{}`$ a gaussian stationary process with variance 1 and correlation function $`\rho ()`$ such that, as $`i+\mathrm{}`$, $`|\rho (i)|L(i)i^\alpha `$, for some $`\alpha >0`$ and some slowly varying function at infinity $`L()`$. Let $`h()L^2\left(d\varphi \right)`$ and denote by $`\tau `$ its Hermite rank. Define
$$\overline{Y}_n=\frac{1}{n}\underset{i=1}{\overset{n}{}}h(Y(i)).$$
Then, for all $`\gamma >0`$, there exists a positive constant $`\kappa _\gamma =\kappa _\gamma (\alpha ,\tau )`$, such that
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)=𝒪\left(n^\gamma \right),$$
(51)
with
$$y_n=y_n(\alpha ,\tau )=\{\begin{array}{ccc}n^{1/2}\mathrm{log}(n)^{1/2}\hfill & \text{if }\hfill & \alpha \tau >1,\hfill \\ n^{1/2}\mathrm{log}(n)^{1/2}L_\tau (n)^{1/2}\hfill & \text{if }\hfill & \alpha \tau =1,\hfill \\ n^{\alpha \tau /2}\mathrm{log}(n)^{\tau /2}L(n)^{\tau /2}\hfill & \text{if }\hfill & 0<\alpha \tau <1.\hfill \end{array}$$
(52)
where $`L_\tau (n)=_{|i|n}|\rho (i)|^\tau `$. In the case $`\alpha \tau =1`$, we assume that for all $`j>\tau `$, the limit, $`lim_{n+\mathrm{}}L_\tau (n)^1_{|i|n}|\rho (i)|^j`$ exists.
Proof. Let $`(y_n)_{n1}`$ be the sequence defined by (52). The proof is splitted into three parts according to the value of $`\alpha \tau `$.
Case $`\alpha \tau \mathbf{<}\mathrm{𝟏}`$ : From Chebyshev’s inequality, we have for all $`q1`$
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)\frac{1}{\kappa _\gamma ^{2q}y_n^{2q}}𝐄\left(\left(\overline{Y}_n\right)^{2q}\right).$$
From Theorem 1 of Breuer and Major (1983) and in particular Equation (2.6), we have, as $`n+\mathrm{}`$
$$𝐄\left(\left(\overline{Y}_n\right)^{2q}\right)\frac{(2q)!}{2^qq!}\frac{1}{n^q}\sigma ^{2q},\text{ with }\sigma ^2=\underset{i}{}\underset{j\tau }{}\frac{(c_j)^2}{j!}\rho (i)^j,$$
(53)
where $`c_j`$ denotes the $`j`$-th Hermite coefficient of $`h()`$. Note that $`\sigma ^2h_{L^2(d\varphi )}^2\rho _\mathrm{}^\tau ^2`$. Thus, for $`n`$ large enough, we have
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)\frac{\lambda }{n^qy_n^{2q}}\frac{(2q)!}{2^qq!}\left(h_{L^2(d\varphi )}^2\rho _\mathrm{}^\tau ^2\kappa _\gamma ^2\right)^q.$$
(54)
From Stirling’s formula, we have as $`q+\mathrm{}`$
$$\frac{(2q)!}{2^qq!}\sqrt{2}q^q(2e^1)^q.$$
(55)
From (52) by choosing $`q=[\mathrm{log}(n)]`$, (54) becomes
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)\lambda \left(2e^1h_{L^2(d\varphi )}^2\rho _\mathrm{}^\tau ^2\kappa _\gamma ^2\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),$$
if $`\kappa _\gamma ^2>2h_{L^2(d\varphi )}^2\rho _\mathrm{}^\tau ^2\mathrm{exp}(\gamma 1)`$.
Case $`\alpha \tau \mathbf{=}\mathrm{𝟏}`$ : Using the proof of Theorem $`1^{}`$ of Breuer and Major (1983), we can prove that for all $`q1`$
$`𝐄\left(\left(n^{1/2}L_\tau (n)^{1/2}\overline{Y}_n\right)^{2q}\right)`$ $``$ $`\lambda {\displaystyle \frac{2q!}{2^qq!}}𝐄\left(\left(n^{1/2}L_\tau (n)^{1/2}\overline{Y}_n\right)^2\right)^q`$ (56)
$``$ $`\lambda {\displaystyle \frac{2q!}{2^qq!}}\left({\displaystyle \underset{j\tau }{}}{\displaystyle \frac{(c_j)^2}{j!}}\underset{n+\mathrm{}}{lim}L_\tau (n)^1{\displaystyle \underset{|i|n}{}}|\rho (i)|^j\right)^q`$
$``$ $`\lambda {\displaystyle \frac{2q!}{2^qq!}}h_{L^2(d\varphi )}^{2q}.`$
Then from Chebyshev’s inequality, we have for all $`q1`$
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)\lambda \frac{L_\tau (n)^q}{n^qy_n^{2q}}\frac{2q!}{2^qq!}\left(h_{L^2(d\varphi )}^2\kappa _\gamma ^2\right)^q.$$
From (52) by choosing $`q=[\mathrm{log}(n)]`$, we obtain
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)\lambda \left(2e^1h_{L^2(d\varphi )}^2\kappa _\gamma ^2\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),$$
if $`\kappa _\gamma ^2>2h_{L^2(d\varphi )}^2\times \mathrm{exp}(\gamma 1)`$.
Case $`\alpha \tau \mathbf{<}\mathrm{𝟏}`$ : Denote by $`k_\alpha `$ the lowest integer satisfying $`k_\alpha \alpha >1`$, that is $`k_\alpha =[1/\alpha ]+1`$, and for $`j\tau `$ denote by $`Z_j`$ the following random variable
$$Z_j=\frac{1}{n}\underset{i=1}{\overset{n}{}}\frac{c_j}{j!}H_j(Y(i)).$$
Denote by $`\kappa _{1,\gamma }`$ and $`\kappa _{2,\gamma }`$ two positive constants such that $`\kappa _\gamma =\mathrm{max}(\kappa _{1,\gamma },\kappa _{2,\gamma })`$. From the triangle inequality,
$$\left(|\overline{Y}_n|\kappa _\gamma y_n\right)\left(|\overline{Y}_n\underset{j=\tau }{\overset{k_\alpha 1}{}}Z_j|\kappa _{1,\gamma }y_n\right)+\underset{j=\tau }{\overset{k_\alpha 1}{}}\left(|Z_j|\kappa _{2,\gamma }y_n\right)$$
(57)
Since
$$\overline{Y}_n\underset{j=\tau }{\overset{k_\alpha 1}{}}Z_j=\frac{1}{n}\underset{i=1}{\overset{n}{}}\underset{jk_\alpha }{}\frac{c_j}{j!}H_j(Y(i))=\frac{1}{n}\underset{i=1}{\overset{n}{}}h^{}(Y(i)),$$
where $`h^{}()`$ is a function with Hermite rank $`k_\alpha `$. Applying Lemma 7 in the case $`\alpha \tau >1`$, it follows that, for all $`\gamma >0`$, there exists a constant $`\kappa _{1,\gamma }`$ such that, for $`n`$ large enough
$$\left(|\overline{Y}_n\underset{j=\tau }{\overset{k_\alpha 1}{}}Z_j|\kappa _{1,\gamma }y_n\right)=𝒪\left(n^\gamma \right).$$
(58)
Now, let $`\tau j<k_\alpha `$ and $`q1`$, from Theorem 3 of Taqqu (1977), we have
$`\left(|Z_j|\kappa _{2,\gamma }y_n\right)`$ $``$ $`{\displaystyle \frac{1}{\kappa _{2,\gamma }^{2q}y_n^{2q}}}\left({\displaystyle \frac{c_j}{j!}}\right)^{2q}n^{2q}𝐄\left({\displaystyle \underset{i_1,\mathrm{},i_{2q}}{}}H_j(Y(i_1))\mathrm{}H_j(Y(i_{2q}))\right)`$ (59)
$``$ $`\lambda {\displaystyle \frac{L(n)^{jq}}{n^{\alpha jq}y_n^{2q}}}\left({\displaystyle \frac{c_j}{j!}}\kappa _{2,\gamma }^1\right)^{2q}\mu _{2q},`$
where $`\mu _{2q}`$ is a constant such that $`\mu _{2q}\left(\frac{2}{1\alpha j}\right)^q𝐄\left(H_j(Y)^{2q}\right)`$. It is also proved in Taqqu (1977) (p. 228), that $`𝐄\left(H_j(Y)^{2q}\right)(2jq)!/(2^{jq}(jq)!),`$ as $`q+\mathrm{}`$. Thus, from Stirling’s formula, we obtain as $`q+\mathrm{}`$
$$(|Z_j|y_n)\lambda \frac{L(n)^{(j\tau )q}}{n^{\alpha (j\tau )q}}\mathrm{log}(n)^{\tau q}q^{jq}\left(\frac{2}{1\alpha j}\left(\frac{c_j}{j!}\right)^2\left(\frac{2j}{e}\right)^j\kappa _{2,\gamma }^1\right)^q.$$
By choosing $`q=[\mathrm{log}(n)]`$, we finally obtain, as $`n+\mathrm{}`$
$$\underset{j=\tau }{\overset{k_\alpha 1}{}}\left(|Z_j|\kappa _{2,\gamma }y_n\right)\lambda \left(\frac{2}{1\alpha \tau }\left(\frac{c_\tau }{\tau !}\right)^2\left(\frac{2\tau }{e}\right)^\tau \kappa _{2,\gamma }^2\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),$$
(60)
if $`\kappa _{2,\gamma }^2>\frac{2}{1\alpha \tau }\left(\frac{c_\tau }{\tau !}\right)^2(2\tau )^\tau \mathrm{exp}(\gamma \tau )`$. From (57), we get the result by combining (58) and (60).
###### Corollary 8
Under conditions of Lemma 7, for all $`\alpha >0`$, $`j1`$ and $`\gamma >0`$, there exists $`q=q(\gamma )1`$ and $`\zeta _\gamma >0`$ such that
$$𝐄\left(\left\{\frac{1}{n}\underset{i=1}{\overset{n}{}}H_j(Y(i))\right\}^{2q}\right)\zeta _\gamma n^\gamma .$$
(61)
Proof. (53), (56) and (59) imply that there exists $`\lambda =\lambda (q)>0`$ such that for all $`q1`$, we have
$`𝐄\left(\left\{{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}H_j(Y(i))\right\}^{2q}\right)\lambda (q)n^q`$ $`=`$ $`\lambda (q)\times \{\begin{array}{cc}n^q\hfill & \text{ if }\alpha j>1\hfill \\ L_{\tau _p}(n)n^q\hfill & \text{ if }\alpha j=1\hfill \\ L(n)^{\alpha jq}n^{\alpha jq}\hfill & \text{ if }\alpha j<1\hfill \end{array}`$ (65)
$`=`$ $`𝒪\left(n^\gamma \right).`$ (66)
Indeed, it is sufficient to choose $`q`$ such that, $`q>\gamma `$ if $`\alpha j1`$ and $`q>\gamma /\alpha j`$ if $`\alpha j<1`$.
###### Lemma 9
Let $`0<p<1`$, denote by $`g()`$ a function satisfying Assumption $`𝐀_\mathrm{𝟒}\mathbf{(}𝛏\mathbf{(}𝐩\mathbf{)}\mathbf{)}`$ and by $`(x_n)_{n1}`$ a sequence with real components, such that $`x_n0`$, as $`n+\mathrm{}`$. Then, for all $`j1`$, there exists a positive constant $`d_j=d_j(\xi (p))<+\mathrm{}`$ such that, for $`n`$ large enough
$$\left|c_j(\xi (p)+x_n)c_j(\xi (p))\right|d_j|x_n|.$$
(67)
Proof. Let $`j1`$, under Assumption $`𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}`$, for $`n`$ large enough, $`\xi (p)+x_n_{i=1}^Lg(U_i)`$. Thus, for $`n`$ large enough,
$`c_j(\xi (p)+x_n)c_j(\xi (p))`$ $`=`$ $`{\displaystyle _{}}\left(h_{\xi (p)+x_n}(t)h_{\xi (p)}(t)\right)H_j(t)\varphi (t)𝑑t`$ (70)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle _{U_i}}\left(\mathrm{𝟏}_{g_i(t)\xi (p)+x_n}\mathrm{𝟏}_{g_i(t)\xi (p)}\right)H_j(t)\varphi (t)𝑑t`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{L}{}}}{\displaystyle _{m_{i,n}}^{M_{i,n}}}(1)^j\varphi ^{(j)}(t)𝑑t,`$
$`=`$ $`\{\begin{array}{cc}_{i=1}^L\left(\varphi (M_{i,n})\varphi (m_{i,n})\right)\hfill & \text{if }j=1,\hfill \\ _{i=1}^L(1)^j\left(\varphi ^{(j1)}(M_{i,n})\varphi ^{(j1)}(m_{i,n})\right)\hfill & \text{if }j>1,\hfill \end{array}`$
where $`g_i()`$ is the restriction of $`g()`$ to $`U_i`$, and where $`m_{i,n}`$ (resp. $`M_{i,n}`$) is the minimum (resp. maximum) between $`g_i^1(\xi (p)+x_n)`$ and $`g_i^1(\xi (p))`$. We leave the reader to check that there exists a positive constant $`d_j`$, such that, for $`n`$ large enough
$$\left|c_j(\xi (p)+x_n)c_j(\xi (p))\right|d_j|x_n|\times \{\begin{array}{cc}_{i=1}^L\left|\varphi ^{(j)}(g_i^{(1)}(u))\left(g_i^{(1)}\right)^{}(u)\right|\hfill & \text{if }j=1,2\hfill \\ _{i=1}^L\left|\varphi ^{(j2)}(g_i^{(1)}(u))\left(g_i^{(1)}\right)^{}(u)\right|\hfill & \text{if }j>2,\hfill \end{array}$$
which is the desired result.
###### Lemma 10
Under conditions of Theorem 2, there exists a constant denoted by $`\kappa _\epsilon =\kappa _\epsilon (\alpha ,\tau _p)`$, such that, we have almost surely, as $`n+\mathrm{}`$
$$\left|\widehat{\xi }(p;𝒈\mathbf{(}𝒀\mathbf{)})\xi _{g(Y)}(p)\right|\epsilon _n,$$
(71)
where $`\epsilon _n=\epsilon _n(\alpha ,\tau (\xi (p)))=\kappa _\epsilon y_n(\alpha ,\tau (\xi (p))`$, $`y_n(,)`$ being defined by (52).
Proof. We have
$$\left(\left|\widehat{\xi }\left(p\right)\xi (p)\right|\epsilon _n\right)=\left(\widehat{\xi }\left(p\right)\xi (p)\epsilon _n\right)+\left(\widehat{\xi }\left(p\right)\xi (p)+\epsilon _n\right).$$
(72)
Using Lemma 1.1.4 $`(iii)`$ of Serfling (1980), we have
$$\left(\widehat{\xi }\left(p\right)\xi (p)\epsilon _n\right)\left(\widehat{F}\left(\xi (p)\epsilon _n\right)p\right).$$
(73)
Under Assumption $`𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}`$, for $`n`$ large enough
$$pF(\xi (p)\epsilon _n)=f(\xi (p))\epsilon _n+o\left(\epsilon _n\right)\frac{f(\xi (p))}{2}\epsilon _n.$$
Consequently, for $`n`$ large enough and from (73)
$$\left(\widehat{\xi }\left(p\right)\xi (p)\epsilon _n\right)\left(\widehat{F}\left(\xi (p)\epsilon _n\right)F(\xi (p)\epsilon _n)\frac{f(\xi (p))}{2}\epsilon _n\right).$$
(74)
Define $`\tau _{p,n}=\tau (\xi (p)\epsilon _n)`$, from Lemma 9, we have for $`n`$ large enough
$$\widehat{F}\left(\xi (p)\epsilon _n\right)F(\xi (p)\epsilon _n)2\left(\widehat{F}\left(\xi (p)\right)F(\xi (p))\right)+2\epsilon _n\underset{jJ_n}{}Z_{n,j},$$
(75)
where
$$J_n=\{\begin{array}{cc}\{\tau _p<j\tau _{p,n}\}\hfill & \text{ if }\tau _{p,n}>\tau _p,\hfill \\ \mathrm{}\hfill & \text{ if }\tau _{p,n}=\tau _p,\hfill \\ \{\tau _{p,n}j<\tau _p\}\hfill & \text{ if }\tau _{p,n}<\tau _p.\hfill \end{array}\text{ and }Z_{n,j}=\frac{1}{n}\underset{i=1}{\overset{n}{}}\frac{d_j}{j!}H_j(Y(i)).$$
Now, define $`c_\epsilon =\kappa _\epsilon f(\xi (p))/4`$. Let $`\gamma >0`$, (61) implies that there exists $`q1`$ such that, for $`n`$ large enough
$`\left(|2\epsilon _nZ_n|{\displaystyle \frac{f(\xi (p))}{2}}\epsilon _n\right)`$ $``$ $`{\displaystyle \underset{jJ_n}{}}\left(|Z_{n,j}|>c_\epsilon \right)`$ (76)
$``$ $`{\displaystyle \underset{jJ_n}{}}{\displaystyle \frac{1}{c_\epsilon ^{2q}}}𝐄\left(Z_{n,j}^{2q}\right)=𝒪\left(n^\gamma \right).`$
Let us fix $`\gamma =2`$. From (74), (75) and (76) and from Lemma 7 (applied to the function $`h_{\xi (p)}()`$), we obtain
$$\left(\widehat{\xi }\left(p\right)\xi (p)\epsilon _n\right)\left(|\widehat{F}\left(\xi (p)\right)F(\xi (p))|c_\epsilon \epsilon _n\right)+𝒪\left(n^2\right)=𝒪\left(n^2\right),$$
if $`c_\epsilon >\kappa _2`$ that is if $`\kappa _\epsilon >4/f(\xi (p))\kappa _2`$.
Let us now focus on the second right-hand term of (72). Following the sketch of this proof, we may also obtain, for $`n`$ large enough
$$\left(\widehat{\xi }\left(p\right)\xi (p)+\epsilon _n\right)=𝒪\left(n^2\right),$$
if $`\kappa _\epsilon >4/f(\xi (p))\kappa _2`$. Thus, for $`n`$ large enough $`\left(\left|\widehat{\xi }\left(p\right)\xi (p)\right|\epsilon _n\right)=𝒪\left(n^2\right)`$, which leads to the result thanks to Borel-Cantelli’s Lemma.
The following Lemma is an analogous result obtained by Bahadur in the i.i.d. framework, see Lemma E p.97 of Serfling (1980).
###### Lemma 11
Under conditions of Theorem 2, denote by $`\mathrm{\Delta }(z)`$ for $`z`$ the random variable, $`\mathrm{\Delta }(z)=\widehat{F}(z;𝐠\mathbf{(}𝐘\mathbf{)})F_{g(Y)}(z)`$. Then, we have almost surely, as $`n+\mathrm{}`$
$$S_n(\xi _{g(Y)}(p),\epsilon _n(\alpha ,\tau _p))=\underset{|x|\epsilon _n}{sup}\left|\mathrm{\Delta }(\xi _{g(Y)}(p)+x)\mathrm{\Delta }(\xi _{g(Y)}(p))\right|=𝒪_{a.s.}\left(r_n(\alpha ,\overline{\tau }_p)\right),$$
(77)
where $`\epsilon _n=\epsilon _n(\alpha ,\tau _p)`$ is defined by (71) and $`r_n(\alpha ,\overline{\tau }_p)`$ is defined by (32).
Proof. Put $`\epsilon _n=\epsilon _n(\alpha ,\tau _p)`$ and $`r_n=r_n(\alpha ,\overline{\tau }_p)`$. Denote by $`(\beta _n)_{n1}`$ and $`(\eta _{b,n})_{n1}`$ the following two sequences
$$\beta _n=\left[n^{3/4}\epsilon _n\right]\text{ and }\eta _{b,n}=\xi (p)+\epsilon _n\frac{b}{\beta _n},$$
for $`b=\beta _n,\mathrm{},\beta _n`$. Using the monotonicity of $`F()`$ and $`\widehat{F}()`$, we have,
$$S_n(\xi (p),\epsilon _n)\underset{\beta _nb\beta _n}{\mathrm{max}}|M_{b,n}|+G_n,$$
(78)
where $`M_{b,n}=\mathrm{\Delta }(\eta _{b,n})\mathrm{\Delta }(\xi (p))`$ and $`G_n=\mathrm{max}_{\beta _nb\beta _n1}\left(F(\eta _{b+1,n})F(\eta _{b,n})\right).`$ Under Assumption $`𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}`$, we have for $`n`$ large enough
$$G_n\left(\eta _{b+1,n}\eta _{b,n}\right)\times \underset{|x|\epsilon _n}{sup}f(\xi (p)+x)=𝒪\left(n^{3/4}\right).$$
(79)
The proof is finished if one can prove that for all $`\gamma >0`$ (in particular $`\gamma =2`$) and for all $`b`$, there exists $`\kappa _\gamma ^{}`$ such that
$$\left(|M_{b,n}|\kappa _\gamma ^{}r_n\right)=𝒪\left(n^\gamma \right).$$
(80)
Indeed, since $`\beta _n=𝒪\left(n^{1/2+\delta }\right)`$ for all $`\delta >0`$, if (80) is true, then we have
$`\left(\underset{\beta _nb\beta _n}{\mathrm{max}}|M_{b,n}|\kappa _2^{}r_n(\alpha ,\tau _p)\right)`$ $``$ $`(2\beta _n+1)\times \underset{\beta _nb\beta _n}{\mathrm{max}}\left(|M_{b,n}|\kappa _2^{}r_n\right)`$
$`=`$ $`𝒪\left(n^{3/2+\delta }\right).`$
Thus, from Borel-Cantelli’s Lemma, we have, almost surely
$$\underset{\beta _nb\beta _n}{\mathrm{max}}|M_{b,n}|=𝒪_{a.s.}\left(r_n\right)$$
And so, from (78) and (79).
$$S_n(\xi (p),\epsilon _n)=𝒪_{a.s.}\left(r_n\right)+𝒪\left(n^{3/4}\right)=𝒪_{a.s.}\left(r_n\right),$$
(81)
which is the stated result.
So, the rest of the proof is devoted to prove (80). For the sake of simplicity, denote by $`h_n^{}()`$ the function $`h_{\eta _{b,n}}()h_{\xi (p)}()`$. For $`n`$ large enough, the Hermite rank of $`h_n^{}()`$ is at least equal to $`\overline{\tau }_p`$, that is defined by (29). In the sequel, we need the following bound for $`h_n^{}_{L^2(d\varphi )}^2`$
$$h_n^{}_{L^2(d\varphi )}^2=𝐄(h_n^{}(Y)^2)=\omega _n(1\omega _n)\text{ with }\omega _n=\left|F_{g(Y)}(\eta _{b,n})F_{g(Y)}(\xi (p))\right|.$$
As previously, we have $`\omega _n=𝒪(\epsilon _n)`$ and so, there exists $`\zeta >0`$, such that
$$h_n^{}_{L^2(d\varphi )}^2\zeta \epsilon _n.$$
(82)
From now on, in order to simplify the proof, we use the following upper-bound
$$\epsilon _n=\epsilon _n(\alpha ,\tau _p)\epsilon _n(\alpha ,\overline{\tau }_p),$$
and with a slight abuse, we still denote $`\epsilon _n=\epsilon _n(\alpha ,\overline{\tau }_p)`$. Note also, that from Lemma 9, the $`j`$-th Hermite coefficient, for some $`j\overline{\tau }_p`$, is given by $`c_j(\eta _{b,n})c_j(\xi (p))`$. And there exists a positive constant $`d_j=d_j(\xi (p))`$ such that for $`n`$ large enough
$$|c_j(\eta _{b,n})c_j(\xi (p)|d_j\epsilon _n\frac{|b|}{\beta _n}d_j\epsilon _n.$$
(83)
We now proceed like in the proof of Lemma 7.
Case $`\alpha \overline{\tau }_p\mathbf{>}\mathrm{𝟏}`$: using Theorem 1 of Breuer and Major (1983) and (54), we can obtain for all $`q1`$
$$\left(|M_{b,n}|\kappa _\gamma ^{}r_n\right)\lambda \frac{1}{n^qr_n^{2q}}\frac{(2q)!}{2^qq!}\frac{1}{(\kappa _\gamma ^{})^{2q}}h_n^{}_{L^2(d\varphi )}^{2q}\rho _{\mathrm{}^{\overline{\tau }_p}}^{2q}.$$
(84)
As $`q+\mathrm{}`$, we get
$$\left(|M_{b,n}|\kappa _\gamma ^{}r_n\right)\lambda \frac{\epsilon _n^q}{n^qr_n^{2q}}q^q\left(2\zeta e^1\rho _{\mathrm{}^{\overline{\tau }_p}}^2\frac{1}{(\kappa _\gamma ^{})^2}\right)^q.$$
From (32), (52) (with $`\tau =\overline{\tau }_p`$) and by choosing $`q=[\mathrm{log}(n)]`$, we have
$$\left(|M_{b,n}|\kappa _\gamma ^{}r_n\right)\lambda \left(2\zeta \kappa _\epsilon e^1\rho _{\mathrm{}^{\overline{\tau }_p}}^2\frac{1}{(\kappa _\gamma ^{})^2}\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),$$
(85)
if $`\kappa _{\gamma }^{}{}_{}{}^{2}>2\zeta \kappa _\epsilon \rho _{\mathrm{}^{\tau _p}}^2\mathrm{exp}(\gamma 1)`$.
Case $`\alpha \overline{\tau }_p\mathbf{=}\mathrm{𝟏}`$ from (56), we can obtain for all $`q1`$
$`𝐄\left(M_{b,n}^{2q}\right)`$ $``$ $`\lambda {\displaystyle \frac{(2q)!}{2^qq!}}{\displaystyle \frac{L_{\overline{\tau }_p}(n)^q}{n^q}}h_n^{}_{L^2(d\varphi )}^{2q}\lambda \zeta ^q{\displaystyle \frac{(2q)!}{2^qq!}}{\displaystyle \frac{L_{\overline{\tau }_p}(n)^q\epsilon _n^q}{n^q}}`$
$``$ $`\lambda {\displaystyle \frac{L_{\overline{\tau }_p}(n)^q\epsilon _n^q}{n^q}}(2\zeta e^1)^qq^q.`$
From (32), (52) (with $`\tau =\overline{\tau }_p`$), by choosing $`q=[\mathrm{log}(n)]`$, we have
$`\left(|M_{b,n}|\kappa _\gamma ^{}r_n\right)`$ $``$ $`{\displaystyle \frac{1}{\kappa _{\gamma }^{}{}_{}{}^{2q}r_n^{2q}}}𝐄\left(M_{b,n}^{2q}\right)`$
$``$ $`\lambda \left(2\zeta \kappa _\epsilon e^1{\displaystyle \frac{d_{\tau _p}^2}{\tau _p!}}{\displaystyle \frac{1}{\kappa _{\gamma }^{}{}_{}{}^{2}}}\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),`$
if $`\kappa _{\gamma }^{}{}_{}{}^{2}>2\zeta \kappa _\epsilon d_{\tau _p}^2/\tau _p!\mathrm{exp}(\gamma 1)`$.
Case $`\alpha \overline{\tau }_p\mathbf{<}\mathrm{𝟏}`$: denote by $`(r_{1,n})_{n1}`$ and by $`(r_{2,n})_{n1}`$ the following two sequences
$$r_{1,n}=n^{1/2\alpha \overline{\tau }_p/4}\mathrm{log}(n)^{\overline{\tau }_p/4+1/2}L(n)^{\overline{\tau }_p/4}\text{ and }r_{2,n}=n^{\alpha \overline{\tau }_p}\mathrm{log}(n)^{\overline{\tau }_p}L(n)^{\overline{\tau }_p}.$$
(86)
Note that $`\mathrm{max}(r_{1,n},r_{2,n})`$ is equal to $`r_{1,n}`$, when $`2/3<\alpha \overline{\tau }_p<1`$ and to $`r_{2,n}`$, when $`0<\alpha \overline{\tau }_p2/3`$. So, in order to obtain (80) in the case $`0<\alpha \overline{\tau }_p<1`$, it is sufficient to prove that there exists $`\kappa _\gamma ^{}`$ such that, for $`n`$ large enough
$$\left(|M_{b,n}|\kappa _\gamma ^{}\mathrm{max}(r_{1,n},r_{2,n})\right)=𝒪\left(n^\gamma \right).$$
Denote by $`k_\alpha `$ the integer $`[1/\alpha ]+1`$ for which $`\alpha k_\alpha >1`$, and by $`Z_{j,n}`$ for $`\overline{\tau }_pj<k_\alpha `$ the random variable defined by
$$Z_{j,n}=\frac{1}{n}\underset{i=1}{\overset{n}{}}\frac{c_j(\eta _{b,n})c_j(\xi (p))}{j!}H_j(Y(i)).$$
From the triangle inequality, we have
$$\left(|M_{b,n}|\kappa _\gamma ^{}\mathrm{max}(r_{1,n},r_{2,n})\right)\left(\left|M_{b,n}\underset{j=\overline{\tau }_p}{\overset{k_\alpha 1}{}}Z_{j,n}\right|\kappa _\gamma ^{}r_{1,n}\right)+\underset{j=\overline{\tau }_p}{\overset{k_\alpha 1}{}}\left(|Z_{j,n}|\kappa _\gamma ^{}r_{2,n}\right).$$
(87)
Since,
$$M_{b,n}\underset{j=\tau _p}{\overset{k_\alpha 1}{}}Z_{j,n}=\frac{1}{n}\underset{i=1}{\overset{n}{}}\underset{jk_\alpha }{}\frac{c_j(\eta _{b,n})c_j(\xi (p))}{j!}H_j(Y(i))=\frac{1}{n}\underset{i=1}{\overset{n}{}}h_n^{\prime \prime }(Y(i)),$$
where $`h_n^{\prime \prime }()`$ is a function with Hermite rank $`k_\alpha `$, such that $`\alpha k_\alpha >1`$, we have from (84)
$$\left(\left|M_{b,n}\underset{j=\overline{\tau }_p}{\overset{k_\alpha 1}{}}Z_{j,n}\right|\kappa _\gamma ^{}r_{1,n}\right)\lambda \frac{1}{n^qr_{1,n}^{2q}}h_n^{}_{L^2(d\varphi )}^{2q}\frac{(2q)!}{2^qq!}\frac{1}{\kappa _{\gamma }^{}{}_{}{}^{2q}}\rho _{\mathrm{}^{k_\alpha }}^{2q}$$
(88)
for all $`q1`$. From (82), we obtain, as $`q+\mathrm{}`$
$$\left(\left|M_{b,n}\underset{j=\tau _p}{\overset{k_\alpha 1}{}}Z_{j,n}\right|\kappa _\gamma ^{}r_{1,n}\right)\lambda \frac{\epsilon _n^q}{n^qr_{1,n}^{2q}}q^q\left(2\zeta e^1\rho _{\mathrm{}^{k_\alpha }}^2\kappa _{\gamma }^{}{}_{}{}^{2}\right)^q.$$
From (52) (with $`\tau =\overline{\tau }_p`$), (86) and by choosing $`q=[\mathrm{log}(n)]`$, we obtain
$$\left(\left|M_{b,n}\underset{j=\tau _p}{\overset{k_\alpha 1}{}}Z_{j,n}\right|\kappa _\gamma ^{}r_{1,n}\right)\lambda \left(2\zeta e^1\rho _{\mathrm{}^{k_\alpha }}^2\kappa _\epsilon \kappa _{\gamma }^{}{}_{}{}^{2}\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),$$
(89)
if $`\kappa _{\gamma }^{}{}_{}{}^{2}>\kappa _{1,\gamma }^{}=2\zeta \rho _{\mathrm{}^{k_\alpha }}^2\kappa _\epsilon \mathrm{exp}(\gamma 1)`$. Now, concerning the last term of (87), from (59), we can prove, for all $`\overline{\tau }_pj<k_\alpha `$
$$\left(Z_{j,n}\kappa _\gamma ^{}r_{2,n}\right)\lambda \frac{L(n)^{jq}}{n^{\alpha jq}r_{2,n}^{2q}}\frac{1}{\kappa _{\gamma }^{}{}_{}{}^{2q}}\left(\frac{c_j(\eta _{b,n})c_j(\xi (p))}{j!}\right)^{2q}\mu _{2q},$$
where $`\mu _{2q}`$ is a constant such that, as $`q+\mathrm{}`$,
$$\mu _{2q}\lambda \left(\frac{2}{1\alpha j}\right)^q\frac{(2jq)!}{2^{jq}(jq)!}.$$
From (83), we have, as $`q+\mathrm{}`$
$$\left(Z_{j,n}\kappa _\gamma ^{}r_{2,n}\right)\lambda \frac{\epsilon _n^{2q}L(n)^{jq}}{n^{\alpha jq}r_{2,n}^{2q}}q^{jq}\left(\frac{2}{1\alpha j}\left(\frac{2j}{e}\right)^jd_j^2\kappa _{\gamma }^{}{}_{}{}^{2}\right)^{2q}.$$
From (32), (52) (with $`\tau =\overline{\tau }_p`$) by choosing $`q=[\mathrm{log}(n)]`$, we have, as $`n+\mathrm{}`$
$$\left(Z_{j,n}\kappa _\gamma ^{}r_{2,n}\right)\lambda \left(\frac{\mathrm{log}(n)L(n)}{n^\alpha }\right)^{(j\overline{\tau }_p)q}\left(\frac{2}{1\alpha j}\left(\frac{2j}{e}\right)^jd_j^2\kappa _\epsilon ^2\kappa _{\gamma }^{}{}_{}{}^{2}\right)^q.$$
Consequently, as $`n+\mathrm{}`$, we finally obtain
$$\underset{j=\overline{\tau }_p}{\overset{k_\alpha 1}{}}\left(Z_{j,n}\kappa _\gamma ^{}r_{2,n}\right)\lambda \left(\frac{2}{1\alpha \overline{\tau }}\left(\frac{2\overline{\tau }}{e}\right)^{\overline{\tau }}d_{\overline{\tau }}^2\kappa _\epsilon ^2\kappa _{\gamma }^{}{}_{}{}^{2}\right)^{\mathrm{log}(n)}=𝒪\left(n^\gamma \right),$$
(90)
if $`\kappa _{\gamma }^{}{}_{}{}^{2}>\kappa _{2,\gamma }^{}=\frac{2}{1\alpha \overline{\tau }}\left(\frac{2\overline{\tau }}{e}\right)^{\overline{\tau }}d_{\overline{\tau }}^2\kappa _\epsilon ^2\mathrm{exp}(\gamma \overline{\tau })`$. Let us choose $`\kappa _\gamma ^{}`$ such that $`\kappa _{\gamma }^{}{}_{}{}^{2}>\mathrm{max}(\kappa _{1,\gamma }^{},\kappa _{2,\gamma }^{})`$. Then, by combining (89) and (90), we deduce from (87) that, for every $`\gamma >0`$
$$\left(|M_{b,n}|\kappa _\gamma ^{}\mathrm{max}(r_{1,n},r_{2,n})\right)=𝒪\left(n^\gamma \right),$$
and so, (80) is proved.
### 6.3 Proof of Theorem 2
Proof. Let us detail the proof presented in Section 6.1. We have
$$\frac{p\widehat{F}\left(\xi (p)\right)}{f(\xi (p))}\left(\widehat{\xi }\left(p\right)\xi (p)\right)=A(p)+B(p)+C(p)$$
with $`A(p)`$, $`B(p)`$ and $`C(p)`$ respectively defined by (48), (49) and (50). Under Assumption $`𝑨_\mathrm{𝟒}\mathbf{(}𝝃\mathbf{(}𝒑\mathbf{)}\mathbf{)}`$, from Lemma 10 and Taylor’s theorem we have almost surely, as $`n+\mathrm{}`$
$$C(p)\underset{|x|\epsilon _n(\alpha ,\tau _p)}{sup}F_{g(Y)}^{\prime \prime }(\xi (p)+x)\left(\widehat{\xi }\left(p\right)\xi (p)\right)^2=𝒪_{a.s.}\left(\epsilon _n(\alpha ,\tau _p)^2\right).$$
From the definition of sample quantile, we have almost surely, see e.g. Serfling (1980), $`A(p)=𝒪_{a.s.}\left(n^1\right)`$. Now, by combining Lemma 10 and Lemma 11, we have almost surely $`B(p)=𝒪_{a.s.}\left(r_n(\alpha ,\overline{\tau }_p)\right)`$. Thus, we finally obtain
$$\widehat{\xi }\left(p\right)\xi (p)=\frac{p\widehat{F}\left(\xi (p)\right)}{f(\xi (p))}+𝒪_{a.s.}\left(n^1\right)+𝒪_{a.s.}\left(r_n(\alpha ,\overline{\tau }_p)\right)+𝒪_{a.s.}\left(\epsilon _n(\alpha ,\tau _p)^2\right),$$
which leads to the result by noticing that $`\epsilon _n(\alpha ,\tau _p)^2=𝒪\left(r_n(\alpha ,\overline{\tau }_p)\right)`$.
### 6.4 Auxiliary Lemmas for the proof of Theorem 4
Let $`0<p_0p_1<1`$.
###### Lemma 12
Under conditions of Theorem 3, there exists a constant denoted by $`\theta =\theta (\alpha ,\tau _{p_0,p_1})`$ such that, we have almost surely, as $`n+\mathrm{}`$
$$T=\underset{p_0pp_1}{sup}\left|\widehat{\xi }(p;𝒈\mathbf{(}𝒀\mathbf{)})\xi _{g(Y)}(p)\right|\epsilon _n(\alpha ,\tau _{p_0,p_1}),$$
(91)
where $`\epsilon _n=\epsilon _n(\alpha ,\tau _{p_0,p_1})=\theta y_n(\alpha ,\tau _{p_0,p_1})`$ and $`y_n`$ is given by (50).
Proof. Define $`p_{j,n}=p_0+\frac{j}{[n^{3/2}]}(p_1p_0)`$ for $`j=0,\mathrm{},[n^{3/2}]`$, and let $`p[p_0,p_1]`$. Using the monotonicity of $`\widehat{\xi }()`$ and $`\xi ()`$, there exists some $`j`$ such that $`p[p_{j,n},p_{j+1,n}]`$ and such that
$`\widehat{\xi }\left(p\right)\xi (p)`$ $``$ $`\widehat{\xi }\left(p\right)\widehat{\xi }(p_{j+1,n}))+\widehat{\xi }(p_{j+1,n}))\xi (p)`$
$``$ $`\widehat{\xi }\left(p_{j+1,n}\right)\xi (p_{j+1,n}))+\xi (p_{j+1,n}))\xi (p_{j,n}))+\xi (p_{j,n}))\xi (p)`$
$``$ $`\widehat{\xi }\left(p_{j+1,n}\right)\xi (p_{j+1,n}))+\xi (p_{j+1,n}))\xi (p_{j,n})).`$
This leads to
$$T\underset{j=0,\mathrm{},[n^{3/2}]}{\mathrm{max}}|\widehat{\xi }\left(p_{j,n}\right)\xi (p_{j,n}))|+\underset{j=0,\mathrm{},[n^{3/2}]1}{\mathrm{max}}|\xi (p_{j+1,n})\xi (p_{j,n}))|.$$
(92)
Under Assumption $`𝑨_\mathrm{𝟓}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{,}𝒑_\mathrm{𝟏}\mathbf{)}`$, it comes
$$\underset{j=0,\mathrm{},[n^{3/2}]1}{\mathrm{max}}|\xi (p_{j+1,n})\xi (p_{j,n}))|=𝒪(n^{3/2}).$$
(93)
Now, following the proof of Lemma 10, one can prove that there exists some constant $`\theta (\alpha ,\tau _{p_0,p_1})`$ such that for all $`j=0,\mathrm{},[n^{3/2}]`$,
$$\left(|\widehat{\xi }\left(p_{j,n}\right)\xi (p_{j,n})|\theta y_n(\alpha ,\tau _{p_0,p_1})\right)=𝒪\left(n^3\right).$$
Therefore, as $`n+\mathrm{}`$,
$`\left(\underset{j=0,\mathrm{},[n^{3/2}]}{\mathrm{max}}|\widehat{\xi }\left(p_{j,n}\right)\xi (p_{j,n})|\epsilon _n\right)`$ $``$ $`([n^{3/2}]+1)\underset{j=0,\mathrm{},[n^{3/2}]}{\mathrm{max}}\left(|\widehat{\xi }\left(p_{j,n}\right)\xi (p_{j,n})|\epsilon _n\right)`$
$`=`$ $`𝒪\left(n^{3/2}\right).`$
which, combined with (92), (93) and Borel-Cantelli’s Lemma, leads to the result.
The following result is an extension of Lemma 11 and Theorem 4.2 obtained by Sen (1971).
###### Lemma 13
Under Assumptions of Theorem 3 and following Lemma 11, we have almost surely, as $`n+\mathrm{}`$
$$S_n^{}=\underset{\begin{array}{c}x,y[\xi (p_0),\xi (p_1)]\\ |xy|\epsilon _n(\alpha ,\tau _{p_0,p_1})\end{array}}{sup}\left|\mathrm{\Delta }(x)\mathrm{\Delta }(y)\right|=𝒪_{a.s.}\left(r_n(\alpha ,\tau _{p_0,p_1})\right)$$
(94)
where $`\tau _{p_0,p_1}`$ is defined by (33).
Proof. Set $`\epsilon _n=\epsilon _n(\alpha ,\tau _{p_0,p_1})`$ and $`r_n=r_n(\alpha ,\tau _{p_0,p_1})`$. Define $`\xi _{j,n}=\xi (p_0)+\frac{j}{p_n}(\xi (p_1)\xi (p_0))`$ for $`j=0,\mathrm{},p_n`$ with $`p_n=\left[\epsilon _n^1\right]`$, and let $`x,y[\xi (p_0),\xi (p_1)]`$ such that $`|xy|\epsilon _n`$. Two cases may occur
* If there exists some $`j`$ such that $`x,y[\xi _{j,n},\xi _{j+1,n}]`$ then
$$|\mathrm{\Delta }(x)\mathrm{\Delta }(y)||\mathrm{\Delta }(x)\mathrm{\Delta }(\xi _{j,n})|+|\mathrm{\Delta }(\xi _{j,n})\mathrm{\Delta }(y)|2\times S_n(\xi _{j,n},\epsilon _n)$$
* Otherwise and witout loss of generality, there exists $`j,k`$ with $`k>j`$ such that $`x[\xi _{j,n},\xi _{j+1,n}]`$ and $`y[\xi _{k,n},\xi _{k+1,n}]`$. Since $`|xy|\epsilon _n`$, it follows that $`|\xi _{k,n}\xi _{j+1,n}|\epsilon _n`$. Then,
$`|\mathrm{\Delta }(x)\mathrm{\Delta }(y)|`$ $``$ $`|\mathrm{\Delta }(x)\mathrm{\Delta }(\xi _{k,n})|+|\mathrm{\Delta }(\xi _{k,n})\mathrm{\Delta }(\xi _{j+1,n})|+|\mathrm{\Delta }(\xi _{j+1,n})\mathrm{\Delta }(y)|`$
$``$ $`S_n(\xi _{k,n},\epsilon _n)+2\times S_n(\xi _{j+1,n},\epsilon _n).`$
In other words, for all $`x,y`$ one may obtain
$$|\mathrm{\Delta }(x)\mathrm{\Delta }(y)|3\times \underset{0jp_n}{\mathrm{max}}S_n(\xi _{j,n},\epsilon _n).$$
Hence, $`S_n^{}3\times \mathrm{max}_{0jp_n}S_n(\xi _{j,n},\epsilon _n)`$. Now, following the proof of Lemma 11, one may prove that there exists some positive constant $`\theta _\gamma `$ such that for $`n`$ large enough and for all $`j=0,\mathrm{},p_n`$,
$$\left(S_n(\xi _{j,n},\epsilon _n)\theta _\gamma r_n\right)=𝒪\left(n^\gamma \right).$$
And in particular for $`\gamma =2`$, it comes
$`\left(\underset{0jp_n}{\mathrm{max}}S_n(\xi _{j,n},\epsilon _n)\theta _2r_n\right)`$ $``$ $`(p_n+1)\underset{j=0,\mathrm{},p_n}{\mathrm{max}}\left(S_n(\xi _{j,n},\epsilon _n)\theta _\gamma r_n\right)`$
$`=`$ $`𝒪\left({\displaystyle \frac{p_n}{n^2}}\right)=𝒪\left(n^{3/2}\right),`$
whatever the value of $`\alpha \tau _{p_0,p_1}`$. This leads to the result by using Borel-Cantelli’s Lemma.
### 6.5 Proof of Theorem 3
Proof. We follow the proof of Theorem 2. Let $`p[p_0,p_1]`$ and let $`\epsilon _n=\epsilon _n(\alpha ,\tau _{p_0,p_1})`$, then
$$\frac{p\widehat{F}\left(\xi (p)\right)}{f(\xi (p))}\left(\widehat{\xi }\left(p\right)\xi (p)\right)=A(p)+B(p)+C(p)$$
where $`A(p),B(p)`$ and $`C(p)`$ are respectively defined by (48), (49) and (50). Similarly to the proof of Theorem 2, one may prove that $`sup_{p_0pp_1}A(p)=𝒪_{a.s.}\left(n^1\right)`$. Under Assumption $`𝑨_\mathrm{𝟓}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{,}𝒑_\mathrm{𝟏}\mathbf{)}`$, $`C(p)\left(sup_{|x|\epsilon _n(\alpha ,\tau _p)}F^{\prime \prime }(x+\xi (p))\right)\frac{\left(\widehat{\xi }\left(p\right)\xi (p)\right)^2}{f(\xi (p))}`$. Therefore, for $`n`$ large enough, $`C(p)\lambda \left(sup_{p_0pp_1}\left(\widehat{\xi }\left(p\right)\xi (p)\right)\right)^2`$. And from Lemma 12, this leads to
$$\underset{p_0pp_1}{sup}C(p)=𝒪_{a.s.}\left(\epsilon _n(\alpha ,\tau _{p_0,p_1})^2\right).$$
In addition, using Lemma 13, one also has $`sup_{p_0pp_1}B(p)=𝒪_{a.s.}\left(r_n(\alpha ,\tau _{p_0,p_1})\right)`$, which ends the proof.
### 6.6 Auxiliary Lemma for the proof of Theorem 4
###### Lemma 14
Consider for $`0<p<1`$ the function $`h_p()`$, given by
$$h_p(t)=\mathrm{𝟏}_{\{|t|\xi _{|Y|}(p)\}}(t)p,$$
(95)
that is the function $`h_{\xi _{g(Y)}(p)}()`$ with $`g()=||`$. Then by denoting $`c_j^{h_p}`$ the $`j`$-th Hermite coefficient of $`h_p()`$, we have for all $`j1`$
$$c_0^{h_p}=c_{2j+1}^{h_p}=0\text{ and }c_{2j}^{h_p}=2H_{2j1}(q)\varphi (q),$$
(96)
where $`q=\xi _{|Y|}(p)=\mathrm{\Phi }^1\left(\frac{1+p}{2}\right)`$.
Proof. Since $`\left(|Y|q\right)=p`$ and $`h_p()`$ is even, we have $`c_0^{h_p}=c_{2j+1}^{h_p}=0`$, for all $`j1`$. Now, (27) implies
$`c_{2j}^{h_p}`$ $`=`$ $`{\displaystyle _{}}h_p(t)H_{2j}(t)\varphi (t)𝑑t=2\times {\displaystyle _0^q}H_{2j}(t)\varphi (t)𝑑t`$
$`=`$ $`2\times \left[\varphi ^{(2j1)}(t)\right]_0^q=2\times \left[H_{2j1}(t)\varphi (t)\right]_0^q`$
$`=`$ $`2H_{2j1}(q)\varphi (q).`$
###### Remark 10
Let $`g()=\stackrel{~}{g}(||)`$, where $`\stackrel{~}{g}()`$ is a strictly increasing function on $`^+`$, then for all $`0<p<1`$, we have
$$\xi _{|Y|}(p)=\stackrel{~}{g}^1\left(\xi _{g(Y)}(p)\right).$$
Consequently, the functions $`h_{\xi _{g(Y)}(p)}()`$ for $`g()=||`$, $`g()=||^\alpha `$ and $`g()=\mathrm{log}||`$ are strictly identical. And so, their Hermite decomposition is given by (96) and their Hermite rank is equal to 2.
### 6.7 Proof of Theorem 4
Proof. $`(i)`$ Define
$$b_n=\frac{1}{2}\underset{m=1}{\overset{M}{}}B_m\mathrm{log}\left(1+\delta _n^{𝒂^𝒎}(0)\right),$$
(97)
where $`\delta _n^{𝒂^𝒎}(0)`$ is given by (4). From (14), (15), and (25), we have almost surely
$`\widehat{H}^\alpha H`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha }}\epsilon _m^\alpha `$ (98)
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha }}\mathrm{log}\left({\displaystyle \frac{\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)}{\xi _{|Y|^\alpha }(𝒑,𝒄)}}\right)+\alpha \times b_n`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha \xi _{|Y|^\alpha }(𝒑,𝒄)}}(\widehat{\xi }(𝒑,𝒄;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)\xi _{|Y|^\alpha }(𝒑,𝒄))(1+o_{a.s.}\left(1\right))+\alpha b_n.`$
and
$`\widehat{H}^{\mathrm{log}}H`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}B_m\epsilon _m^{\mathrm{log}}`$ (99)
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}B_m\left(\widehat{\xi }(𝒑,𝒄;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{\mathrm{log}|Y|}(𝒑,𝒄)\right)+b_n`$
Under Assumption $`𝑨_\mathrm{𝟔}\mathbf{(}𝜼\mathbf{)}`$, we have
$$b_n=𝒪\left(n^\eta \right).$$
(100)
Moreover, let $`i,j1`$, under Assumption $`𝑨_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}𝝂\mathbf{)}`$, we have, from Lemma 1
$$𝐄(Y^{𝒂^𝒎}(i)Y^{𝒂^𝒎}(i+j))=\rho ^{𝒂^𝒎}(j)=𝒪\left(|j|^{2H2\nu }\right).$$
(101)
Then, for all $`m=1,\mathrm{},M`$ and for all $`k=1,\mathrm{},K`$, from Lemma 10 and Remark 10, we obtain, that almost surely
$`\widehat{\xi }(p_k;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)\xi _{|Y|^\alpha }(p_k)`$ $`=`$ $`𝒪_{a.s.}\left(y_n(2\nu 2H,\tau _{p_k})\right),`$
$`\widehat{\xi }(p_k;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{\mathrm{log}|Y|}(p_k)`$ $`=`$ $`𝒪_{a.s.}\left(y_n(2\nu 2H,\tau _{p_k})\right),`$
where the sequence $`y_n(,)`$ is defined by (52) with $`L()=1`$. The result (37) is obtained by combining (98), (99) and (100).
$`(ii)`$ Let us apply Theorem 2 to the sequence $`𝒈\mathbf{(}𝒀^{𝒂^𝒎}\mathbf{)}`$, for some $`m=1,\mathrm{},M`$, with $`g()=||`$, $`g()=||^\alpha `$ and $`g()=\mathrm{log}||`$. For all $`k=1,\mathrm{},K`$, we have almost surely
$`\widehat{\xi }(p_k;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{|Y|}(p_k)`$ $`=`$ $`{\displaystyle \frac{p_k\widehat{F}(\xi _{|Y|}(p_k);\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})}{f_{|Y|^\alpha }(\xi _{|Y|}(p_k))}}+𝒪_{a.s.}\left(r_n\right)`$
$`\widehat{\xi }(p_k;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)\xi _{|Y|^\alpha }(p_k)`$ $`=`$ $`{\displaystyle \frac{p_k\widehat{F}(\xi _{|Y|^\alpha }(p_k);\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)}{f_{|Y|^\alpha }(\xi _{|Y|^\alpha }(p_k))}}+𝒪_{a.s.}\left(r_n\right)`$
$`\widehat{\xi }(p_k;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{\mathrm{log}|Y|}(p_k)`$ $`=`$ $`{\displaystyle \frac{p_k\widehat{F}(\xi _{\mathrm{log}|Y|}(p_k);\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})}{f_{\mathrm{log}|Y|}(\xi _{\mathrm{log}|Y|}(p_k))}}+𝒪_{a.s.}\left(r_n\right),`$
where, for the sake of simplicity, $`r_n=r_n(2\nu 2H,\overline{\tau }_{p_k})`$ defined by (35) and (36). Note that from Remark 10 $`\overline{\tau }_{p_k}=2`$ for all $`k=1,\mathrm{},K`$.
With some little computation, we can obtain, almost surely
$$\widehat{\xi }(p_k;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)\xi _{|Y|^\alpha }(p_k)=\alpha \xi _{|Y|}(p_k)^{\alpha 1}\left(\widehat{\xi }(p_k;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{|Y|}(p_k)\right)+𝒪_{a.s.}\left(r_n\right),$$
(102)
and
$$\widehat{\xi }(p_k;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{\mathrm{log}|Y|}(p_k)=\xi _{|Y|}(p_k)^1\left(\widehat{\xi }(p_k;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{|Y|}(p_k)\right)+𝒪_{a.s.}\left(r_n\right).$$
(103)
From (98), (99), (102), (103) and properties of Gaussian variables, the following results hold almost surely
$$\widehat{H}^\alpha H=\underset{m=1}{\overset{M}{}}\underset{k=1}{\overset{K}{}}\frac{B_mc_k}{2q_k\varphi (q_k)}\pi _{k,\alpha }\left(\widehat{F}(q_k;\mathbf{|}𝒀\mathbf{|})p_k\right)+𝒪_{a.s.}\left(r_n\right)+𝒪\left(b_n\right),$$
(104)
and
$$\widehat{H}^{\mathrm{log}}H=\underset{m=1}{\overset{M}{}}\underset{k=1}{\overset{K}{}}\frac{B_mc_k}{2q_k\varphi (q_k)}\left(\widehat{F}(q_k;\mathbf{|}𝒀\mathbf{|})p_k\right)+𝒪_{a.s.}\left(r_n\right)+𝒪\left(b_n\right),$$
(105)
where $`q_k`$ and $`\pi _k^\alpha `$ are defined by (42). Denote by $`\theta _{m,k}^\alpha `$ the following constant
$$\theta _{m,k}^\alpha =\frac{B_mc_k}{2q_k\varphi (q_k)}\pi _k^\alpha .$$
Since $`\pi _k^0=1`$, (104) and (105) can be rewritten as
$`\widehat{H}^\alpha H`$ $`=`$ $`Z_n^\alpha +𝒪_{a.s.}\left(r_n\right)+𝒪\left(b_n\right)`$ (106)
$`\widehat{H}^{\mathrm{log}}H`$ $`=`$ $`Z_n^0+𝒪_{a.s.}\left(r_n\right)+𝒪\left(b_n\right),`$ (107)
where for $`\alpha 0`$,
$$Z_n^\alpha =\underset{m=1}{\overset{M}{}}\underset{k=1}{\overset{K}{}}\theta _{m,k}^\alpha \left(\widehat{F}(q_k;\mathbf{|}𝒀\mathbf{|})p_k\right).$$
(108)
Thus, under Assumption $`𝑨_\mathrm{𝟔}\mathbf{(}𝜼\mathbf{)}`$, we have, as $`n+\mathrm{}`$,
$`MSE(\widehat{H}^\alpha H)`$ $`=`$ $`𝒪\left(𝐄\left((Z_n^\alpha )^2\right)\right)+𝒪\left(r_n(2\nu 2H,2)^2\right)+𝒪\left(n^{2\eta }\right),`$ (109)
$`MSE(\widehat{H}^{\mathrm{log}}H)`$ $`=`$ $`𝒪\left(𝐄\left((Z_n^0)^2\right)\right)+𝒪\left(r_n(2\nu 2H,2)^2\right)+𝒪\left(n^{2\eta }\right).`$ (110)
Now,
$$𝐄\left((Z_n^\alpha )^2\right)=\frac{1}{n^2}\underset{m_1,m_2=1}{\overset{M}{}}\underset{k_1,k_2=1}{\overset{K}{}}\underset{i_1,i_2=1}{\overset{n}{}}\theta _{m_1,k_1}^\alpha \theta _{m_2,k_2}^\alpha 𝐄\left(h_{q_{k_1}}(Y^{𝒂^{m_1}}(i_1))h_{q_{k_2}}(Y^{𝒂^{m_2}}(i_2))\right).$$
For $`k_1,k_2=1,\mathrm{},K`$, $`m_1,m_2=1,\mathrm{},M`$ and $`i_1,i_2=1,\mathrm{},n`$, we have from Lemma 14,
$`𝐄\left(h_{q_{k_1}}(Y^{𝒂^{m_1}}(i_1))h_{q_{k_2}}(Y^{𝒂^{m_2}}(i_2))\right)`$ $`=`$ $`{\displaystyle \underset{j_1\tau _{p_{k_1}}/2}{}}{\displaystyle \underset{j_2\tau _{p_{k_2}}/2}{}}{\displaystyle \frac{c_{2j_1}^{h_{p_{k_1}}}c_{2j_2}^{h_{p_{k_2}}}}{(2j_1)!(2j_2)!}}`$ (111)
$`\times 𝐄\left(H_{2j_1}(Y^{𝒂^{m_1}}(i_1))H_{2j_2}(Y^{𝒂^{m_2}}(i_2))\right)`$
$`=`$ $`{\displaystyle \underset{j1}{}}{\displaystyle \frac{c_{2j}^{h_{p_{k_1}}}c_{2j}^{h_{p_{k_2}}}}{(2j)!}}\rho ^{𝒂^{m_1},𝒂^{m_2}}(i_2i_1)^{2j}.`$
Under Assumption $`𝑨_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}𝝂\mathbf{)}`$, we have from Lemma 1, $`\rho ^{𝒂^{m_1},𝒂^{m_2}}(i)=𝒪\left(|i|^{2H2\nu }\right)`$. Now, we leave the reader to check that, as $`n+\mathrm{}`$
$$\frac{1}{n^2}\underset{i_1,i_2=1}{\overset{n}{}}\rho ^{𝒂^{m_1},𝒂^{m_2}}(i_2i_1)^2=𝒪(\frac{1}{n}\underset{|i|n}{}|i|^{2(2H2\nu )})=𝒪(v_n(2\nu 2H))),$$
where the sequence $`v_n()`$ is given by (39). Thus, we have, as $`n+\mathrm{}`$, $`𝐄\left((Z_n^\alpha )^2\right)=𝒪\left(v_n(2\nu 2H)\right)`$, which leads to the result from (109) and (110).
$`(iii)`$ Assume $`\nu >H+1/4`$ and $`\eta >1/2`$, then from (106) and (107), the following equivalences in distribution hold
$$\sqrt{n}\left(\widehat{H}_n^\alpha H\right)\sqrt{n}Z_n^\alpha \text{ and }\sqrt{n}\left(\widehat{H}_n^{\mathrm{log}}H\right)\sqrt{n}Z_n^0.$$
(112)
Now, decompose $`Z_n^\alpha =T_n^1+T_n^2`$, where
$$T_n^1=\frac{1}{\sqrt{n}}\underset{m=1}{\overset{M}{}}\underset{k=1}{\overset{K}{}}\theta _{m,k}^\alpha \underset{i=\mathrm{}+1}{\overset{M\mathrm{}+1}{}}h_{q_k}(Y^{𝒂^𝒎}(i))$$
and
$$T_n^2=\sqrt{n}\underset{m=1}{\overset{M}{}}\underset{k=1}{\overset{K}{}}\theta _{m,k}^\alpha \left\{\frac{1}{n}\underset{i=M\mathrm{}+1}{\overset{n}{}}h_{q_k}(Y^{𝒂^𝒎}(i))\right\},$$
Clearly, $`T_n^1`$ converges to 0 in probability, as $`n+\mathrm{}`$. Therefore, we have, as $`n+\mathrm{}`$
$$Z_n^\alpha \sqrt{n}\left\{\frac{1}{n}\underset{i=M\mathrm{}+1}{\overset{n}{}}G^\alpha (Y^{𝒂^\mathrm{𝟏}}(i),\mathrm{},Y^{𝒂^𝑴}(i))\right\}$$
(113)
where $`G^\alpha `$ is the function from $`^M`$ to $``$ defined for $`\alpha 0`$ and $`t_1,\mathrm{},t_M`$ by:
$$G^\alpha (t_1,\mathrm{},t_M)=\underset{m=1}{\overset{M}{}}\underset{k=1}{\overset{K}{}}\theta _{m,k}h_{q_k}(t_m).$$
(114)
Denote by $`\stackrel{\mathbf{~}}{𝒀}^𝒂(i)`$, the vector defined for $`i=M\mathrm{}+1,\mathrm{},n`$ by
$$\stackrel{\mathbf{~}}{𝒀}^𝒂(i)=(Y^{𝒂^\mathrm{𝟏}}(i),\mathrm{},Y^{𝒂^𝑴}(i)).$$
We obviously have $`𝐄\left(G^\alpha (\stackrel{\mathbf{~}}{𝒀}^𝒂(i))^2\right)<+\mathrm{}`$. Since, for all $`k=1,\mathrm{},K`$, the functions $`h_{q_k}`$ have Hermite rank $`\tau _{p_k}`$, the function $`G^\alpha `$ has Hermite rank $`2`$ (see e.g. Arcones (1994) for the definition of the Hermite rank of multivariate functions). Moreover under Assumption $`𝑨_\mathrm{𝟏}\mathbf{(}\mathrm{𝟐}𝝂\mathbf{)}`$, we have from Lemma 1, as $`j+\mathrm{}`$
$`𝐄\left(Y^{𝒂^{𝒎_\mathrm{𝟏}}}(i)Y^{𝒂^{𝒎_\mathrm{𝟐}}}(i+j)\right)`$ $`=`$ $`\rho ^{𝒂^{𝒎_\mathrm{𝟏}}\mathbf{,}𝒂^{𝒎_\mathrm{𝟐}}}(j)=𝒪\left(|j|^{2H2\nu }\right)\mathrm{}^2(),`$
as soon as $`\nu >H+1/4`$. Thus, from Theorem 4 of Arcones (1994), there exists $`\sigma _\alpha ^2`$ (defined for $`\alpha 0`$) such that, as $`n+\mathrm{}`$, the following convergence in distribution holds
$$Z_n^\alpha 𝒩(0,\sigma _\alpha ^2)$$
with
$$\sigma _\alpha ^2=\underset{i}{}𝐄\left(G^\alpha \left(\stackrel{\mathbf{~}}{𝒀}^𝒂(i^{})\right)G^\alpha \left(\stackrel{\mathbf{~}}{𝒀}^𝒂(i^{}+i)\right)\right).$$
With previous notations, we have
$`\sigma _\alpha ^2`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{m_1,m_2=1}{\overset{M}{}}}{\displaystyle \underset{k_1,k_2=1}{\overset{K}{}}}\theta _{m_1,k_1}^\alpha \theta _{m_2,k_2}^\alpha 𝐄\left(h_{p_{k_1}}(Y^{𝒂^{𝒎_\mathrm{𝟏}}}(i^{}))h_{p_{k_2}}(Y^{𝒂^{𝒎_\mathrm{𝟐}}}(i^{}+i))\right)`$ (115)
$`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{m_1,m_2=1}{\overset{M}{}}}{\displaystyle \underset{k_1,k_2=1}{\overset{K}{}}}{\displaystyle \underset{jr}{}}{\displaystyle \frac{c_{2j}^{h_{p_{k_1}}}c_{2j}^{h_{p_{k_2}}}}{(2j)!}}\theta _{m_1,k_1}^\alpha \theta _{m_2,k_2}^\alpha \rho ^{𝒂^{𝒎_\mathrm{𝟏}}\mathbf{,}𝒂^{𝒎_\mathrm{𝟐}}}(i)^{2j}.`$
From (96), we can see that formula (115) is equivalent to (41), which ends the proof from (112).
### 6.8 Proof of Corollary 5
Proof. Equation (113) is still available for a sequence $`\alpha _n`$ such that $`\alpha _n0`$ as $`n+\mathrm{}`$, that is
$$\sqrt{n}\left(\widehat{H}_n^{\alpha _n}H\right)\sqrt{n}\left\{\frac{1}{n}\underset{i=M\mathrm{}+1}{\overset{n}{}}G^{\alpha _n}(Y^{𝒂^\mathrm{𝟏}}(i),\mathrm{},Y^{𝒂^𝑴}(i))\right\}$$
From (114) and since $`\pi _k^{\alpha _n}1`$, as $`n+\mathrm{}`$, we have $`G^{\alpha _n}()G^0()`$. Therefore, the following equivalence in distribution holds, as $`n+\mathrm{}`$
$$\sqrt{n}\left(\widehat{H}_n^{\alpha _n}H\right)\sqrt{n}\left(\widehat{H}_n^{\mathrm{log}}H\right),$$
which ends the proof.
### 6.9 Auxiliary Lemma for the proof of Theorem 6
###### Lemma 15
Let $`0<\beta _1\beta _2<1`$ and let $`𝐙=(Z_1,\mathrm{},Z_n)`$ $`n`$ random variables identically distributed, such that $`sup_{\beta _1p\beta _2}\widehat{\xi }_Z(p;𝐙)=𝒪_{a.s.}\left(1\right)`$, then
$$\overline{𝒁}^{(𝜷)}\frac{1}{1\beta _2\beta _1}_{\beta _1}^{1\beta _2}\widehat{\xi }_Z(p;𝒁)𝑑p=𝒪_{a.s.}\left(n^1\right).$$
Proof. It is sufficient to notice that for $`i=1,\mathrm{},n`$
$$n_{\frac{i1}{n}}^{\frac{i}{n}}\widehat{\xi }(p;𝒁)𝑑pZ_{(i),n}n_{\frac{i}{n}}^{\frac{i+1}{n}}\widehat{\xi }(p;𝒁)𝑑p.$$
which leads to
$$\frac{n}{n[n\beta _2][n\beta _1]}_{\frac{[n\beta _1]}{n}}^{\frac{n[n\beta _2]}{n}}\widehat{\xi }(p;𝒁)𝑑p\overline{𝒁}^{(𝜷)}\frac{n}{n[n\beta _2][n\beta _1]}_{\frac{[n\beta _1]}{n}+\frac{1}{n}}^{\frac{n[n\beta _2]}{n}+\frac{1}{n}}\widehat{\xi }(p;𝒁)𝑑p.$$
The end is omitted.
### 6.10 Proof of Theorem 6
Proof. $`(i)`$ From (20), (21), and (26), we have
$`\widehat{H}^{\alpha ,tm}H`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha }}\epsilon _m^{\alpha ,tm}`$ (116)
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha }}\mathrm{log}\left(\overline{\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶}^{(𝜷)}/\overline{|Y|^\alpha }^{(𝜷)}\right)+\alpha \times b_n`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha \overline{|Y|^\alpha }^{(𝜷)}}}(\overline{\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶}^{(𝜷)}\overline{|Y|^\alpha }^{(𝜷)})(1+o_{a.s.}\left(1\right))+\alpha b_n.`$
and
$`\widehat{H}^{\mathrm{log},tm}H`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}B_m\epsilon _m^{\mathrm{log},tm}`$ (117)
$`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}B_m\left(\overline{\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}}^{(𝜷)}\overline{\mathrm{log}|Y|}^{(𝜷)}\right)+b_n`$
Let us notice that from Lemma 12, one can apply Lemma 15 for the vectors $`\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶`$ and $`\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}`$. Then it comes
$`\overline{\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶}^{(𝜷)}\overline{|Y|^\alpha }^{(𝜷)}`$ $`=`$ $`{\displaystyle \frac{1}{1\beta _2\beta _1}}{\displaystyle _{\beta _1}^{1\beta _2}}\left(\widehat{\xi }(p;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}^𝜶)\xi _{|Y|^\alpha }(p)\right)𝑑p+𝒪_{a.s.}\left(n^1\right),`$
$`\overline{\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|}}^{(𝜷)}\overline{\mathrm{log}|Y|}^{(𝜷)}`$ $`=`$ $`{\displaystyle \frac{1}{1\beta _2\beta _1}}{\displaystyle _{\beta _1}^{1\beta _2}}\left(\widehat{\xi }(p;\mathrm{𝐥𝐨𝐠}\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})\xi _{\mathrm{log}|Y|}(p)\right)𝑑p+𝒪_{a.s.}\left(n^1\right).`$
Hence, from (101) and Lemma 12 and Remark 10 and under Assumption $`\mathbf{(}𝑨_\mathrm{𝟔}\mathbf{(}𝜼\mathbf{)}\mathbf{)}`$, we obtain
$`\widehat{H}^{\alpha ,tm}H`$ $`=`$ $`𝒪_{a.s.}\left(y_n(2H2,2)\right)+𝒪\left(n^\eta \right)+𝒪_{a.s.}\left(n^1\right)`$
$`\widehat{H}^{\mathrm{log},tm}H`$ $`=`$ $`𝒪_{a.s.}\left(y_n(2H2,2)\right)+𝒪\left(n^\eta \right)+𝒪_{a.s.}\left(n^1\right),`$
where the sequence $`y_n(,)`$ is defined by (52) with $`L()=1`$. This leads to the result by noticing that $`n^1=𝒪\left(y_n(2H2,2)\right)`$.
$`(ii)`$ By following the proof of Theorem 4 $`(ii)`$ and from Theorem 3, we may obtain the following representation
$`\widehat{H}^{\alpha ,tm}H`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \frac{B_m}{\alpha \overline{|Y|^\alpha }^{(𝜷)}}}\times {\displaystyle \frac{1}{1\beta _2\beta _1}}{\displaystyle _{\beta _1}^{1\beta _2}}{\displaystyle \frac{\widehat{F}(q;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})p}{2\frac{1}{\alpha }q^{1\alpha }\varphi (q)}}𝑑p`$
$`+𝒪_{a.s.}\left(r_n\right)+𝒪_{a.s.}\left(n^1\right)+𝒪\left(n^\eta \right),`$
$`\widehat{H}^{\mathrm{log},tm}H`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}B_m\times {\displaystyle \frac{1}{1\beta _2\beta _1}}{\displaystyle _{\beta _1}^{1\beta _2}}{\displaystyle \frac{\widehat{F}(q;\mathbf{|}𝒀^{𝒂^𝒎}\mathbf{|})p}{2q\varphi (q)}}𝑑p`$
$`+𝒪_{a.s.}\left(r_n\right)+𝒪_{a.s.}\left(n^1\right)+𝒪\left(n^\eta \right).`$
With such a representation, we observe that the result $`(ii)`$ can be proved similarly to the one of Theorem 4.
$`(iii)`$ By assuming that $`\eta >1/2`$ and $`\nu >H+1/4`$, one may obtain the asymptotic normality of $`\widehat{H}^{\alpha ,tm}`$ and $`\widehat{H}^{\mathrm{log},tm}`$ by using the same tools as the one presented in the proof of Theorem 4 $`(iii)`$. Therefore, let us just explicit the asymptotic variance of estimators $`\widehat{H}^{\alpha ,tm}`$ and $`\widehat{H}^{\mathrm{log},tm}`$. If $`\nu >H+1/4`$ and $`\eta >1/2`$, then from previous representations and from 111 we obtain as $`n+\mathrm{}`$
$`Var\left(\sqrt{n}\left(\widehat{H}^{\alpha ,tm}H\right)\right)`$ $``$ $`{\displaystyle \frac{n}{n^2}}{\displaystyle \underset{m_1,m_2=1}{\overset{M}{}}}{\displaystyle \underset{i_1,i_2=1}{\overset{n}{}}}B_{m_1}B_{m_2}{\displaystyle \frac{1}{(\overline{|Y|^\alpha }^{(𝜷)})^2}}\times {\displaystyle \frac{1}{(1\beta _2\beta _1)^2}}\times `$
$`{\displaystyle _{\beta _1}^{1\beta _2}}{\displaystyle _{\beta _1}^{1\beta _2}}{\displaystyle \underset{j1}{}}{\displaystyle \frac{c_{2j}^{h_{p_1}}c_{2j}^{h_{p_2}}}{(2j)!}}{\displaystyle \frac{q_1^{\alpha 1}q_2^{\alpha 1}}{4\varphi (q_1)\varphi (q_2)}}dp_1dp_2\rho ^{a^{m_1},a^{m_2}}(i)^{2j},`$
with $`q_k=\mathrm{\Phi }^1\left(\frac{1+p_k}{2}\right)`$ for $`k=1,2`$. Due to (96) and since $`\overline{|Y|^\alpha }^{(𝜷)}=\frac{1}{1\beta _2\beta _1}_{\beta _1}^{1\beta _2}q^\alpha 𝑑p`$, this variance converges towards $`\sigma _{\alpha ,tm}^2`$ given by (46), as $`n+\mathrm{}`$.
We leave the reader to check that the asymptotic variance of $`\sqrt{n}\left(\widehat{H}^{\mathrm{log},tm}H\right)`$ is given by $`\sigma _{0,tm}^2`$.
Acknowledgement. The author is very grateful to Anestis Antoniadis and Rémy Drouilhet for helpful comments and to Kinga Sipos for a careful reading of the present paper.
J.-F. Coeurjolly
LJK, SAGAG Team, Université Grenoble 2
1251 Av. Centrale BP 47
38040 GRENOBLE Cedex 09
France
E-mail: Jean-Francois.Coeurjolly@upmf-grenoble.fr
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# Accuracy of Semiclassics: Comparative Analysis of WKB and Instanton Approaches.
## I Introduction
It is a textbook wisdom that if the de Broglie wavelengths $`\lambda `$ of particles are small in comparison with the characteristic space scales of a given problem, then the problem can be treated semiclassically. The commonly used WKB method (phase integral approach) LL65 \- HE62 is intended for the conditions of ”geometrical optics”, in which the gradient of the action $`\sigma `$ is large, but slowly variable (this is suggested also by it containing the factor $`\mathrm{}`$, since we are dealing with a semiclassical approximation, in which $`\mathrm{}`$ is taken as small). The corresponding condition can be formulated more quantitatively LL65 as follows: $`\lambda `$ must vary only slightly over distances of the order of itself
$`\left|{\displaystyle \frac{d(\lambda (x)/2\pi )}{dx}}\right|1,`$ (1)
where $`\lambda (x)=2\pi \mathrm{}/p(x)`$, and $`p(x)`$ is a classical momentum. However, this simple criterion is not a practical tool to estimate how accurate could be found semiclassical solutions of particular problems, since nothing is specified regarding the convergence of the semiclassical series. The criterion (1) does not work to estimate the magnitude of the error involved in the approximate calculation of physical quantities (e.g., matrix elements), neither to find the domain of validity in the complex plane in which the semiclassical solutions are defined. Indeed from (1) one can conclude only that higher order corrections to semiclassical wave functions are small in the asymptotic regions, but this mathematical criterion has almost nothing to do with say the physical accuracy of semiclassical matrix elements which depends on the wave function accuracy in space regions providing main contributions into the matrix elements under consideration. For example, the energy eigenvalues are determined by the asymptotical region of the linear turning points (i.e., the region distant from these points), and as well by the proximity region to the second order turning points, since in the both regions the wave functions possess the largest values.
Within the WKB method such kind of a physical accuracy estimation has been performed long ago by N. and P.O. Fröman FF65 . They analyzed higher order corrections to the semiclassical wave functions and found that although those are really small over $`1/\gamma ^2`$ ($`\gamma 1`$ is semiclassical parameter), the corrections are proportional to the factor $`[(E/\gamma )U]^2`$, (where $`E`$ is energy and $`U`$ is potential), and thus the function has non-integrable singularity at the linear turning points where $`[(E/\gamma )U]=0`$ (or, within the alternative to WKB semiclassical formalism so-called extreme tunneling trajectory or instanton instanton approach PO77 \- BV02 , the corrections are singular in the second order turning points). To surmount this problem in FF65 (see also OL59 , and OL74 ) the analytical continuation of the correction function into the complex plane has been proposed, and it gives impractically bulky expressions even for simple model potentials. Since this problem has relevance far beyond WKB treatment of a particular model potential, this is an issue of general interest to develop a simple and convenient in practice quantitative method to study the accuracy of the semiclassical approach, and it is the immediate motivation of the present paper to develop a systematic procedure how to do it.
The idea of our approach is to construct two linearly independent continuous (with continuous first derivatives) approximate solutions to the Schrödinger equation, which in the asymptotic region coincide with semiclassic solutions, and in the vicinity of the turning points - with the exact solutions of the so-called comparison equation (i.e. the exact solution of the Schrödinger equation for the chosen appropriately approximate near the turning points potentials $`V_c(X)`$, henceforth will be referred to as the comparison potential). Although, scanning the literature we found one rather old paper LA37 with a similar comparison equation approach, but our accuracy criterion is formulated as the majorant inequalities for a certain matrix (which we find in the explicit analytical form and calculated numerically) connecting our approximate and exact solutions in the finite space interval (not only in the vicinity of the turning points). Since this method has largely gone unnoticed in the study of semiclassics, we found it worthwhile to present its derivation in a short and explicit form, and also to point out its practical usability.
The remainder of this paper is organized as follows. In section II we present the basic expressions necessary for our investigation. In this section we also present the main steps and qualitative idea of our method. Section III contains our results. We derive the inequalities which enable us to find the finite space interval (not at the isolated points) where the solutions have to be matched, and calculate the $`2\times 2`$ coordinate dependent matrix connecting the approximate and exact solutions. Since the semiclassical solutions of the harmonic potential coincide with the exact solutions, the accuracy of any semiclassically treated problem depends crucially on its potential energy anharmonicity. That is why as the touchstone to test our method the results presented in the section III are applied to an anharmonic oscillator in section IV. We end with some brief conclusions in the same section.
## II Semiclassical equations in the WKB and instanton forms
Technically the basic idea how to overcome the difficulty of the semiclassical solutions in the vicinity of the turning points is reduced to an appropriate (admitting exact analytic solutions) approximation of the potential near the turning points. After that step one has to match the asymptotics of this exact solution to the Schrödinger equation for an approximate potential with the semiclassical solutions to the Schrödinger equation for the potential under consideration (i.e. approximate solutions of the exact potential) far from the turning points. To illustrate main ideas of any semiclassic method (and to retain compactness and transparency of expressions) we discuss here a one dimensional case. As it is well known LL65 in the WKB method solutions to the Schrödinger equation are sought in the form
$`\psi =A\mathrm{exp}\left({\displaystyle \frac{i\sigma }{\mathrm{}}}\right),`$ (2)
where for the function $`\sigma `$ called action the one particle Schrödinger equation (traditionally termed as Hamilton Jacoby equation) reads as
$`{\displaystyle \frac{1}{2m}}\left({\displaystyle \frac{\sigma }{x}}\right)^2=EU,`$ (3)
where $`m`$ is a particle mass, $`E`$ is its energy, and $`U`$ is external field potential. Since the system is supposed quasi-classical in its properties, we seek $`\sigma `$ in the form of a series expanded in powers of $`\mathrm{}`$. Depending on normalization prefactor $`A(x)`$ entering (2) can be also found but the corresponding equation (referred traditionally as transport equation) plays a pure passive role since it is fully determined by the action $`\sigma `$ found as the solution of the Hamilton - Jacoby equation
$`{\displaystyle \frac{i\mathrm{}}{m}}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\sigma }{x^2}}A+{\displaystyle \frac{A}{x}}{\displaystyle \frac{\sigma }{x}}\right]+{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{^2A}{x^2}}=0,`$ (4)
where in the spirit of the semiclassical approximation the last term ($`\mathrm{}^2`$) is neglected. Technically of course more convenient to use instead of $`\mathrm{}`$ an expansion over equivalent but dimensionless parameter $`\gamma ^11`$ we will call in what follows as semiclassical parameter and define as
$`\gamma {\displaystyle \frac{m\mathrm{\Omega }_0a_0^2}{\mathrm{}}}1,`$ (5)
where $`a_0`$ is a characteristic length of the problem, e.g. the tunneling distance, $`\mathrm{\Omega }_0`$ is a characteristic frequency, e.g. the oscillation frequency around the potential minimum. Evidently the semiclassical parameter $`\gamma 1`$, and by its physical meaning it is determined by the ratio of the characteristic potential scale over the zero oscillation energy. We put $`\mathrm{}=1`$, and use $`\mathrm{\Omega }_0`$ and $`a_0`$ to set corresponding dimensionless scales, i.e. we introduce dimensionless energy $`ϵE/\gamma \mathrm{\Omega }_0`$, dimensionless coordinate $`Xx/a_0`$, dimensionless potential $`VU/\mathrm{\Omega }_0`$ (except where explicitely stated to the contrary and dimensions are necessary for understanding or numerical estimations).
The analogous to (2), (3) procedure for the Schrödinger equation in the imaginary time (instanton formalism, corresponding to the Wick rotation in the phase space, when coordinates remain real valued $`xx`$ but conjugated momenta become imaginary $`p_xip_x`$) can be formulated as the following substitution for the wave function (cf. to (2))
$`\psi =A_E(X)\mathrm{exp}(\gamma \sigma _E),`$ (6)
where the action $`\sigma _E`$ and we use the subscript $`E`$ to denote so-called Euclidean action obtained from the WKB action $`\sigma `$ after the Wick rotation. Performed above rotation is not a harmless change of variables. The deep meaning of this transformation within the instanton approach is related to redistribution of different terms between the Hamilton - Jacoby and the transport equations. Indeed, like that is in the WKB method, eigenvalues for the ground and for the low-lying states are of the order of $`\gamma ^0`$, while all other terms in (3) are of the order of $`\gamma ^1`$. Therefore to perform a regular expansion over $`\gamma ^1`$ for the substitution (6) one has to remove the energy term from the Hamilton - Jacoby equation, and to include this term into the transport equation. Besides in the first order over $`\gamma ^1`$ one can neglect the term with the second derivative of the prefactor. As a result of this redistributions the both equations are presented as
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sigma _E}{X}}\right)^2=V(X),`$ (7)
instead of the WKB Hamilton - Jacoby equation (3), and the transport equation is
$`{\displaystyle \frac{A_E}{X}}{\displaystyle \frac{\sigma _E}{X}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\sigma _E}{X^2}}A_E=ϵA_E.`$ (8)
One can easily note by a simple inspection of the WKB (3), (4) and of the instanton (7), (8) equations that although the both semiclassical methods can be formulated neglecting terms of the order of $`\gamma ^1`$, therefore possessing the same accuracy over $`\gamma ^1`$, the solutions evidently coincide in the asymptotic classically forbidden region $`V(X)ϵ/\gamma `$, but their behavior, number and type of turning points (where any semiclassic approximation does not work) are quite different. For example in the WKB formalism there are two turning points where $`V(X)(ϵ/\gamma )=0`$ around each minimum of the potential, while in the instanton approach since the energy does not enter the Hamilton - Jacoby equation (thus one can say that no classically accessible regions at all) the turning points are extremal points (minima for the case) of the potential. Furthermore as a consequence of this difference, in the WKB method all turning points are linear, whereas in the instanton approach they are second order (quadratic over $`X`$).
## III Accuracy of semiclassical approximation
Armed with this knowledge we are in the position now to construct our approximants. Let us introduce besides the comparison potential $`V_c(X)`$, one more specially chosen potential $`V_{sc}`$ (henceforth will be referred to as the semiclassical potential). This potential is chosen by the requirement that the exact solutions to the Schrödinger equation with $`V_{sc}`$ coincide asymptotically with the semiclassical solutions to the Schrödinger equation with the potential $`V(X)`$ the problem under study. Thus according to the construction, the semiclassical wave function $`\mathrm{\Psi }_{sc}`$ satisfies the equation
$`\mathrm{\Psi }_{sc}^1{\displaystyle \frac{d^2\mathrm{\Psi }_{sc}}{dX^2}}=2\gamma ^2\left(V_{sc}(X){\displaystyle \frac{ϵ}{\gamma }}\right),`$ (9)
and from here we can relate the semiclassical potential ($`V_{sc}(X)`$) with the bare one ($`V(X)`$)
$`V_{sc}^{(1,2)}=V(X){\displaystyle \frac{1}{2\gamma ^2}}A^1\left({\displaystyle \frac{d^2A}{dX^2}}\right),`$ (10)
in the vicinity of the first order or of the second order (superscripts 1 or 2) turning points.
Since near the turning point $`X_0`$ the prefactors $`A^{(1)}|XX_0|^{1/4}`$, and $`A^{(2)}(XX_0)^n`$ (where $`n`$ is an integer number which occurs from the transport equation (8) solution at the energy $`ϵ=n+(1/2)`$) the potential $`V_{sc}^{(1)}`$ at $`XX_0`$ is singular and negative, and $`V_{sc}^{(2)}`$ has the same singularity ($`(XX_0)^2`$) but positive. The difference is due to the fact that near the WKB linear turning points we have deal with the $`V_{sc}^{(1)}`$ well, whereas near the second order instanton turning points one has to treat the potential barrier $`V_{sc}^{(2)}`$. It might be useful to illustrate the essential features of the introduced above potentials $`V_c`$ and $`V_{sc}`$ applying the definition (10) to a simple (but the generic touchstone) example of the following anharmonic oscillator
$`V(X)={\displaystyle \frac{1}{2}}[X^2+\alpha X^3+\beta X^4].`$ (11)
We show in Fig. 1 the semiclassical and the comparison potentials associated with (11) for the WKB (Fig. 1a) and instanton (Fig. 1b) methods ($`\alpha =1.25`$, $`\beta =0.5`$, and the energy window corresponds to $`n=3`$ excited state of the potential (11)).
The key elements to construct our approximants are the following combinations related to probability flows to and from the turning points
$`J(X_1)=\mathrm{\Psi }_{sc}^1\left({\displaystyle \frac{d\mathrm{\Psi }_{sc}}{dX}}\mathrm{\Psi }_c^1{\displaystyle \frac{d\mathrm{\Psi }_c}{dX}}\right)_{X=X_1}=2\gamma ^2\mathrm{\Psi }_{sc}^1(X_1)\mathrm{\Psi }_c^1(X_1){\displaystyle _{\mathrm{}}^{X_1}}\mathrm{\Psi }_{sc}(X)\mathrm{\Psi }_c(X)(V_{sc}(X)V_c(X))𝑑X,`$ (12)
where $`X_1<X_0`$ and analogously for $`X_2>X_0`$ the flow function $`J(X_2)`$ is given by (12) where the integration limits are from $`X_2`$ to $`+\mathrm{}`$. Since the exact wave functions are continuous with continuous first derivatives (providing due to these features the continuity of the density probability currents), the idea of our procedure is to require the same from the approximate wave functions.
The integrals entering $`J(X_1)`$ and $`J(X_2)`$ can be calculated easily for any form of the potential, and the maximum accuracy of the any semiclassical approach can be achieved upon the matching of the approximate solutions at the characteristic points $`X_{1,2}^\mathrm{\#}`$ where $`J(X_{1,2}^\mathrm{\#})=0`$. The points $`X_{1,2}^\mathrm{\#}`$ do exist in the case when the potentials $`V_c`$ and $`V_{sc}`$ intersect in the region where the approximate wave functions $`\mathrm{\Psi }_{sc}`$ and $`\mathrm{\Psi }_c`$ are monotone ones. It is easy to realize (see e.g., Fig. 1) that the both points occur in the vicinity of the linear turning points for the potentials with $`d^2V/dX^2>0`$. One such a point disappears when the potential turning point becomes the inflection point, and there are no $`X_{1,2}^\mathrm{\#}`$ points at all for $`d^2V/dX^2<0`$. In the vicinity of the second order turning point the comparison potential $`V_c`$ is a parabolic one. The curvature of the latter potential can be always chosen to guarantee the two intersection points always exist. The choice of the comparison potential corresponds to a certain renormalization ($`\gamma ^2`$) of the characteristic oscillation frequency $`(d^2V/dX^2)_{X=X_0}`$. Note in passing that the approach we are advocating here conceptually close (although not identical) to the scale transformation proposed by Miller and Good MG53 and further developed inPE71 .
Thus the conditions $`J(X_{1,2}^\mathrm{\#})=0`$ allow us to construct well controlled approximate solutions to the Schrödinger equation. The accuracy of the approximation depends on the deviation of the approximate wave functions from the exact ones in the vicinity of the characteristic points $`X_{1,2}^\mathrm{\#}`$. In own turn, the deviation is determined by the higher over $`(XX_0)`$ terms of the potential $`V(X)`$ which are not included in the harmonic comparison potential $`V_c`$. Include explicitely the corresponding higher order terms to distinguish $`V_c`$ and $`V_{sc}`$ potentials, we find in the vicinity of the linear turning points
$`V_{sc}^{(1)}V_c{\displaystyle \frac{c_1}{\gamma ^2}}\left(XX_0\right)^2+{\displaystyle \frac{\omega ^2}{2}}\left(XX_0\right)^2,`$ (13)
where the universal numerical constant $`c_1=5/32`$, and the second term in the r.h.s. is related to deviation of the bare potential from the linear one. The same manner near the second order turning points
$`V_{sc}^{(2)}V_c{\displaystyle \frac{c_2}{\gamma ^2}}\left(XX_0\right)^2+\alpha (XX_0)^3,`$ (14)
where the universal constant $`c_2=n(n1)/4`$ is zero for the lowest vibrational states $`n=0,\mathrm{\hspace{0.17em}1}`$, and the last term describes non-parabolicity of the potential. Note that unlike the semiclassical action which within the instanton method is independent of quantum numbers $`n`$, the position of the characteristic points does depend on $`n`$, and the $`X_{1,2}^\mathrm{\#}`$ points are placed near the boundaries of the classically accessible region.
Now we are in the position to construct the approximant wave functions
$`\stackrel{~}{\mathrm{\Psi }}(X)=\{\begin{array}{c}\mathrm{\Psi }_c(X),X_1^\mathrm{\#}<X<X_2^\mathrm{\#}\\ \mathrm{\Psi }_{sc},X<X_1^\mathrm{\#},X>X_2^\mathrm{\#}\end{array},`$ (17)
which are the solutions to the Schrödinger equation with the following piecewise smooth approximating potential
$`\stackrel{~}{V}(X)=\{\begin{array}{c}V_c(X),X_1^\mathrm{\#}<X<X_2^\mathrm{\#}\\ V_{sc},X<X_1^\mathrm{\#},X>X_2^\mathrm{\#}\end{array}.`$ (20)
The wave functions calculated according to (17) in the framework of the instanton approach close to the Weber functions in the classically accessible regions, but their exponentially decaying tails in the classically forbidden regions correspond to the exact (bare) potential, not to its harmonic approximant. Analogously in the WKB method these functions (17) coincide with the semiclassical ones out of the interval $`(X_1^\mathrm{\#},X_2^\mathrm{\#})`$, and with the Airy functions in this interval.
To proceed further on we have to relate our approximant wave functions (17) and two linearly independent solutions to the bare Schrödinger equation $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$. It can be written down formally as
$`\mathrm{\Psi }(X)=\stackrel{~}{\mathrm{\Psi }}(X)+{\displaystyle _{X_0}^X}𝑑X^{}v(X^{})G(X,X^{})\mathrm{\Psi }(X^{}),`$ (21)
where $`v=V(X)\stackrel{~}{V}(X)`$, $`G(X,X^{})`$ is the Green function for the Schrödinger equation with the potential (20),
$`G(X,X_1)const[\stackrel{~}{\mathrm{\Psi }}_1(X)\stackrel{~}{\mathrm{\Psi }}_2(X_1)\stackrel{~}{\mathrm{\Psi }}_1(X_1)\stackrel{~}{\mathrm{\Psi }}_2(X)];X_1X,`$ (22)
where the constant in (22) is the Wronskian equal to $`(2\gamma )^1`$ in the instanton method, and $`i(2\gamma )^1`$ within the WKB approach. In (21) $`\mathrm{\Psi }(\mathrm{\Psi }_1,\mathrm{\Psi }_2)`$ (the same definition for $`\stackrel{~}{\mathrm{\Psi }}`$), and we take the turning point $`X_0`$, where the functions $`\mathrm{\Psi }`$ and $`\stackrel{~}{\mathrm{\Psi }}`$ are close to each other at the lower integration limit.
The solution to the integral equation (21) is expressed as the Neumann series expansion,
$`\mathrm{\Psi }(X)=\stackrel{~}{\mathrm{\Psi }}(X)+{\displaystyle _{X_0}^X}𝑑X_1v(X_1)G(X,X_1)\stackrel{~}{\mathrm{\Psi }}(X_1)+\mathrm{},`$ (23)
and the $`m`$-th order term can be factorized and estimated as
$`{\displaystyle \frac{1}{m!}}\left({\displaystyle _{X_0}^X}𝑑X_1v(X_1)G(X,X_1)\stackrel{~}{\mathrm{\Psi }}(X_1)\right)^m.`$ (24)
The integrals entering this estimation
$`{\displaystyle _{X_0}^X}𝑑X_1v(X_1)G(X,X_1)\stackrel{~}{\mathrm{\Psi }}(X_1)=L_{12/22}(X)\stackrel{~}{\mathrm{\Psi }}_1(X)+L11/21(X)\stackrel{~}{\mathrm{\Psi }}_2(X)`$ (25)
contain the $`2\times 2`$ matrix with the following matrix elements
$`L_{ij}={\displaystyle _{X_0}^X}𝑑X^{}\stackrel{~}{\mathrm{\Psi }}_i(X^{})v(X^{})\stackrel{~}{\mathrm{\Psi }}_j(X^{}).`$ (26)
It is convenient to introduce the matrix $`\widehat{C}^{(n)}`$ relating the $`n`$-th order wave function correction $`\delta \mathrm{\Psi }^{(n)}`$ with the wave function $`\stackrel{~}{\mathrm{\Psi }}`$
$`\left(\begin{array}{c}\delta \mathrm{\Psi }_1^{(n)}\\ \delta \mathrm{\Psi }_2^{(n)}\end{array}\right)=\widehat{C}^{(n)}\left(\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}_1\\ \stackrel{~}{\mathrm{\Psi }}_2\end{array}\right),`$ (31)
and the full connection matrix between the exact and approximate wave functions
$`\left(\begin{array}{c}\mathrm{\Psi }_1\\ \mathrm{\Psi }_2\end{array}\right)=\widehat{C}\left(\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}_1\\ \stackrel{~}{\mathrm{\Psi }}_2\end{array}\right),`$ (36)
is
$`\widehat{C}={\displaystyle \underset{n}{}}\widehat{C}^{(n)}.`$ (37)
According to the inequality (24)
$`\widehat{C}^{(n)}{\displaystyle \frac{1}{n!}}\widehat{C}_0^n,`$ (38)
where
$`\widehat{C}_0=\left(\begin{array}{cc}L_{12}& L_{11}\\ L_{22}& L_{12}\end{array}\right).`$ (41)
Combining finally the expressions (37) - (41) we end up with the upper and lower bounds for the correction matrix $`\widehat{C}`$ estimation
$`1+\widehat{C}_0\widehat{C}\mathrm{exp}(\widehat{C}_0).`$ (42)
We conclude from (24), (26) that the Neumann series posses an absolute convergence if all the matrix elements $`L_{ij}`$ are finite. Besides, unlike the correction functions introduced within the Fröman approach FF65 , the integrals in (26) have no singularities on the real axis. Evidently the integrals (26) are finite with the oscillating WKB functions, since the perturbation potential $`v(X)`$ is not zero only in the close proximity to the characteristic points $`X_{1,2}^\mathrm{\#}`$. However in the instanton method due to mixing of increasing and decreasing exponents, the matrix elements $`L_{22}`$ is divergent. Despite of this divergency the product $`L_{22}\stackrel{~}{\mathrm{\Psi }}_1`$, we are only interested in, is finite, and it is convenient to perform one more transformation to exclude explicitely this divergency.
Technically one can easily eliminate the both off-diagonal elements of the matrix $`\widehat{C}_0`$ and thus to get rid of the divergency of the (exponentially decreasing solution $`\stackrel{~}{\mathrm{\Psi }}_1`$) amplitude due to the contribution to the $`\stackrel{~}{\mathrm{\Psi }}_1`$ the exponentially increasing solution $`\stackrel{~}{\mathrm{\Psi }}_2`$. These linear transformations renormalize the correction matrix elements $`L_{22}`$ and $`L_{11}`$ as follows
$`L_{22}(X)\stackrel{~}{\mathrm{\Psi }}_1(X)L_{22}^{}(X)\stackrel{~}{\mathrm{\Psi }}_2(X);L_{11}(X)\stackrel{~}{\mathrm{\Psi }}_2(X)L_{11}^{}(X)\stackrel{~}{\mathrm{\Psi }}_1(X),`$ (43)
where the renormalized matrix elements $`L_{22}^{}`$ and $`L_{11}^{}`$ read as
$`L_{22}^{}(X)={\displaystyle _{X_0}^X}𝑑X^{}{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_1(X)\stackrel{~}{\mathrm{\Psi }}_2(X^{})}{\stackrel{~}{\mathrm{\Psi }}_1(X^{})\stackrel{~}{\mathrm{\Psi }}_2(X)}}\stackrel{~}{\mathrm{\Psi }}_1(X^{})v(X^{})\stackrel{~}{\mathrm{\Psi }}_2(X^{})L_{12}(X),`$ (44)
and
$`L_{11}^{}(X)={\displaystyle _{X_0}^X}𝑑X^{}{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_1(X^{})\stackrel{~}{\mathrm{\Psi }}_2(X)}{\stackrel{~}{\mathrm{\Psi }}_1(X)\stackrel{~}{\mathrm{\Psi }}_2(X^{})}}\stackrel{~}{\mathrm{\Psi }}_1(X^{})v(X^{})\stackrel{~}{\mathrm{\Psi }}_2(X^{})L_{12}(X).`$ (45)
Now all the integrals entering $`L_{11}^{}`$, $`L_{22}^{}`$, and $`L_{12}`$ are convergent and finite at any $`X`$, the correction matrix $`\widehat{C}_0`$ is transformed into the diagonal and positively defined matrix $`\widehat{C}_0^{}`$
$`\widehat{C}_0^{}=\left(\begin{array}{cc}L_{11}^{}L_{12}& 0\\ 0& L_{12}L_{22}^{}\end{array}\right).`$ (48)
The $`n`$-th order corrections in the instanton approach satisfy the inequality
$`0\delta \mathrm{\Psi }_1^{(n)}(n!)^1(L_{11}L_{12})\stackrel{~}{\mathrm{\Psi }}_1,\mathrm{\hspace{0.17em}0}\delta \mathrm{\Psi }_2^{(n)}(n!)^1(L_{12}L_{22}^{})\stackrel{~}{\mathrm{\Psi }}_2.`$ (49)
Explicit summation of r.h.s in (49) gives us the upper and the lower bound limits for the solutions of the initial Schrödinger equation, i.e. the stripe where increasing and decreasing solutions are confined
$`|\stackrel{~}{\mathrm{\Psi }}_1(X)||\mathrm{\Psi }_1(X)||\stackrel{~}{\mathrm{\Psi }}_1(X)|\mathrm{exp}(L_{11}^{}L_{12}),|\stackrel{~}{\mathrm{\Psi }}_2(X)||\mathrm{\Psi }_2(X)||\stackrel{~}{\mathrm{\Psi }}_2(X)|\mathrm{exp}(L_{12}L_{22}^{}).`$ (50)
It is our main result in this paper, and the stripe (50) gives the accuracy of the semiclassical instanton method. Besides we are in the position now to estimate the contribution of increasing semiclassical solutions into decreasing ones (what is relevant to solve eigenvalue problems). The summation convergent majorant series enables us to estimate the upper bound for this contribution
$`{\displaystyle \frac{L_{11}^{}}{L_{22}^{}}}\left(1\mathrm{exp}(L_{12})\right).`$ (51)
Therefore at $`L_{12}1`$ the summation of all order perturbation terms enhances the 1-st order correction by the factor $`L_{12}/L_{22}^{}`$. Analogously the majorant estimates described above can be used to construct the connection matrices linking the semiclassical solutions through the turning points. The comparison of the bare connection matrices (see e.g., HE62 , BM94 , BV02 ) with the matrices calculated accordingly to (26) - (44) provides the estimates for the eigenvalue accuracy.
As it was mentioned already, for the WKB method the procedure is even more simple, since no any divergency and therefore no need to perform the transformation (43). The similar to (24) - (25) factorization gives the $`\widehat{C}_0`$ matrix (cf. with (41) for the instanton approach)
$`\widehat{C}_0=i\left(\begin{array}{cc}L_{12}& L_{11}\\ L_{22}& L_{12}\end{array}\right),`$ (54)
and the estimations for the $`n`$-th order contribution (cf. with (49) can be formulated now as
$`|\widehat{C}^{(2n)}|{\displaystyle \frac{1}{(2n)!}}\left(\begin{array}{cc}\mathrm{\Delta }& 0\\ 0& \mathrm{\Delta }\end{array}\right)^n,`$ (57)
and
$`|\widehat{C}^{(2n+1)}|{\displaystyle \frac{1}{(2n+1)!}}\mathrm{\Delta }^n\widehat{C}_0,`$ (58)
where we denote $`\mathrm{\Delta }L_{12}^2L_{11}L_{22}`$. Correspondingly to (57), (58) the diagonal and off-diagonal correction matrix elements are bounded from above
$`|C_{11/22}|\left|\mathrm{cos}\sqrt{\mathrm{\Delta }}+iL_{12}{\displaystyle \frac{\mathrm{sin}\sqrt{\mathrm{\Delta }}}{\sqrt{\mathrm{\Delta }}}}\right|,`$ (59)
and
$`|C_{12/21}|\left|L_{11}{\displaystyle \frac{\mathrm{sin}\sqrt{\mathrm{\Delta }}}{\sqrt{\mathrm{\Delta }}}}\right|.`$ (60)
The whole procedure we employed is rationalized in the Fig. 2, where we compare the solutions to the comparison equation with the anharmonic oscillator semiclassical wave functions computed within the instanton and WKB approaches and indicate the optimal matching points $`X^\mathrm{\#}`$ found accordingly to the condition $`J(X^\mathrm{\#})=0`$.
## IV Anharmonic oscillator
In closing let us illustrate how our estimations (50), (58) work for a strongly anharmonic potential (11). Although it is not great triumph to re-derive the known results, our derivation illustrates several characteristic features of the correction matrix techniques derived in the section III: better accuracy, rapid convergence, simple disposal of divergences, and ease of computation in particular. The main message of our consideration in the precedent section III is that the quantitative accuracy of the semiclassics depends crucially on the proximity of the semiclassical wave functions to the solutions of the comparison equation in the region of the asymptotically smooth matching. Therefore, it is tempting to improve the accuracy by taking into account the anharmonic corrections to the comparison potential $`V_c(X)`$. However, since the eigenvalues and the normalization of the wave functions are almost independent of the detailed behavior in the vicinity of the linear turning points (because near these points, situated at the boundaries of the classically accessible region, the probability density (i.e. $`|\mathrm{\Psi }|^2`$) is exponentially small) this idea is useless for the WKB approach. In contrast with this, for the instanton method, the accuracy can be improved considerably upon including the anharmonic corrections into the comparison potential. Indeed, within the instanton method the accuracy is determined by the vicinity of the second order turning points where the wave functions acquire the largest values (and just in this region the smooth matching described above has to be performed).
Let us remind first the traditional (but formulated within the semiclassical framework) perturbation theory. Keeping in the transport equation (4) the second derivative of the prefactor $`A`$, the system of the equations of Hamilton - Jacoby (3) and the exact transport equation
$`{\displaystyle \frac{d^2A}{dX^2}}2\gamma \left[{\displaystyle \frac{dA}{dX}}{\displaystyle \frac{d\sigma }{dX}}+\left({\displaystyle \frac{d^2\sigma }{dX^2}}ϵ\right)A\right]=0`$ (61)
are exactly equivalent to the Schrödinger equation under consideration. For the second order turning points the anharmonic corrections
$`V_p={\displaystyle \frac{1}{2}}\left(\alpha X^3+\beta X^4\right)`$ (62)
can be considered as a perturbation and it is convenient to include this perturbation $`V_p`$ into the transport equation (61). Then the comparison equation is reduced to the inhomogeneous Weber equation, and its solutions can be expanded over the Weber functions $`D_\nu (X)`$ EM53 (see also OL59 , OL74 )
$`\mathrm{\Psi }_\nu (X)=N_\nu \left(D_\nu (X)+{\displaystyle \underset{k}{}}b_{k\nu }D_{\nu +k}(X)\right),`$ (63)
where $`N_\nu ^2=1+_kb_{\nu k}^2`$ is the wave function normalization factor, the expansion coefficients are proportional to the small parameters $`\alpha /\sqrt{\gamma }`$, and $`\beta /\gamma `$, and the Weber function index $`\nu `$ is related to the energy eigenvalue $`ϵ=\nu +(1/2)`$. This expansion (63) looks like a conventional perturbation series, but it does not. In the comparison equation we are keeping the both (decreasing and increasing) waves, and as a result of it, the indices of the Weber functions $`\nu +k`$ are not integer numbers. In the first order over the perturbation $`V_p`$ the only non-zero coefficients in (63) correspond to the following selection rules
$`k=\pm 1,\pm 3;\mathrm{and}k=0,\pm 2,\pm 4`$ (64)
for the cubic and fourth order anharmonic corrections respectively. Explicitely these non-zero expansion coefficients can be found by straitforward calculations, and they are
$`b_{0\nu }={\displaystyle \frac{3}{2}}\beta \left(\nu ^2+\nu +{\displaystyle \frac{1}{2}}\right),b_{1\nu }=3\alpha \nu ^2,b_{1\nu }=\alpha (\nu +1),b_{2\nu }{\displaystyle \frac{\beta }{2}}(\nu 1)\left(\nu {\displaystyle \frac{1}{2}}\right),b_{2\nu }{\displaystyle \frac{\beta }{2}}\left(\nu +{\displaystyle \frac{3}{2}}\right),`$ (65)
$`b_{3\nu }={\displaystyle \frac{1}{3}}\alpha \nu (\nu 1)(\nu 2)(\nu 3),b_{3\nu }={\displaystyle \frac{1}{3}}\alpha ,b_{4\nu }={\displaystyle \frac{1}{4}}\beta \nu (\nu 1)(\nu 2)(\nu 3)(\nu 4),b_{4\nu }={\displaystyle \frac{1}{4}}\beta .`$
On equal footing we can find the perturbative corrections to the Bohr-Sommerfeld quantization rules, and therefore the eigenvalues. The calculation is straightforward, though deserves some precaution and rather tedious. Skipping a large amount of tedious algebra we end up with the fractional part of the quantum number $`\nu `$
$`\nu n+\chi _n,`$ (66)
and up to the second order over the anharmonic perturbation $`V_p`$ we find
$`\chi _n^{(2)}={\displaystyle \frac{15\alpha ^2}{2\gamma }}\left(n^2+n+{\displaystyle \frac{11}{30}}\right)+{\displaystyle \frac{3\beta }{\gamma }}\left(n^2+n+{\displaystyle \frac{1}{2}}\right).`$ (67)
However the described standard perturbative approach leads to qualitatively wrong features of the solutions. For example, the wave functions (63), (65) are represented as a product of $`\nu `$ independent exponential factors and dependent of $`\nu `$ polynomials. As it is well known LL65 in one dimension the $`n`$-th excited state wave function must have $`n`$ zeros (and the number of zeros may not be changed by any perturbation). However in the $`m`$-th order perturbation theory approximation, the wave function (63) corresponding to a certain excited state $`n`$ contains Hermitian polynomials up to the order $`n+3m`$ or $`n+4m`$ for the cubic or quartic anharmonic perturbations, respectively. Therefore some false zeros of the wave function appears in the standard perturbation theory, and the region where the function oscillates becomes more and more wide in the higher order over perturbations approximation. The contributions of these qualitatively and quantitatively incorrect higher order terms become dominating in the asymptotic region at $`|\alpha \gamma |1`$, $`\beta \gamma 1`$. It conforms with the classical results due to Bender and Wu BW69 , BW73 who have shown that for the quartic anharmonic potential ($`\alpha =0`$, $`\beta >0`$ in (11)) the convergency radius is zero. Moreover, for $`\beta =0`$ and for an arbitrary small $`\alpha `$ (11) is the cubic anharmonic potential, i.e., the decay one. Thus it should have only complex eigenvalues, what is not the case for the eigenvalues calculated within the perturbation theory. The method we developed in section III enables us not only to estimate more accurate the anharmonic corrections to the eigenvalues, and to bring the whole schema of the calculations in a more elegant form. Our finding of the correction matrices is not merely to surpass a technical difficulty of the standard perturbative method, it is more one of principle, and we will show that the method has no drawbacks of the perturbation theory.
The proof proceeds as follows. Let us consider first the instanton method for the anharmonic potential (61) possessing one second order turning point $`X=0`$. As it was shown in the section III one has to find also two other characteristic points which are the roots of the equation (12). We denote the points as $`X_L^\mathrm{\#}`$, and $`X_R^\mathrm{\#}`$ (to refer by the self-explanatory subscripts $`L`$ and $`R`$ to the left and to the right from the turning point $`X=0`$). At the next step using the correction matrices introduced in the section III, we can define formally the transformation of our approximate wave functions (17) into the unknown exact wave functions $`\mathrm{\Psi }`$ as
$`\left(\begin{array}{c}\mathrm{\Psi }_L\\ \mathrm{\Psi }_R\end{array}\right)=\widehat{C}\left(\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}_L\\ \stackrel{~}{\mathrm{\Psi }}_R\end{array}\right),`$ (72)
where as above the subscripts $`L`$ and $`R`$ refer to the wave functions in the regions to the left and to the right from the turning point $`X=0`$. We do not know the correction matrix $`\widehat{C}`$ but we do know (see (31) - (45)) the boundary estimations for the matrix.
The Eq. (72) can be used also to correct the known at the second order turning point the connection matrix $`\widehat{M}`$ HE62 , BV02 . Indeed the connection matrices link the semiclassical solutions in the $`X`$-regions to the left and to the right from the turning points. For the isolated second order turning point (we are dealing within the instanton method), the connection matrix $`\widehat{M}`$ links the exponentially increasing and decreasing solutions in the space regions separated by the turning point. The condition ensuring the correct asymptotic behavior is the quantization rule for this case which can be formulated as $`M_{11}=0`$ ($`M_{ij}`$ are the matrix elements of the connection matrix $`\widehat{M}`$). Since in the regions to the left and to the right from the turning point our approximate solutions $`\stackrel{~}{\mathrm{\Psi }}`$ coincide by their definition (17) with the semiclassical ones, the correction matrix method enables us to correct the quantization rule too. Namely, the quantization rules are formulated within the connection matrix technique read now as
$`C_{22}^RT_2C_{22}^LC_{21}^RC_{21}^L{\displaystyle \frac{\mathrm{sin}^2(\pi \nu )}{T_2}}+(C_{22}^LC_{21}^R+C_{21}^LC_{22}^R)\mathrm{cos}(\pi \nu )=0,`$ (73)
where the Stokes constant for the second order turning point HE62 is
$`T_2={\displaystyle \frac{\sqrt{2\pi }}{\mathrm{\Gamma }(\nu )}}),`$ (74)
and $`C_{ij}^{R,L}`$ are the correction matrices at the $`X_R^\mathrm{\#}`$ or $`X_L^\mathrm{\#}`$ characteristic points respectively. Expanding the Gamma function entering (73) around the integer numbers, i.e., as above (66), $`\nu =n+\chi _n`$ one can find from the equation (73) the fractional part of the quantum number. If we were known the correction matrix $`\widehat{C}`$ the solution of (73) would provide the exact eigenvalues. But we do know only the estimations from below and from above for the $`\widehat{C}`$ matrix. In the same spirit we can calculate the estimations for the correction matrices (and therefore for the eigenvalues) within the WKB approach.
Luckily it turns out that the mathematical nature of the semiclassical problem is on our side here, and, in fact, even the first order estimation from below $`\widehat{C}^{(0)}`$ (31) gives already the accuracy comparable with the standard perturbation procedure, and the estimation from above (37) gives the eigenvalues almost indistinguishable from the ”exact” ones obtained by the numerical diagonalization of the Hamiltonian. The same true for the wave functions found by the correction matrix technique. We show in Fig. 3 $`|\mathrm{\Psi }_3|^2`$ for the same anharmonic potential (11) with $`\alpha =1.25`$, $`\beta =0.5`$. Clearly the exact numerical results and those obtained by our correction matrix techniques are correct qualitatively and in the very good quantitative agreement (indistinguishable starting from the second order approximation) unlike the situation with the standard perturbation theory. Besides we present in the table the eigenvalues of the anharmonic potential. We take the anharmonic coefficients $`\alpha `$ and $`\beta `$ in (11) so large, that corresponding perturbations of the eigen values are of the order of the bare harmonic frequency (one in our dimensionless units $`\alpha =1.2,\beta =0.5`$). In the table the eigenvalues found by the numerical diagonalization are presented in the column $`I`$. The column $`II`$ contains the harmonic approximation results, the column $`III`$ is the second order perturbation theory (67), and the columns $`IV`$ and $`V`$ results are obtained by applying our correction matrix technique: estimation from below with the first order correction matrix (31) in the column $`IV`$, and the estimation from above with the matrix (37) in the column $`V`$.
To conclude, as we have shown how to estimate the corrections to the main technical tool for the semiclassical approach, the connection matrices linking the solutions to the left and to the right from the turning points. Everything (e.g., the upper and the lower bounds for $`\chi _n`$) is determined by the matrices $`L_{12}`$ and $`L_{22}^{}`$. Approximating to (11) potential is found from (13) and after that straightforward computing according to (26) the matrices $`L_{ij}`$ and their renormalization (44) leads to the corrections we are looking for, presented in the Fig. 4, which allow us to estimate the accuracy of the semiclassical eigenstates and eigenfunctions. We conclude that even for a strongly anharmonic potential the both methods (WKB and instanton) are fairly accurate ones (about $`5\%`$) up to the energy close to the potential barrier top (in the region of negative curvature, we already discussed above). It is worth noting that in the frame work of the conventional perturbation theory (due to zero convergency radius with respect to $`\beta `$ coefficient in (11)) pure computational problems to get the same accuracy become nearly unsurmountable, see e.g., BW69 , TU84 . Our findings show that to estimate quantitatively the semiclassical accuracy it is enough to compare two linearly independent (with the same quantum number) solutions of the initial potential under study, and of the approximating piecewise smooth potential. The main advantage of the approach is related to the appropriate (13), (20) choice of the approximating potential, providing absolutely convergent majorant series (21) for the solutions. Actually our correction matrix technique is a fairly universal one and enables to estimate (and improve!) the semiclassical accuracy for arbitrary one dimensional potentials with any combination of the turning and of the crossing points.
###### Acknowledgements.
The research described in this publication was made possible in part by RFFR Grants. One of us (E.K.) is indebted to INTAS Grant (under No. 01-0105) for partial support, and V.B. and E.V. are thankful to CRDF Grant RU-C1-2575-MO-04.
Figure Caption
Fig. 1
The characteristic semiclassical potentials $`V_{sc}`$ (dot-dashed lines) and $`V_c`$ (dashed lines) for the bare anharmonic potential (11): $`\alpha =1.25`$, $`\beta =0.5`$ ($`\gamma =33`$, and the energy window corresponds to $`n=3`$):
(a) instanton approach;
(b) WKB method.
Fig. 2
Semiclassical wave function $`\mathrm{\Psi }_3`$ ($`n=3`$) for the anharmonic potential (11) with $`\alpha =1.25`$, $`\beta =0.5`$ ($`\gamma =33`$). Stars indicate the matching points, dashed lines show the solutions to the comparison equations, and:
(a) - solid line traces the instanton solution;
(b) - solid line shows the WKB wave function.
Fig. 3
Comparison of the exact $`|\mathrm{\Psi }_3|^2`$ (solid line) for the anharmonic potential (11) with $`\alpha =1.25`$, $`\beta =0.5`$ ($`\gamma =33`$) with two lowest (zero and first order) approximations of the correction matrix method (dashed and dot-dashed lines). Note that the second order approximation with the relative accuracy $`10^2`$ is indistinguishable from the exact numerical results.
Fig. 4
Corrections to the decreasing solutions in for the anharmonic potential (11);
$`\beta =0.5,n=3,\gamma =33`$;
(1, 3) - instanton method, (2, 4) - WKB, (1, 2) - the first order corrections $`L_{11}`$, (3 , 4) - the upper bound estimation summing up all terms. The vertical dashed lines indicate the values of the cubic anharmonic term ($`\alpha `$) where the inflection point and new extrema of the potential are appeared.
Table
Eigenvalues of the anharmonic potential (11) ($`\alpha =1.2,\beta =0.5`$).
$`I`$ the eigenvalues found by the numerical diagonalization;
$`II`$ the harmonic oscillator eigenvalues;
$`III`$ the eigenvalues in the second order perturbation theory (67);
$`IV`$ the eigenvalues estimated from below by the correction matrix (31);
$`V`$ the estimation from above with the matrix (37).
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# 1 Introduction
## 1 Introduction
Unlike the standard model, which contains many free parameters, it is expected that string theory, should be able to do a much better job, as regards explaining things like quark and lepton masses and the gauge hierarchy, etc, from the vacuum expectation values (vev’s) of the moduli fields. In addition to giving vevs to the moduli (), string compactifications in the presence of fluxes, can also give a positive cosmological constant , making the study even more phenomenologically interesting. The inclusion of fluxes necessarily requires the introduction of warp factors in the metric , which can be argued to provide a mechanism for generating the large hierarchy of scales - the (not-so) recent TeV-scale quantum gravity proposal being particularly significant. The study of fluxes is also of much importance from the point of view of open/closed (topological) string dualities . Thus, turning on fluxes is of great importance for establishing connection between string theory and the observable universe.
Turning on $`NSNS`$ fluxes in, e.g., Heterotic theory, the internal manifold cannot be Kähler anymore implying that one then has torsion. As shown in , Iwasawa manifold, an example of a “half-flat” manifold (see also and references therein) satisfies the conditions for $`𝒩=1`$ supersymmetry in the presence of $`NSNS`$ fluxes. Based on , we study the explicit uplifts of the Iwasawa metric at the “standard complex structure” point of the moduli space of almost complex structures, to seven-folds with either $`G_2`$ holonomy or $`SU(3)`$ structure.
The plan of the paper is as follows. In Section 2, we give a brief review of six-folds with $`SU(3)`$ structure, the relevant torsion classes and the Iwasawa manifold. In Section 3, we discuss the uplift of the Iwasawa manifold to seven-folds of $`G_2`$ holonomy using Hitchin’s flow equations using “size” deformations and we show the impossibility of the uplift using “size” and “shape” deformations of the Iwasawa via Hitchin’s flow equations. We then give a brief review of seven-folds with $`SU(3)`$ structure and the relevant torsion classes and then discuss the uplift of the Iwasawa (at the “standard complex structure” point) to seven-folds with $`SU(3)`$ structure using “size” and “shape” deformations via generalizations of the Hitchin’s flow equations. Section 4 has the interpretation of the non-vanishing Bianchi identity for $`M`$-theory compactified on seven-folds with $`SU(3)`$ structure, in the presence of $`G`$-flux sourced by $`M5`$-branes wrapping two-cycles embedded in the seven-fold (analagous to ) and the generic movement in the moduli space of almost complex structures on the Iwasawa by “size” and “shape” deformations of the Iwasawa metric. Using the results of , we show that one can associate a Riemann surface with a pair of “shape deformation” functions and the dilaton, via the Weierstra$`ß`$ representation for the conformal immersion of a surface in $`𝐑^l`$ for a suitable $`l`$. Section 5 has the conclusion. There are two appendices.
## 2 Six-Folds with $`SU(3)`$ Structure
In this section, we give a brief review of six-folds with $`SU(3)`$ structure, relevant to, e.g., Heterotic theory compactifications in the presence of $`NSNS`$ flux, and a short review of the “half-flat” Iwasawa manifold.
### 2.1 Brief review of torsion classes
As per the work of , the requirements for $`𝒩=1`$ supersymmetric compactification of heterotic string theory are:
(a) The internal 6-dim manifold has to be complex. That means that the Nijenhuis tensor $`N_{mnp}`$ has to vanish.
(b) Up to a constant factor, there is exactly one holomorphic (3,0)-form $`\omega `$ whose norm is related to the complex structure $`J`$ by $`dJ=i(\overline{})\mathrm{log}||\omega ||`$.
(c) The Yang Mills background field strength must be a (1,1)-form and must satisfy $`trFF=tr\stackrel{~}{R}\stackrel{~}{R}i\overline{}J`$ as well as $`F_{mn}J^{mn}=0`$ (the Donaldson-Uhlenbeck-Yau condition), $`\stackrel{~}{R}`$ being the modified curvature two-form in the presence of torsion.
(d) The Warp factor is given by $`(y)=\varphi (y)+\mathrm{const}`$ ; the dilaton by $`\varphi (y)=\frac{1}{8}\mathrm{log}\omega +\mathrm{const}`$.
(e) The background 3-form $`H`$ is determined in terms of $`J`$ by $`H=\frac{i}{2}(\overline{})J.`$
It is possible to reformulate these conditions in terms of torsional constraints . NS-NS flux $`0`$ requires $`\stackrel{~}{Y}`$ to be a manifold with $`SU(3)`$ structure but not $`SU(3)`$ holonomy. $`SU(3)`$ structure implies: $``$ J, $`\mathrm{\Omega }`$ such that $`dJ0`$ ; $`d\mathrm{\Omega }0`$ which means that the Manifold is not Kähler and then can also not be CY. $`SU(3)`$-structure can now be determined in terms of torsion classes.
The difference between any two metric-compatible connections is a tensor, known as the contorsion $`\kappa _{mnp}`$ defined via:
$$_m^{(T)}\eta =_m\eta \frac{1}{4}\kappa _{mnp}\mathrm{\Gamma }^{np}\eta =0,$$
$`\eta `$ being the globally defined spinor that is covariantly constant w.r.t. the connection modified by the the three-form flux $`H`$. The contorsion can be related to torsion $`T_{mnp}`$ through
$$T_{mnp}=\frac{1}{2}(\kappa _{mnp}\kappa _{nmp}).$$
The torsion classes can be defined in terms of J, $`\mathrm{\Omega }`$, dJ, $`d\mathrm{\Omega }`$ and the contraction operator $`\mathrm{\_}|:\mathrm{\Lambda }^kT^{}\mathrm{\Lambda }^nT^{}\mathrm{\Lambda }^{nk}T^{}`$ where $`J`$ is given by:
$$J=e^1e^2+e^3e^4+e^5e^6,$$
and the (3,0)-form $`\mathrm{\Omega }`$ is given by
$$\mathrm{\Omega }=(e^1+ie^2)(e^3+ie^4)(e^5+ie^6).$$
The basis of $`(1,0)`$-forms is given by
$$e^i+iJe^i\mathrm{\Lambda }^{(1,0)},$$
where $`Je^a=J_{}^{a}{}_{b}{}^{}e^b`$ and consequently $`JJ=1`$.
The torsion classes are defined in the following way: $`W_1[dJ]^{(3,0)}`$, given by real numbers $`W_1=W_1^++W_1^{}`$ with $`d\mathrm{\Omega }_+J=\mathrm{\Omega }_+dJ=W_1^+JJJ`$ and $`d\mathrm{\Omega }_{}J=\mathrm{\Omega }_{}dJ=W_1^{}JJJ`$; similarly for $`W_2[d\mathrm{\Omega }]_0^{(2,2)}`$ : $`(d\mathrm{\Omega }_+)^{(2,2)}=W_1^+JJ+W_2^+J`$ and $`(d\mathrm{\Omega }_{})^{(2,2)}=W_1^{}JJ+W_2^{}J`$; $`W_3[dJ]_0^{(2,1)}`$ is defined as $`W_3=dJ^{(2,1)}[JW_4]^{(2,1)}`$; $`W_4JdJ`$ : $`W_4=\frac{1}{2}J\mathrm{\_}|dJ`$; $`W_5[d\mathrm{\Omega }]_0^{(3,1)}`$: $`W_5=\frac{1}{2}\mathrm{\Omega }_+\mathrm{\_}|d\mathrm{\Omega }_+`$ (the subscript 0 indicative of the primitivity of the respective forms). Depending on the classes of torsion one can obtain different types of manifolds:
* (complex) special-hermitian manifolds with $`W_1=W_2=W_4=W_5=0`$ which means that $`\tau W_3`$;
* (complex) Kähler manifolds with $`W_1=W_2=W_3=W_4=0`$ which means $`\tau W_5`$;
* (complex) balanced Manifolds with $`W_1=W_2=W_4=0`$ which means $`\tau W_3W_5`$;
* (complex) Calabi-Yau manifolds with $`W_1=W_2=W_3=W_4=W_5=0`$ which means $`\tau =0`$.
* half-flat manifolds (may or may not be complex) with $`W_1=W_2=0`$ whch means $`\tau W_3W_4W_5`$.
### 2.2 Heterotic string on the Iwasawa manifold
The torsional constraints are: $`\tau W_3W_4W_5`$$`2W_4+W_5=0`$ and $`W_{4,5}`$ real, exact. We consider now manifolds fullfilling the torsional constraints and satisfying supersymmetry requirements. From this we obtain, that we need special-hermitian manifolds for which the torsion is $`\tau W_3`$; they are complex and half-flat. Since $`W_4=W_5=0`$ it follows that the dilaton is constant. We consider nilmanifolds (6-dim) which are special-hermitian manifolds which are constructed from simply-connected nilpotent Lie group G by quotienting with discrete subgroup $`\mathrm{\Gamma }`$ of G for which $`G\mathrm{\Gamma }`$ is compact.
There are 34 classes of such manifold and they do not admit a Kähler metric. 18 of these admit complex structure. An example is the Iwasawa manifold. For any of 18 classes one can choose complex structure, compatible with metric and $`W_4=0`$. Iwasawa manifold is a nilmanifold obtained as the compact quotient space $`M=\mathrm{\Gamma }G`$, where $`G`$ the complex Heisenberg group is given by a set of matrices under multiplication
$$G=\{\left(\begin{array}{ccc}1& z& u\\ 0& 1& v\\ 0& 0& 1\end{array}\right):u,v,zC\}.$$
The discrete subgroup $`\mathrm{\Gamma }`$ is defined by restricting $`u,v,z`$ to Gaussian integers:
$`v`$ $``$ $`v+m`$
$`z`$ $``$ $`z+m`$
$`u`$ $``$ $`u+p+nv`$
where $`m,n,pZiZ`$. To find an Iwasawa manifold solution of torsional constraints the parameters have to be choosen such that the torsion lies in $`\tau W_3`$ with $`2W_4+W_5=0`$ and $`W_4=W_5=0`$. It turns out that the moduli space of complex structures has two disconnected components referred as the standard complex structure $`J_0`$ and ’edge’. It can be shown that in both components $`W_4=W_5=0`$.
The Standard complex stucture is given by: $`J_0=e^1e^2+e^3e^4+e^5e^6`$ with the (1,0)-forms are given by
$$\alpha =e^1+ie^2,\beta =e^3+ie^4,\gamma =e^5ie^6.$$
Complex coordinates $`(z,v,u)`$ can be introduced
$`\alpha `$ $`=`$ $`dz,`$
$`\beta `$ $`=`$ $`dv,`$
$`\gamma `$ $`=`$ $`i(duzdv).`$
These are holomorphic left-invariant (1,0)-forms with respect to the standard complex structure.
The two-form in the standard complex structure limit, is given by
$$J_0=\frac{i}{2}[dzd\overline{z}+dvd\widehat{v}+(duzdv)(d\widehat{u}\overline{z}d\overline{v})],$$
thereby implying that the metric is $`ds^2=|dz|^2+|dv|^2+|duzdv|^2`$. The Iwasawa manifold can thus be viewed as $`T^2`$ fibration over a $`T^2\times T^2`$ base. The Euler characteristic for the same is hence zero. The 3-form is given by $`\mathrm{\Omega }=\alpha \beta \gamma =idzdvdu.`$
## 3 Uplift
$`𝒩=1`$ theories in four dimensions starting from M-theory in the presence of G ($``$ 4-form) fluxes require 7-folds with either $`G_2`$ holonomy or $`SU(3)`$ structure. We first discuss the uplift of the Iwasawa to seven-folds with $`G_2`$ holonomy.
### 3.1 Hitchin’s construction of 7-folds with $`G_2`$ Holonomy from 6-folds that are half-flat
From a six dimensional manifold $`M`$ with $`SU(3)`$ structure ($`M,J,\mathrm{\Omega }`$) one can construct a 7-dim manifold as a warped product
$$X_7=M\times I;I𝐑.$$
A $`G_2`$ -Manifold is defined by: $`\varphi =Jdt+\mathrm{\Omega }_+`$ where the calibration $`\varphi `$ is closed and coclosed and $`(J,\mathrm{\Omega }_+,\mathrm{\Omega }_{})`$ are $`t`$-dependent. This implies:
$`d\varphi =\left(\widehat{d}J{\displaystyle \frac{\mathrm{\Omega }_+}{t}}\right)dt+\widehat{d}\mathrm{\Omega }_+=0;`$
$`d\varphi =\left(\widehat{d}\mathrm{\Omega }_{}J{\displaystyle \frac{J}{t}}\right)dt+J\widehat{d}J=0.`$ (1)
Through this the forms have now been promoted to seven dimensions. $`M`$ is half-flat which means that its $`SU(3)`$ structure is such that:
$$\widehat{d}\mathrm{\Omega }_+=J\widehat{d}J=0.$$
(2)
The conditions (3.1) for the seven-manifold $`X_7`$ to have $`G_2`$ holonomy , i.e. $`d\varphi =d\varphi =0`$, yield:
$`\widehat{d}J`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_+}{t}},`$ (3)
$`\widehat{d}\mathrm{\Omega }_{}`$ $`=`$ $`J{\displaystyle \frac{J}{t}}`$ (4)
called Hitchin’s Flow equations.
### 3.2 Uplift of the Iwasawa manifold to $`G_2`$-holonomy manifold via Hitchin’s Flow equation
In this subsection, we discuss the uplift of the Iwasawa to a seven-fold of $`G_2`$-holonomy using “size”(implying $`t`$-dependent) deformation of the Iwasawa metric. We later show that “size” alongwith nontrivial “shape” (implying $`z,\overline{z};v,\overline{v}`$-dependent) deformations of the Iwasawa metric can not be used to uplift to a $`G_2`$-holonomy manifold.
We use the following ansatz for $`J`$ with “size” deformation:
$`J`$ $`=`$ $`e^{a(t)}e^{12}+e^{b(t)}e^{34}+e^{c(t)}e^{56}`$ (5)
$`\mathrm{\Omega }`$ $`=`$ $`e^{\frac{1}{2}(a(t)+b(t)+c(t))}(e^1+ie^2)(e^3+ie^4)(e^5+ie^6)`$ (6)
We now use the following two-step algebra:
$`\widehat{d}e^1=\widehat{d}e^2=\widehat{d}e^3=\widehat{d}e^4=0;\widehat{d}e^5=e^{14}+e^{23};\widehat{d}e^6=e^{13}+e^{42};`$
with the representations:
$`e^1+ie^2=dz;e^3+ie^4=dv;e^5ie^6=i(duzdv).`$
$`\widehat{d}J`$ $`=`$ $`e^{c(t)}(e^{146}+e^{236}e^{513}e^{542})`$ (7)
The condition $`\widehat{d}\mathrm{\Omega }_+=J\widehat{d}J=0`$ is identically fullfilled. We solve $`\widehat{d}\mathrm{\Omega }_{}=J\frac{J}{t}`$:
$`\dot{a}=\dot{b}=\dot{c};`$
$`{\displaystyle \frac{d}{dt}}(e^{a+b})=4e^{\frac{a+b+c}{2}}`$ (8)
These are consistent with $`\widehat{d}J=\frac{\mathrm{\Omega }_+}{t}`$. For the metric we then obtain a one-parameter ($`\lambda `$) family of singular $`G_2`$-metrics:
$$ds_{7}^{}{}_{}{}^{2}=dt^2+(1+\lambda t)^{\frac{2}{3}}|dz|^2+(1+\lambda t)^{\frac{2}{3}}|dv|^2+(1+\lambda t)^{\frac{2}{3}}|duzdv|^2.$$
(9)
By specifying suitable boundary conditions at the end-points of the interval, one could perhaps use such manifolds in (heterotic-)$`M`$-theory (like) compactifications.
It is not possible to uplift the Iwasawa to $`G_2`$-manifold via both “size” and (non-trivial) “shape” deformations of the Iwasawa metric. We show the same now. Using e.g. the following ansatz:
$`J=e^{a(t)}e^{12}+A(z,\overline{z},v,\overline{v})e^{b(t)}e^{34}+B(z,\overline{z},v,\overline{v})e^{c(t)}e^{56}`$
$`\mathrm{\Omega }=e^{\frac{a+b+c}{2}}(e^{A_1}e^1+ie^{A_2}e^2)(e^{A_3}e^3+ie^{A_4}e^4)(e^{A_5}e^5+ie^{A_6}e^6),`$ (10)
with ”shape” and ”size” deformation we will obtain from the Hitchin’s Flow equations, the conditions for halfflatness and the additional condition $`\mathrm{\Sigma }A_i=A+B`$. From the halfflatness condition $`\widehat{d}\mathrm{\Omega }_+=0`$ we obtain:
$`A_{135}^{}{}_{z}{}^{}=A_{135}^{}{}_{\overline{z}}{}^{}`$
$`A_{135}^{}{}_{v}{}^{}=A_{135}^{}{}_{\overline{v}}{}^{}`$
$`A_{146}^{}{}_{z}{}^{}=A_{146}^{}{}_{\overline{z}}{}^{}`$
$`A_{146}^{}{}_{v}{}^{}=A_{146}^{}{}_{\overline{v}}{}^{}`$
$`A_{236}^{}{}_{z}{}^{}=A_{236}^{}{}_{\overline{z}}{}^{}`$
$`A_{236}^{}{}_{v}{}^{}=A_{236}^{}{}_{\overline{v}}{}^{}`$
$`A_{245}^{}{}_{z}{}^{}=A_{245}^{}{}_{\overline{z}}{}^{}`$
$`A_{245}^{}{}_{v}{}^{}=A_{245}^{}{}_{\overline{v}}{}^{},`$ (11)
where $`A_{ijk}^{}{}_{z\mathrm{or}\overline{z}}{}^{}\frac{(A_i+a_j+A_l)}{z(\mathrm{or}\overline{z})}`$. From the halfflatness condition $`J\widehat{d}J=0`$ we obtain:
$`A_z=B_z`$
$`A_{\overline{z}}=B_{\overline{z}}.`$ (12)
Considering the Hitchin’s Flow equation $`\frac{\mathrm{\Omega }_+}{t}=\widehat{d}J`$ it follows:
$`{\displaystyle \frac{\dot{a}+\dot{b}+\dot{c}}{2}}e^{\frac{a+b+c}{2}}e^{A_{ijk}}=e^ce^B`$
$`(ijk=135;146;236;245)`$
$`A_{135}=A_{146}=A_{236}=A_{245}.`$ (13)
The above equations give us
$$A_{135}=A_{146}=A_{236}=A_{245}=\mathrm{const}$$
(14)
From the other Hitchin’s Flow equation $`\widehat{\mathrm{\Omega }_{}}=J\frac{J}{t}`$ we obtain:
$`\dot{a}=\dot{b}=\dot{c}`$
$`e^Ae^{a+b}(\dot{a}+\dot{b})=e^{A_{136}}+e^{A_{145}}+e^{A_{235}}+e^{A_{246}}`$
$`A_{136}^{}{}_{z}{}^{}=A_{136}^{}{}_{\overline{z}}{}^{}`$
$`A_{136}^{}{}_{v}{}^{}=A_{136}^{}{}_{\overline{v}}{}^{}`$
$`A_{145}^{}{}_{z}{}^{}=A_{145}^{}{}_{\overline{z}}{}^{}`$
$`A_{145}^{}{}_{v}{}^{}=A_{145}^{}{}_{\overline{v}}{}^{}`$
$`A_{235}^{}{}_{z}{}^{}=A_{235}^{}{}_{\overline{z}}{}^{}`$
$`A_{235}^{}{}_{v}{}^{}=A_{235}^{}{}_{\overline{v}}{}^{}`$
$`A_{246}^{}{}_{z}{}^{}=A_{246}^{}{}_{\overline{z}}{}^{}`$
$`A_{246}^{}{}_{v}{}^{}=A_{246}^{}{}_{\overline{v}}{}^{}.`$ (15)
We can now use the additional condition:
$$A_{135}+A_{146}+A_{236}+A_{245}=2\mathrm{\Sigma }A_i=2(A+B)=\mathrm{const}$$
(16)
Thus $`A`$ has to be a constant what tells us, that there cannot be a “shape” deformation in $`J`$. Now we can use $`\mathrm{\Sigma }A_i=\mathrm{const}`$ again and write:
$`A_{123456}^{}{}_{z;\overline{z};v;\overline{v}}{}^{}=(A_{136}+A_{245})_{z;\overline{z};v;\overline{v}}=(A_{145}+A_{236})_{z;\overline{z};v;\overline{v}}`$
$`=(A_{235}+A_{146})_{z;\overline{z};v;\overline{v}}=(A_{246}+A_{135})_{z;\overline{z};v;\overline{v}}=0.`$ (17)
Using equation (14), it follows:
$$A_{136}=A_{145}=A_{235}=A_{246}=\mathrm{const}$$
(18)
That means that we also cannot have “shape” deformation in $`\mathrm{\Omega }`$. The uplift of Iwasawa manifold to $`G_2`$ holonomy can only be done with “size” deformation.
### 3.3 Intrinsic torsion classes for M-theory with fluxes in 7-dimension with $`SU(3)`$ structure
The difference compared to 6-dimensional case is the existence of a globally defined vector $`v`$. $`SU(3)`$ structure in d=7 is then described by a triplet $`v,J,\mathrm{\Omega }`$. The 2-form $`J`$ now satisfies $`J_a^bJ_b^c=\delta _a^c+v_av^c`$. The metric of seven-dim space can be written as $`ds_{7}^{}{}_{}{}^{2}(x,t)=ds_{6}^{}{}_{}{}^{2}(x,t)+vv`$ with $`v=e^{q\varphi (x)}dt`$. In 7 dimensions the decomposition of torsion gives us 15 classes:
$$\tau RC_{1,2}V_{1,2}W_{1,2}S_{1,2}𝒜_{1,2}T$$
where $`C_{1,2},W_{1,2}`$ and T are complex.
After using necessary and sufficient conditions for obtaining $`N=1`$ supersymmetry solutions of M theory with fluxes, and comparing with the decomposition of the 4-form $`G`$-flux given in terms of irreducible $`SU(3)`$ representations:
$$G=\frac{Q}{6}JJ+J𝒜+\mathrm{\Omega }_{}V+v(\alpha _1\mathrm{\Omega }_++\alpha _2\mathrm{\Omega }_{}+JW+U),$$
(19)
and the corresponding expression for $`_7G`$ (that involves an additional flux component “$`S`$”) we obtain:
$`R`$ $`=`$ $`C_1=W_1=W_2=𝒜_1=T=S_2=0;`$
$`C_2`$ $`=`$ $`\overline{C}_2=0;`$
$`V_1`$ $`=`$ $`{\displaystyle \frac{2}{3}}V_3=\sigma ;`$
$`𝒜_2`$ $`=`$ $`𝒜;`$
$`S_1`$ $`=`$ $`2S.`$
### 3.4 Irreducible 7-manifold
A naturally induced $`SU(3)`$-structure on seven-dimensional manifold without a killing isometry is given by:
$$v=e^{q\varphi }dt,J=\widehat{J}(t),\mathrm{\Omega }=\widehat{\mathrm{\Omega }}(t)$$
where $`t`$ is variable parameterizing the interval $`I`$. For simplicity we first discuss $`q=0`$:
$$ds_{7}^{}{}_{}{}^{2}(y,t)=ds_{6}^{}{}_{}{}^{2}(y,t)+dt^2.$$
The conditions $`\widehat{d}J=2S0,\widehat{d}\mathrm{\Omega }=0`$ give us that $`M_6`$ is a special-hermitian manifold,i.e. it is a complex non-Kähler manifold <sup>5</sup><sup>5</sup>5As the torsion classes $`W_1`$ and $`W_2`$ are defined via: $`(\widehat{d}\widehat{\mathrm{\Omega }}_\pm )^{(2,2)}=W_1^\pm \widehat{}\widehat{J}+W_2^\pm \widehat{J}`$, this implies $`W_1=W_2=0`$, implying the vanishing of the Nijenhuis tensor. Hence the manifold is a complex manifold. Further, $`W_5=\frac{1}{2}\widehat{\mathrm{\Omega }_+}\mathrm{\_}|\widehat{d}\widehat{\mathrm{\Omega }_+}=0`$. As $`W_1`$ is also identified with $`[\widehat{d}\widehat{J}]^{(3,0)}`$ and given that $`W_1=0`$ this implies that $`dJ\mathrm{\Lambda }^{(2,1)}\mathrm{\Lambda }^{(1,2)}`$\[As $`\widehat{d}\widehat{J}0`$, writing $`\alpha ^{(p,q)}`$ in terms of a real $`\alpha ^{p+q}`$: $`\alpha _{m_1\mathrm{}..m_{p+q}}^{(p,q)}=(P^+)_{m_1}^{n_1}\mathrm{}(P^+)_{m_p}^{n_p}(P^{})_{m_{p+1}}^{n_{p+1}}\mathrm{}.(P^{})_{m_{p+q}}^{n_{p+q}}\alpha _{n_1\mathrm{}n_{p+q}}^{p+q}`$, where $`(P^\pm )_m^n\frac{1}{2}(\delta _m^n\pm iJ_m^n)`$, this implies that $`\widehat{d}\alpha ^{(p,q)}\mathrm{\Lambda }^{(p+2,q1)}\mathrm{\Lambda }^{(p+1,q)}\mathrm{\Lambda }^{(p,q+1)}\mathrm{\Lambda }^{(p1,q+2)}`$, where $`\alpha ^{(p,q)}\mathrm{\Lambda }^{(p,q)}`$.\]. Further, as $`W_3`$ is identified with $`[\widehat{d}\widehat{J}]^{(2,1)}`$ and $`\widehat{d}\widehat{J}^{(2,1)}=(\widehat{J}W_4)^{(2,1)}+W_3`$, one concludes that $`W_4=0`$. Thus, the six-fold is a special hermitian manifold.. One can now build a 7-manifold, that can be used in flux solutions of M-theory by solving diff̃erential equations in $`t`$. This construction generalizes the Hitchin’s Flow equations for construction of $`G_2`$-holonomy manĩfold:
$`{\displaystyle \frac{J}{t}}={\displaystyle \frac{2}{3}}QJ2𝒜`$
$`{\displaystyle \frac{\mathrm{\Omega }}{t}}=Q\mathrm{\Omega }.`$ (20)
We now use an ansatz for $`J`$ with “shape” and “size” deformations:
$$J=e^{a(t)}e^{12}+A(z,\overline{z},v,\overline{v})e^{34}+B(z,\overline{z},v,\overline{v})e^{56}$$
(21)
where $`Im(A)=Im(B)=0`$.
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2}}\widehat{d}J={\displaystyle \frac{1}{2}}[(e^{134}+ie^{234}){\displaystyle \frac{A}{z}}+(e^{134}ie^{234}){\displaystyle \frac{A}{\overline{z}}}`$ (22)
$`+(e^{156}+ie^{256}){\displaystyle \frac{B}{z}}+(e^{156}ie^{256}){\displaystyle \frac{B}{\overline{z}}}+(e^{356}+ie^{456}){\displaystyle \frac{B}{v}}`$
$`+(e^{356}ie^{456}){\displaystyle \frac{B}{\overline{v}}}+B(e^{135}+e^{245}+e^{146}+e^{236})]`$
The Bianchi identity reads:
$`{\displaystyle \frac{dQ}{6}}JJ+{\displaystyle \frac{2}{9}}Q^2JJv{\displaystyle \frac{4}{3}}QvJ𝒜{\displaystyle \frac{Q}{3}}JJ\sigma `$
$`+2S𝒜+2v𝒜𝒜+J\sigma 𝒜3v\sigma JW`$
$`2vSW+vJdW+dS\stackrel{\mathrm{?}}{=}0`$ (23)
where the question mark over the equality sign is indicative of the possibility that one can not satisfy a source-free Bianchi identity (as will be the case).
With $`QJ=\frac{3}{2}[\dot{J}+2A]Q=\frac{3}{2}\dot{a}(t)`$ it follows:
$`{\displaystyle \frac{dQ}{6}}JJ={\displaystyle \frac{\ddot{a}}{2}}dt[e^aAe^{1234}+e^aBe^{1256}+ABe^{3456}]`$ (24)
$`{\displaystyle \frac{2}{9}}Q^2JJv=\dot{a}^2dt[e^aAe^{1234}+e^aBe^{1256}+ABe^{3456}]`$ (25)
And with $`𝒜\frac{1}{2}\dot{a}(t)[A(z,\overline{z},v,\overline{v})e^{34}+B(z,\overline{z},v,\overline{v})e^{56}]`$ it follows:
$`{\displaystyle \frac{4}{3}}QvJA=\dot{a}^2dt[e^aAe^{1234}+e^aBe^{1256}+2ABe^{3456}]`$ (26)
$`2S𝒜={\displaystyle \frac{\dot{a}(t)}{2}}[e^{13456}({\displaystyle \frac{}{z}}+{\displaystyle \frac{}{\overline{z}}})(AB)+ie^{23456}({\displaystyle \frac{}{z}}{\displaystyle \frac{}{\overline{z}}})(AB)]`$ (27)
$`2v𝒜𝒜=\dot{a}^2ABdte^{3456}`$ (28)
$`d_7S={\displaystyle \frac{1}{2}}dt[2e^{1256}({\displaystyle \frac{^2A}{z^2}}+{\displaystyle \frac{^2A}{\overline{z}^2}})`$
$`+2e^{1234}({\displaystyle \frac{^2B}{z^2}}+{\displaystyle \frac{^2B}{\overline{z}^2}}+{\displaystyle \frac{^2B}{v^2}}+{\displaystyle \frac{^2B}{\overline{v}^2}}+2B)`$
$`e^{2356}({\displaystyle \frac{^2A}{zv}}+{\displaystyle \frac{^2A}{\overline{z}v}}+c.c.)ie^{2456}({\displaystyle \frac{^2A}{zv}}+{\displaystyle \frac{^2A}{\overline{z}v}}c.c.)`$
$`ie^{1356}({\displaystyle \frac{^2A}{zv}}{\displaystyle \frac{^2A}{\overline{z}v}}c.c.)+e^{1456}({\displaystyle \frac{^2A}{zv}}{\displaystyle \frac{^2A}{\overline{z}v}}+c.c.)`$
$`i(e^{1346}+e^{2345})({\displaystyle \frac{B}{v}}{\displaystyle \frac{B}{\overline{v}}})(e^{1345}e^{2346})({\displaystyle \frac{B}{v}}+{\displaystyle \frac{B}{\overline{v}}})`$
$`i(e^{1245}+e^{1236})[({\displaystyle \frac{B}{z}}{\displaystyle \frac{B}{\overline{z}}})+({\displaystyle \frac{A}{z}}{\displaystyle \frac{A}{\overline{z}}})]`$
$`+i(e^{1246}+e^{1235})[({\displaystyle \frac{B}{z}}+{\displaystyle \frac{B}{\overline{z}}})({\displaystyle \frac{A}{z}}+{\displaystyle \frac{A}{\overline{z}}})]]`$ (29)
The Hodge dual of $`dG`$ then reads:
$`_{11}dG=e^{4\mathrm{\Delta }}dx^0dx^1dx^2dx^3`$ (30)
$`[({\displaystyle \frac{\ddot{a}}{2}}e^aA({\displaystyle \frac{^2B}{z^2}}+{\displaystyle \frac{^2B}{\overline{z}^2}}+{\displaystyle \frac{^2B}{v^2}}+{\displaystyle \frac{^2B}{\overline{v}^2}}+2B))e^{56}`$
$`+`$ $`({\displaystyle \frac{\ddot{a}}{2}}e^aB({\displaystyle \frac{^2A}{z^2}}+{\displaystyle \frac{^2A}{\overline{z}^2}}))e^{34}+{\displaystyle \frac{\ddot{a}}{2}}ABe^{12}{\displaystyle \frac{\dot{a}(t)}{2}}({\displaystyle \frac{}{z}}+{\displaystyle \frac{}{\overline{z}}})ABdte^2`$
$`i{\displaystyle \frac{\dot{a}(t)}{2}}({\displaystyle \frac{}{z}}{\displaystyle \frac{}{\overline{z}}})ABdte^1`$
$`{\displaystyle \frac{1}{2}}[e^{14}({\displaystyle \frac{^2A}{zv}}+{\displaystyle \frac{^2A}{\overline{z}v}}+c.c.)ie^{13}({\displaystyle \frac{^2A}{zv}}+{\displaystyle \frac{^2A}{\overline{z}v}}c.c.)`$
$`ie^{24}({\displaystyle \frac{^2A}{zv}}{\displaystyle \frac{^2A}{\overline{z}v}}c.c.)+e^{23}({\displaystyle \frac{^2A}{zv}}{\displaystyle \frac{^2A}{\overline{z}v}}+c.c.)`$
$`i(e^{25}+e^{16})({\displaystyle \frac{B}{v}}{\displaystyle \frac{B}{\overline{v}}})(e^{26}e^{15})({\displaystyle \frac{B}{v}}+{\displaystyle \frac{B}{\overline{v}}})`$
$`i(e^{36}+e^{45})[({\displaystyle \frac{B}{z}}{\displaystyle \frac{B}{\overline{z}}})+({\displaystyle \frac{A}{z}}{\displaystyle \frac{A}{\overline{z}}})]`$
$`+i(e^{35}+e^{46})[({\displaystyle \frac{B}{z}}+{\displaystyle \frac{B}{\overline{z}}})({\displaystyle \frac{A}{z}}+{\displaystyle \frac{A}{\overline{z}}})]]]`$
We use this form because $`dG=_{11}_6`$, the six-form $`_6`$ specifying the position of the $`M5`$-branes transverse to the world volume, when one interprets the metric to represent $`M5`$-branes wrapped around two-cycles (with densities $`\rho _{ij}`$) (See ). Computing $`_{11}dG`$ we obtain informations about $`_6`$.
We now simplify in the following way: $`A,B`$ do not depend on $`v`$ and $`\overline{v}`$; $`A=B`$; $`\frac{A}{z}=\frac{A}{\overline{z}}`$. One parameter family of solution satisfying the assumptions and the requirement of periodicity in the $`T^2`$-valued $`z`$, is given by $`A(z,\overline{z})=cos[m\pi (z+\overline{z})]`$ where $`m𝐙`$.
The Hodge dual of $`dG`$ then reads:
$`_{11}dG=e^{4\mathrm{\Delta }}dx^0dx^1dx^2dx^3`$
$`\left[\left({\displaystyle \frac{\ddot{a}}{2}}e^a+2(m\pi )^22\right)Ae^{56}+\left({\displaystyle \frac{\ddot{a}}{2}}e^a+2(m\pi )^2\right)Ae^{34}+{\displaystyle \frac{\ddot{a}}{2}}A^2e^{12}2\alpha \dot{a}{\displaystyle \frac{A}{z}}e^{2t}\right]`$
$`e^{4\mathrm{\Delta }}dx^0dx^1dx^2dx^3(\rho _{12}e^{12}+\rho _{34}e^{34}+\rho _{56}e^{56}+\rho _{2t}e^{2t})`$ (31)
Interestingly, along hypersurfaces given by: $`z+\overline{z}=\frac{(2k+1)}{2m}`$, one gets a source-free Bianchi identity. For other points, we consider now different cases for $`_{11}dG`$:
* $`a=\mathrm{const}`$, which would mean that we do not have ”size” deformation. $`_{11}dG=2\alpha ^2Ae^{34}+2(m\pi )^22)Ae^{56}`$ That means the densities $`\rho _{34}`$ and $`\rho _{56}`$ are not zero; it is the only posibility that two terms vanish.
* $`\dot{a}(t)=\mathrm{const}`$: $`\rho _{12}=0`$
* We can also require that $`\rho _{34}=0`$ \- we then have non-zero density $`\rho _{12}`$ with $`a(t)`$ given by: $`a(t)=ln[\frac{e^{\gamma _1t+\gamma _2}(32(m\pi )^2+e^{\gamma _1t+\gamma _2})^2}{16\gamma _1^2}]`$, $`\gamma _{1,2}`$ being constants of integration. One can get a similar expression for $`a(t)`$ if one sets $`\rho _{56}=0`$.
The Warp-factor can now be calculated to:
$$(t)=\frac{1}{2}a(t)+\mathrm{const};\text{set}\mathrm{const}=0$$
For the metric it follows:
$$ds_{11}^{}{}_{}{}^{2}=e^{a(t)}\eta _{\mu \nu }dx^\mu dx^\nu +ds_{7}^{}{}_{}{}^{2}.$$
### 3.5 Discussion for $`q0`$
For the 7-dimensional metric we assume $`ds_{7}^{}{}_{}{}^{2}(y,t)=e^{p\varphi }ds_{6}^{}{}_{}{}^{2}(y,t)+e^{2\varphi }dt^2`$. One can also show the following relations:
$`\sigma \widehat{d}={\displaystyle \frac{1}{2}}\widehat{d}\varphi ;\dot{}={\displaystyle \frac{1}{3}}Qe^\varphi .`$ (32)
Unlike what is shown in , the following equation is identically fullfilled:
$`3d\sigma =dQv+2Qv\sigma .`$
The two- and the three-form of SU(3) structure on the six-fold are now:
$`J=e^{p\varphi }\widehat{J};\mathrm{\Omega }=e^{\frac{3}{2}p\varphi }\widehat{\mathrm{\Omega }}.`$ (34)
Thus, one gets a non-Kähler complex manifold referred to as a balanced manifold<sup>6</sup><sup>6</sup>6This time, $`W_5=\frac{1}{2}\widehat{\mathrm{\Omega }_+}\mathrm{\_}|\widehat{d}\widehat{\mathrm{\Omega }_+}=\frac{3}{8}\widehat{d}\varphi 0`$. Hence, the torsion $`\tau W_3W_5`$ \- a balanced manifold.. We obtain:
$`\widehat{d}\widehat{J}=2e^{\frac{1}{2}\varphi }S;\widehat{d}\widehat{\mathrm{\Omega }}={\displaystyle \frac{3}{4}}\widehat{d}\varphi \widehat{\mathrm{\Omega }}`$ (35)
and for the Flow equations it follows:
$`{\displaystyle \frac{\widehat{J}}{t}}={\displaystyle \frac{2}{3}}e^\varphi Q\widehat{J}2e^{\frac{\varphi }{2}}𝒜{\displaystyle \frac{1}{2}}\dot{\varphi }\widehat{J}`$
$`{\displaystyle \frac{\widehat{\mathrm{\Omega }}}{t}}=Qe^\varphi \widehat{\mathrm{\Omega }}{\displaystyle \frac{3}{4}}\dot{\varphi }\widehat{\mathrm{\Omega }}.`$ (36)
For simplicity we choose $`p=\frac{1}{2}`$. We assume $`A=\frac{1}{2}e^\varphi \dot{a}(Ae^{34}+Be^{56})`$ from where it follows $`Q=\frac{3}{2}\dot{a}e^\varphi `$. Hence, $`\mathrm{\Delta }=\frac{1}{2}a(t)\frac{1}{2}\varphi `$. Through comparison of the coefficients using $`\sigma =\frac{1}{2}\widehat{d}\varphi =\frac{2}{3}W\mathrm{\_}|J`$ we obtain :
$`W={\displaystyle \frac{3}{4}}[e^1i(\varphi _z\varphi _{\overline{z}})e^ae^2(\varphi _z+\varphi _{\overline{z}})e^a+e^3i{\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}e^4{\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}].`$ (37)
Now we can calculate the terms for the Bianchi identity. The same is done in Appendix A. The $`D=11`$-Hodge dual of $`dG`$ can then be calculated, as done towards the end of Appendix A. Now, using the simplifications we already used for q=0: $`A,B`$ do not depend on $`v`$ and $`\overline{v}`$; $`A=B`$; $`\frac{A}{z}=\frac{A}{\overline{z}}`$, thereby obtaining a one-parameter family of solution satisfying the assumptions $`A(z,\overline{z})=cos[m\pi (z+\overline{z})],m𝐙`$, and additionally assuming $`\varphi `$ does not depend on $`v`$ and $`\overline{v}`$, the eleven-dimensional Hodge dual of $`dG`$ then reads:
$`_{11}dG=e^{4\mathrm{\Delta }}dx^0dx^1dx^2dx^3`$
$`[(e^\varphi {\displaystyle \frac{\ddot{a}}{2}}e^a2\alpha ^22+{\displaystyle \frac{\varphi _{zz}+\varphi _{\overline{z}\overline{z}}}{2}}+\alpha (\varphi _z+\varphi _{\overline{z}})+{\displaystyle \frac{3}{2}}e^\varphi e^a\alpha (\varphi _z+\varphi _{\overline{z}}){\displaystyle \frac{1}{2}}\dot{a}\dot{\varphi }e^\varphi )Ae^{56}`$
$`+`$ $`(e^\varphi {\displaystyle \frac{\ddot{a}}{2}}e^a2\alpha ^2+{\displaystyle \frac{\varphi _{zz}+\varphi _{\overline{z}\overline{z}}}{2}}+\alpha (\varphi _z+\varphi _{\overline{z}})+{\displaystyle \frac{3}{2}}e^\varphi e^a\alpha (\varphi _z+\varphi _{\overline{z}}){\displaystyle \frac{1}{2}}\dot{a}\dot{\varphi }e^\varphi )Ae^{34}`$
$`+`$ $`({\displaystyle \frac{\ddot{a}}{2}}{\displaystyle \frac{1}{2}}\dot{a}\dot{\varphi }e^\varphi )A^2e^{12}+{\displaystyle \frac{A}{4}}(\varphi _z+\varphi _{\overline{z}})[3e^\varphi e^a1](e^{35}e^{46})`$
$`+`$ $`{\displaystyle \frac{A}{4}}(\varphi _z\varphi _{\overline{z}})[3e^\varphi e^a1]i(e^{36}+e^{45})2\alpha \dot{a}e^\varphi A^2dte^2]`$
$`\mathrm{vol}\mathrm{form}(𝐑^{3,1})(\rho _{12}e^{12}+\rho _{34}e^{34}+\rho _{56}e^{56}+\rho _{35}e^{35}+\rho _{36}e^{36}+\rho _{45}e^{45}+\rho _{46}e^{46}+\rho _{2t}e^{2t}).`$
Generically, one thus sees, that the deformed Iwasawa is a balanced manifold. However, if one assumes that $`\widehat{d}\varphi =0`$ and $`\varphi =\varphi (t)`$, then the balanced six-fold becomes a special-hermitian six-fold. In such a situation, $`\rho _{35}=\rho _{36}=\rho _{45}=\rho _{46}=0`$. Again, along hypersurfaces given by: $`z+\overline{z}=\frac{(2k+1)}{2m}`$, one gets a source-free Bianchi identity.
For a balanced manifold, assuming $`\dot{\varphi }=0,\varphi _z=\varphi _{\overline{z}}`$, which is satisfied by, e.g., $`\varphi =ln[cos[in\pi (z\overline{z})]],n𝐙`$ (once again ensuring periodicity w.r.t. the $`T^2`$-valued $`z`$), one gets the following metrics for seven-folds with $`SU(3)`$ structure:
* having frozen the “size” deformation of the Iwasawa, the seven-fold is a fibration of an interval over a balanced six-fold:
$$ds_7^2=(cos[2n\pi \mathrm{Im}(z)])^2dt^2+|dz|^2+cos[2m\pi \mathrm{Re}(z)]|dv|^2+cos[2m\pi \mathrm{Re}z]|duvdz|^2,$$
(39)
which gives
$$ds_{11}^2=sec[2n\pi \mathrm{Im}(z)]ds_{𝐑^{3,1}}^2+ds_7^2;$$
(40)
this implies that $`\rho _{12}=\rho _{t2}=0`$. The interval, for this case, could also be replaced by $`S^1`$ as periodicity w.r.t. this $`S^1`$ will then not be a problem.
* For the six-folds to be special-hermitian manifolds by assuming $`\varphi =\varphi (t)=a(t)`$ (i.e., $`\widehat{d}\varphi =0`$), and further $`\ddot{a}(\dot{a})^2=22(m\pi )^2`$, solved by $`a(t)=\pm \left(\sqrt{2(m\pi )^22}\right)t+\mathrm{constant}`$, implying $`\rho _{56}=0`$, one gets the seven-fold to be a warped product of a balanced manifold and an interval:
$`ds_7^2`$ $`=`$ $`e^{2\left(\sqrt{2(m\pi )^22}\right)t+\gamma _3}dt^2+e^{\left(\sqrt{2(m\pi )^22}\right)t+\frac{\gamma _3}{2}}|dz|^2+cos[2m\pi \mathrm{Re}(z)]|dv|^2`$ (41)
$`+cos[2m\pi \mathrm{Re}(z)]|duvdz|^2,`$
which gives a two($`m𝐙,\gamma _3\mathrm{R}`$))-parameter family of solutions:
$$ds_{11}^2=e^{\left(\sqrt{2(m\pi )^22}\right)t\frac{\gamma _3}{2}}ds_{𝐑^{3,1}}^2+ds_7^2.$$
(42)
## 4 Interpretation of the Uplifts
For uplifts to seven-folds with $`SU(3)`$ structure, as $`dG`$, generically, is not zero, we only have solutions with sources. The non-zero piece of $`dG`$, for $`q=0`$, in the ansatz we worked with, could be interpreted due to $`M5`$-branes wrapping 2-cycles with respect to the standard complex structure in the internal seven-fold viewed as a warped product of a (special hermitian or balanced) six-fold and the internal. For $`q0`$, in general, one ends up with balanced manifolds. However, the issue of supersymmetry of the wrapped M5-branes should be looked into more carefully, by , e.g., looking at the world-volume theory of these branes<sup>7</sup><sup>7</sup>7AM thanks O.Lunin for a discussion on this point. to ensure that the wrapped $`M5`$-branes that minimize the energy functional, belong to $`M`$-theory.
The two-cycle is calibrated w.r.t. the generalized calibration $`J`$ thereby guaranteeing the minimization of the corresponding energy functional, but away from the standard-complex-structure point in the moduli space of almost complex structures on the Iwasawa. At any generic point in the moduli space of almost complex structures on the Iwasawa manifold, the two form $`J`$ can be described by using a basis constructed of $`SO(4)`$ matrices $`P^i`$ and the one forms $`e_i`$ of the orthonormal basis :
$$J=\frac{i}{2}(\alpha \overline{\alpha }+\beta \overline{\beta }+\gamma \overline{\gamma })$$
where:
$`\alpha \mathrm{cos}\theta f^1\mathrm{sin}\theta e^6+if^2`$
$`\beta f^3+if^4`$
$`\gamma +i(\mathrm{cos}\theta e^6\mathrm{sin}\theta f^1+if^2),`$
and $`f^i=P_{}^{i}{}_{j}{}^{}e^j,PSO(4)`$.
The wedge product of two different $`f`$s is defined as:
$$f^if^j=\frac{1}{2}P_{}^{i}{}_{[k}{}^{}P_{}^{j}{}_{l]}{}^{}e^ke^l.$$
This means that we cannot have a non-trivial ”shape” (and ”size”) deformation using the standard complex structure. By turning on fluxes one moves away from the standard complex structure. The moduli space is then neither ”edge” nore standard complex structure with respect to the one forms $`e_i`$ of the orthonormal basis.
If one had used a singular ansatz for $`J`$ implying localization w.r.t. directions transverse to the $`M5`$-brane world volume, and such that $`e^{2\mathrm{\Delta }}`$ vanished at the location of the $`M5`$-branes <sup>8</sup><sup>8</sup>8AM thanks J.Maldacena for emphasizing this point to him., then using the results of , one could interpret the 11-dimensional metric after uplifting the balanced manifold to a seven-fold with $`SU(3)`$ structure, as corresponding to $`M5`$ brane wrapped around a Riemann surface in $`𝐂^3`$:
$$ds_{11}^2=H_1^2ds_{𝐑^{3,1}}^2+g_{m\overline{n}}dz^md\overline{z}^{\overline{n}}+H_2^2dt^2,$$
(43)
one gets the required interpretation by setting $`H_1=\mathrm{\Delta }`$ and $`H_2=\varphi `$ (See ) and regarding $`z,v,u`$ as the complex coordinates on $`𝐂^3`$.
Now, as an interesting mathematical curiosity, using the results of , we will show that one can associate a Riemann surface to a choice of the “shape” deformation functions. Consider the system of partial differential equations:
$$\frac{A}{z}=\varphi B;\frac{B}{\overline{z}}=\varphi A,$$
(44)
which is equivalent to:
$$\frac{^2B}{z\overline{z}}=\frac{1}{\varphi }\frac{\varphi }{z}\frac{B}{\overline{z}}\varphi ^2B$$
(45)
(or one could switch $`B`$ with $`A`$ and $`z\overline{z}`$). We see that the following choice solves (44):
$`A_1(z,\overline{z})=sin[m\pi (z+\overline{z})],B_1(z,\overline{z})=cos[m\pi (z+\overline{z})],\varphi (z,\overline{z})=m\pi ;`$
$`A_1(z,\overline{z})=cos[m\pi (z+\overline{z})],B_1(z,\overline{z})=sin[m\pi (z+\overline{z})],\varphi (z,\overline{z})=m\pi .`$ (46)
One can construct an $`𝐑^3`$ with coordinates::
$`X^1+iX^2=i{\displaystyle _\mathrm{\Gamma }}(A_1^2dz^{}B_1^2d\overline{z}^{}),`$
$`X^1iX^2=i{\displaystyle _\mathrm{\Gamma }}(B_1^2dz^{}A_1^2d\overline{z}^{}),`$
$`X^3={\displaystyle _\mathrm{\Gamma }}(A_1B_1dz^{}+A_1B_1d\overline{z}^{}),`$ (47)
and for the contour $`\mathrm{\Gamma }=(x,y)|x=y,0<x<1`$,
$`X^1+iX^2=x{\displaystyle \frac{i}{4m\pi }}sin(4m\pi x),`$
$`X^3={\displaystyle \frac{cos(4m\pi x)1}{2m\pi }},`$ (48)
which would be the Weierstra$`ß`$ representation for the conformal immersion in $`𝐑^3`$ with the condition of conformal immersion (implying $`g_{zz}=g_{\overline{z}\overline{z}}=0`$):
$$\left(\frac{X^1}{z}\right)^2+\left(\frac{X^2}{z}\right)^2+\left(\frac{X^3}{z}\right)^2=0,$$
(49)
a quadric in $`\mathrm{𝐂𝐏}^2`$. The induced metric on the Riemann surface is given by:
$$ds^2=(A_1^2+B_1^2)^2|dz|^2=|dz|^2,$$
(50)
which is just a $`T^2`$. In appendix B, we discuss a Weierstra$`ß`$ representation of conformal immersion in $`𝐂^3`$ for the linear dilaton ansatz but for solutions that are valid only locally because of lack of periodicity w.r.t. the $`T^2`$-valued coordinates.
## 5 Conclusion
We uplifted the Iwasawa to a $`G_2`$-manifold via Hitchin’s Flow equations using ”size”-deformations, and found that an uplift via Hitchin’s flow equations (to a seven-fold with $`G_2`$ holonomy) with non-trivial ”shape” deformation is not possible. Furthermore we uplifted the Iwasawa to an irreducible 7 manifold with $`SU(3)`$-structure through deforming the metric using ”shape” and ”size”-deformations via generalization of the Hitchin’s Flow equations. Without simplification the seven-fold of $`SU(3)`$ structure turns out to be a warped product of either a special hermtian or a balanced six-fold and an interval (which for cases when one freezes the “size deformations”, could also be replaced by an $`S^1`$). The uplifted metric could be interpreted as $`M5`$-branes wrapping two-cycles calibrated by the generalized calibration (like ), the two-form corresponding to the almost complex structure, but away from the standard complex structure point (in the moduli space of almost complex structures on the Iwasawa). The supersymmetry of these wrapped membranes needs to be further looked into. We also showed how to associate a Riemann surface (via Weierstra$`ß`$ representation for conformal immersion in $`𝐑^l`$) with the “shape deformation” functions of seven-folds with $`SU(3)`$ structure and the dilaton. The relationship of the same to, e.g., supersymmetric two-cycles or Riemann surfaces embedded in seven-folds of $`SU(3)`$ structure around which $`M5`$-branes wrap, needs to be understood.
However, if one allows for sources, as was the case, then there seems to be the possibility of considering singular uplifts of the following type. Given that it is possible to impose just the requirement of (anti-)analyticity on the (anti-)holomorphic parts of the functions $`A(z,\overline{z};v,\overline{v}),B(z,\overline{z};v,\overline{v})`$ and $`\varphi (z,\overline{z};v,\overline{v})`$, remembering that the complex coordinates $`z,v`$ are $`T^2`$-valued, one can hence introduce the doubly periodic functions: the Weierstra$`ß`$ elliptic function with double poles and the Jacobi elliptic function with simple poles, into the metric of the seven-fold. The elliptic surface corresponding to the Weierstra$`ß`$ elliptic function could perhaps be related to the Riemann surface relevant to the world volume of the $`M5`$-brane relevant to the uplifts. See for some connections relevant to this study.
## Acknowledgements
AF would like to thank the physics department of IIT Roorkee for the hospitality during her stay there as part of her semester-long study leave from RWTH, Aachen, and the organizers of THEP-I (2005) at IIT Roorkee where some preliminary results of this work were presented. One of us (AM) would like to thank S.Chiossi, G.Papadopoulos and J.Gauntlett for useful communications, and especially J.Maldacena for critical comments on the material of the first version of this paper and O.Lunin for discussions, the Harvard high energy theory group (and S.Minwalla in particular) and IAS, Princeton for the hospitality during his stay there where part of this work was completed, the Department of Atomic Energy (Board of Research in Nuclear Sciences), Govt. of India for a research grant (under the DAE young scientist award scheme) and the organizers of PASCOS05, Gyeongju, Korea, where some preliminary results of this work were presented. We thank H.S.Solanki for participation in the initial stages of this project.
## Appendix A The Bianchi Identity for $`q0`$
In this appendix, we discuss the evaluation of the various terms in the Bianchi identity.
$`{\displaystyle \frac{dQ}{6}}JJ={\displaystyle \frac{e^\varphi }{2}}[\ddot{a}(e^aAe^{1234t}+e^aBe^{1256t}+ABe^{3456t})`$
$`\dot{a}[(\varphi _z+\varphi _{\overline{z}})ABe^{13456}+i(\varphi _z\varphi _{\overline{z}})ABe^{23456}+(\varphi _v+\varphi _{\overline{v}})e^aBe^{12356}`$
$`+i(\varphi _v\varphi _{\overline{v}})e^aBe^{12456}]+\dot{\varphi }[e^aAe^{1234t}+e^aBe1256t+ABe^{3456t}]`$ (A1)
$$\frac{2}{9}Q^2JJv=\dot{a}^2e^\varphi dt[e^aAe^{1234}+e^aBe^{1256}+ABe^{3456}]$$
(A2)
$$\frac{4}{3}QvJA=\dot{a}^2e^\varphi dt[e^aAe^{1234}+e^aBe^{1256}+2ABe^{3456}]$$
(A3)
$`{\displaystyle \frac{Q}{3}}JJ\sigma ={\displaystyle \frac{\dot{a}}{2}}e^\varphi [(\varphi _z+\varphi _{\overline{z}})ABe^{13456}+i(\varphi _z\varphi _{\overline{z}})ABe^{23456}`$
$`+(\varphi _v+\varphi _{\overline{v}})e^aBe^{12356}+i(\varphi _v\varphi _{\overline{v}})e^aBe^{12456}]]`$ (A4)
$`2S𝒜={\displaystyle \frac{\dot{a}}{2}}e^\varphi [e^{13456}[({\displaystyle \frac{}{z}}+{\displaystyle \frac{}{\overline{z}}})AB(\varphi _z+\varphi _{\overline{z}})AB]`$
$`+ie^{23456}[({\displaystyle \frac{}{z}}{\displaystyle \frac{}{\overline{z}}})AB(\varphi _z\varphi _{\overline{z}})AB]`$
$`e^{12356}{\displaystyle \frac{e^a}{2}}B(\varphi _v+\varphi _{\overline{v}})e^{12456}i{\displaystyle \frac{e^a}{2}}B(\varphi _v\varphi _{\overline{v}})]`$ (A5)
$$2v𝒜𝒜=e^\varphi \dot{a}^2ABdte^{3456}$$
(A6)
$`J\sigma 𝒜={\displaystyle \frac{\dot{a}}{2}}e^\varphi [(\varphi _z+\varphi _{\overline{z}})ABe^{13456}+i(\varphi _z\varphi _{\overline{z}})ABe^{23456}+(\varphi _v+\varphi _{\overline{v}}){\displaystyle \frac{e^a}{2}}Be^{12356}`$
$`+i(\varphi _v\varphi _{\overline{v}}){\displaystyle \frac{e^a}{2}}Be^{12456}]`$ (A7)
$`3v\sigma JW={\displaystyle \frac{9}{8}}e^\varphi [e^{1234t}[Ae^a(4\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{A}}(4\varphi _v\varphi _{\overline{v}})]e^{1256t}Be^a4\varphi _z\varphi _{\overline{z}}`$
$`+e^{1356t}iB[(\varphi _z\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})e^a+{\displaystyle \frac{1}{A}}(\varphi _z+\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})]`$
$`+e^{1456t}B[(\varphi _z\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})e^a{\displaystyle \frac{1}{A}}(\varphi _z+\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})]`$
$`+e^{2356t}B[(\varphi _z+\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})e^a{\displaystyle \frac{1}{A}}(\varphi _z\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})]`$
$`+e^{2456t}iB[(\varphi _z+\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})e^a{\displaystyle \frac{1}{A}}(\varphi _z\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})]`$
$`e^{3456t}{\displaystyle \frac{B}{A}}4\varphi _v\varphi _{\overline{v}}]`$ (A8)
$`2vSW={\displaystyle \frac{3}{4}}e^\varphi dt[e^{1234}[2e^a[A_z\varphi _{\overline{z}}A_{\overline{z}}\varphi _z+A\varphi _z\varphi _{\overline{z}}]+{\displaystyle \frac{e^a}{A}}2\varphi _v\varphi _{\overline{v}}]`$
$`+e^{1256}2e^a[B_z\varphi _{\overline{z}}B_{\overline{z}}\varphi _z+B\varphi _z\varphi _{\overline{z}}]+e^{3456}{\displaystyle \frac{2}{A}}[B_v\varphi _{\overline{v}}B_{\overline{v}}\varphi _v+B\varphi _v\varphi _{\overline{v}}]`$
$`+ie^{1356}[{\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}[(B_z+B_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z+\varphi _{\overline{z}})]e^a(\varphi _z\varphi _{\overline{z}})[(B_v+B_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v+\varphi _{\overline{v}})]]`$
$`e^{1456}[{\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}[(B_z+B_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z+\varphi _{\overline{z}})]e^a(\varphi _z\varphi _{\overline{z}})[(B_vB_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v\varphi _{\overline{v}})]]`$
$`e^{2356}[{\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}[(B_zB_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z\varphi _{\overline{z}})]e^a(\varphi _z+\varphi _{\overline{z}})[(B_v+B_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v+\varphi _{\overline{v}})]]`$
$`ie^{2456}[{\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}[(B_zB_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z\varphi _{\overline{z}})]e^a(\varphi _z+\varphi _{\overline{z}})[(B_vB_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v\varphi _{\overline{v}})]]`$
$`+B[ie^a(e^{1236}+e^{1245})(\varphi _z\varphi _{\overline{z}})e^a(e^{1246}e^{1235})(\varphi _z+\varphi _{\overline{z}})`$
$`+i{\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}(e^{1346}+e^{2345})+{\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}(e^{2346}e^{1345})]]`$ (A9)
$`dW`$ can be calculated as follows:
$`dW={\displaystyle \frac{3}{4}}[\dot{a}[i(\varphi _z\varphi _{\overline{z}})e^ae^{1t}(\varphi _z+\varphi _{\overline{z}})e^ae^{2t}]`$
$`+[e^{1t}i(\dot{\varphi _z}\dot{\varphi _{\overline{z}}})e^a+e^{2t}(\dot{\varphi _z}+\dot{\varphi _{\overline{z}}})e^a+e^{3t}i{\displaystyle \frac{(\dot{\varphi _v}\dot{\varphi _{\overline{v}}})}{A}}e^{4t}{\displaystyle \frac{(\dot{\varphi _v}+\dot{\varphi _{\overline{v}}})}{A}}]4e^{12}\varphi _{\overline{z}z}e^a]`$
$`ie^{13}[(\varphi _{zv}\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}\varphi _{\overline{v}z}+\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_z+A_{\overline{z}})}{A^2}}(\varphi _v\varphi _{\overline{v}})]`$
$`+e^{14}[(\varphi _{zv}\varphi _{\overline{z}v}\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}+\varphi _{\overline{v}z}+\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_z+A_{\overline{z}})}{A^2}}(\varphi _v+\varphi _{\overline{v}})]`$
$`+e^{23}[(\varphi _{zv}+\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_zA_{\overline{z}})}{A^2}}(\varphi _v\varphi _{\overline{v}})]`$
$`+ie^{24}[(\varphi _{zv}+\varphi _{\overline{z}v}\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}+\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_zA_{\overline{z}})}{A^2}}(\varphi _v+\varphi _{\overline{v}})]`$
$`+e^{34}[4\varphi _{\overline{v}v}+{\displaystyle \frac{2A_v\varphi _{\overline{v}}+2A_{\overline{z}}\varphi _v}{A^2}}]`$ (A10)
$`vJdW`$
$`={\displaystyle \frac{3}{4}}e^\varphi dt[B[e^{1256}4\varphi _{\overline{z}z}e^a]`$
$`ie^{1356}[(\varphi _{zv}\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}\varphi _{\overline{v}z}+\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_z+A_{\overline{z}})}{A^2}}(\varphi _v\varphi _{\overline{v}})]`$
$`+e^{1456}[(\varphi _{zv}\varphi _{\overline{z}v}\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}+\varphi _{\overline{v}z}+\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_z+A_{\overline{z}})}{A^2}}(\varphi _v+\varphi _{\overline{v}})]`$
$`+e^{2356}[(\varphi _{zv}+\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_zA_{\overline{z}})}{A^2}}(\varphi _v\varphi _{\overline{v}})]`$
$`+ie^{2456}[(\varphi _{zv}+\varphi _{\overline{z}v}\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}+\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}+{\displaystyle \frac{(A_zA_{\overline{z}})}{A^2}}(\varphi _v+\varphi _{\overline{v}})]`$
$`+e^{3456}[4\varphi _{\overline{v}v}+{\displaystyle \frac{2A_v\varphi _{\overline{v}}+2A_{\overline{z}}\varphi _v}{A^2}}]`$
$`e^{1234}[A4\varphi _{\overline{z}z}e^a+e^a[4\varphi _{\overline{v}v}+{\displaystyle \frac{2A_v\varphi _{\overline{v}}+2A_{\overline{z}}\varphi _v}{A^2}}]]`$ (A11)
$`d_7S={\displaystyle \frac{1}{2}}dt[e^{1256}[2(A_{zz}+A_{\overline{z}\overline{z}})A(\varphi _{zz}+\varphi _{\overline{z}\overline{z}})A_z\varphi _zA_{\overline{z}}\varphi _{\overline{z}}]`$
$`e^{2356}[A_{zv}+A_{\overline{z}\overline{v}}+A_{z\overline{v}}+A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}+\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})`$
$`{\displaystyle \frac{(A_v+A_{\overline{v}})}{2}}(\varphi _z+\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}\varphi _{\overline{v}z}\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}})]`$
$`ie^{2456}[A_{zv}+A_{\overline{z}\overline{v}}A_{z\overline{v}}A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}+\varphi _{\overline{z}v}\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})`$
$`{\displaystyle \frac{(A_vA_{\overline{v}})}{2}}(\varphi _z+\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}+\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}})]`$
$`ie^{1356}[A_{zv}A_{\overline{z}\overline{v}}+A_{z\overline{v}}A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}\varphi _{\overline{z}v}+\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})`$
$`{\displaystyle \frac{(A_v+A_{\overline{v}})}{2}}(\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}\varphi _{\overline{v}z}+\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}})]`$
$`+e^{1456}[A_{zv}A_{\overline{z}\overline{v}}A_{z\overline{v}}+A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}\varphi _{\overline{z}v}\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})`$
$`{\displaystyle \frac{(A_vA_{\overline{v}})}{2}}(\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}+\varphi _{\overline{v}z}+\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}})]`$
$`+e^{1234}[2(B_{zz}+B_{\overline{z}\overline{z}}+B_{vv}+B_{\overline{v}\overline{v}})B(\varphi _{zz}+\varphi _{\overline{z}\overline{z}}+\varphi _{vv}+\varphi _{\overline{v}\overline{v}})B_z\varphi _zB_{\overline{z}}\varphi _{\overline{z}}`$
$`B_v\varphi _vB_{\overline{v}}\varphi _{\overline{v}}+4B]e^{3456}e^a(\varphi _{vv}+\varphi _{\overline{v}\overline{v}})`$
$`+(e^{1235}e^{1246})[B_z+B_{\overline{z}}A_zA_{\overline{z}}{\displaystyle \frac{A}{2}}(\varphi _z+\varphi _{\overline{z}})]`$
$`+i(e^{2345}+e^{1346})[B_vB_{\overline{v}}{\displaystyle \frac{e^a}{2}}(\varphi _v\varphi _{\overline{v}})]`$
$`i(e^{1245}+e^{1236})[B_zB_{\overline{z}}+A_zA_{\overline{z}}{\displaystyle \frac{A}{2}}(\varphi _z\varphi _{\overline{z}})]`$
$`+(e^{1345}+e^{2346})[B_v+B_{\overline{v}}+{\displaystyle \frac{e^a}{2}}(\varphi _v+\varphi _{\overline{v}})]]+`$ (A12)
From the above, one can calculate $`_{11}dG`$ as follows:
$`_{11}dG`$ $`=`$ $`e^{4\mathrm{\Delta }}dx^0dx^1dx^2dx^3`$
$`[(\ddot{a}{\displaystyle \frac{e^\varphi }{2}}e^aA+{\displaystyle \frac{9}{8}}e^\varphi [Ae^a(4\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{A}}(4\varphi _v\varphi _{\overline{v}})]`$
$``$ $`{\displaystyle \frac{3}{2}}e^\varphi [e^a[A_z\varphi _{\overline{z}}A_{\overline{z}}\varphi _z+A\varphi _z\varphi _{\overline{z}}]+{\displaystyle \frac{e^a}{A}}\varphi _v\varphi _{\overline{v}}]`$
$`+`$ $`{\displaystyle \frac{3}{2}}e^\varphi [A2\varphi _{\overline{z}z}e^a+e^a[2\varphi _{\overline{v}v}+{\displaystyle \frac{A_v\varphi _{\overline{v}}+A_{\overline{z}}\varphi _v}{A^2}}]]`$
$``$ $`{\displaystyle \frac{1}{2}}[2(B_{zz}+B_{\overline{z}\overline{z}}+B_{vv}+B_{\overline{v}\overline{v}})B(\varphi _{zz}+\varphi _{\overline{z}\overline{z}}+\varphi _{vv}+\varphi _{\overline{v}\overline{v}})B_z\varphi _zB_{\overline{z}}\varphi _{\overline{z}}`$
$`B_v\varphi _vB_{\overline{v}}\varphi _{\overline{v}}+4B]{\displaystyle \frac{1}{2}}\dot{a}\dot{\varphi }e^\varphi e^aA)e^{56}`$
$`+`$ $`(\ddot{a}{\displaystyle \frac{e^\varphi }{2}}e^aB+{\displaystyle \frac{9}{8}}e^\varphi Be^a4\varphi _z\varphi _{\overline{z}}{\displaystyle \frac{3}{2}}e^\varphi e^a[B_z\varphi _{\overline{z}}B_{\overline{z}}\varphi _z+B\varphi _z\varphi _{\overline{z}}]`$
$`+`$ $`3e^\varphi B[\varphi _{\overline{z}z}e^a]{\displaystyle \frac{1}{2}}[2(A_{zz}+A_{\overline{z}\overline{z}})A(\varphi _{zz}+\varphi _{\overline{z}\overline{z}})A_z\varphi _zA_{\overline{z}}\varphi _{\overline{z}}]{\displaystyle \frac{1}{2}}\dot{a}\dot{\varphi }e^\varphi e^aB)e^{34}`$
$`+`$ $`(\ddot{a}{\displaystyle \frac{e^\varphi }{2}}AB+{\displaystyle \frac{9}{8}}e^\varphi {\displaystyle \frac{B}{A}}4\varphi _v\varphi _{\overline{v}}]{\displaystyle \frac{3}{2}}e^\varphi {\displaystyle \frac{1}{A}}[B_v\varphi _{\overline{v}}B_{\overline{v}}\varphi _v+B\varphi _v\varphi _{\overline{v}}]`$
$`+`$ $`{\displaystyle \frac{3}{2}}e^\varphi [2\varphi _{\overline{v}v}+{\displaystyle \frac{A_v\varphi _{\overline{v}}+A_{\overline{z}}\varphi _v}{A^2}}]+{\displaystyle \frac{1}{2}}e^a(\varphi _{vv}+\varphi _{\overline{v}\overline{v}}){\displaystyle \frac{1}{2}}\dot{a}\dot{\varphi }e^\varphi e^aAB)e^{12}`$
$``$ $`{\displaystyle \frac{\dot{a}}{2}}e^\varphi [({\displaystyle \frac{}{z}}+{\displaystyle \frac{}{\overline{z}}})AB]e^{2t}{\displaystyle \frac{\dot{a}}{2}}e^\varphi [({\displaystyle \frac{}{z}}{\displaystyle \frac{}{\overline{z}}})AB]ie^{1t}`$
$`+`$ $`{\displaystyle \frac{\dot{a}}{4}}e^\varphi e^a(\varphi _v+\varphi _{\overline{v}})Be^{4t}+{\displaystyle \frac{\dot{a}}{4}}e^\varphi e^a(\varphi _v\varphi _{\overline{v}})Be^{3t}`$
$`+`$ $`({\displaystyle \frac{9}{8}}e^\varphi B[(\varphi _z\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})e^a+{\displaystyle \frac{1}{A}}(\varphi _z+\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})]`$
$``$ $`{\displaystyle \frac{3}{4}}e^\varphi [{\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}[(B_z+B_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z+\varphi _{\overline{z}})]`$
$`e^a(\varphi _z\varphi _{\overline{z}})[(B_v+B_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v+\varphi _{\overline{v}})]]`$
$``$ $`{\displaystyle \frac{3}{4}}e^\varphi [(\varphi _{zv}\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _{vz}\varphi _{\overline{v}z}+\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}`$
$`+`$ $`{\displaystyle \frac{(A_z+A_{\overline{z}})}{A^2}}(\varphi _v\varphi _{\overline{v}})]+{\displaystyle \frac{1}{2}}[A_{zv}A_{\overline{z}\overline{v}}+A_{z\overline{v}}A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}\varphi _{\overline{z}v}+\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})`$
$``$ $`{\displaystyle \frac{(A_v+A_{\overline{v}})}{2}}(\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}\varphi _{\overline{v}z}+\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}})])ie^{24}`$
$`+`$ $`({\displaystyle \frac{9}{8}}e^\varphi B[(\varphi _z\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})e^a{\displaystyle \frac{1}{A}}(\varphi _z+\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})]`$
$`+`$ $`{\displaystyle \frac{3}{4}}e^\varphi [{\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}[(B_z+B_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z+\varphi _{\overline{z}})]`$
$``$ $`e^a(\varphi _z\varphi _{\overline{z}})[(B_vB_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v\varphi _{\overline{v}})]]`$
$`+`$ $`{\displaystyle \frac{3}{4}}e^\varphi [(\varphi _{zv}\varphi _{\overline{z}v}\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _vz+\varphi _{\overline{v}z}+\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}`$
$`+`$ $`{\displaystyle \frac{(A_z+A_{\overline{z}})}{A^2}}(\varphi _v+\varphi _{\overline{v}})]{\displaystyle \frac{1}{2}}[A_{zv}A_{\overline{z}\overline{v}}A_{z\overline{v}}+A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}\varphi _{\overline{z}v}\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})`$
$``$ $`{\displaystyle \frac{(A_vA_{\overline{v}})}{2}}(\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}+\varphi _{\overline{v}z}+\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}})])e^{23}`$
$`+`$ $`({\displaystyle \frac{9}{8}}e^\varphi B[(\varphi _z+\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})e^a{\displaystyle \frac{1}{A}}(\varphi _z\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})]`$
$`+`$ $`{\displaystyle \frac{3}{4}}e^\varphi [{\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}[(B_zB_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z\varphi _{\overline{z}})]`$
$``$ $`e^a(\varphi _z+\varphi _{\overline{z}})[(B_v+B_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v+\varphi _{\overline{v}})]]`$
$`+`$ $`{\displaystyle \frac{3}{4}}e^\varphi [(\varphi _{zv}+\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})e^a(\varphi _vz\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}`$
$`+`$ $`{\displaystyle \frac{(A_zA_{\overline{z}})}{A^2}}(\varphi _v\varphi _{\overline{v}})]+{\displaystyle \frac{1}{2}}[A_{zv}+A_{\overline{z}\overline{v}}+A_{z\overline{v}}+A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}+\varphi _{\overline{z}v}+\varphi _{z\overline{v}}+\varphi _{\overline{z}\overline{v}})`$
$``$ $`{\displaystyle \frac{(A_v+A_{\overline{v}})}{2}}(\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}\varphi _{\overline{v}z}\varphi _{v\overline{z}}+\varphi _{\overline{v}\overline{z}})])e^{14}`$
$`+`$ $`({\displaystyle \frac{9}{8}}e^\varphi B[(\varphi _z+\varphi _{\overline{z}})(\varphi _v\varphi _{\overline{v}})e^a{\displaystyle \frac{1}{A}}(\varphi _z\varphi _{\overline{z}})(\varphi _v+\varphi _{\overline{v}})]`$
$`+`$ $`{\displaystyle \frac{3}{4}}e^\varphi [{\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}[(B_zB_{\overline{z}}){\displaystyle \frac{B}{2}}(\varphi _z\varphi _{\overline{z}})]`$
$``$ $`e^a(\varphi _z+\varphi _{\overline{z}})[(B_vB_{\overline{v}}){\displaystyle \frac{B}{2}}(\varphi _v\varphi _{\overline{v}})]]`$
$`+`$ $`{\displaystyle \frac{3}{4}}e^\varphi [(\varphi _{zv}+\varphi _{\overline{z}v}\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})e^a(\varphi _vz+\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}}){\displaystyle \frac{1}{A}}`$
$`+`$ $`{\displaystyle \frac{(A_zA_{\overline{z}})}{A^2}}(\varphi _v+\varphi _{\overline{v}})]+{\displaystyle \frac{1}{2}}[A_{zv}+A_{\overline{z}\overline{v}}A_{z\overline{v}}A_{\overline{z}\overline{v}}{\displaystyle \frac{A}{2}}(\varphi _{zv}+\varphi _{\overline{z}v}\varphi _{z\overline{v}}\varphi _{\overline{z}\overline{v}})`$
$``$ $`{\displaystyle \frac{(A_vA_{\overline{v}})}{2}}(\varphi _z+\varphi _{\overline{z}})+{\displaystyle \frac{e^a}{2}}(\varphi _{vz}+\varphi _{\overline{v}z}\varphi _{v\overline{z}}\varphi _{\overline{v}\overline{z}})])ie^{13}`$
$`+`$ $`({\displaystyle \frac{3}{4}}e^\varphi e^aB(\varphi _z\varphi _{\overline{z}})+{\displaystyle \frac{1}{2}}[B_zB_{\overline{z}}+A_zA_{\overline{z}}{\displaystyle \frac{A}{2}}(\varphi _z\varphi _{\overline{z}})])i(e^{45}+e^{36})`$
$`+`$ $`({\displaystyle \frac{3}{4}}e^\varphi e^aB(\varphi _z+\varphi _{\overline{z}})+{\displaystyle \frac{1}{2}}[B_z+B_{\overline{z}}A_zA_{\overline{z}}+{\displaystyle \frac{A}{2}}(\varphi _z+\varphi _{\overline{z}})])(e^{35}e^{46})`$
$`+`$ $`({\displaystyle \frac{3}{4}}e^\varphi {\displaystyle \frac{(\varphi _v\varphi _{\overline{v}})}{A}}B{\displaystyle \frac{1}{2}}[B_vB_{\overline{v}}{\displaystyle \frac{e^a}{2}}(\varphi _v\varphi _{\overline{v}})])i(e^{25}+e^{16})`$
$`+`$ $`({\displaystyle \frac{3}{4}}e^\varphi {\displaystyle \frac{(\varphi _v+\varphi _{\overline{v}})}{A}}B{\displaystyle \frac{1}{2}}[B_v+B_{\overline{v}}+{\displaystyle \frac{e^a}{2}}(\varphi _v+\varphi _{\overline{v}})])(e^{15}e^{26})]`$
This is relevant to the interpretation of the uplift as $`M5`$-branes wrapped around supersymmetric two-cycles calibrated by a generalized calibration $`J`$, but generically, away from the standard complex structure point in the moduli space of almost complex structures. In the absence of $`M5`$-branes, $`_{11}dG=0`$ would lead to a system of coupled non-linear partial differential equations which will, if the solutions exists, be very difficult to solve.
## Appendix B The conformal immersion in $`𝐂^3`$
In this appendix, we give an example of a Riemann surface obtained as a Weierstra$`ß`$ representation of conformal immersion of a surface in $`𝐂^3`$, using $`\varphi =x`$ Re(z) (the “linear dilaton” ansatz) and $`A,B𝐑`$. As the solutions we get will not involve periodic $`A,B,\varphi `$, they are relevant only locally to the (non-compact) seven-folds of $`SU(3)`$-structure.
One needs to solve the following differential equation (use $`\frac{A,B}{y}=0`$)
$$(\frac{1}{4}\frac{d^2}{dx^2}\frac{1}{2x}\frac{d}{dx}+x^2)B=0.$$
(B1)
The two solutions, $`B_{1,2}`$ are given by:
$$B_1(x)=x^{\frac{3}{2}}J_{\frac{3}{4}}(x^2),B_2(x)=x^{\frac{3}{2}}J_{\frac{3}{4}}(x^2).$$
(B2)
Hence, the corresponding two values for $`A`$ would be:
$`A_1(x)={\displaystyle \frac{1}{2\sqrt{x}}}(2x^2J_{\frac{1}{4}}(x^2)3J_{\frac{3}{4}}(x^2)+2x^2J_{\frac{7}{4}}(x^2)),`$
$`A_2(x)={\displaystyle \frac{1}{2\sqrt{x}}}(2x^2J_{\frac{7}{4}}(x^2)3J_{\frac{3}{4}}(x^2)+2x^2J_{\frac{1}{4}}(x^2)).`$ (B3)
Then the Weierstra$`ß`$ representation for conformal immersion of a surface ($``$ Riemann surface) in $`𝐂^3`$ (See and references therein) is given by defining the following coordinates $`X^{i=1,\mathrm{},6}(z,\overline{z})`$, of the Rieman surface in $`𝐂^3`$:
$`X^1+iX^2=i{\displaystyle _\mathrm{\Gamma }}(A_1^2dz^{}B_1^2d\overline{z}^{}),`$
$`X^1iX^2=i{\displaystyle _\mathrm{\Gamma }}(B_1^2dz^{}A_1^2d\overline{z}^{}),`$
$`X^3={\displaystyle _\mathrm{\Gamma }}(A_1B_1dz^{}+A_1B_1d\overline{z}^{}),`$
$`X^4+iX^5=i{\displaystyle _\mathrm{\Gamma }}(A_2^2dz^{}B_2^2d\overline{z}^{}),`$
$`X^4iX^5=i{\displaystyle _\mathrm{\Gamma }}(B_2^2dz^{}A_2^2d\overline{z}^{}),`$
$`X^6={\displaystyle _\mathrm{\Gamma }}(A_2B_2dz^{}+A_2B_2d\overline{z}^{}),`$ (B4)
where $`\mathrm{\Gamma }`$ is a contour in $`𝐂(z,\overline{z})`$. The $`𝐂^3`$ coordinates for values of $`B_{1,2}`$ and $`A_{1,2}`$ as given in (B2) and (B), for the path $`\mathrm{\Gamma }=\{(x,y)|x=y\}`$, is given in appendix B. The induced metric on the Riemann surface is then given by:
$$ds^2=(A_1^2+A_2^2+B_1^2+B_2^2)^2|dz|^2.$$
(B5)
The condition for a conformal immersion (implying $`g_{zz}=g_{\overline{z}\overline{z}}=0`$):
$$(\frac{X^1}{z})^2+(\frac{X^2}{z})^2+(\frac{X^3}{z})^2+(\frac{X^4}{z})^2+(\frac{X^5}{z})^2+(\frac{X^6}{z})^2=0,$$
(B6)
is a quadric in $`\mathrm{𝐂𝐏}^5`$:
$$w_1^2+w_2^2+w_3^2+w_4^2+w_5^2+w_6^2=0,$$
(B7)
$`w_i`$ being homogeneous coordinates on $`\mathrm{𝐂𝐏}^5`$.
Along the path $`\mathrm{\Gamma }:\{(x^{},y^{})|y^{}=x^{},0x^{}x0,x<1\}`$, one sees that:
$$X^1+iX^2=iI_2I_1,X^3=I_3;X^4+iX^5=iI_5I_4,X^6=I_6,$$
(B8)
where
$`I_1{\displaystyle _\mathrm{\Gamma }}(A_1^2+B_1^2),I_2{\displaystyle _\mathrm{\Gamma }}(A_1^2B_1^2),I_3{\displaystyle _\mathrm{\Gamma }}A_1B_1;`$
$`I4_{}{\displaystyle _\mathrm{\Gamma }}(A_1^2+B_1^2),I_5{\displaystyle _\mathrm{\Gamma }}(A_1^2B_1^2),I_6{\displaystyle _\mathrm{\Gamma }}A_1B_1.`$ (B9)
The integrals $`I_i,i=1,\mathrm{},6`$ are given as:
$`I_1={\displaystyle \frac{x^4\left((J_{\frac{1}{4}}(x^2))^2J_{\frac{5}{4}}(x^2)J_{\frac{3}{4}}(x^2)\right)}{4}}+{\displaystyle \frac{x^4\left((J_{\frac{3}{4}}(x^2))^2J_{\frac{1}{4}}(x^2)J_{\frac{7}{4}}(x^2)\right)}{4}}`$
$`+{\displaystyle \frac{x^4\left((J_{\frac{7}{4}}(x^2))^2J_{\frac{3}{4}}(x^2)J_{\frac{11}{4}}(x^2)\right)}{4}}+{\displaystyle \frac{\left(x^2\right)_2^{\frac{3}{2}}F_3(\{\frac{3}{4},\frac{5}{4}\},\{\frac{3}{2},\frac{7}{4},\frac{7}{4}\},x^4)}{\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{7}{4})}}+`$
$`{\displaystyle \frac{3\left(x^2\right)_2^{\frac{3}{2}}F_3(\{\frac{3}{4},\frac{5}{4}\},\{\frac{7}{4},\frac{7}{4},\frac{5}{2}\},x^4)}{8\sqrt{2}\mathrm{\Gamma }(\frac{7}{4})}}^2{\displaystyle \frac{3\left(x^2\right)_2^{\frac{7}{2}}F_3(\{\frac{7}{4},\frac{9}{4}\},\{\frac{11}{4},\frac{11}{4},\frac{7}{2}\},x^4)}{28\sqrt{2}\mathrm{\Gamma }(\frac{7}{4})\mathrm{\Gamma }(\frac{11}{4})}}`$
$`{\displaystyle \frac{\left(x^2\right)^{\frac{7}{2}}\left(363_2F_3(\{\frac{5}{4},\frac{7}{4}\},\{\frac{5}{2},\frac{11}{4},\frac{11}{4}\},x^4)56x_2^4F_3(\{\frac{9}{4},\frac{11}{4}\},\{\frac{7}{2},\frac{15}{4},\frac{15}{4}\},x^4)\right)}{2541\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{11}{4})}}`$ (B10)
$`I_2={\displaystyle \frac{x^4\left((J_{\frac{1}{4}}(x^2))^2J_{\frac{5}{4}}(x^2)J_{\frac{3}{4}}(x^2)\right)}{4}}{\displaystyle \frac{x^4\left((J_{\frac{3}{4}}(x^2))^2J_{\frac{1}{4}}(x^2)J_{\frac{7}{4}}(x^2)\right)}{4}}`$
$`+{\displaystyle \frac{x^4\left((J_{\frac{7}{4}}(x^2))^2J_{\frac{3}{4}}(x^2)J_{\frac{11}{4}}(x^2)\right)}{4}}+{\displaystyle \frac{\left(x^2\right)_2^{\frac{3}{2}}F_3(\{\frac{3}{4},\frac{5}{4}\},\{\frac{3}{2},\frac{7}{4},\frac{7}{4}\},x^4)}{\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{7}{4})}}+`$
$`{\displaystyle \frac{3\left(x^2\right)_2^{\frac{3}{2}}F_3(\{\frac{3}{4},\frac{5}{4}\},\{\frac{7}{4},\frac{7}{4},\frac{5}{2}\},x^4)}{8\sqrt{2}(\mathrm{\Gamma }(\frac{7}{4}))^2}}{\displaystyle \frac{3\left(x^2\right)_2^{\frac{7}{2}}F_3(\{\frac{7}{4},\frac{9}{4}\},\{\frac{11}{4},\frac{11}{4},\frac{7}{2}\},x^4)}{28\sqrt{2}\mathrm{\Gamma }(\frac{7}{4})\mathrm{\Gamma }(\frac{11}{4})}}`$
$`{\displaystyle \frac{\left(x^2\right)^{\frac{7}{2}}\left(363_2F_3(\{\frac{5}{4},\frac{7}{4}\},\{\frac{5}{2},\frac{11}{4},\frac{11}{4}\},x^4)56x_2^4F_3(\{\frac{9}{4},\frac{11}{4}\},\{\frac{7}{2},\frac{15}{4},\frac{15}{4}\},x^4)\right)}{2541\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{11}{4})}}`$ (B11)
$`I_3={\displaystyle \frac{x^4}{180\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{7}{4})^2\mathrm{\Gamma }(\frac{11}{4})}}[36\sqrt{x^2}\mathrm{\Gamma }({\displaystyle \frac{7}{4}})\mathrm{\Gamma }({\displaystyle \frac{11}{4}})_2F_3(\{{\displaystyle \frac{5}{4}},{\displaystyle \frac{5}{4}}\},\{{\displaystyle \frac{3}{2}},{\displaystyle \frac{7}{4}},{\displaystyle \frac{9}{4}}\},x^4)`$
$`+\mathrm{\Gamma }({\displaystyle \frac{3}{4}})(27\mathrm{\Gamma }({\displaystyle \frac{11}{4}})_2F_3(\{{\displaystyle \frac{5}{4}},{\displaystyle \frac{5}{4}}\},\{{\displaystyle \frac{7}{4}},{\displaystyle \frac{9}{4}},{\displaystyle \frac{5}{2}}\},x^4)+5x^4\mathrm{\Gamma }({\displaystyle \frac{7}{4}})_2F_3(\{{\displaystyle \frac{9}{4}},{\displaystyle \frac{9}{4}}\},\{{\displaystyle \frac{11}{4}},{\displaystyle \frac{13}{4}},{\displaystyle \frac{7}{2}}\},x^4)))]`$
$`I_4={\displaystyle \frac{x^4\left((J_{\frac{7}{4}}(x^2))^2J_{\frac{11}{4}}(x^2)J_{\frac{3}{4}}(x^2)\right)}{4}}+{\displaystyle \frac{x^4((J_{\frac{3}{4}}(x^2))^2J_{\frac{7}{4}}(x^2)J_{\frac{1}{4}}(x^2))}{4}}+`$
$`{\displaystyle \frac{x^4\left((J_{\frac{1}{4}}(x^2))^2J_{\frac{3}{4}}(x^2)J_{\frac{5}{4}}(x^2)\right)}{4}}{\displaystyle \frac{4\sqrt{2}_2F_3(\{\frac{3}{4},\left(\frac{1}{4}\right)\},\{\frac{3}{2},\frac{1}{4},\frac{1}{4}\},x^4)}{\left(x^2\right)^{\frac{3}{2}}\mathrm{\Gamma }(\frac{3}{4}),\mathrm{\Gamma }(\frac{1}{4})}}`$
$`{\displaystyle \frac{3_2F_3(\{\left(\frac{3}{4}\right),\left(\frac{1}{4}\right)\},\{\left(\frac{1}{2}\right),\frac{1}{4},\frac{1}{4}\},x^4)}{\sqrt{2}\left(x^2\right)^{\frac{3}{2}}(\mathrm{\Gamma }(\frac{1}{4})^2}}{\displaystyle \frac{3\sqrt{2}\sqrt{x^2}_2F_3(\{\frac{1}{4},\frac{3}{4}\},\{\frac{1}{2},\frac{5}{4},\frac{5}{4}\},x^4)}{\mathrm{\Gamma }(\frac{1}{4})\mathrm{\Gamma }(\frac{5}{4})}}`$
$`{\displaystyle \frac{4\sqrt{2}\sqrt{x^2}\left(75_2F_3(\{\frac{1}{4},\frac{1}{4}\},\{\left(\frac{1}{2}\right),\frac{5}{4},\frac{5}{4}\},x^4)+8x_2^4F_3(\{\frac{3}{4},\frac{5}{4}\},\{\frac{1}{2},\frac{9}{4},\frac{9}{4}\},x^4)\right)}{75\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{5}{4})}}`$
$`I_5={\displaystyle \frac{x^4\left((J_{\frac{7}{4}}(x^2))^2J_{\frac{11}{4}}(x^2)J_{\frac{3}{4}}(x^2)\right)}{4}}{\displaystyle \frac{x^4\left((J_{\frac{3}{4}}(x^2))^2J_{\frac{7}{4}}(x^2)J_{\frac{1}{4}}(x^2)\right)}{4}}+`$
$`{\displaystyle \frac{x^4\left((J_{\frac{1}{4}}(x^2))^2J_{\frac{3}{4}}(x^2)J_{\frac{5}{4}}(x^2)\right)}{4}}{\displaystyle \frac{4\sqrt{2}_2F_3(\{\left(\frac{3}{4}\right),\left(\frac{1}{4}\right)\},\{\left(\frac{3}{2}\right),\frac{1}{4},\frac{1}{4}\},x^4)}{\left(x^2\right)^{\frac{3}{2}}\mathrm{\Gamma }(\left(\frac{3}{4}\right))\mathrm{\Gamma }(\frac{1}{4})}}`$
$`{\displaystyle \frac{3_2F_3(\{\frac{3}{4},\frac{1}{4}\},\{\frac{1}{2},\frac{1}{4},\frac{1}{4}\},x^4)}{\sqrt{2}\left(x^2\right)^{\frac{3}{2}}\mathrm{\Gamma }(\frac{1}{4})^2}}{\displaystyle \frac{3\sqrt{2}\sqrt{x^2}_2F_3(\{\frac{1}{4},\frac{3}{4}\},\{\frac{1}{2},\frac{5}{4},\frac{5}{4}\},x^4)}{\mathrm{\Gamma }(\frac{1}{4})\mathrm{\Gamma }(\frac{5}{4})}}`$
$`{\displaystyle \frac{4\sqrt{2}\sqrt{x^2}\left(75_2F_3(\{\frac{1}{4},\frac{1}{4}\},\{\left(\frac{1}{2}\right),\frac{5}{4},\frac{5}{4}\},x^4)+8x_2^4F_3(\{\frac{3}{4},\frac{5}{4}\},\{\frac{1}{2},\frac{9}{4},\frac{9}{4}\},x^4)\right)}{75\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{5}{4})}}`$
$`I_6={\displaystyle \frac{\sqrt{2}}{3\sqrt{x^2}\mathrm{\Gamma }(\frac{3}{4})\mathrm{\Gamma }(\frac{1}{4})^2\mathrm{\Gamma }(\frac{5}{4})}}[12\mathrm{\Gamma }({\displaystyle \frac{1}{4}})\mathrm{\Gamma }({\displaystyle \frac{5}{4}})_2F_3(\{{\displaystyle \frac{1}{4}},{\displaystyle \frac{1}{4}}\},\{{\displaystyle \frac{3}{2}},{\displaystyle \frac{1}{4}},{\displaystyle \frac{3}{4}}\},x^4)`$
$`+\mathrm{\Gamma }({\displaystyle \frac{3}{4}})(9\mathrm{\Gamma }({\displaystyle \frac{5}{4}})_2F_3(\{{\displaystyle \frac{1}{4}}),{\displaystyle \frac{1}{4}}\},\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{4}},{\displaystyle \frac{3}{4}}\},x^4)+x^4\mathrm{\Gamma }({\displaystyle \frac{1}{4}})_2F_3(\{{\displaystyle \frac{3}{4}},{\displaystyle \frac{3}{4}}\},\{{\displaystyle \frac{1}{2}},{\displaystyle \frac{5}{4}},{\displaystyle \frac{7}{4}}\},x^4))].`$
We thus see that the immersion is severely non-linear in nature for the “linear dilaton background”.
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# Entangling power of permutation invariant quantum states
## I Introduction
Recently it has been argued that for critical (gapless) quantum spin systems the entanglement entropy of a block of $`n`$ spins diverges logarithmically as $`\gamma \mathrm{log}_2n`$, while for non-critical systems it converges to a constant finite value Fazio ; Vidal ; Latorre\_condmat . This property was interpreted in the framework of conformal field theory Korepin associated with the corresponding quantum phase transition and the prefactor $`\gamma `$ of the logarithm was related to the central charge $`c`$ of the theory as $`c=3\gamma `$. This was shown explicitly for the exactly solvable antiferromagnetic Heisenberg spin $`1/2`$ chain, i.e., the $`XXZ`$ model where the anisotropy parameter $`\mathrm{\Delta }`$ belongs in the critical regime to the interval $`(1,1)`$ and $`\gamma =1/3`$. Rather surprisingly, at the transition point $`\mathrm{\Delta }=1`$ from gapless to noncritical behaviour, the entanglement (von Neumann) entropy of a block of spins in the ground state was found to grow faster than in the critical domain $`1<\mathrm{\Delta }1`$, namely as $`\gamma \mathrm{log}_2n`$ with the logarithmic prefactor satisfying $`\frac{1}{2}\gamma 1`$ Entanglement\_Heisenberg . At $`\mathrm{\Delta }=1`$ the ground state of the $`XXZ`$ Hamiltonian has permutational invariance (up to a gauge transformation), and is degenerate with respect to the total $`z`$-magnetization, so that the whole system is generically described via a density matrix which can be written as a (weighted) sum of projectors on the multiplet components. The lower bound $`\gamma =\frac{1}{2}`$ is attained if the state of the whole system is a pure state (only one component of the multiplet is present), while the upper bound $`\gamma =1`$ is reached for a mixed state where all the components of the multiplet have equal weights. Note that the pure state is generically not factorizable, see (4) below, thus producing a mixed state after partial tracing.
In the present paper we generalize the approach of Entanglement\_Heisenberg to arbitrary permutation invariant quantum spin states. In particular, we consider the case of a ferromagnetic spin chain with arbitrary spin $`\sigma `$ on every site and show that the entanglement entropy of a block of $`n`$ spins in the ferromagnetic ground state generically diverges as $`\sigma \mathrm{log}_2n`$. Our approach uses the invariance of the ground state under site permutations, thus allowing us to compute the entanglement entropy exactly for blocks of arbitrary size and for systems of arbitrary length.
The paper is organized as follows. In section 2 we introduce the permutation invariant states and list physical systems whose ground states have this symmetry. In section 3 we formulate a theorem which gives the analytical expression of the eigenvalues of the reduced density matrix for arbitrary spin $`\sigma `$. Using this theorem we compute the entanglement entropy of a block of size $`n`$ in the finite system of total length $`L`$ in specific ground state sectors. Taking the limit of large subsystem sizes, we derive analytic expressions for the entanglement entropy $`S_{(n)}`$ both for $`n,L1`$ and for $`n1,L=\mathrm{}`$. As a result, we find that in the ground state sector with a fixed value of $`S^z`$ the block entanglement entropy diverges for large $`n`$ as $`S_{(n)}=\sigma \mathrm{log}_2[2\pi en(Ln)/L]+C`$. In section 4 the case of spin $`\sigma 1`$ is treated in more detail. A discussion and some further remarks close the paper.
## II Permutation invariant states
Let us consider states in a Hilbert space $`(\text{ }\mathrm{C}^d)^L`$ of a quantum spin chain of local spin $`\sigma `$, where $`d=2\sigma +1`$ is the dimension of the spin space at every site $`i`$ and $`d^L`$ the dimension of the whole Hilbert space. Permutation invariant states constitute a subspace $`Q`$ of substantially smaller dimension
$$\kappa (L)=dimQ=\left(\genfrac{}{}{0pt}{}{L+2\sigma }{2\sigma }\right),$$
(1)
this being the number of possible ways to distribute $`L`$ indistinguishable objects among $`2\sigma +1`$ boxes. Denoting by $`N_j`$ the number of objects in the block $`j`$, the state of the whole system is completely characterized by
$`|\mathrm{\Psi }(L,N_0,N_1,`$ $`\mathrm{},N_{2\sigma })=\sqrt{{\displaystyle \frac{N_0!N_1!\mathrm{}N_{2\sigma }!}{L!}}}\times `$ (2)
$`{\displaystyle \underset{P}{}}|\underset{N_0}{\underset{}{\mathrm{}}}\mathrm{}\underset{N_i}{\underset{}{\mathrm{}}}\mathrm{}\underset{N_{2\sigma }}{\underset{}{\mathrm{}}},`$
where $`N_0`$ is the number of spins pointing down ($`\sigma _z=\sigma `$), $`N_j`$ is the number of spins with $`\sigma _z=\sigma +j`$ , up to $`N_{2\sigma }`$ spins with maximal $`\sigma _z=\sigma `$, occupying in total $`L=N_0+N_1+\mathrm{}+N_{2\sigma }`$ sites. The sum is taken over all possible permutations, the total number of which is $`\frac{L!}{N_0!N_1!\mathrm{}N_{2\sigma }!}`$, with the prefactor in (2) taking care of the normalization.
We will be interested in the entanglement of a block of $`n`$ spins with the remaining $`Ln`$ spins (playing here the role of the environment), characterized by the von Neumann entropy
$$S_{(n)}=tr(\rho _n\mathrm{log}_2\rho _n)=\underset{k}{}\lambda _k\mathrm{log}_2\lambda _k,$$
(3)
where $`\rho _n`$ is the reduced density matrix of the block, obtained from the density matrix $`\rho `$ of the whole system by tracing out the degrees of freedom of the environment $`\rho _{(n)}=tr_{(Ln)}\rho `$, and $`\lambda _k`$ are its eigenvalues. The density matrix of the whole system is a projector on the pure state in (2), i.e. $`\rho =|\mathrm{\Psi }(L,N_0,N_1,\mathrm{}N_{2\sigma })\mathrm{\Psi }(L,N_0,N_1,\mathrm{}N_{2\sigma })|`$. Notice that due to the permutational symmetry, $`S_{(n)}`$ does not depend on the particular choice of the sites in the block but only on its size $`n`$. The eigenvalues $`\lambda _k`$ of the reduced density matrix are by construction all real, nonnegative, and sum up to one: $`_k\lambda _k=1`$.
Before giving the general expression for $`\lambda _k`$ we discuss separately the two-states case of $`\sigma =1/2`$. In this case the state of the system (2) reduces to
$$|\mathrm{\Psi }(L,N)=\left(\genfrac{}{}{0pt}{}{L}{N}\right)^{1/2}\underset{P}{}|\underset{N}{\underset{}{\mathrm{}}}\underset{LN}{\underset{}{\mathrm{}}}$$
(4)
where $``$ and $``$ denote spin up and spin down respectively. This state appears in the literature in different physical situations. Since the entanglement properties do not depend on the underlying model but only on the form of the wavefunction, we list some models for which (4) is an exact ground state to show the diversity of applications.
i) The isotropic Heisenberg ferromagnet, $`H=J_{i=1}^L\stackrel{}{\sigma _i}\stackrel{}{\sigma _{i+1}}`$ where $`\sigma _i`$ are Pauli matrices, $`J<0`$ denotes the exchange constant and $`L`$ is the total number of spins. (We assume periodic boundary conditions $`L+11`$.) The ground state belongs to a multiplet of total spin $`S=L/2`$ and is $`(L+1)`$-fold degenerate with $`S^z=\frac{L}{2},\frac{L}{2}+1,\mathrm{}\frac{L}{2}`$. The state (4) is a pure state, corresponding to the multiplet component with a fixed number $`N`$ of spins up.
ii) The Heisenberg antiferromagnet at $`\mathrm{\Delta }=1`$. In this case, the state (4) is obtained from the ground state by a unitary transformation, inverting every other spin along the chain (note that the von Neumann entropy is invariant under unitary transformations of the state of the system).
iii) The generalized Hubbard model in the limit of strong attraction (the so-called eta-pairing states Albertini ).
iv) Hardcore bosons on a complete graph Toth , Penrose ,Mario .
v) The Lipkin-Meshkov-Glick model LMG .
The state (4) also appears in a classical context as stationary distribution of the asymmetric exclusion process (ASEP) describing a Markov process of nonequilibrium stochastic motion of $`N`$ particles on a ring with $`L`$ sites with hardcore exclusion. In this case (4) means that all particle configurations have equal probabilities in the stationary state Gunter\_Hamiltonian\_approach . The block entanglement entropy for the case of spin $`\sigma =1/2`$ was obtained in Entanglement\_Heisenberg . In particular, it was shown that the eigenvalues $`\lambda _k`$ of the reduced density matrix $`\rho _n`$ are $`\lambda _k(L,n)=C_k^nC_{Nk}^{Ln}/C_N^L,`$ where $`C_m^n`$ denotes the binomial coefficient $`n!/(m!(nm)!)`$ and $`k=0,1,\mathrm{}\mathrm{min}(n,N)`$. In the limit of large $`n`$, the von-Neumann entropy was found to be
$$S_{(n)}(p)\frac{1}{2}\mathrm{log}_2(pq)+\frac{1}{2}\mathrm{log}_22\pi e\frac{n(Ln)}{L}.$$
(5)
where $`p=\frac{N}{L},q=1p`$. In the next section we generalize these results to the case of arbitrary spins.
## III Entanglement entropy for arbitrary on-site spin $`\sigma `$
We formulate the main result of this section in the following
Theorem: The eigenvalues of the reduced density matrix $`\rho _n(N_0,N_1,\mathrm{}N_{2\sigma })`$ of a block of $`n`$ spins in the permutation invariant state (2) of the whole system are given by
$$\lambda _𝐤(L,n,\sigma )=\frac{1}{C_n^L}\underset{i=0}{\overset{2\sigma }{}}C_{k_i}^{N_i}$$
(6)
where $`k_i,N_i`$ satisfy $`k_0+\mathrm{}+k_{2\sigma }=n`$, $`N_0+\mathrm{}+N_{2\sigma }=L`$. Here and below we denote the set of $`k_i`$ by the bold $`𝐤`$. To prove the theorem we note that the reduced density matrix $`\rho _n`$ is decomposed with respect to symmetric orthogonal subspaces of the system of $`n`$ spins, classified by the numbers $`k_j=0,1,\mathrm{}\mathrm{min}(n,N_j)`$ of spins with $`\sigma _z=\sigma +j`$ in the block
$$\rho _n(N)=\underset{𝐤}{}w_𝐤|\psi (n,𝐤)\psi (n,𝐤)|.$$
(7)
Here the state $`|\psi (n,𝐤)`$ has the same structure as (2),
$$|\psi (n,𝐤)=\underset{P}{}|\underset{k_0}{\underset{}{...}}\mathrm{}\underset{k_i}{\underset{}{\mathrm{}}}\mathrm{}\underset{k_{2\sigma }}{\underset{}{\mathrm{}}}$$
(8)
and $`w_k`$ is the corresponding probability, given by the number of ways one can distribute $`N_j`$ spins (of values $`\sigma +j`$) on $`k_j`$ sites for all possible values of $`N_j,k_j`$, $`j=0,\mathrm{},2\sigma `$, divided by the number of ways one can distribute $`N=N_0+\mathrm{}+N_{2\sigma }`$ spins on $`n=k_0+\mathrm{}+k_{2\sigma }`$ sites (the total number of states), i.e.
$$w_𝐤=\frac{1}{C_n^L}\underset{j=0}{\overset{2\sigma }{}}C_{k_j}^{N_j}.$$
(9)
Since the states $`|\psi (n,𝐤)`$ are orthogonal, the representation in (7) is diagonal and the $`w_𝐤`$ coincide with the eigenvalues of $`\rho _n(N)`$. This proves of the theorem. Notice that for the case $`\sigma =1/2`$ Eq. (6) reproduces the results derived in Entanglement\_Heisenberg (use $`N_0=LN_1`$, $`k_0=nk_1`$ and the invariance under the exchange of $`N_1`$ and $`n`$).
Having found the eigenvalues of $`\rho _n(N)`$ one can compute the entanglement entropy $`S_{(n)}`$ for arbitrary $`L,n`$ and $`N`$ . In the large-$`n`$ limit one obtains, see the Appendix, in analogy with the case of $`\sigma =1/2`$ Entanglement\_Heisenberg , that
$`S_{(n)}(L,\sigma )`$ $`=C+\sigma \mathrm{log}_2\left(2\pi e{\displaystyle \frac{n(Ln)}{L}}\right),\text{ }n1\text{,}`$ (10)
$`C`$ $`={\displaystyle \frac{1}{2}}\mathrm{log}_2\left({\displaystyle \underset{k=0}{\overset{2\sigma }{}}}p_k\right),`$ (11)
where $`p_k=\frac{N_k}{L}`$. In Figs. 1 ,2 we compare the exact entropy of finite systems, as computed from the exact expressions Eqs. (3, 6), with the analytical expression (10), from which we see that there is an excellent agreement also for small values of $`n{\displaystyle \underset{k=0}{\overset{2\sigma }{}}}p_k`$. In the thermodynamic limit $`L\mathrm{},\frac{N_i}{L}p_i`$, the eigenvalues of the reduced density matrix (6) simplify to
$$\lambda _𝐤(\mathrm{},n,\{p_i\}_{i=0}^{2\sigma })=n!\underset{i=0}{\overset{2\sigma }{}}\frac{(p_i)^{k_i}}{k_i!}$$
The entanglement entropy then becomes
$$S_{(n)}(\mathrm{},\{p_i\}_{i=0}^{2\sigma })=C(\{p_i\})+\sigma \mathrm{log}_2\left(2\pi en\right)$$
(12)
where the constant $`C`$ is given by (11).
Notice that while the general formulae (3,6) are valid for an arbitrary choice of parameters, the analytical result (10) is valid for $`n{\displaystyle \underset{k=0}{\overset{2\sigma }{}}}p_k1`$. If in the limit of an infinite system $`L\mathrm{}`$ some of the $`p_k`$ vanish (say, $`z`$ coefficients $`p_k0`$ out of a total number of coefficients $`2\sigma `$), the entangled state of $`n`$ spins will effectively behave like the one with effective site spin $`\stackrel{~}{\sigma }=\sigma \frac{z}{2}`$, and the von Neumann entropy will respectively grow as $`S_n(\sigma \frac{z}{2})\mathrm{log}_2n`$ instead of $`S_n\sigma \mathrm{log}_2n`$, see (12).
It is instructive to recall that Eq. (10) is derived under the assumption of a pure state of the global system. If the state of the whole system is mixed in the ensemble of states with the same symmetry (2), the resulting von-Neumann entropy $`S_{(n)}`$ will become larger. An upper bound of $`S_{(n)}`$ can be derived by noting that the number of nonzero eigenvalues of the reduced density matrix is equal to the number of terms in the decomposition (7), which is $`\kappa (n)`$, see (1). The upper limit of the entanglement entropy is thus reached when all $`\lambda _k`$ in (3) are equal, implying
$$sup(S_{(n)})=\mathrm{log}_2\left(\genfrac{}{}{0pt}{}{n+2\sigma }{2\sigma }\right).$$
(13)
This limit is attained for a mixed state of the global system with all components of the ensemble (2) equally weighted. In this case the von-Neumann entropy is given by the above expression even for finite $`n,L`$ (see Entanglement\_Heisenberg for an example). For large $`n`$, $`sup(S_{(n)})2\sigma \mathrm{log}_2n`$.
We also remark that in the critical spin models where spin-spin correlation decay algebraically in the ground state (e.g., the region $`1<\mathrm{\Delta }1`$ of the antiferromagnetic Heisenberg chain), there are three distinct physical properties contributing to the entanglement: (a) on-site correlations due to single-site quantum fluctuations, (b) algebraically decaying spin-spin correlations which survive in the thermodynamic limit and (c) the correlations due to the constraint of fixed magnetization which vanish in every domain of finite size in the thermodynamic limit. For permutation invariant states (2) only contribution (a) is left in the thermodynamic limit for finite domains, but for domains of finite volume fraction also the correlations (c) due to the constraint of fixed magnetization remain relevant.
## IV Models with higher spin
As remarked above the entanglement entropy of a quantum state does not depend on any underlying model but only on the properties of that state. Nevertheless it is of interest to have some insight for which systems the permutation invariant states (2) considered here are the ground state of that quantum system. First of all, the generalizations of the Sutherland model Sutherland describing quantum spin chains with an interaction given by the permutation operator in $`SU(N)`$,
$$H=J_{ij}P_{i,j}$$
(14)
are obviously invariant under the permutation group $`S_N`$. Here the set of $`J_{ki}`$ is defined on any connected graph, an example being nearest neighbor interaction. For ferromagnetic interaction, all $`J_{ij}<0`$, the states (2) span its ground state, which can be proved along the lines of Pratt . For $`SU(2)`$ and $`J_{ij}=J\delta _{i,j+1}`$ the Hamiltonian (14) reduces to the isotropic Heisenberg Hamiltonian. For the general case we recall that the permutation operator is written in terms of spin operators as
$$P=\underset{i=0}{\overset{2S}{}}(1)^{2S+i}\underset{ki}{\overset{2S}{}}\frac{2\left(\stackrel{}{𝐒}\stackrel{}{𝐒}\right)k(k+1)+2S(S+1)}{i(i+1)k(k+1)}$$
where $`N=2S+1`$, (see, e.g., BatchelorMaslen ). For the above example, all states (2) are eigenstates of (14) with the lowest energy.
The increase of entanglement entropy compared to the $`s=1/2`$ case appears to result from the larger local state space of the $`SU(N)`$ chain, but not from the $`SU(N)`$-symmetry itself. This picture is supported by the generalized disordered SU(2)-symmetric spin-$`\sigma `$ Heisenberg ferromagnet
$$H=\underset{i,j}{}J_{i,j}g(\stackrel{}{\sigma _i}\stackrel{}{\sigma _j}1)$$
(15)
where the exchange energies $`J_{i,j}0`$ may be non-zero between any pair of lattice sites on an arbitrary lattice and $`\stackrel{}{\sigma _i}`$ are local $`SU(2)`$ generators in the spin-$`\sigma `$ representation. For any polynomial function $`g`$ with positive expansion coefficients linear combinations of the permutation invariant states (2) with fixed total magnetization are ground states of the Hamiltonian, see Schu94 for a detailed discussion of the ground states of this model in a probabilistic setting. Since these are not pure states in the sense discussed above (with all quantum numbers $`N_k`$ fixed) the entanglement entropy of these ground states is higher than those of the pure ground states of the $`SU(N)`$ spin chain.
Other models with pair interaction, but no symmetry: Using the Perron-Frobenius theorem it is straightforward to construct quantum Hamiltonians of the structure
$$H=\underset{i,j}{}J_{i,j}g_{i,j}$$
(16)
where $`g_{i,j}`$ is a hermitian pair interaction matrix satisfying $`g_{i,j}|s=0`$ for all $`i,j`$ and where the wave function $`|s`$ which is constant for all spin configurations in $`(\text{ }\mathrm{C}^d)^L`$ is the ground state of $`g_{i,j}`$. Then, if all $`J_{i,j}0`$, the vector $`|s`$ is also the ground state of $`H`$, and if furthermore $`g_{i,j}`$ has no invariant subspaces, it is the unique ground state with maximal entanglement entropy (13). Such quantum systems also have a probabilistic interpretation as generator of some irreducible Markov chain Gunter\_Hamiltonian\_approach .
## V Summary and conclusions
We have obtained exact eigenvalues of the reduced density matrix for permutation invariant states, for arbitrary length of the system. In the thermodynamic limit, it was shown that the von Neumann entropy of entanglement of a block of $`n`$ spins $`\sigma `$ with the environment grows logarithmically fast $`S_{(n)}\gamma \mathrm{log}_2n`$ , with a prefactor $`\gamma =\sigma `$ for a pure global state and with $`\gamma =2\sigma `$ for homogenously (maximally) mixed global state. Various models, the ground states of which are permutation invariant, are given, for spin $`\sigma =1/2`$ and higher. We note that the logarithmic growth of entanglement entropy due to permutation invariance is faster than the one of critical (conformally invariant) models, where $`S_{(n)}\gamma \mathrm{log}_2n`$ with $`\gamma =1/3`$ was observed Vidal .
It is also interesting to compare the finite size corrections of the entropy of the permutation invariant states (10) with those of critical spin chains, $`S_{(n)}^{cr}\frac{c}{3}\mathrm{log}_2(\frac{L}{\pi }\mathrm{sin}(\frac{\pi n}{L}))`$, obtained in Cardy , see Eq.(3.8) in this paper. Expanding in the first nonvanishing order of $`1/L`$ both expressions, we get the finite size corrections $`\mathrm{\Delta }^{cr}(n)=S_{(n)}^{cr}(L)S_{(n)}^{cr}(\mathrm{})`$, $`\mathrm{\Delta }^{per}(n)=S_{(n)}^{per}(L)S_{(n)}^{per}(\mathrm{})`$ for permutation invariant and critical models as
$`\mathrm{\Delta }^{per}(n)`$ $`=\sigma \mathrm{log}_2(1{\displaystyle \frac{n}{L}})\sigma {\displaystyle \frac{1}{\mathrm{ln}2}}{\displaystyle \frac{n}{L}}+O\left({\displaystyle \frac{n}{L}}\right)^2`$
$`\mathrm{\Delta }^{cr}(n)`$ $`{\displaystyle \frac{c}{3}}\mathrm{log}_2(1{\displaystyle \frac{1}{3}}({\displaystyle \frac{\pi n}{L}})^2){\displaystyle \frac{c}{9}}{\displaystyle \frac{1}{\mathrm{ln}2}}\left({\displaystyle \frac{\pi n}{L}}\right)^2+O({\displaystyle \frac{n}{L}})^4`$
Thus, besides the difference of the coefficient of the principal logarithmic divergence ($`\gamma \mathrm{log}_2n`$ with $`\gamma =1/3`$ for critical $`XXZ`$ and $`\gamma =\sigma `$ for permutation invariant states resp., a substantial difference in the finite size corrections (linear in $`n`$ and of order $`1/L`$ for permutation invariance) and quadratic in $`n`$ and of order $`1/L^2`$ resp. (for conformal invariant critical states) is also observed.
###### Acknowledgements.
VP acknowledges the INFM, Unitá di Salerno, for providing a three months grant during which this work was initiated, and the Department of Physics of the University of Salerno for the hospitality. MS acknowledges partial financial support from Forschungszentrum Jülich, from the University of Cologne within SFB/TR12 project ”Symmetries and Universality in Mesoscopic Systems” and from a MURST-PRIN-2003 Initiative.
## VI Appendix
To prove Eq. (10) we first give the general scheme and then demonstrate the result on a particular example. For simplicity we shall discuss the limit when the size of the global system tends to infinity, $`L\mathrm{}`$. The eigenvalues (6) $`\lambda _𝐤\lambda (k_0,k_1,\mathrm{}k_{2\sigma })`$ for large values of $`n`$ and $`k_i`$, can be approximated by a multi-dimensional Gaussian distribution with mean and moments given by
$$k_i=np_i$$
(17)
$$\left(k_ik_i\right)^2=np_i\left(1p_i\right)$$
(18)
$$\left(k_ik_i\right)\left(k_jk_j\right)=np_ip_j\text{}ij$$
(19)
where $`p_i=\frac{N_i}{L}`$, $`i=0,1,\mathrm{}2\sigma `$, is the average density of spins $`\sigma +i`$ in the system. To prove Eqs. (17)-(19) we first observe that, using the Stirling approximation $`m!=m^mexp(m)\sqrt{2\pi m}`$, in the limit $`L,N_i\mathrm{}`$ the eigenvalues (6) of the reduced density matrix can be approximated as
$$\lambda _𝐤c_𝐤(n)\underset{i=0}{\overset{2\sigma }{}}p_i^{k_i}$$
(20)
with $`c_𝐤(n)=n!/(k_0!k_1!\mathrm{}k_{2\sigma }!)`$. Notice that in this limit the distribution of the $`\lambda _𝐤`$ coincides with the multinomial distribution
$`\left({\displaystyle \underset{i=0}{\overset{2\sigma }{}}}p_i\right)^n`$ $`={\displaystyle \underset{𝐤}{}}c_𝐤(n){\displaystyle \underset{i=0}{\overset{2\sigma }{}}}p_i^{k_i}={\displaystyle \underset{𝐤}{}}\lambda _𝐤=1.`$
Using this expression, the mean value of $`k_\alpha `$, $`\alpha =0,1,\mathrm{},2\sigma `$, can be found as
$`k_\alpha ={\displaystyle \underset{𝐤}{}}k_\alpha c_𝐤(n){\displaystyle \underset{i=0}{\overset{2\sigma }{}}}p_i^{k_i}=p_\alpha {\displaystyle \frac{}{p_\alpha }}({\displaystyle \underset{i=0}{\overset{2\sigma }{}}}p_i)^n=np_\alpha .`$
In a similar manner we obtain
$`k_\alpha \left(k_\alpha 1\right)=p_\alpha ^2{\displaystyle \frac{^2}{p_\alpha ^2}}\left({\displaystyle \underset{i=0}{\overset{2\sigma }{}}}p_i\right)^n=n(n1)p_\alpha ^2,`$
$`k_\alpha k_\beta =p_\alpha p_\beta {\displaystyle \frac{^2}{p_\alpha p_\beta }}\left({\displaystyle \underset{i=0}{\overset{2\sigma }{}}}p_i\right)^n=n(n1)p_\alpha p_\beta ,`$
from which Eqs. (18), (19) readily follow. Denoting by $`x_i=k_i/n`$, so that $`0<x_i<1`$, the eigenvalues $`\lambda _𝐤`$ are approximated as
$$\lambda _𝐤\frac{\sqrt{detA}}{(2\pi )^\sigma }\mathrm{exp}\left(\frac{1}{2}\underset{i,j=1}{\overset{2\sigma }{}}A_{ij}\left(x_ip_i\right)\left(x_jp_j\right)\right)$$
(21)
where the shifting of $`x_i,x_j`$ is introduced to account for the nonzero mean (17), and the coefficients of the symmetric matrix $`a_{ij}=a_{ji}`$ are fixed using the moments (18), (19). On the other hand, the computation of the moments from the distribution (21) gives
$$\left(x_ip_i\right)\left(x_jp_j\right)=\frac{M_{ij}}{detA}$$
(22)
where $`M_{ij}`$ are the minors of the matrix $`A`$. Comparing Eq. (22) with Eqs. (18 ),(19), we can fix the elements $`A_{ij}`$ of the matrix. To illustrate the method let us consider the case of spin $`\sigma =1`$. Then,
$$\lambda _𝐤\lambda (x,y)\frac{\sqrt{D}}{2\pi }\mathrm{exp}\left(\frac{1}{2}\left(ax^2+2bxy+cy^2\right)\right)$$
where we introduced $`x=x_1p_1,`$ $`y=x_2p_2`$ and $`D=detA=acb^2`$ for brevity of notation. From (18),(19) we have $`x^2=np_1(1p_1)`$, $`y^2=np_2(1p_2)`$, and $`xy=np_1p_2`$. On the other hand, from (22) we obtain $`x^2=c/D`$, $`y^2=a/D`$, and $`xy=b/D`$. Computing the determinant
$$D=acb^2=D^2\left(x^2y^2xy^2\right)$$
and substituting the moments, we get
$$D^1=n^2p_1p_2p_3\text{.}$$
The von Neumann entropy is then computed as
$$S𝑑x𝑑y\left(\lambda (x,y)\mathrm{log}_2\lambda (x,y)\right)$$
Notice that for finite $`p_i`$ the larger contribution to the integral comes from the neighbor of the origin $`x=y=0`$ so that we can extend the limits of integration to the whole real axis and perform the integral exactly. This leads to
$$S(n)\mathrm{log}_2\frac{2n\pi e}{\sqrt{D}}=\mathrm{log}_22\pi ne+\frac{1}{2}\mathrm{log}_2\left(p_1p_2p_3\right),$$
which coincides with the expression (10) in the limit $`L=\mathrm{}`$ for the case of spin $`\sigma =1`$. A more detailed analysis, analogous to the one done in Entanglement\_Heisenberg , restores the finite-size dependence on $`L`$ of the von-Neumann entropy as
$$S(n)\mathrm{log}_2\left(2\pi e\frac{n(Ln)}{L}\right)+\frac{1}{2}\mathrm{log}_2\left(p_1p_2p_3\right)\text{ for }\sigma =1$$
Working out the same procedure for an arbitrary spin $`\sigma `$ we obtain (10).
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# IFT-UAM/CSIC-05-24 A New Mechanism of Kähler Moduli Stabilization in Type IIB Theory
## 1 Introduction
The problem of stabilization of moduli fields in string theory (scalar fields with flat directions in the potential) has been extensively studied, see for example , since it has important theoretical and phenomenological consequences. From the experimental point of view the existence of such massless fields should be observed in fifth force experiments <sup>1</sup><sup>1</sup>1I would like to thank J. Conlon for this remark., but since it is not, the consistency of the theory requires such fields to acquire a mass. Cosmological observations also impose constraints on the moduli fixing scale to reproduce reheating at the inflationary epoch, having as a lower bound 100 TeV. The moduli fields come from two different sectors: closed and open string sectors. Moduli associated to the closed string sector give information about the size (Kähler moduli) and the shape (complex structure moduli) of the compact manifold, and also the dilaton. Moduli associated to the open string sector correspond to the presence of Wilson lines and to the parametrization of the position of D-branes in those cases when they are present, in the transverse dimensions.
A relevant result found in showed that a linear combination of RR, and NSNS three-form fluxes on IIB theory were able to stabilize the dilaton and the complex structure moduli. This mechanism presented a great advance in the resolution of the problem although the Kähler moduli remained unfixed. The basic reason for this is due to the superpotential. It does not contain any dependence on the Kähler moduli, leading to a no-scale scalar potential at leading order in $`\alpha ^{^{}}`$. In KKLT found a way to fix one overall Kähler modulus by using non perturbative mechanisms such as condensation of gauginos or instanton effects. This led to an Anti De-Sitter (AdS) vacuum in a particular compactification manifold. By explicitly breaking supersymmetry through the introduction of anti-D3 branes and by fine tuning fluxes, they were able to lift it to a de Sitter vacuum. de Sitter vacua have acquired great importance due to the recent data that suggest the acceleration of the universe and also because of their close relationship with the inflationary scenario .
New advances in the context of the landscape have been achieved by Douglas et al. obtaining manifolds with all moduli fixed through non-perturbative mechanisms and able to lead to a static cosmology. These last advances, although significant, have not been able to provide for realistic compactification manifolds. A potential problem in generating non-perturbative superpotentials from strong infrared dynamics is that it is model dependent. It can also generate too much massless charged matter.
In the model of KKLT was improved. They induce a supersymmetric model where one Kähler modulus is present by introducing magnetic fluxes contained on $`D7`$ branes.
Model building in this context has several fine-tuning and stability issues. In particular, the superpotential induced by the fluxes must be hierarchically small $`(<10^4)`$ in order to obtain solutions with large volume in which the effective field theory approximation can be trusted. On the other hand the fluxes must fix the dilaton at small string coupling to suppress loop effects. Very recently some papers have appeared , see also , that find a way to stabilize all moduli by considering the combination of $`\alpha ^{}_{}{}^{}3`$ effects and non perturbative contributions to the superpotential, giving rise to a large volume AdS vacuum. Its minimum is independent of the value of the flux superpotential. In this model supersymmetry is broken by the Kähler moduli and the gravitino mass is not flux dependent through the superpotential.
Here we propose a perturbative mechanism of moduli fixing which is fullfilled at the supersymmetric minimum of the theory and it is able to lead to realistic descriptions with all of the moduli fixed. In principle one could think that it is possible to induce any other dependence in the superpotential by introducing branes in the model. However it has been conjectured by that B-type branes (D-branes which wrap 2n-cycles with magnetic fluxes and the type of branes interesting in IIB models) can couple to the Kähler moduli only through Fayet-Illiopoulos (FI) terms. Several other works have also studied this problem. We will consider a IIB theory compactified on Calabi-Yau (CY) orientifolds, with RR, NSNS 3-form fluxes and magnetic fluxes. Three-form fluxes couple to the three cycles via the superpotential associated to the F-term, which does not depend on the Kähler moduli. We show that taking into account the coupling of the 3-form fluxes to the open string sector of magnetized D-branes which leads to flux induced mass terms, i.e.$`\mu `$-terms or soft breaking terms with magnetic fields, together with FI terms gives a scalar potential which stabilizes the Kähler moduli. This mechanism is model independent and we think that it is a generic procedure for stabilizing moduli in a manifold. One advantage of this method is that by being supersymmetric or breaking spontaneously supersymmetry, it can be put in the 4D standard supergravity form and it keeps control over the types of interactions that can be induced. In a previous paper a method was proposed to stabilize some Kähler moduli through the coupling of fluxes to the FI-term in order to fix the blow-up moduli of the model. In that case the expected soft breaking term contribution did not include magnetic fields and failed in the attempt of stabilizing the full Kähler moduli. This idea is an extension of that one including magnetic fields in the configurations. In , in a different approach, they also consider mass terms as a mechanism to stabilize moduli.
In a type I theory was proposed with three-form fluxes and magnetic fluxes. They claim that they are able to stabilize Kähler moduli just through FI-terms. We argue that what they find does not constitute a true stabilization of the moduli since the moduli are free to acquire any other vev without any energy cost. To stabilize Kähler moduli it is necessary to also have a coupling of fluxes to the open string sector.
The paper is organized as follows: in section 2 we give a brief summary of three-form flux stabilization. We show how a linear combination of three form fluxes stabilizes complex structure moduli and the dilaton under very general assumptions. Non perturbative mechanisms are used to stabilize Kähler moduli. In section 3 we review how soft terms appear. We particularly focus on the flux induced soft breaking terms with magnetic fields and we find a general expression for these terms in a toroidal orientifold generalizing to the case when all Kähler moduli are different. In section 4 we describe D-term supersymmetry breaking. FI terms that couple to B-type branes represent a deviation from the supersymmetry conditions for the branes and a shift in the value at which the moduli have a supersymmetric value. However a proper recombination mechanism can restore the supersymmetry. In section 5 we propose the mechanism for stabilizing moduli without introducing non perturbative mechanisms and we describe in detail the minimization of the scalar potential. In section 6 we provide a concrete realization of this mechanism (of phenomenological interest) for the case of ISD three-form fluxes. We perform IIB compactified on $`T^6/Z_2\times Z_2\times \mathrm{\Omega }R`$ moduli stabilization and we indicate the explicit values at which Kähler moduli get fixed. We conclude in section 7 with a discussion, summarizing the main results.
## 2 Review of moduli stabilization
In type IIB theory on a Calabi-Yau manifold the closed string moduli content associated to the geometry are: an axion-dilaton $`S`$, $`h_{11}`$ Kḧaler moduli $`T_i`$ parametrizing the $`CY_3`$ size and $`h_{21}`$ complex structure moduli $`U_i`$ parametrizing the $`CY_3`$ shape, where $`h_{11},h_{21}`$ are the Hodge numbers characterizing the Calabi-Yau three-fold. The Kähler potential associated to the closed string sector has, for toroidal compactifications where the metric factors are into three two-by-two blocks, the following expression,
$`k^2𝒦(S,U_i,T_i)=ln[S+\overline{S}]{\displaystyle \underset{j=1}{\overset{h_{11}}{}}}ln[T_j+\overline{T}_j]{\displaystyle \underset{j=1}{\overset{h_{21}}{}}}ln[U_j+\overline{U}_j].`$ (2.1)
Giddings, Kachru and Polchinski introduced the methods to stabilize the dilaton and complex structure moduli by turning on 3-form fluxes. In type IIB theory, strings can have RR and NSNS 3-form field strengths, which can wrap dual 3-cycles labeled by A and B leading to quantized background fluxes,
$`{\displaystyle \frac{1}{4\pi ^2\alpha ^{^{}}}}{\displaystyle _A}F_3=M{\displaystyle \frac{1}{4\pi ^2\alpha ^{^{}}}}{\displaystyle _B}H_3=K`$ (2.2)
where K and M are arbitrary integers. These forms allow consistent string compactifications of generic orientifold $`CY_3`$ manifolds. Fluxes also have some other interesting consequences: they induce a warp factor in the metric that deforms the manifold and generates hierarchies , fix the moduli partially, and also give a mechanism to break supersymmetry in a controlled manner by inducing soft breaking terms.
To understand the mechanism of the partial fixing of moduli we have to remark that these fluxes generate a superpotential found by
$`W={\displaystyle _{CY_3}}G_3\mathrm{\Omega }_3,`$ (2.3)
where $`G_3=F_3SH_3`$, with $`S`$ the complex axion-dilaton of type IIB theory. $`\mathrm{\Omega }_3`$ is the unique (3,0) form of the Calabi-Yau threefold. The holomorphic three form $`\mathrm{\Omega }_3`$ has a non-trivial dependence on the complex structure moduli $`U_i`$. This superpotential is independent of the Kähler moduli $`T_i`$. It gives the $`D=4,N=1`$ scalar potential
$`V_F=\mathrm{exp}^{k^2𝒦}(K^{I\overline{J}}D_IW\overline{D_JW}3k^2|W|^2)`$ (2.4)
Here $`I,J`$ label all the above geometric moduli of the manifold.
The covariant derivatives are defined as $`D_IW=_IW+k^2_I𝒦W`$ where $`𝒦`$ denotes the Kähler potential and $`K^{I\overline{J}}`$ the inverse of the Kähler metric defined in terms of the Kähler potential $`K_{I\overline{J}}=_I_{\overline{J}}𝒦`$.
This potential is of no-scale type ($`\mathrm{\Lambda }=0`$ at tree level in $`\alpha ^{^{}}`$) since the Kähler dependence on $`T_i`$ cancels exactly the $`3W^2`$ contribution. Since the potential is positive definite, the global minimum of this potential lies at zero. The values of complex structure moduli and the dilaton get fixed for the particular values at which the superpotential is minimized, $`D_iW=0`$, where $`i`$ runs over all fields except the Kähler moduli. The Kähler moduli however remain unfixed as the superpotential has no dependence on them. The minimum of the potential is supersymmetric if $`D_TW=_I𝒦W=0`$ and non-susy otherwise.
The Kähler moduli have been stabilized by a different mechanism, namely non-perturbative contributions that introduce an exponential dependence. This contribution together with the one induced by fluxes gives the total superpotential
$`W=W_{flux}+A\mathrm{exp}(aT)`$ (2.5)
and generate a scalar potential which typically has an AdS minimum at a finite value of $`T`$ and a run-away behaviour at infinity. Several mechanisms for breaking supersymmetry allow in principle this minimum to be lifted to a de Sitter vacuum, i.e. .
## 3 Soft breaking terms with magnetic fields
The low energy effective action can have susy-breaking soft terms coming from magnetized or non magnetized configurations. Soft terms are operators of $`dim<4`$ which do not induce quadratic divergences that spoil the good properties of the $`N=1`$ supersymmetry. They give mass to the superpartners of the SM fields, when spontaneous supersymmetry breaking occurs. Their scale is expected to be not much above the electroweak scale. In fact when they are induced through fluxes they are of the order of the flux scale. These soft terms arise from the interaction of low dimensional D-brane charges induced on the D7 branes with the background flux (see ). The soft terms contain gaugino masses, trilinear terms, and scalar masses of MSSM fields. Regarding the stabilization of Kähler moduli we are only going to be interested in scalar masses since they are the dominant terms in the expression. In this section we want to extend the results found in , (see also ) to more general settings. Unless specified otherwise we will follow the notation of but extending their results to be valid for a general configuration of D-branes compactified on toroidal orientifolds of type IIB theory on $`T^6/Z_2\times Z_2\times \mathrm{\Omega }R`$ .
We consider a six torus factorized as $`T^6=_{i=1}^3T_i^2`$, performing the quotient by $`Z_2\times Z_2\times \mathrm{\Omega }R`$ symmetry as in , filled with D7 branes some of which are magnetized and D9 branes carrying anti-D3 brane charges in the hidden sector to cancel RR and NSNS tadpole conditions while preserving $`N=1`$ supersymmetry. The magnetic constant field is defined in terms of the wrapping number $`m`$ of each stack of $`D7_a`$ branes transverse to the torus $`i`$, and $`n`$ which represents the units of magnetic flux.
$`{\displaystyle \frac{m_a^i}{2\pi }}{\displaystyle _{T_i^2}}F_a^i=n_a^i`$ (3.1)
$`F_a^i`$ is the world-volume magnetic field. For later convenience we define the following angles,
$`\mathrm{\Psi }_a^i=arctan(2\pi \alpha ^{^{}}F_a^i)=arctan({\displaystyle \frac{\alpha ^{^{}}n_a^i}{m_a^iA_i}})`$ (3.2)
$`A_i`$ represents the area of the two torus. Open strings give rise to charged fields. There are two types of states living on stacks of Dp-branes. Untwisted states are chiral fields coming from open strings whose ends are attached to the same stack of branes and twisted sector chiral fields those lying at two different stacks of magnetized Dp-branes. We are only going to consider the twisted sector fields as they are the ones that have bosonic soft breaking terms. From D7 brane twisted sectors $`D7_iD7_j`$ there are only chiral massless multiplets denoted by $`C^{7i7_j}`$ transforming in the bifundamentals of $`G_i\times G_j`$, with $`G_i`$ being the gauge group associated to the enhanced symmetry of each stack of D7 branes. The low energy dynamics of the massless fields is governed by a $`D=4,N=1`$ supergravity action that depends on the Kähler potential, the gauge kinetic functions and the superpotential. The moduli in this type of constructions, as already explained in , are $`\{M,C_I\}`$ with $`M`$ the closed string moduli and $`C_I=\{C^{7_i},C^{7_i7_j}\}`$, where $`C_i^{7_i}`$ are the moduli fields representing the position of the $`D7_i`$ branes in the transverse $`T_i^2`$ whereas $`C_j^{7_i}`$ correspond to the presence of Wilson lines turned on, in the two complex directions parallel to the $`D7_i`$ brane. In the following we will not care about open string moduli since there are models of phenomenological interest free of them, such as the one we propose in the last section. Closed geometric moduli that appear in the 4D $`N=1`$ supergravity action consist of the complex dilaton
$`S=e^{\varphi _{10}}+ia_0`$ (3.3)
where $`a_0`$ is the R-R 0-form and $`e^{\varphi _{10}}=g_s`$ the string coupling constant; the complex structure moduli $`U_j,j=1,2,3`$, which are equivalent to the following geometrical moduli for the case of toroidal compactifications,
$`\tau _j={\displaystyle \frac{1}{e_{jx}^2}}(A_j+ie_{jx}.e_{jy})`$ (3.4)
where $`e_{jx},e_{jy}`$ are the $`T_j^2`$ lattice vectors, $`A_j`$ is the area of the two-torus, which for the particular case of the square $`T_j^2`$ is equal to $`A_j=R_{jx}R_{jy}`$, and the dual angle is $`\phi _a^i=\mathrm{arctan}(\frac{n_a^iR_{ix}}{m_a^iR_{iy}})`$. The geometric Kähler moduli $`\rho _i`$ are described by,
$`\rho _j=A_j+ia_j,`$ (3.5)
where the axions $`a_j`$ arise from the R-R 4-form. However, as explained in the Kähler moduli field denoted by $`T_i,i=1,2,3`$ are not equivalent to the geometric moduli $`\rho _i`$, since its correct expression is
$`T_i=\mathrm{exp}^{\varphi _{10}}{\displaystyle \frac{A_jA_k}{\alpha ^{}_{}{}^{}2}}+ia_ijik`$ (3.6)
and the gravitational coupling in $`D=4`$ is $`G_N=k^2/8\pi `$ with
$`k^2={\displaystyle \frac{M_{pl}^2}{8\pi }}=\mathrm{exp}^{2\varphi _{10}}{\displaystyle \frac{A_1A_2A_3}{\pi \alpha ^{}_{}{}^{}4}}.`$ (3.7)
The string scale is defined as $`M_s=\frac{1}{\sqrt{\alpha ^{^{}}}}`$. We have chosen this type of compactification for its simplicity: The moduli in this type of compactification are just the dilaton $`S`$, three complex structure moduli associated to the relation between the radius of each tori $`U_i,i=1,.,3`$ and three Kähler moduli associated to the size of each torus $`T_i,i=1,2,3`$.
The Kähler potential contains the Kähler part $`\widehat{K}`$ associated to the closed string moduli, $`M=\{S,T_i,U_i\}`$, and the Kähler part $`\stackrel{~}{K}`$ which is associated to the matter fields $`C_I=\mathrm{\Phi }_{aa},\mathrm{\Phi }_{ab}`$ of the open string sector .
$`k^2𝒦`$ $`=`$ $`k^2\widehat{K}+k^2{\displaystyle \underset{IJ}{}}\stackrel{~}{K}_{IJ}C_{\overline{I}}C_J+k^2{\displaystyle \underset{IJ}{}}Z_{IJ}(C_IC_J+c.c)+\mathrm{}`$ (3.8)
The standard expression for the soft breaking terms found in is
$`m_I^2=m_{3/2}^2+V_0{\displaystyle \underset{M,N}{}}\overline{F}^{\overline{M}}F^N_{\overline{M}}_Nlog(\stackrel{~}{K_{IJ}})`$ (3.9)
in supergravity conventions with all quantities measured in Planck units. $`I=aa,ab,ba,bb`$ labels the different stacks of magnetized D7 branes. $`F^M`$ are the auxiliary fields of the corresponding chiral multiplet $`\mathrm{\Phi }_A`$ which in general have the following expression
$`\overline{F}^{\overline{A}}=k^2\mathrm{exp}^{k^2K/2}K^{\overline{A}B}D_BW`$ (3.10)
where $`\mathrm{\Phi }_A=\{M,C_I\}`$ and as before $`C_I=\{C_j^{7_i},C^{7_i7_j}\}`$. When a spontaneous breaking of supersymmetry occurs, the auxiliary vevs acquire a vacuum expectation value. In this case the supersymmetry breaking is produced by the presence of a non supersymmetric 3-form flux, ($`G_3`$ containing a $`(0,3)`$ piece). The auxiliary fields are parametrized as
$`F^S=\sqrt{3}Csm_{3/2}\mathrm{sin}(\theta )\mathrm{exp}^{i\gamma _s},`$ (3.11)
$`F^{T_i}=\sqrt{3}Ct_i\eta _im_{3/2}\mathrm{cos}(\theta )\mathrm{exp}^{i\gamma _i},`$ (3.12)
such that $`_i\eta _i^2=1`$ and $`\eta `$ and $`\theta `$, the goldstino angle, control whether $`S`$ or $`T_i`$ dominate the SUSY breaking. Here
$`C^2`$ $`=`$ $`1+{\displaystyle \frac{V_0}{3m_{3/2}^2}}`$ (3.13)
$`m_{3/2}`$ $`=`$ $`\mathrm{exp}^{1/2k^2𝒦}W`$ (3.14)
$`V_0`$ $`=`$ $`F^m\widehat{K}_{mn}F^n3m_{3/2}^2`$ (3.15)
and $`m_{3/2}`$ denotes the gravitino mass and $`V_0`$ the cosmological constant.
For the case of soft breaking terms induced by 3-form and magnetic fluxes, the expression becomes
$`m_I^2`$ $`=m_{3/2}^2+V_01/4\overline{F}^{\overline{M}}F^N_{\overline{M}}_N\mathrm{ln}(st_1t_2t_3)+{\displaystyle \underset{i=1}{\overset{3}{}}}(_{\overline{M}}_N\nu _i)\overline{F}^{\overline{M}}F^N(ln(u_i))`$
$`1/2{\displaystyle \underset{i=1}{\overset{3}{}}}\overline{F}^{\overline{M}}F^N_{\overline{M}}_N(\mathrm{ln}(\mathrm{\Gamma }(1\nu _i))\mathrm{ln}(\mathrm{\Gamma }(\nu _i))).`$
where the Kähler potential has the following form,
$`k^2𝒦`$ $`=`$ $`k^2\widehat{K}+k^2{\displaystyle \underset{IJ}{}}\stackrel{~}{K}_{IJ}C_{\overline{I}}C_J+k^2{\displaystyle \underset{IJ}{}}Z_{IJ}(C_IC_J+c.c)+\mathrm{}`$ (3.17)
$`\stackrel{~}{K}`$ $`=`$ $`{\displaystyle \frac{(st_1t_2t_3)^{\frac{1}{4}}}{2\pi \alpha ^{^{}}}}\mathrm{\Pi }_{j=1}^3u_j^{\nu _j}\sqrt{{\displaystyle \frac{\mathrm{\Gamma }(1\nu _j)}{\mathrm{\Gamma }(\nu _j)}}}`$ (3.18)
$`k^2\widehat{K}`$ $`=`$ $`\mathrm{ln}(s){\displaystyle \underset{i}{}}\mathrm{ln}t_i{\displaystyle \underset{i}{}}\mathrm{ln}(u_i)`$ (3.19)
with the conventions of
$`s=S+\overline{S};t_i=T_i+\overline{T_i};u_i=U_i+\overline{U_i}.`$ (3.20)
$`k^2=4\pi \alpha ^{^{}}(st_1t_2t_3)^{1/4}`$, and $`\widehat{\nu }_i=\frac{1}{\pi }(\mathrm{\Psi }_{ab}^i)=\frac{1}{\pi }(\mathrm{\Psi }_b^i\mathrm{\Psi }_a^i)`$, where $`\mathrm{\Psi }_a^i=\mathrm{arctan}(g_a\beta _i)`$ with $`\beta _i=\sqrt{\frac{st_i}{t_{jt_k}}}`$ and $`g_a^i=\frac{n_a^i}{m_a^i}`$. $`\widehat{\nu }_i=0`$ trivially, which is the condition for two stack of D-branes to preserve the same supersymmetry. The $`\nu _i`$ are computed in terms of $`\widehat{\nu }_i/\pi `$, such that $`0<\nu _i<1`$ and $`_{i=1}\nu _i=2`$ as in . Here $`\nu _i=1+\widehat{\nu }_i/\pi `$ iff $`\widehat{\nu }_i0`$ and $`\nu _i=\widehat{\nu }_i/\pi `$ otherwise. This last expression is closely related to the one in , although it has been generalized to the case when all of the Kähler moduli are different.
Regrouping terms,
$`m_{ab}^2`$ $`=m_{3/2}^2+V_01/4\overline{F}^{\overline{M}}F^N_{\overline{M}}_N\mathrm{ln}(st_1t_2t_3)+{\displaystyle \underset{i=1}{\overset{3}{}}}(_{\overline{M}}_N\nu _i)\overline{F}^{\overline{M}}F^N`$
$`(ln(u_i)1/2B_0^i(\nu _i))1/2{\displaystyle \underset{i=1}{\overset{3}{}}}B_1^i(\nu _i)\overline{F}^{\overline{M}}F^N(_{\overline{M}}\nu _i_N\nu _i).`$
The $`B_p`$ are defined in terms of polygamma functions as,
$`B_0^i(\nu _i)`$ $`=`$ $`{\displaystyle \frac{_\nu \mathrm{\Gamma }(1\nu _i)}{\mathrm{\Gamma }(1\nu _i)}}{\displaystyle \frac{_\nu \mathrm{\Gamma }(\nu _i)}{\mathrm{\Gamma }(\nu _i)}},`$ (3.22)
$`B_p^i`$ $`=`$ $`_\nu B_{p1}(\nu _i),`$ (3.23)
and in terms of a useful analytical expression is
$`B_0(z)=\pi cotan(\pi z)`$ (3.24)
and its derivatives. Let us remark that this definition of $`B_p(\nu )^i`$ is slightly different of the one used in . A special value of this function is at $`\nu _i=1/2`$ where $`B_k^i(1/2)=0`$.
We finally obtain the following expression,
$`m_{ab}^2`$ $`=`$ $`m_{3/2}^2+V_03/4C^2m_{3/2}^2[1+1/4\pi ^2{\displaystyle \underset{i=1}{\overset{3}{}}}(ln(u_i)1/2B_0^i(\nu _i))`$
$`[4𝐏_i\mathrm{sin}(2\pi \mathrm{\Psi }_i)_{ab}𝐐_i\mathrm{sin}(4\pi \mathrm{\Psi }_i)_{ab}]+1/8\pi ^3{\displaystyle \underset{i}{\overset{3}{}}}B_1^i(\nu _i)𝐐_i(\mathrm{sin}(2\pi \mathrm{\Psi }_i)_{ab})^2]`$
where $`\mathrm{sin}(\pi \mathrm{\Psi }_i)_{ab}=\mathrm{sin}(\pi \mathrm{\Psi }_i)_b\mathrm{sin}(\pi \mathrm{\Psi }_i)_a`$. and the above variables $`𝐏,𝐐`$ are defined in terms of goldstino angles,
$`𝐏_i`$ $`=`$ $`\mathrm{sin}^2\theta +\mathrm{cos}^2\theta (\eta _i^2\eta _j^2\eta _k^2)`$ (3.26)
$`𝐐_i`$ $`=`$ $`12\mathrm{cos}^2\theta \{\eta _i\eta _jcos(\gamma _i\gamma _j)+\eta _i\eta _kcos(\gamma _i\gamma _k)\eta _j\eta _kcos(\gamma _j\gamma _k)\}`$
$`+\mathrm{sin}(2\theta )\{\eta _icos(\gamma _i\gamma _s)\eta _jcos(\gamma _j\gamma _s)\eta _kcos(\gamma _k\gamma _s\}`$
so they are sensitive to the particular choice of 3-form fluxes.
The expression for soft breaking mass terms can be finally expressed as,
$`m_{ab}^2={\displaystyle \frac{1}{4}}(1+3A)m_{3/2}^2+{\displaystyle \frac{1}{4}}(3+A)V_0`$ (3.28)
where
$`A`$ $``$ $`1/4\pi ^2{\displaystyle \underset{i=1}{\overset{3}{}}}(ln(u_i)1/2B_0(\nu _i))[4𝐏_i\mathrm{sin}(2\pi \mathrm{\Psi }_i)_{ab}𝐐_i\mathrm{sin}(4\pi \mathrm{\Psi }_i)_{ab}]`$ (3.30)
$`1/8\pi ^3{\displaystyle \underset{i}{\overset{3}{}}}B_1(\nu _i)𝐐_i(\mathrm{sin}(2\pi \mathrm{\Psi }_i)_{ab})^2.`$
These results recover the ones of by making appropriate substitutions for their particular configuration and imposing $`t_2=t_3`$. Although the Kähler part of the potential $`\stackrel{~}{K}_{C\overline{C}}`$ correspond to the twisted magnetized sectors, with the above definitions, the soft terms at twisted, unmagnetized $`D7_2D7_3`$ stacks are also correctly found due to appropriate cancellations.
## 4 D-term supersymmetry breaking
In this section we are interested in characterizing D-term behaviour in the presence of B-type branes. The motivation is the following: in order to solve the problem of Kähler moduli stabilization, naively we can think of finding a perturbative generalization of the superpotential that contains Kähler moduli in its definition. However this seems not to be possible. Kähler moduli are the moduli associated to the $`(1,1)`$ forms that naturally can be thought of as being stabilized by magnetic fluxes living on D-branes. In general there are $`h_{11}`$ of them, being the $`h_{11}`$ Hodge number counting the number of 2-cycles present in the compactified manifold. The Decoupling Statement of Douglas e al. establishes that Kähler moduli on Type IIB can only couple to B-type branes (i.e. Branes wrapping even cycles) through Fayet Iliopoulos (FI) terms in the scalar potential. The mirror statement of this on type IIA says that complex structure moduli can only couple to A-branes (branes wrapping odd cycles (3-cycles)) through FI terms. This fact seems disappointing from the perspective of finding a perturbative superpotential for Kähler moduli. Moreover this is a common feature for the case of D-branes at singularities .
In an $`𝒩=1`$ 4D type IIA theory on $`\frac{T^6}{Z_2\times Z_2}`$ orientifold with D6 branes at angles was studied. In this section we will review the main properties. This model is dual to magnetized D9 branes on type IIB theory with discrete torsion. Requiring supersymmetric models imposes a condition between orientifold planes and D branes. On type IIA, in a supersymmetric model each stack of D6 branes is related to the orientifold planes O6 by a rotation in $`SU(3)`$. The supersymmetry configurations imposes a condition on the angles $`\theta _i`$ of the D6 branes with respect to the orientifold plane in the i-th direction of the two-torus, which for the case they considered was,
$`{\displaystyle \underset{i=1}{\overset{3}{}}}\theta _i=0.`$ (4.1)
This condition is equivalent to
$`{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{arctan}(\chi _i{\displaystyle \frac{m_i}{n_i}})=0,`$ (4.2)
with $`R_{xi},R_{yi}`$ the radius of the i-th two-torus and $`\chi _i=\frac{R_{yi}}{R_{xi}}`$ the untwisted complex structure moduli. This model is T-dual to a type IIB orientifolded theory with discrete torsion and twisted Kähler moduli with the condition
$`{\displaystyle \underset{i}{}}\mathrm{\Psi }_a^i`$ $`=`$ $`{\displaystyle \frac{3\pi }{2}}mod(2\pi ),`$ (4.3)
where we are using the convention for angles introduced in the preceding section, $`\mathrm{\Psi }_a^i=arctan(\frac{\alpha ^{^{}}n_a^i}{m_a^iA_i})`$ in terms of the tori areas $`A_i`$. $`A_i`$ are expected to couple open string modes in the D9 branes on 2-cycles (B-branes) as Fayet-Illiopoulos (FI) term. This is the dual version to the one indicated in . In it is explained that FI terms are proportional to the deviation from the susy condition, giving an effective action proportional to the deviation,
$`{\displaystyle \underset{a}{}}{\displaystyle 𝑑x^4\xi _aD_a}`$ (4.4)
where for small FI terms , , $`\xi _a=\delta _i\mathrm{\Psi }_a^i=(_i\mathrm{\Psi }_a^isusycond.)`$ and vanishes for supersymmetric configurations of D-branes. Since $`D_a=_b^{N_b}\left|\varphi _{ab}\right|^2\left|\varphi _{ba}\right|^2+\xi _a`$, and $`V_D=\frac{1}{2}D_aD^a`$ then, the D-term piece of the scalar potential is then
$`V_D={\displaystyle \underset{a}{\overset{N_a}{}}}({\displaystyle \underset{b}{\overset{N_b}{}}}\left|\varphi _{ab}\right|^2{\displaystyle \underset{b}{}}\left|\varphi _{ba}\right|^2+\xi _a)^2,`$ (4.5)
where $`\left|\varphi _{ab}\right|`$ represents the charged matter fields lying on an oriented string with ends attached at two different stacks $`a`$ and $`b`$. $`\left|\varphi _{ba}\right|`$ represents the matter fields at the intersection of the same two stacks $`a`$ and $`b`$ with reversal orientation and has also to be considered. $`N_a,N_b`$ are the number of parallel Dp-branes at each stack.
The dependence of the scalar potential on a $`\varphi _{ab}`$ matter field is
$`V_D(\varphi _{ab})=(\left|\varphi _{ab}\right|^2+1/2\xi _a)^2+(\left|\varphi _{ab}\right|^2+1/2\xi _b)^2`$ (4.6)
where we have renamed for the second part of the R.H.S $`ab`$ and we have used that $`\left|\varphi _{ba}\right|^2=\left|\varphi _{ab}\right|^2`$ and the mixed term which gives a mass term for $`\left|\varphi _{ab}\right|`$
$`\left|\varphi _{ab}\right|^2(\xi _a\xi _b)=\left|\varphi _{ab}\right|^2\delta \mathrm{\Psi }_{ab},`$ (4.7)
defining $`\delta \mathrm{\Psi }_{ab}=\xi _b\xi _a`$ and approaching $`(\delta \mathrm{\Psi }_{ab})^2\frac{1}{2}(\xi _a^2+\xi _b^2)`$, we arrive to the familiar expression for D-terms appearing from twisted sector, :
$`V_D={\displaystyle \underset{I}{}}(q_I\left|C_I\right|^2\delta \mathrm{\Psi }_I)^2`$ (4.8)
where $`I`$ runs over all the indices of the matter field , i.e, $`aa,ab,ba,bb`$. $`q_I`$ are the charges of the matter fields under the $`U(1)`$ gauge group of the D7 branes, with $`q_{aa}=0,q_{ab}=q_{ba},q_{ab}=+1,0,1`$; and $`C_I=\varphi _{aa},\varphi _{ab}`$ are all of the open string moduli. The FI terms appear always when there are supersymmetry breaking coming from the 2-form magnetic fluxes. In standard type IIB orientifolded actions the FI terms can be properly tuned by adjusting the twisted moduli. Physically this term reproduces at leading order the splitting between scalar and fermion masses . Namely in the $`ab`$ sector, chiral fermions remain massless at tree level while their scalar partners obtain a mass proportional to $`\delta \nu _{ab}=\frac{1}{\pi }_i(\mathrm{\Psi }_a^i\mathrm{\Psi }_b^i)`$.
An important remark regarding supersymmetry is the following: In it has been pointed out that since some of these scalar fields can acquire vevs, the existence of a FI terms by itself does not automatically imply a susy breaking since they may acquire them so as to make the D-term vanish. Physically it is due to a recombination process of the D6 branes -in type IIA picture- (which now are not supersymmetric since the angles have changed) into supersymmetric smooth 3-cycles. This process gives a vev to some of the scalar fields $`\varphi _{ab}`$, at the intersection.
In the type IIB picture, the argument remains valid. D9 branes have magnetic fields at supersymmetric values. In the presence of non-vanishing FI terms the D9 branes become non susy, but if some scalars acquire mass cancelling D-term contribution, then they change their magnetic values on the 2-cycles restoring supersymmetry.
It has also been argued that although the existence of a FI-term in local compactifications allows for supersymmetry breaking, for the case of global compactifications recombination processes seem to be a generic mechanism to restore supersymmetry.
As already mentioned in the introduction, in a mechanism was proposed to stabilize Kähler moduli through three-form fluxes and magnetic fluxes. The argument was that the Kähler moduli get fixed at its supersymmetric value since the supersymmetry condition for magnetized D-branes depends explicitly on them. However, because of the argument above, this is not a true stabilization since matter fields can acquire a vev without any energy cost to cancel the D-term contribution, modifying the value of the twisted moduli. Then the scalar potential has a flat direction with $`V=0`$.
## 5 A new mechanism of moduli stabilization
We are going to consider a decoupling approach where complex structure moduli and the dilaton have been stabilized by the standard mechanism. Turning on suitable three-form fluxes stabilizes their vevs and allows their dynamics to be integrated out as in . This approach is valid in the regime when the mass scale of the Kähler moduli is much less than that of the dilaton and the complex structure moduli. The warping induced by the fluxes will be negligible in the approximation of large volume. We propose the following mechanism to stabilize the Kähler moduli perturbatively: We consider the coupling of fluxes (3-form fluxes and magnetic ones) to the open string sector that induce bosonic flux induced masses. These mass terms in principle can be soft breaking terms or supersymmetric. Supersymmetric mass terms are the $`\mu `$ terms appearing in the superpotential when supersymmetric 3-form fluxes (2,1) are turned on and the D-brane stacks share at least one parallel direction . These flux induced masses $`m_I^2`$ for magnetized D-branes combined with the FI contribution represent a new coupling in the scalar potential that lift the flat directions of the potential giving masses to all the moduli. This fact will happen generically when soft masses are present but as remarked above, for supersymmetric masses to exist on susy flux-backgrounds, only very specific configurations will be allowed. The effective scalar potential in $`D=4`$ taking into account the F-piece and the D-piece of the potential is equal to
$`V_T`$ $`=`$ $`V_F^{back}+\text{flux induced mass terms}+FI=`$ (5.1)
$`=`$ $`V_F^{back}+{\displaystyle \underset{I}{}}m_I^2|C_I|^2+{\displaystyle \underset{I}{}}(q_I|C_I|^2{\displaystyle \underset{i}{}}\delta \mathrm{\Psi }_I^i)^2,`$ (5.2)
where $`V_F^{back}`$ is the background no-scale potential induced by the superpotential. Flux induced mass terms have the effect of lifting the flat directions in the potential. These terms can be positive or negative. For the case of ISD $`G_3`$ the mass terms are positive and $`V`$ is also positive definite. We will see that there exists a global minimum of the potential $`V`$ that fixes all of the moduli. In order to minimize this potential we should consider the minimization with respect to every moduli $`C_I,M`$. It is important to ask whether the resulting critical point is a saddle point or a minimum. The answer is that iff there exist a critical point such that $`V_T=0`$, being $`V_T`$ definite positive, then this point is a global minimum of the theory and in this case it corresponds to its supersymmetric value. (The supersymmetric minimum in fact corresponding to F-flatness and D-flatness condition) since it is bounded from below to a value greater or equal than $`0`$. We find the value of this minimum at $`|C_I|=0,_i\delta \mathrm{\Psi }_I^i=0`$.
Let see it in more detail. To see the stabilization we need to minimize the potential with respect to the full moduli $`C_I,M`$, since the flat directions were associated in to the presence of matter fields which did not acquire a vev. As explained before, the non-vanishing of the FI term represents a shifting from the supersymmetric condition associated to the D-brane configuration $`_i\delta \mathrm{\Psi }_I^i0`$. As we are interested in stabilizing Kähler moduli on type IIB, then we are going to impose that the flux induced terms have to be associated to the presence not only of 3-form fluxes but also of magnetic fields which depend explicitly on the Kähler moduli. The presence of magnetic fluxes gives an extra dependence of the intersecting angles $`_i\delta \mathrm{\Psi }_I^i`$ on the Kähler moduli.
The superpotential is not renormalized at any order in perturbation theory and receives no $`\alpha ^{^{}}`$ corrections from the bulk, however it could receive from brane contributions . The D-term, by an appropriate selection of fluxes, can leave the model supersymmetric and so will not receive corrections changing its structure. Moreover, if we analyse the equation we can see that this holds generically given $`m_I^2>0`$ (as we will see this is the case for the interesting ISD 3-form fluxes).
We have in principle two kinds of flux contributions, ISD fluxes and IASD fluxes. We do not consider interaction terms between them. Since IASD fluxes do not solve type IIB 10D equations of motion unless non-perturbative effects will be taken into account, we will not consider them in the analysis. In the presence of ISD fluxes only, as pointed out in , soft terms are positive. The argument relies in the following : they can be regarded as geometric moduli of the F/M-theory fourfold and generate positive definite scalar potential. In it is argued that this is a general property of ISD fluxes also valid in the case of magnetic fluxes at the intersections <sup>2</sup><sup>2</sup>2We thank L. Ibanez for helpful clarifying explanation at this point. A way of illustrating this point is that for the case of soft breaking terms with magnetic fluxes, the superpotential
$`W={\displaystyle _{CY_4}}G_4\mathrm{\Omega }_4`$ (5.3)
leads to a positive definite scalar potential. $`W`$ in particular includes D7-brane geometric moduli , which in this case are described in terms of the homological charges also coming from the magnetized D-branes.
For the case of ISD fluxes, the $`V_F^{TT}=0`$ since it is a no-scale potential ($`D_TW\overline{D_T}W=3\left|W\right|^2`$), then extremizing with respect to the moduli $`T_N`$ gives the equation
$`_NV=0{\displaystyle \underset{I}{}}(2q_I|C_I|^2_N\mathrm{\Psi }_I+2(\delta \mathrm{\Psi }_I)_N\mathrm{\Psi }_i+(_Nm_I^2)|C_I|^2)=0.`$ (5.4)
Minimizing with respect matter field gives as solution
$`|C_I|`$ $`=`$ $`0,V(C_I)=|\delta \mathrm{\Psi }_I|^2`$ (5.5)
$`|C_I|`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{2q_I\delta \mathrm{\Psi }_Im_I^2}{2q_I^2}}},V_D(|C_I^2|_{min})={\displaystyle \frac{m_I^2}{4q_I^2}}(4q_I(\delta \mathrm{\Psi }_I)m_I^2).`$ (5.6)
The first solution imposes in this model that the global minimum lies at $`\delta \mathrm{\Psi }_I=0`$, fixing the moduli through the dependence of $`\mathrm{\Psi }_I`$ to its supersymmetric value iff $`\text{soft terms}>0`$ for all $`I`$. The value of the minimum of the potential corresponds for supersymmetric configuration of D-branes, to a no-scale potential ($`V=0`$), and in those cases when supersymmetry can be consistently broken through the FI term to a de Sitter minimum, in the same spirit as . The other two extrema lead to an AdS vacua and they are only possible in those cases when $`\left|C_I\right|`$ is real, that correspond to have negative squared mass terms for the supersymmetric case, $`m_I^2<0`$, so it is not possible for ISD $`G_3`$ fluxes or to have $`0<m_I^2<2_Iq_I\delta \mathrm{\Psi }_I`$ for a non trivial FI term. Assuming that $`_Iq_I\delta \mathrm{\Psi }_I>0`$ one can see that $`m_I^22_Iq_I\delta \mathrm{\Psi }_I`$ always. The bound is never saturated unless $`\nu =1/2,u_i=1,s10^3`$ <sup>3</sup><sup>3</sup>3I would like to thank G. Tasinato for his comments regarding this point. and these values do not allow to have consistent flux compactifications on this background compatible with tadpole cancellation conditions, so these other two possibilities are never achieved nor by making fine tuning of fluxes. The unique minimum is at $`C_I=0`$.
To be sure that critical points are not saddle points if we were interested in these cases associated to IASD fluxes, ( that we are not), we should perform the Hessian calculation -very involved due to the higly non-trivial structure of the $`V_T`$\- , including also the F-part of the potential that is no longer vanishing.
One could also ask whether trilinear couplings of the soft terms could change this behaviour. The answer is not for the particular case we are considering, i.e. ISD fluxes. let see it in more detail. The full structure of the soft and susy terms is rather complicated. Generically,
$`\text{soft terms}={\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}(M_a\lambda ^a\lambda ^a+h.c.)+{\displaystyle \underset{I}{}}m_I^2|C_I|^2+`$ (5.7)
$`{\displaystyle \underset{I,J,K}{}}A_{IJK}C_IC_JC_K+1/2B_{IJ}C_IC_J+h.c.`$ (5.8)
where $`\lambda ^a`$ represents gaugino mass terms, $`m_I`$ are the scalar mass terms and there are bilinear and trilinear couplings each of which with a highly nontrivial dependence on all of the moduli. However one can see that all of the induced mass that appear in the scalar potential are polynomial of lower bound $`2`$ in the matter fields $`C_I`$, hence, as before $`C_I=0`$ is still an extrema, and as before $`\delta _I\mathrm{\Psi }_i=0`$ corresponds to $`V_T=0`$ which is the global minimum of the theory. Moreover, the expression for the trilinear terms in this case corresponding to $`W=W_{flux}`$ is the following,
$`A_{IJK}=F^M\widehat{K_M}_Mlog(\stackrel{~}{K_{II}}\stackrel{~}{K_{JJ}}\stackrel{~}{K_{KK}})`$ (5.9)
trilinear terms can be estimated as $`A_{IJK}1/2(m_{3/2}_i1/\nu _i)1/2m_{3/2}`$ and scalar mass terms $`m_I^2m_{3/2}^2(1/4+3/4(1/_i1/\nu _i+_i\frac{1}{\nu _i^2}))1/4m_{3/2}^2`$. As before one can ask if there exists an appropiate fine tuning in the configuration of D-branes and in the choice of 3-form fluxes in such a way that for a particular configuration of brane case the other two minima will be $`C_I=0`$ with a $`FI0`$, giving $`V_T=0`$ and the answer is that again it seems not possible. Both bounds are saturated for $`\nu _i=1/2`$ which corresponds to the case cosidered above that leads to non physical solution in these set-ups. Scalar soft terms are the dominant contribution. An important clue in this result is the fact that the whole mass terms are positive defined for ISD $`G_3`$. Otherwise there could also other solutions corresponding to have $`V_T<0`$, and in those cases a careful analysis should be performed.
In the absence of magnetic fluxes moduli can not be fixed since $`\mathrm{\Psi }_I`$ is independent of them. However as already shown in the preceding section the presence of magnetic fluxes does not by itself imply the stabilization of Kähler moduli, as we already explained, as matter fields are free to acquire any vev to cancel D-term contribution leaving Kähler moduli unfixed. We want to emphasize that only when there are magnetic fluxes in the 3-form flux induced soft (susy) <sup>4</sup><sup>4</sup>4 This mechanism is also valid for supersymmetric masses generated by similar mechanisms whenever appropiate configurations are considered. breaking terms combined with the FI term, the Kähler moduli get fixed. This contribution coupling to the FI-term, prevents matter fields from acquiring a vev cancelling this contribution as it is energetically disfavoured.
This is then a generic mechanism, in the same way that 3-form fluxes $`H_3,F_3`$ serve to stabilize the complex structure moduli $`(2,1)`$ and the dilaton through the superpotential. The magnetic fluxes given by a particular configurations of the magnetized D-branes fix the Kähler moduli (1,1) at their supersymmetric values, only once the potential has no flat directions. These directions are lifted by the combined action of these 3-form flux-induced soft terms with magnetic fields, and the FI-terms. We think that this solution gives a final answer to the question of perturbative stabilization of moduli. Both types of moduli need to be present to achieve a model without moduli. The scale of stabilization of the moduli vevs is at supersymmetric values. The mass scale of the moduli however depends on the model considered, and is given by the scale of the flux induced terms.
All of the previous discussion in principle remains valid for the case of susy flux induced terms, i.e. mass terms with the same masses for fermions and their scalar superpartners (it has been obtained in D-brane action calculations, see for example, , and F-theory .). These susy mass terms, however do not appear always that there are three form fluxes ISD (2,1) , it is needed a D-brane configuration with at least two chiral multiplets such that they are at least parallel in one direction to generate $`\mu `$-terms. To have a true $`N=1`$ model it is also needed masses for all of the moduli present in the configuration, allowing to obtain at the same time non-trivial areas. If such a model is provided then it constitutes a model with Kähler moduli stabilized. This implies that for a supersymmetric model, whenever the induced flux mass terms are all positive, which is the case for ISD fluxes, the matter fields are going to be fixed at zero. Spontaneous susy breaking will lift this value to the lowest metastable de Sitter vacua once the Kähler moduli have been stabilized by the supersymmetric condition. In both cases ( non-susy and susy) this result is very appealing since it is valid for the ISD fluxes which are the ones we are interested in since they solve the 10D equations of motion in type IIB theory. From the calculational point of view, they stabilize the Kähler moduli for a given Dp-brane content without the need to specify a detailed expression for the highly complicated soft terms. We provide an explicit example in the last section of the paper.
### 5.1 Beyond imaginary self-duality condition
In general for IASD fluxes or a combination of both, the contribution of $`V_F`$ has to be taken into account but this does not change the behaviour of the stabilization. The two possible extra extrema are
$`|C_I|^2={\displaystyle \frac{2q_I\delta \mathrm{\Psi }_Im_I^2}{2q_I^2}}`$ (5.10)
For a susy model this reduces to:
$`V(|C_I|_{min}^2)`$ $`=`$ $`{\displaystyle \frac{|m_I|^4}{4q_I^2}}<0iffm_I^2<0\text{I}.`$ (5.11)
The minima of the potential correspond to a AdS vacua. For the case of negative bosonic soft terms only very rough approximations to the stabilized values can be done because of the complicated equation to solve. Moreover the Kähler potential cannot be simplified so much since the bilinear $`Z_{ij}`$ has to be calculated. This case can appear when IASD fluxes are taken into account. The explicit value at which Kähler moduli get fixed now it is going to be determined by the explicit expression of the flux induced mass terms, that one for the soft terms ( see section 4) is highly nonlinear. On the other hand IASD fluxes do not solve the equations of motion in this description since it is purely done at perturbative level ( the situation changes if non-perturbative mechanisms are taken into account), and they generate poorly understood run-away potentials. We will not perform the calculation.
## 6 An example of IIB on $`\frac{T^6}{Z_2\times Z_2\times \mathrm{\Omega }R}`$
In this section we provide a concrete example of phenomenological interest which realizes the new mechanism proposed in the preceding section. Some examples of flux compactifications of phenomenological interest are . See also , in which some of the flux compactifications of phenomenological interest are able to lead to KKLT models.
In the following we will review the main features of the constructions of which are a global embedding of the local model proposed in and used in . We will then construct our explicit realization.
The model of is based on type IIB string theory compactified on a $`\frac{T^6}{Z_2\times Z_2}`$ modded out by the orientifold action, which has also been examined in. We consider $`T^6=\mathrm{\Pi }_{i=1}^3T_i^2`$. The generators of the orbifold symmetries $`Z_2\times Z_2`$ are $`\theta ,\omega `$ which act on the complex coordinates of the tori as
$`\theta :(z_1,z_2,z_3)(z_1,z_2,z_3),`$ (6.1)
$`\omega :(z_1,z_2,z_3)(z_1,z_2,z_3).`$ (6.2)
The orientifold action is given by $`\mathrm{\Omega }R`$ where $`\mathrm{\Omega }`$ is the usual world-sheet parity and
$`R:(z_1,z_2,z_3)(z_1,z_2,z_3).`$ This model contains 64 $`O3`$ planes each one on a fixed point of $``$ and 4 $`O7_i`$ planes, located at the $`Z_2`$ fixed points of the i-th $`T^2`$ and wrapping the other tori. We consider the case of intrinsic torsion as in . The open string sector contains D9 branes with non trivial magnetic fluxes. The non-trivial gauge bundle generically reduces the rank of the group and upon KK reduction leads to D=4 chiral fermions. The magnetic fluxes also induce D-brane charges of lower dimension that contribute to the tadpoles. Magnetized D9 branes usually have D7, D5 and D3 charges. As explained in the topological information of these models is encoded in three numbers. $`(N_a,(n_a^i,m_a^i))`$ where $`N_a`$ is the number of D9 branes contained on the $`a`$-stack, $`m_a^i`$ is the number of times that a-stack of D-branes wrap the i-th $`T^2`$ and $`n_a^i`$ is the units of magnetic flux in that torus induced by D-branes. The unit of magnetic flux of the D-branes in that torus, as seen in section 3, is
$`{\displaystyle \frac{m_a^i}{2\pi }}{\displaystyle _{T_i^2}}F_a^i=n_a^i.`$ (6.3)
The $`D9_a`$ branes preserve the same supersymmetry of the orientifold planes provided that
$`{\displaystyle \underset{i=1}{\overset{3}{}}}\mathrm{\Psi }_a^i`$ $`=`$ $`{\displaystyle \frac{3\pi }{2}}mod(2\pi ),`$ (6.4)
$`\text{with}\pi \mathrm{\Psi }_a^i`$ $`=`$ $`arctan({\displaystyle \frac{n_a^i\beta _i}{m_a^i}}),`$ (6.5)
where no summation over index $`i`$ is performed. This condition also guarantees that any two sets of branes preserve a common supersymmetry since the relative angles
$`\theta _{ab}^i=\mathrm{\Psi }_b^i\mathrm{\Psi }_a^i`$ (6.6)
trivially satisfy
$`{\displaystyle \underset{i}{}}\theta _{ab}^i=0mod(2\pi ).`$ (6.7)
In the case of a global analysis we need to add new branes which are expected to be in a hidden sector in order to cancel global tadpoles . The conditions are,
$`1)`$ $`{\displaystyle \underset{\alpha }{}}N_\alpha n_\alpha ^1n_\alpha ^2n_\alpha ^3+1/2N_{fluxes}=16N_{min}=8\text{with torsion}`$ (6.8)
$`2)`$ $`{\displaystyle \underset{\alpha }{}}N_\alpha m_\alpha ^in_\alpha ^jm_\alpha ^k=16ijk\text{and}i,j,k=1,2,3`$ (6.9)
with $`N_{flux}=64.nnZ`$. Global cancellation of $`𝐙_\mathrm{𝟐}`$ RR charges must also be imposed to cancel the contribution of $`D5_i\overline{D5_i}`$ and $`D9_i\overline{D9_i}`$ pairs, which is equivalent to satisfying the following conditions for a case with torsion,
$`{\displaystyle \underset{\alpha }{}}N_\alpha m_\alpha ^1m_\alpha ^2m_\alpha ^3\mathrm{𝟒}𝐙`$ (6.10)
$`{\displaystyle \underset{\alpha }{}}N_\alpha n_\alpha ^in_\alpha ^jm_\alpha ^k\mathrm{𝟒}𝐙ijk\text{and}i,j,k=1,2,3`$ (6.11)
and for the case without torsion the conditions are the same by making these changes,
$`{\displaystyle \underset{\alpha }{}}N_\alpha n_\alpha ^1n_\alpha ^2n_\alpha ^3+1/2N_{fluxes}=16N_{min}=4`$ (6.12)
$`{\displaystyle \underset{\alpha }{}}N_\alpha m_\alpha ^1m_\alpha ^2m_\alpha ^3\mathrm{𝟖}𝐙`$ (6.13)
$`{\displaystyle \underset{\alpha }{}}N_\alpha n_\alpha ^in_\alpha ^jm_\alpha ^k\mathrm{𝟖}𝐙ijk\text{and}i,j,k=1,2,3`$ (6.14)
Let us focus on the case with torsion. The following model is a concrete realization that serves our purpose.
#### MSSM
(6.24)
This includes a slight modification to the proposal of since their configuration was not able to fix all of the Kähler moduli. In fact the spectrum keeps its interesting phenomenological properties except for the number of families which depend on the particular homological charges chosen. To satisfy the consistency conditions the following equation has to be solved,
$`4n+g^2+N_f=8.`$ (6.25)
It gives different flux vacua for the different possible replication of families $`I_{ab}=\mathrm{\Pi }_i(n_a^im_b^in_b^im_a^i)`$. Searching for a realistic scenario it is necessary to obtain three family replication, unfortunately a simple examination reveals that this model does not contain it. The possible vacua are:
For $`g=2`$
$`n`$ $`=`$ $`0N_f=4`$ (6.26)
$`n`$ $`=`$ $`1N_f=0`$ (6.27)
For $`g=1`$
$`n`$ $`=`$ $`0N_f=7`$ (6.28)
$`n`$ $`=`$ $`1N_f=3`$ (6.29)
To illustrate our example we choose the vacua for $`g=1`$, containing fluxes, i.e. $`n=1,N_f=3`$.
#### Spectrum
Locally the spectrum of this model can contain the MSSM since $`g`$ is not constrained. This has been studied in . This sector is called the visible sector. We show it below as a remainder,
(6.39)
Table 1: Chiral spectrum of the MSSM-like model.
However the global completion imposes constraints in such a way that $`g3`$ so it is not possible to obtain 3-family replication and in that sense, we are proposing is a toy model although it keeps the nice properties of chiral matter, correct quantum numbers, etc.. It contains the intersection of the so-called, visible sector and the hidden sector and hidden-hidden sector. We analyze the full spectrum in terms of a Pati-Salam model.
(6.58)
Table: Chiral spectrum of a one generation Pati-Salam $`N=1`$ chiral model of table 1. The abelian generator of the unique massless U(1) is given by $`Q^{^{}}=\frac{1}{3}Q_a2(Q_{h_1}Q_{h_2})`$.
Some linear combinations of U(1) will leave $`U(1)`$ fields massive in a Green-Schwarz mechanism but since our purpose is to show the perturbative stabilization of moduli we have not included them in the spectrum calculation. We show that a realization of string compactification on a $`CY_3`$ orientifold with some phenomenological properties render the moduli fixed.
When we substitute for the supersymmetry conditions,by taking the values of table $`(6.15)`$ in $`(6.4)`$, we obtain the following equations
$`\delta _2=\delta _3t_2=t_3,`$ (6.59)
$`\pi (\phi _1+\phi _2+\phi _3)={\displaystyle \frac{3\pi }{2}},`$ (6.60)
$`\pi (\varphi _2+\varphi _3)={\displaystyle \frac{\pi }{2}}.`$ (6.61)
Using $`(6.5)`$ a straightforward calculation gives
$`\beta _1=0.157,\beta _2=\beta _3={\displaystyle \frac{1}{\sqrt{20}}}`$ (6.62)
which means that the vacuum expectation value at which Kähler moduli gets fixed is,
$`t_1=20s,t_2=t_3=127.4s,`$ (6.63)
where $`s`$ denotes the vacuum expectation value of the dilaton which is the string coupling constant. Substituting in the value of areas ($`3.6`$) this means that
(6.66)
This mechanism implies that for this toroidal compactification, are needed three and only three stacks of D-branes with different magnetic fluxes in such a way that lead to different equations in order to have a determined compatible system of equations. This is the subtlety that does not allow us to use Marchesano et al.’s model in our mechanism as an example of D-brane configuration since they have just two different equations. This construction is in no way unique, and represents a toy model to illustrate the mechanism of stabilization. We expect that more complicated models can be constructed with realistic spectrum and with all of the moduli fixed.
Regarding the blow-up moduli associated to the orbifold fixed points, they are going to be also stabilized with this mechanism since for models with intrinsic torsion, the moduli are Kähler, and can get fixed in the same way as it was done in . With the above results we can perform the explicit calculation of soft mass terms for a given three-form flux configuration with complex-structure moduli and dilaton stabilized. We choose the one for $`n=1,N_{flux}=64`$ with $`𝒩=0`$ supersymmetry
$`G_3`$ $`=`$ $`2(d\overline{z}_1dz_2dz_3+dz_1d\overline{z}_2dz_3+dz_1dz_2d\overline{z}_3+d\overline{z}_1d\overline{z}_2d\overline{z}_3),`$ (6.67)
$`W`$ $`=`$ $`8(u_1u_2u_3s).`$ (6.68)
The complex-structure moduli and the dilaton are stabilized at $`u_i=s=i`$, which in our notation is equivalent to $`u_i=s=1`$. This values lead to a non-acceptable value for perturbative analysis since $`g_s=1`$ as explained in . However we will use this to give an example of explicit calculation of the soft terms. The soft terms from T-dominance (ISD) fluxes correspond to $`V_0=0`$. The gravitino vev is
$`m_{3/2}^2={\displaystyle \frac{|W|^2}{s\mathrm{\Pi }_it_iu_i}}`$ (6.69)
and we have made $`C=1,cos(\theta )=1,\eta _i=1/\sqrt{3},\gamma _i=\gamma _T`$ as is explained in detailed in .
The soft breaking terms for this model are:
(6.80)
in terms of gravitino mass $`m_{3/2}`$. Since the microscopic source of SUSY-breaking is the above ISD flux in a toroidal setting, by using the above definitions, $`|W|^2=256`$, and the gravitino mass is $`m_{3/2}=0.019`$, so the soft terms have unrealistic values as they are extremely high.
For general toroidal/orbifold models with intersecting D6-branes the flux-induced soft terms are typically of the order of the string scale $`(1/\alpha ^{^{}})`$, which is only slightly smaller than $`M_{Pl}`$ and not able to solve hierarchy problems . This fact is due to the simplicity of the compactification manifolds as well as the fact that fluxes are distributed uniformly.
## 7 Discussion and conclusions
We have shown a new method to dynamically stabilize, with fluxes, all moduli in supersymmetric and non supersymmetric models. In supersymmetric models we have explained that the mechanism is restricted to particular configurations able to generate $`\mu `$ terms for all of the moduli, i.e. configurations of stacks parallel to each other in at least one direction in which a suitable $`(2,1)`$ $`G_3`$ flux has been turned on. Three-form fluxes generically stabilize the dilaton and the complex structure moduli. They generate an F-term part of the scalar potential. Magnetic fluxes are two forms that fix Kähler moduli at their supersymmetric value once the potential has no flat directions. This goal is achieved by inducing through 3-form flux a non-susy (susy) breaking (flux-induced) mass terms with magnetic fields in the scalar potential, that combined with the FI term lift the flat directions. We think this is a generic mechanism that gives a final answer to the problem of perturbative moduli stabilization. These mass terms avoid the possibility of cancelling the D-term by the consistent adjustment of the matter fields vev, since it requires an energy cost. To leading order, the F-part of the potential is a non-scale potential and together with D-term are responsible for the stabilization of the Kähler moduli. ISD fluxes induce positive mass terms and fix the Kähler value of the D-term to its supersymmetric value. In the supersymmetric case it can be possible to implement a similar mechanism of the one found in to generate a de Sitter space (by spontaneous supersymmetry breaking ). However, possibly toroidal models will be unable to generate adequate fine tuning to guarantee that the different approximations are still valid (small cosmological constant, small quantum corrections, a big potential barrier). Maybe this mechanism combined with the one of could be implemented in a more complicated model (it is necessary that $`U(1)`$’s of the FI term will not be charged or the matter will have some special properties that are not present in the case considered). For the case $`m^2<0`$ both parts of the potential, F- and D-, including the supergravity potential $`V_F^{background}`$ contribute and an explicit calculation in terms of the particular soft breaking terms is needed. They can be associated to IASD or a combination of (IASD and ISD) fluxes. The analysis of IASD contribution has not been performed as they do not lead to solutions of physical interest. The scalar potential generically has two AdS minima.
Clearly the case with $`m^2>0`$ which correspond to turning on ISD three form fluxes is much more interesting since it solves the equation of motion and stabilizes the moduli at a value of the order of the string scale which means that it does not depend on the scale of supersymmetry breaking and is higher enough to induce reheating processes at early stages of the universes. The mass scale of the moduli presumably is of the order of the soft breaking mass scale (which is of the order of flux scale $`\frac{\alpha ^{^{}}}{R^3}`$). We have given a concrete Kähler moduli-free realization of phenomenological interest of Type IIB on $`\frac{T^6}{Z_2\times Z_2\times \mathrm{\Omega }R}`$. However this model represents a toy model. We expect improvements in the search for realistic compactifications moduli-free (i.e. three family generation, lower soft masses) in more complicated scenarios as those with warped metrics due to throats, that are also able to explain the hierarchy problem, as well as, address inflation.
## 8 Acknowledgements
I would like to thank to the referee for his/her comments that have helped me very much to improve the paper. I want to thank J.L.F.Barbón, D. Cremades, J. Conlon, E. López, L. Ibañez, F.Quevedo and G. Tasinato for useful conversations. I am also very grateful to A. Font for many comments and clarifying explanations about her work. Finally, I am specially indebted to Angel Uranga for his continuous orientation and discussions all throughout this work. M.P.G.M. is supported by a postdoctoral grant of the Consejería de Educación, Cultura, Juventud y Deportes de la Comunidad Autónoma de La Rioja (Spain).
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# Towards a global theory for the high 𝑻_𝒄 cuprates: Explanation of the puzzling optical properties
## I Introduction
The unusual physics of the cuprates, and specifically the occurrence of high-$`T_c`$ superconductivity (SC) in them Muller1 , continues to be one of the forefront problems in physics. Even though numerous mechanisms have been proposed to explain this puzzling system (e.g., Refs. Brusov ; Anderson ; Laughlin1 ; Varma ; Lee ; Emery ; Kresin ; Laughlin ; Varma1 ; Annette ; Mihailovic1 ; Fine ), a consensus has not been reached yet. The cuprates, as well as related systems of interest, are characterized by such a complexity that it may make it unrealistic, at present, to predict the different aspects of their behavior by a precisely solvable model or numerical calculations. Their complexity involves single-particle Pickett ; Andersen ; Bansil , many-body Macridin ; Anisimov1 ; Munoz , and lattice effects Egami ; Lanzara1 , as well as the occurrence of nanoscale inhomogeneity Bian ; Tran1 ; Muller2 ; Gorkov ; Kapitul ; Davis1 ; Davis2 .
Even though approximate numerical calculations, and solutions of simple models, have been helpful Macridin ; Anisimov1 ; Munoz ; Zaanen1 ; Machida ; Emery1 ; Castro ; Maki ; Dagotto ; Scalapino ; Auerbach ; Norman ; Hanke ; Carbotte1 ; Lichten ; Anisimov2 ; Ovchi ; Neto ; Eremin2 ; Eremin1 ; Mihailovic2 ; Markiew , an approach incorporating results of different schemes may be necessary for the global understanding of the cuprates. In spite that a rigorous analytical or numerical solution of the incorporated scheme, in its globality, is beyond reach at present, it provides an insight to the understanding of the physics of the cuprates, which is missing when models which describe the system partially are applied.
Such an approach, including a minimal framework within which different aspects of the physics of the cuprates could be described correctly, has been worked out by the author Ashk01 ; Ashk03 ; Ashk04 . Within this emerging ‘global theory’ of the cuprates (GTC) their puzzling physics, including the occurrence of high-$`T_c`$ SC, can be understood as the result of a behavior typical of their structure within the regime of a Mott transition Lee . A solution has been studied, corresponding to a state where translational symmetry within the CuO<sub>2</sub> planes is disturbed by dynamical stripe-like inhomogeneities. Signatures of such inhomogeneities have been observed, e.g., in Refs. Tran1 ; Kapitul . They accommodate the competing effects of hopping and antiferromagnetic (AF) exchange on large-$`U`$ electrons, and their existence had been predicted theoretically Zaanen1 ; Machida ; Emery1 ; Castro .
These inhomogeneities provide quasi-one-dimensional segments in which the behavior of the electrons is described in terms of separate carriers of charge and spin Anderson . Since these carriers are strongly coupled to electrons in the regions where spin and charge are inseparable, a two-component scenario is obtained of heavy and light charge carriers, coupled through spin carriers. The speed of the dynamics of the stripe-like inhomogeneities is self-consistently determined by the width of spin excitations around the AF wave vector $`𝐐=(\frac{\pi }{a},\frac{\pi }{a})`$. Small-width excitations, and slow dynamics, are obtained in the SC state, and to some degree in the pseudogap (PG) state Ashk03 ; Ashk04 .
This theory was found to provide a highly-plausible explanation to all the anomalous properties of the cuprates that were studied by it. This includes the systematic anomalous behavior of the resistivity, Hall constant, and thermoelectric power (TEP) Ashk01 , and of spectroscopic anomalies Ashk03 ; Ashk04 . Also were explained the low- and high-energy spin excitations around $`𝐐`$, including the neutron-resonance mode Ashk03 ; Ashk04 , and its connection to the peak-dip-hump structure observed in tunneling and ARPES Ashk04 . Pairing was shown Ashk03 to be induced by the energy gain due to the hopping of pair states perpendicular to the stripe-like inhomogeneities. Pairing symmetry was predicted to be of the $`d_{x^2y^2}`$ type, but to include features not characteristic of this symmetry Ashk03 .
The phase diagram of the cuprates was found Ashk04 to result from an interplay between pairing and coherence within the regime of the Mott transition. Both pairing and coherence are necessary for SC to occur. Coherence without pairing results in a metallic Fermi-liquid (FL) state. Incoherent pairing results in the PG state, consisting of electrons and localized electron pairs. At low temperatures ($`T`$), the stripe-like inhomogeneities partially freeze in the PG state into a glassy “checkerboard” structure Davis1 , and the electrons become localized. This results in the opening of localization minigaps on the Fermi surface (FS); they contribute to the PG, and their size determines the lower doping limit of the SC phase. If SC is suppressed, the borderline between the FL and the PG states persists down to $`T=0`$, where a metal-insulator-transition (MIT) quantum critical point (QCP) occurs Boeb ; Tallon .
A nanoscale heterogeneity was predicted Ashk04 , especially in the underdoped (UD) regime, consisting of “perfectly” SC regions which become completely paired for $`T0`$, and “PG-like” SC regions, where pairing remains partial as in the PG state. A distribution of such regions of different degrees of pairing has been observed by STM Davis2 . This heterogeneity is expected to be intrinsic, and set in below $`T_c`$ even in very pure samples. Its scale is larger than that of the dynamical stripe-like inhomogeneities, which are also intrinsic and essential for high-$`T_c`$ SC Ashk03 . Injection of pairs, similar to the one occurring in p–n junctions in semiconductors, is expected in junctions including slices of a cuprate in the SC state and in the PG state, resulting in the observed “giant proximity effect” (GPE) Bozovic .
Similarly to transport, optical properties detect the electrons within the crystal, with no transfer of electrons into, or out of it (as occurs e.g. in tunneling or ARPES). Thus, even though their theoretical evaluation involves an integration over the Brillouin zone (BZ), they still could be very sensitive to fine many-body effects. The relevance of optical results to the theoretical predictions of Refs. Ashk01 ; Ashk03 ; Ashk04 has been mentioned there just in passing. Here it is demonstrated how this GTC naturally explains a variety of optical results (whichever were tested by it), detailed below, part of which have been lacking a satisfactory understanding so far.
The partial Glover–Ferrell–Tinkham sum rule (f–sum rule) Kubo1 ; Tinkham ; Norman ; Hanke ; Marel3 , over the conduction band, is studied. The GTC is applied to understand observed “violations” of the sum rule through $`T_c`$, due to the transfer of spectral weight from energies $`\mathrm{}>2`$eV to the vicinity of the Fermi level ($`E__\mathrm{F}`$) Basov3 ; Marel1 ; Bontemps ; Homes2 . Also the optical signatures of $`c`$-axis collective modes Marel3 ; Marel2 ; Dulic1 ; Marel4 ; Marel5 ; Dordevic2 are understood.
The physical interpretation of the optical density to mass ($`n/m^{}`$) ratio, derived, e.g., through the f-sum rule is clarified. The GTC is shown to explain “Tanner’s law” Tanner1 , under which the above ratio is about 4–5 times the contribution to it from the Drude part of the optical conductivity $`\sigma `$, which happens to be just a little greater than the ratio $`n_s/m^{}`$ based on the low-$`T`$ superfluid density $`n_s`$. The approximate factor of four is connected to the periodicity within the stripe-like inhomogeneities. The increase in this factor in the heavily UD regime, as well as the occurrence of a constant effective mass of carriers through the transition to the AF regime Basov5 , are also understood.
The optical quantity $`\rho _s=4\pi \mathrm{e}^2n_s/m^{}`$ is shown here to be not identical with the quantity obtained through the relation $`\rho _s=(c/\lambda )^2`$ from measurements of the penetration-depth $`\lambda `$ by methods like $`\mu SR`$ (as has been observed Tajima ), due to a difference between the many-body effects on them. In the second case the GTC predicts Ashk01 ; Ashk03 ; Ashk04 a boomerang-type behavior in the overdoped (OD) regime, in agreement with experiment Niedermayer . But, as is shown here, no such behavior is expected for the optical $`\rho _s`$, also in agreement with experiment Timusk1 .
Within the GTC, the dynamical stripe-like inhomogeneities are intertwined with low-energy spin excitations around $`𝐐`$ (including the resonance mode), which contribute narrow peaks only in the PG and SC states. It is demonstrated here that their optical signatures have been observed Lupi1 ; Lupi2 ; Kim ; Carbotte2 ; Timusk in these states at energies on the edge of the gap, and within it.
This theory is also shown to explain the different optical signatures of the PG in the $`c`$-direction and in the $`ab`$-plane. Namely the observation of a depression, within the PG energy range, of the optical conductivity $`\sigma (\omega )`$, in the $`c`$-direction Homes3 ; Puchkov , but of the optical scattering rate $`1/\tau (\omega )`$ in the $`ab`$-plane Puchkov ; Basov .
The sharp drop in this scattering rate below $`T_c`$ Bonn ; Puchkov is also predicted by the GTC. The existence of the QCP in the phase diagram Ashk04 results in marginal-Fermi-liquid (MFL) Varma behavior of the scattering rate above $`T_c`$, close to the QCP Timusk2 , and to critical behavior of optical quantities Marel6 .
“Uemura’s law”, under which $`T_c\rho _s`$ in the UD regime Uemura , has been understood Ashk01 ; Ashk03 ; Ashk04 on the basis of the pairing phase “stiffness” Emery . Also the optically-derived “Homes’ law” Homes1 , under which $`\rho _s35\sigma (T_c)T_c`$, for cuprates in the entire doping ($`x`$) regime, both in the $`ab`$-plane, and in the $`c`$-direction, and also for low-$`T_c`$ SC’s in the dirty limit, is shown to be consistent with the GTC. The validity of this law in the cuprates is related to the existence of quantum criticality there Zaanen .
Even though Homes’ law is an approximate one (presented in a log–log scale), its coexistence with Uemura’s law in the UD regime implies that the DC conductivity at $`T_c`$ varies considerably less with $`x`$ (in this regime) than its variation at high temperatures Takagi . This behavior is predicted by the GTC as well, and also are understood apparent deviations from Uemura’s law in the heavily UD regime Zuev .
## II Scheme of the Theory Ashk01 ; Ashk03 ; Ashk04
### II.1 Large-$`U`$-limit formalism
Ab-initio calculations Andersen in the cuprates indicate that the electrons in the vicinity of $`E__\mathrm{F}`$ could be analyzed in terms of a band (corresponding dominantly to copper and oxygen orbitals within the CuO<sub>2</sub> planes) for which large-$`U`$-limit approximations are adequate, which is somewhat hybridized to other bands for which small-$`U`$-limit approximations may be suitable. A perturbation expansion in $`U`$ is inadequate for the large-$`U`$ orbitals, and they are treated by the auxiliary-particles approach Barnes . Thus, within the CuO<sub>2</sub> planes, a large-$`U`$ electron in site $`i`$ and spin $`\sigma `$ (which is assigned numbers $`\pm 1`$, corresponding to $``$ and $``$, respectively) is created by $`d_{i\sigma }^{}=e_i^{}s_{i,\sigma }^{}`$, if it is in the “upper-Hubbard-band”, and by $`d_{i\sigma }^{}=\sigma s_{i\sigma }^{}h_i^{}`$, if it is in a Zhang-Rice-type “lower-Hubbard-band”.
Here $`e_i^{}`$ and $`h_i^{}`$ are creation operators of “excessions” and “holons”, and $`s_{i\sigma }^{}`$ are creation operators of “spinons”. If either the excessions, or the holons are ignored, than the large-$`U`$ band could be treated within the $`t`$$`t^{}`$$`J`$ (or $`t`$$`t^{}`$$`t^{\prime \prime }`$$`J`$) model. The auxiliary-particle approach is applied here using the “slave fermion” method Barnes , within which the holons/excessions are fermions and the spinons are bosons. This method had been successful treating antiferromagnetic (AF) systems, and implies taking good account of the effect of AF correlations. For rigorous treatment, one should in principle impose in each site the constraint: $`e_i^{}e_i^{}+h_i^{}h_i^{}+_\sigma s_{i\sigma }^{}s_{i\sigma }^{}=1`$.
In order to treat this constraint properly, an auxiliary Hilbert space is introduced within which a chemical-potential-like Lagrange multiplier is used to impose the constraint on the average. But since the physical observables are projected into the physical space as combinations of Green’s functions, whose time evolution is determined by the Hamiltonian which obeys the constraint rigorously, it is expected that it would not be violated as long as justifiable approximations were used.
The dynamics of the stripe-like inhomogeneities is approached adiabatically, treating them statically with respect to the electrons dynamics. The striped structure Tran1 consists of narrow charged stripes forming antiphase domain walls between wider AF stripes. Since the spin-charge separation approximation (under which two-auxiliary-particle spinon–holon/excession Green’s functions are decoupled into single-auxiliary-particle Green’s functions) is valid in one-dimension, it is justified to assume the existence of effective spinless charge carriers within the narrow charged stripes, but not within the whole CuO<sub>2</sub> plane (as is assumed in RVB theory Anderson ).
### II.2 “Bare” auxiliary particles
The spinons are diagonalized by applying the Bogoliubov transformation for bosons Ashk94 : $`s_\sigma ^{}(𝐤)=\mathrm{cosh}(\xi _{\sigma 𝐤})\zeta _\sigma ^{}(𝐤)+\mathrm{sinh}(\xi _{\sigma 𝐤})\zeta _\sigma ^{}(𝐤).`$ Spinon states, created by $`\zeta _\sigma ^{}(𝐤)`$, have bare energies $`ϵ^\zeta (𝐤)`$ with a V-shape zero minimum at $`𝐤=𝐤_0`$. Bose condensation results in an AF order of wave vector $`𝐐=2𝐤_0`$. Within the lattice BZ there are four inequivalent possibilities for $`𝐤_0`$: $`\pm (\frac{\pi }{2a},\frac{\pi }{2a})`$ and $`\pm (\frac{\pi }{2a},\frac{\pi }{2a})`$, thus introducing a broken symmetry. One has Ashk94 :
$`\mathrm{cosh}(\xi _𝐤)`$ $``$ $`\{\begin{array}{cc}+\mathrm{},\hfill & \text{for }𝐤𝐤_0\text{,}\hfill \\ 1,\hfill & \text{for }𝐤\text{ far from }𝐤_0\text{,}\hfill \end{array}`$
$`\mathrm{sinh}(\xi _𝐤)`$ $``$ $`\{\begin{array}{cc}\mathrm{cosh}(\xi _𝐤),\hfill & \text{for }𝐤𝐤_0\text{,}\hfill \\ 0,\hfill & \text{for }𝐤\text{ far from }𝐤_0\text{.}\hfill \end{array}`$ (1)
Holons (excessions) within the charged stripes are referred to as “stripons”; they carry charge $`\mathrm{e}`$ but no spin, and are created by fermion operators $`p_\mu ^{}(𝐤)`$. A starting point of localized stripon states is assumed, due to the fatal effect of imperfections in the striped structure on itineracy in one dimension. This assumption is supported by the atomic-scale structure observed recently by STM Davis1 , which is consistent with a fluctuating domino-type two-dimensional arrangement of stripe-like inhomogeneities. Such a two-dimensional arrangement has been predicted by the author Ashk03 . The $`𝐤`$ wave vectors of the stripon states present $`𝐤`$-symmetrized combinations of localized states to be treated in a perturbation expansion when coupling to the other fields is considered.
Away from the charged stripes, creation operators of approximate fermion basis states of spinon–holon and spinon–excession pairs are constructed Ashk01 ; Ashk03 . Together with the small-$`U`$ states they form, within the auxiliary space, a basis to “quasi-electron” (QE) states, carrying charge $`\mathrm{e}`$ and spin $`{\scriptscriptstyle \frac{1}{2}}`$, and created by $`q_{\iota \sigma }^{}(𝐤)`$. The bare QE energies $`ϵ_\iota ^q(𝐤)`$ form quasi-continuous ranges of bands within the BZ.
Atomic doping in the cuprates is taking place in inter-planar layers (such as in the chains in YBCO) between CuO<sub>2</sub> planes, and orbitals of the doped atoms contribute states close to $`E__\mathrm{F}`$. These states hybridize with the QE bands, and provide charge transfer to the planes with doping. The QE bands are spanned between the upper and lower Hubbard bands and the vicinity of $`E__\mathrm{F}`$, with more holon-based states closer to the lower Hubbard band (relevant mainly for “p-type” cuprates), and more excession-based states closer to the upper Hubbard band (relevant mainly for “n-type” cuprates). As the doping level $`x`$ is increasing, more QE states are moving from the Hubbard bands to the vicinity of $`E__\mathrm{F}`$, and the system is moving from the insulating to the metallic side of the Mott transition regime.
As was mentioned above, the treatment of the constraint, within the auxiliary space, introduces an additional chemical-potential-like Lagrange multiplier. Since the stripons carry only charge, while the QE’s carry both charge and spin, their treatment involves two “chemical potentials”, $`\mu ^q`$ and $`\mu ^p`$ (corresponding to QE’s and stripons, respectively) whose values are determined by the correct charge, and averaged constraint.
Hopping and hybridization terms introduce strong coupling between the QE, stripon, and spinon fields, which is expressed by a Hamiltonian term of the form (for p-type cuprates):
$`^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{\iota \lambda \mu \sigma }{}}{\displaystyle \underset{𝐤,𝐤^{}}{}}\{\sigma ϵ_{\iota \lambda \mu }^{qp}(\sigma 𝐤,\sigma 𝐤^{})q_{\iota \sigma }^{}(𝐤)p_\mu ^{}(𝐤^{})`$ (2)
$`\times [\mathrm{cosh}(\xi _{\lambda ,\sigma (𝐤𝐤^{})})\zeta _{\lambda \sigma }^{}(𝐤𝐤^{})`$
$`+\mathrm{sinh}(\xi _{\lambda ,\sigma (𝐤𝐤^{})})\zeta _{\lambda ,\sigma }^{}(𝐤^{}𝐤)]+h.c.\},`$
introducing a vertex between their propagators Ashk99 ).
The stripe-like inhomogeneities are strongly coupled to the lattice Bian , and it is the presence of stripons which creates the charged stripes within them. Also the matrix elements in $`^{}`$, are sensitive to the atomic positions Andersen . Consequently Ashk04 , the spinons are renormalized, being “dressed” by phonons, and thus carry some lattice distortion (but no charge) in addition to spin $`{\scriptscriptstyle \frac{1}{2}}`$. Such phonon-dressed spinons are referred to as “svivons”, and they replace the spinons in the $`^{}`$ vertex. This results in strong coupling between electronic spin excitations and Cu–O optical phonon modes Egami , and the existence of an anomalous isotope effect Lanzara1 . The effect of spin-lattice coupling has been studied by Eremin et al. Eremin2 .
The physical signature of the auxiliary fields, considering only the large-$`U`$ band, within the $`t`$$`t^{}`$$`J`$ model, is demonstrated in Fig. 1. An adiabatic “snapshot” of a section of a CuO<sub>2</sub> plane, including a stripe-like inhomogeneity, is shown. Within the adiabatic time scale a site is “spinless” either if it is “charged”, removing the spinned electron/hole on it (as in “stripon sites” in Fig. 1), or if the spin is fluctuating on a shorter time scale (due to, e.g., being in a singlet spin pair). In this description, a site (bare) stripon excitation represents a transition between these two types of a spinless site within the charged stripes, a site (bare) svivon excitation represents a transition between a spinned site and a fluctuating-spin spinless site, and a site (bare) QE excitation represents a transition between a spinned site and a charged spinless site within the AF stripes.
### II.3 Renormalized auxiliary particles
The $`^{}`$ vertex introduces self-energy corrections to the QE, stripon, and svivon fields Ashk99 . Since the renormalized stripon bandwidth is considerably smaller than the QE and svivon bandwidths, a phase-space argument could be used, as in the Migdal theorem, to ignore vertex corrections. Self-consistent expressions were derived Ashk01 ; Ashk03 for the self-energy corrections, and spectral functions $`A^q`$, $`A^p`$, and $`A^\zeta `$, for the QE, stripon, and svivon fields, respectively. The auxiliary-particle energies $`ϵ`$ are renormalized to: $`\overline{ϵ}=ϵ+\mathrm{}\mathrm{\Sigma }(\overline{ϵ})`$, where $`\mathrm{\Sigma }`$ is the self energy. Due to the quasi-continuous range of QE bands, the bandwidth renormalization is particularly strong for the stripon energies, resulting in a very small bandwidth, and limited itineracy due to hopping via intermediary QE–svivon states.
The small stripon bandwidth introduces Ashk01 ; Ashk03 a low-energy scale of $`0.02`$eV, and an apparent “zero-energy” non-analytic behavior of the QE and svivon self energies within a higher energy range (analyticity is restored in the low-energy range). This behavior results in QE and svivon scattering rates with a term $`\omega `$, as in the MFL approach Varma , but also with a constant term, resulting in a logarithmic singularity in $`\mathrm{}\mathrm{\Sigma }(\omega )`$ at $`\omega =0`$ (which is truncated by analyticity in the low-energy range). Consequently the QE self-energy has a kink-like behavior of the renormalized QE energies $`\overline{ϵ}^q`$ around zero energy Ashk03 . A typical renormalization of the svivon energies, around the V-shape zero minimum of $`ϵ^\zeta `$ at $`𝐤_0`$, is shown in Fig. 2, where the major effect is due to the truncated logarithmic singularity in the svivon self energy. The svivon spectral functions $`A^\zeta (\omega )`$ vary continuously around $`\omega =0`$, being positive for $`\omega >0`$ and negative for $`\omega <0`$.
The renormalization of the svivon energies changes the physical signature of their Bose condensation from an AF order to the observed stripe-like inhomogeneities (including both their spin and lattice aspects). The structure of $`A^\zeta `$ and $`\overline{ϵ}^\zeta `$ around the minimum at $`𝐤_0`$ (see Fig. 2) determines the structure of the inhomogeneities. Striped structure of the type shown in Fig. 1 results from a direction-dependent slope of $`\overline{ϵ}^\zeta (𝐤)`$ at small negative energies. The speed of the dynamics of these inhomogeneities depends on the linewidth of $`\overline{ϵ}^\zeta `$ at small negative energies, which is large, unless the system is in a pairing state (thus the SC or PG state), where svivon scattering at such energies is suppressed due to the gap Ashk03 (see below). Consequently, the dynamics of the stripe-like inhomogeneities becomes sufficiently slow for them to be detected only in a pairing state, in agreement with experiment.
For the renormalized svivons, $`\mathrm{cosh}(\xi _𝐤)`$ and $`\mathrm{sinh}(\xi _𝐤)`$ do not diverge at $`𝐤_0`$, as in Eq. (1). But still the values of both of them are large in the negative-$`\overline{ϵ}^\zeta (𝐤)`$ region around this point (see Fig. 2), and thus this region contributes significantly to processes involving svivons. By Eq. (2), the coupling between QE’s, stripons, and svivons of this region, are particularly strong when $`\overline{ϵ}^q\overline{ϵ}^p\pm \overline{ϵ}^\zeta `$. On the other hand, a relatively smaller contribution is obtained to such processes from the positive-$`\overline{ϵ}^\zeta (𝐤)`$ regions, especially when they are sufficiently away from $`𝐤_0`$, where by Eq. (1) $`\mathrm{cosh}(\xi _𝐤)1`$, and $`\mathrm{sinh}(\xi _𝐤)0`$.
Since the stripons are based on states in the charged stripe-like inhomogeneities, which occupy about a quarter of the CuO<sub>2</sub> plane Tran1 (see Fig. 1), the number of stripon states is about a quarter of the number of states in the BZ. The $`𝐤`$ values mostly contributing to these states reflect, on one hand, the structure of the stripe-like inhomogeneities, and on the other hand, the minimization of free energy, achieved when they reside mainly at BZ areas where their coupling to the QE’s and svivons, is optimal. This occurs Ashk03 for stripon coupling with svivons around $`𝐤_0`$ (see Fig. 2), and with QE’s at BZ areas of highest density of states (DOS) close to $`E__\mathrm{F}`$, which are found in most of the cuprates around the “antinodal” points $`(\frac{\pi }{a},0)`$ and $`(0,\frac{\pi }{a})`$. If (from its four possibilities) $`𝐤_0`$ were chosen at $`(\frac{\pi }{2a},\frac{\pi }{2a})`$, then the BZ areas in those cuprates, which the $`𝐤`$ values contributing to the stripon states mostly come from, would be Ashk03 at about a quarter of the BZ around $`\pm 𝐤^p=\pm (\frac{\pi }{2a},\frac{\pi }{2a})`$. Creation operators $`p_\mathrm{e}^{}(\pm 𝐤^p)`$ and $`p_\mathrm{o}^{}(\pm 𝐤^p)`$ of stripon states, which are an even and an odd combination of states at $`𝐤^p`$ and $`𝐤^p`$, were demonstrated Ashk04 to be compatible with the striped structure shown in Fig. 1, in areas where the stripes are directed along either the $`a`$ or the $`b`$ direction.
### II.4 Hopping-induced pairing
The electronic structure obtained here for the cuprates provides a hopping-energy-driven pairing mechanism. Diagrams for such pairing, based on transitions between pair states of stripons and QE’s, through the exchange of svivons, were presented in Ref. Ashk99 . It has been demonstrated Ashk03 , within the $`t`$$`t^{}`$$`J`$ model, that there is an energy gain in inter-stripe hopping of pairs of neighboring stripons through intermediary states of pairs of opposite-spin QE’s (where a svivon is exchanged when one pair is switched to the other), compared to the hopping of two uncorrelated stripons, through intermediary QE–svivon states (since the intermediary svivon excitations are avoided). The contribution of orbitals beyond the $`t`$$`t^{}`$$`J`$ model to the QE states results in further gain in stripon pairing energy, due to both intra-plane and inter-plane pair hopping.
This pairing scheme provides Eliasherg-type equations, of coupled stripon and QE pairing order parameters, which can be combined to give BCS-like equations (though including strong-coupling effects) for both of them in the second order Ashk03 . To study these equations, a domino-type two-dimensional arrangement of stripe-like inhomogeneities, including crossover between stripe segments directed in the $`a`$ and the $`b`$ directions, was assumed Ashk03 , and found later to be consistent with STM results Davis1 . An overall $`d_{x^2y^2}`$-type pairing symmetry was obtained Ashk03 , under which pair correlations are maximal between opposite-spin nearest-neighbor QE sites, and vanish between same-spin next-nearest-neighbor QE sites (see Fig. 1). Sign reversal is obtained Ashk03 for the QE order parameter through the charged stripes. Thus, the lack of long-range coherence in the details of the stripe-like inhomogeneities, especially between different CuO<sub>2</sub> planes, results in features different from those of a simple $`d_{x^2y^2}`$-wave pairing (especially when $`c`$-direction hopping is involved), as has been observed Dynes .
### II.5 Pairing and coherence
Within the GTC Ashk03 ; Ashk04 , the phase diagram of the cuprates is largely the consequence of interplay between pairing and coherence within the regime of a Mott transition. The pairing mechanism, which depends on the stripe-like inhomogeneities, is stronger when the AF/stripes effects are stronger, thus closer to the insulating side of the Mott transition regime. Consequently, the temperature $`T_{\mathrm{pair}}`$, below which pairing occurs, decreases with the doping level $`x`$, as is sketched in the pairing line in Fig. 3. On the other hand, phase coherence, be it of single electrons or of pairs, requires the energetic advantage of itineracy around $`E__\mathrm{F}`$, which is easier to achieve closer to the metallic side of the Mott transition regime. Thus the temperature $`T_{\mathrm{coh}}`$, below which coherence occurs, increases with $`x`$, as is sketched in the coherence line in Fig. 3. Thus $`T_c\mathrm{min}(T_{\mathrm{pair}},T_{\mathrm{coh}})`$.
For $`x\mathrm{}>0.19`$, one has $`T_{\mathrm{pair}}<T_{\mathrm{coh}}`$, and $`T_c`$ is determined by $`T_{\mathrm{pair}}`$. Single-electron coherence, which means the existence of an FL state, then exists for $`T_c<T<T_{\mathrm{coh}}`$ (see Fig. 3), as indicated by ARPES results Ashk03 ; Ashk04 . Within the non-FL approach used here, the stripe-like inhomogeneities are treated adiabatically; but the absence of a pairing gap results in a large svivon linewidth around $`𝐤_0`$, and thus fast stripes dynamics, which is consistent with the existence of an FL state below $`T_{\mathrm{coh}}`$.
For $`x\mathrm{}<0.19`$, one has $`T_{\mathrm{coh}}<T_{\mathrm{pair}}`$, and the normal-state PG, observed in the cuprates in this regime, is Ashk03 ; Ashk04 (at least partly) a pair-breaking gap at $`T_c<T<T_{\mathrm{pair}}`$ (see Fig. 3). In this regime $`T_{\mathrm{pair}}`$ is generally referred to as $`T^{}`$, and $`T_c`$ is determined by $`T_{\mathrm{coh}}`$. Its value is of the order of the phase stiffness, estimated through the treatment of “classical” phase fluctuations by the “X–Y” model Emery . It is given in this regime by:
$`k__\mathrm{B}T_c`$ $``$ $`k__\mathrm{B}T_{\mathrm{coh}}\left({\displaystyle \frac{\mathrm{}}{\mathrm{e}}}\right)^2{\displaystyle \frac{a\rho _s}{16\pi }},`$ (3)
$`\rho _s`$ $`=`$ $`{\displaystyle \frac{4\pi \mathrm{e}^2n_s^{}}{m_s^{}}}=\left({\displaystyle \frac{c}{\lambda }}\right)^2,`$ (4)
in agreement with Uemura’s law Uemura (based on the determination of the penetration-depth $`\lambda `$ from $`\mu SR`$ results).
In order to clarify somewhat the phase stiffness energy in Eq. (3), let us write it as: $`N_s(\mathrm{}K_1)^2/2m_s^{}`$, where $`N_s=n_s^{}(Na)^2a/8\pi ^2`$ is the average number of pairs in the volume of a slice of thickness $`a/8\pi ^2`$ around a plane parallel to the CuO<sub>2</sub> planes (assumed to be squares of length $`Na`$), $`m_s^{}`$ is a pair’s mass, and $`K_1=2\pi /Na`$ is the closest wave-vector point (within a plane) to the pairs’ $`𝐊=0`$ ground-state point. Thus the phase stiffness is related to the energy needed (due to lattice discreteness) to excite a macroscopic number of pairs from the single-pair ground state to the first excited state. Thus, in the PG state, single-pair states are occupied (though not macroscopically) up to energies $`k__\mathrm{B}T>k__\mathrm{B}T_{\mathrm{coh}}`$ (which are a considerable fraction of the pairs bandwidth). Consequently moderate interactions, mixing between $`𝐊`$ pair states (such as induced by phonons), are likely to mix them, resulting in bipolaron-like localized pair states.
The application of a compressive strain, in this regime, results in the increase of $`a\rho _s`$, and thus also of $`T_{\mathrm{coh}}`$ and $`T_c`$ \[see Eq. (3)\]. And indeed, ARPES results in LSCO thin films Pavuna show that $`T_c`$ rises under such a strain. As was discussed above, the increase in $`T_{\mathrm{coh}}`$ is consistent with a move towards the metallic side of the Mott transition regime, which is reflected in the increase of the width of the bands around $`E__\mathrm{F}`$ under this strain Pavuna .
When a junction is made, including slices of a cuprate in the SC state and in the PG state, an injection of pairs is expected to occur between them, as in p–n junctions in semiconductors. Consequently, the density of pairs in a range within the PG side of the junction could exceed the phase-stiffness limit of Eq. (3) for the occurrence of SC there. This explains the observation of a GPE in trilayer junctions of cuprate thin films Bozovic , oriented both in the $`ab`$ plane, and in the $`c`$ direction. Unlike the regular proximity effect, where the range is determined by the coherence length, the range of the GPE is determined by the range where injection of carriers between the SC and PG slices occurs. This range is not related to the coherence length, and could be larger by more than an order of magnitude from it, as is observed Bozovic .
### II.6 Gap equations and heterogeneity
In order to understand the natures of the SC and PG states, let us regard the QE and stripon (Bogoliubov) energy bands obtained through their BCS-like equations Ashk04 :
$`E_\pm ^q(𝐤)`$ $`=`$ $`\pm \sqrt{\overline{ϵ}^q(𝐤)^2+\mathrm{\Delta }^q(𝐤)^2},`$ (5)
$`E_\pm ^p(𝐤)`$ $`=`$ $`\pm \sqrt{\overline{ϵ}^p(𝐤)^2+\mathrm{\Delta }^p(𝐤)^2}.`$ (6)
$`2\mathrm{\Delta }^q`$ and $`2\mathrm{\Delta }^p`$ are related to the observed pairing gap Ashk04 , as will be discussed below. They scale with $`T_{\mathrm{pair}}`$, approximately according to the BCS factors, with an increase due to the effect of strong coupling. A $`d`$-wave pairing factor Maki is relevant for $`\mathrm{\Delta }^q`$, which has its maximum $`\mathrm{\Delta }_{\mathrm{max}}^q`$ at the antinodal points. Since the stripons reside in about a quarter of the BZ around $`\pm 𝐤^p`$, $`|\mathrm{\Delta }^p|`$ does not vary much from its mean value $`\overline{\mathrm{\Delta }}^p`$, and an $`s`$-wave pairing factor is relevant for it. Thus, at low $`T`$:
$$2\mathrm{\Delta }_{\mathrm{max}}^q\mathrm{}>4.3k__\mathrm{B}T_{\mathrm{pair}},2\overline{\mathrm{\Delta }}^p\mathrm{}>3.5k__\mathrm{B}T_{\mathrm{pair}}.$$
(7)
The stripon bandwidth $`\omega ^p`$ has been estimated from TEP results Ashk01 (see below) to be $`0.02`$eV, and by the measured $`T_{\mathrm{pair}}`$ (see Fig. 3) and the above expression, it is considerably smaller than $`\overline{\mathrm{\Delta }}^p`$ in the UD regime, and exceeds its value only in the heavily OD regime. Consequently, in the UD regime, $`E_\pm ^p(𝐤)\pm |\mathrm{\Delta }^p(𝐤)|`$, and the Bogoliubov transformation dictates Ashk04 an approximate half filling of the band of the paired stripons. Thus, if all the stripons were paired, the stripon band would have been approximately half filled (thus $`n^p={\scriptscriptstyle \frac{1}{2}})`$ at low $`T`$. This is inconsistent with TEP results Ashk01 (see below) according to which $`n^p>{\scriptscriptstyle \frac{1}{2}}`$ in the UD regime, and becomes $`{\scriptscriptstyle \frac{1}{2}}`$ for $`x0.19`$.
Consequently Ashk04 , only a part of the stripons could be paired in the UD regime, while complete QE pairing (except for the nodal and localized states discussed below) in expected for $`T0`$. Thus the PG state consists of both paired and unpaired stripons. Since an SC ground state is normally characterized by complete pairing for $`T0`$, an intrinsically heterogenous SC state is obtained there Ashk04 , with nanoscale perfectly SC regions, where, locally, $`n^p{\scriptscriptstyle \frac{1}{2}}`$, and partial-pairing PG-like regions, where $`n^p>{\scriptscriptstyle \frac{1}{2}}`$ locally, but SC still occurs in proximity to the $`n^p{\scriptscriptstyle \frac{1}{2}}`$ regions. This effect is expected to be weaker, or absent, in the OD regime, where $`\overline{\mathrm{\Delta }}^p`$ becomes comparable, and even smaller than $`\omega ^p`$.
Such a nanoscale heterogeneity, was indeed observed in STM data in the SC phase Davis2 , and it is also supported by optical results Dordevic1 . Its features are consistent with the above prediction Ashk04 about the existence of perfectly SC regions, and (especially in the UD regime) partial-pairing PG-like regions. The size of these regions is comparable to the SC coherence length, so that SC is maintained also in the partially-paired regions.
### II.7 Quantum critical point
The degeneracy of the paired states in a perfectly SC state is maintained by the dynamics of the stripe-like inhomogeneities Ashk04 . Thus the free-energy gain in this state keeps them dynamical to $`T0`$. This is not the case in the PG state (and also within PG-like regions in an heterogenous SC state), where these inhomogeneities partially freeze at low $`T`$ into a glassy checkerboard structure Davis1 ; Davis2 . Within this structure one can observe Davis1 $`a`$\- and $`b`$-directed stripe segments, and also $`(4a)\times (4a)`$ patterns obtained due to switching between $`a`$\- and $`b`$-directed segments, when fluctuations occur between domino-type two-dimensional arrangements of the inhomogeneities Ashk03 .
The formation of this glassy structure is of a similar nature to CDW/SDW transitions, and orbital effects of the type of the DDW Laughlin ; Varma1 may also play a role. Energy is gained by creating partial or complete minigaps for the unpaired carriers (which may result in an additional superstructure Davis1 ); but even if there is no real gap within the whole BZ, the unpaired carriers would become Anderson-localized due to their low DOS in the low-$`T`$ disordered glassy structure. And indeed, low-$`T`$ upturns are observed in the electrical resistivity in the PG states, either for low $`x`$, or if the SC state is suppressed, for $`x\mathrm{}<0.19`$, by applying a magnetic field Boeb , or by doping Tallon . When such a suppression occurs, the pairing and coherence lines (see Fig. 3) meet at $`T=0`$ around $`x0.19`$, where an MIT occurs Ashk04 between the FL metallic phase, and the PG non-metallic phase.
The $`T=0`$ MIT point in Fig. 3 satisfies the conditions of a QCP Marel6 . The stoichiometry $`x`$ where it occurs is close to $`x_c0.19`$, where the fractional stripon occupancy, as was determined from the TEP results Ashk01 (see below), is $`n^p={\scriptscriptstyle \frac{1}{2}}`$. The existence of the MIT close to this stoichiometry is plausible Ashk04 , because for higher doping levels the bare stripons become too packed within the charged stripes (see Fig. 1), and inter-atomic Coulomb repulsion between them is likely to destabilize the stripe-like inhomogeneities in the PG state (though the energy gain in the SC state helps maintaining them for higher $`x`$) and stabilize the homogeneous FL state.
## III Previous Applications of the GTC
### III.1 Electron spectrum
Spectroscopic measurements (as in tunneling and ARPES) based on the transfer of electrons into, or out of, the crystal, are determined by the electron’s spectral function $`A_e`$, obtained by projecting the auxiliary spectral functions $`A^q`$, $`A^p`$, and $`A^\zeta `$ to the physical space. Such an expression was derived for $`A_e`$ Ashk03 , and it includes a QE ($`A^q`$) term, and a convoluted stripon–svivon ($`A^pA^\zeta `$) term. From the quasi-continuum of QE bands, only few bands, which are closely related to those of physical electrons, contribute “coherent” bands, while the other QE bands contribute an “incoherent” background to $`A_e`$. Both the bands and the background include hybridized $`A^q`$ and $`A^pA^\zeta `$ contributions, having widths including Ashk03 an MFL-type term $`\omega `$, and a constant term, in agreement with experiment. $`A_e`$ is spanned between the upper and lower Hubbard bands, and the vicinity of $`E__\mathrm{F}`$, and its weight around $`E__\mathrm{F}`$ is increasing with $`x`$.
Since the stripon states reside (in most cuprates) mainly in a quarter of the BZ around points $`\pm 𝐤^p`$ Ashk03 ; Ashk04 (see above), a significant $`A^pA^\zeta `$ contribution to $`A_e`$ close to $`E__\mathrm{F}`$ (at energies around $`\overline{ϵ}^p\pm \overline{ϵ}^\zeta `$) is obtained, with svivons around their energy minimum at $`𝐤_0`$, in BZ areas around the antinodal points. Thus the $`A^pA^\zeta `$ contribution to $`A_e`$ is not significant close to “nodal” FS crossing points, in the vicinity of $`\pm (\frac{\pi }{2a},\pm \frac{\pi }{2a})`$, where $`A_e`$ is determined primarily by $`A^q`$.
Thus the shape of the electron bands around the nodal points is similar to that of $`\overline{ϵ}^q`$, and the (almost) $`T`$-independent “nodal kink” observed by ARPES Lanzara2 ; Johnson , closely below $`E__\mathrm{F}`$, in p-type cuprates, corresponds Ashk04 to the $`T`$-independent kink-like behavior obtained for $`\overline{ϵ}^q`$ there (see above). The absence of such a kink in ARPES measurements in the n-type cuprate NCCO Armitage1 is also consistent with the GTC, which predicts it to occur there closely above $`E__\mathrm{F}`$ Ashk04 (and thus out of the range of ARPES).
On the other hand, the $`T`$-dependent “antinodal kink”, observed by ARPES around the antinodal points Gromko ; Sato ; Lanzara1 , where its major part appears only below $`T_c`$, is due to the $`A^pA^\zeta `$ contribution to the electron bands there Ashk04 . Since (see below) the opening of a pairing gap causes a decrease in the svivon linewidth around the energy minimum at $`𝐤_0`$ (see Fig. 2), this contribution narrows down as $`T`$ is decreased below $`T_{\mathrm{pair}}`$, and especially below $`T_c`$ Ashk04 , as is observed in this kink Gromko ; Sato ; Lanzara1 .
### III.2 Pair-breaking excitations
Thus the antinodal kink is a spectroscopic signature of the pairing gap. Since (in a pairing state) the bandwidth of the Bogoliubov bands $`E_+^p`$ and $`E_{}^p`$, in Eq. (6), is small, the convoluted stripon–svivon states of energies $`E_+^p\pm \overline{ϵ}^\zeta `$ and $`E_{}^p\pm \overline{ϵ}^\zeta `$ form Ashk04 (for svivons around $`𝐤_0`$) spectral peaks, centered at $`E_+^p`$ and $`E_{}^p`$, around the antinodal points. The size of the SC gap is experimentally determined by the spacing between the closest spectral maxima on its two sides. Thus it is given by:
$`2|\mathrm{\Delta }^{\mathrm{SC}}(𝐤)|`$ $`=`$ $`2\mathrm{min}[|\mathrm{\Delta }^q(𝐤)|,E_{\mathrm{peak}}(𝐤)],`$
$`E_{\mathrm{peak}}(𝐤)`$ $`=`$ $`|E_\pm ^p(𝐤\pm 𝐤_0)|.`$ (8)
By Eq. (7), $`\mathrm{\Delta }_{\mathrm{max}}^q>\overline{\mathrm{\Delta }}^p`$; consequently $`\mathrm{\Delta }^{\mathrm{SC}}`$ is determined by $`\mathrm{\Delta }^q`$ around its zeroes at the nodal points, and by $`\mathrm{\Delta }^p`$ around its maxima at the antinodal points \[where $`E_{\mathrm{peak}}(𝐤)`$ exist\]. Actually, since this convoluted stripon–svivon peak lies on the slope of the QE gap, its maximum is shifted to an energy slightly above $`E_{\mathrm{peak}}(𝐤)`$.
Since the QE and convoluted stripon–svivon states hybridize with each other, the states at energies $`E_\pm ^q(𝐤)`$, around the antinodal points, are scattered to stripon–svivon states at energies $`E_\pm ^p\pm \overline{ϵ}^\zeta `$ of magnitudes above $`E_{\mathrm{peak}}(𝐤)`$ (see Fig. 2 and the discussion following it). This results in the widening of the QE coherence peak \[due to Eq. (5)\], at the QE gap edge, to a hump Ashk04 .
In the PG state Ashk04 , the pairs lack phase coherence, and thus Eq. (5) does not yield a coherence peak in the QE gap edge. Furthermore, in this state unpaired convoluted stripon–svivon states exist Ashk04 (see discussion above) within the gap, resulting in the widening of the low-energy svivon states, and thus of $`\pm E_{\mathrm{peak}}(𝐤)`$, due to scattering. Consequently, the gap-edge stripon–svivon peak is smeared, and at temperatures well above $`T_c`$ the PG becomes a depression of width: $`2|\mathrm{\Delta }^{\mathrm{PG}}(𝐤)|=2|\mathrm{\Delta }^q(𝐤)|`$ in the DOS Ashk04 , in agreement with tunneling results Renner ; Kugler .
Thus the pair-breaking excitations in the SC state are characterized by Ashk04 a peak-dip-hump structure (on both sides of the gap) in agreement with tunneling Renner ; Kugler and ARPES results Gromko ; Feng1 ; Borisenko ; Janowitz ; Plate . The peak is largely contributed by the convoluted stripon–svivon states around $`E_{\mathrm{peak}}(𝐤)`$, the dip results from the sharp descent at the upper side of this peak, and the hump above them is of the QE gap edge and other states, widened due to the scattering to stripon–svivon states above the peak, discussed above. By Eqs. (7), and (8), $`\mathrm{\Delta }^{\mathrm{SC}}`$ and $`\mathrm{\Delta }^{\mathrm{PG}}`$ scale with $`T_{\mathrm{pair}}`$, and thus decrease with $`x`$, following the pairing line in Fig. 3, as has been observed Ashk04 .
The heterogeneous existence, in an SC state, of perfectly SC and PG-like partial-pairing regions, discussed above, has been observed by STM Davis2 through the distribution in the heights and widths of the gap-edge peak (widened due to scattering of svivons to unpaired stripons and QE’s). These STM results also show Davis2 that, unlike the gap-edge peak, the low energy excitations near the SC gap minimum are not affected by this heterogeneity. This is consistent with the prediction Ashk04 \[see Eq. (8)\] that the SC gap is determined, around its zeroes at the nodal points, by the QE gap $`\mathrm{\Delta }^q`$. The magnitudes of the QE energies $`E_\pm ^q(𝐤)`$ around the nodal points, are below the range of $`E_{\mathrm{peak}}`$, and convoluted stripon–svivon states, which may hybridize with them, correspond to svivons which are not at the vicinity of $`𝐤_0`$, and to energies $`E_\pm ^p\pm \overline{ϵ}^\zeta `$ of magnitudes above $`E_{\mathrm{peak}}`$. Thus the hybridization between them is insignificant (see Fig. 2 and the discussion following it), and the linewidths of the QE states around the nodal points are small in the SC state, and do not vary between the perfectly SC and PG-like regions, as the gap-edge states do.
### III.3 Localization gaps, and nodal FS arcs
As was mentioned above, the formation of the glassy (checkerboard) structure in the PG state, and the PG-like SC regions, is self-consistently intertwined with the formation of (at least partial) minigaps and the localization of the unpaired stripons. Such gaps are formed there also in QE states around the antinodal points, which are strongly hybridized with these unpaired stripon states (convoluted with svivons around $`𝐤_0`$), and they become localized too. QE states are extended over a larger range in space than the size of the SC regions, and those with the above localization gaps can participate in the pairing process only if these gaps are (approximately) smaller than their pairing gaps.
Since the QE bandwidth is much larger than $`|\mathrm{\Delta }^q|`$ \[see Eq. (5)\], they become almost completely paired at low $`T`$ (also in the PG state), except for the QE states which have too big localization gaps to participate in the pairing process, and those at the vicinity of the nodal points. The minimal doping level $`x_00.05`$ Basov5 , for which SC pairing occurs, is determined by the condition that for $`xx_0`$ the number of QE states, which have sufficiently small localization gaps to be coupled to stripon states in the pairing process Ashk03 , drops below the minimum necessary for the pairing to occur.
Consequently, ARPES measurements in the SC as well as the PG state Marshall ; Norman1 , show that parts of FS around the antinodal points disappear, due to the formation of the QE gap (including parts which are due to localization, and parts which are due to pairing). On the other hand, arcs of the FS remain around the zero-$`\mathrm{\Delta }^q`$ nodal points, where unpaired QE states persist at low $`T`$, and become localized for $`T0`$. Such arcs continue to exist also for $`x<x_0`$ Yoshida2 . The $`x<x_0`$ regime is characterized by “diagonal stripes”, where the stripon states do not contribute at $`E__\mathrm{F}`$ Yoshida2 , and transport is due to the QE’s on the nodal arcs (with the rest of the FS missing). The modulation caused by the diagonal stripes, in a direction perpendicular to them, could be the reason for the minigap observed Shen by ARPES in the nodal direction in this regime.
Resistivity ($`\rho `$) measurements through $`x=x_0`$ Ando confirm the low-$`T`$ localization on the nodal FS arcs (whether it is Anderson localization, or due to a minigap), and show a monotonous variation of $`\rho (T)`$ with $`x`$, which is not affected (except for the occurrence of SC) by the change in the striped structure or the AF transition. This indicates that the QE’s on the nodal FS arcs are hardly affected by the change in the striped structure. Their hopping is likely to be dominated by $`t^{}`$ processes Ashk03 , which do not disturb the AF order. One could distinguish between the heavily UD regime Basov5 of $`x<x_0^{}0.09(>x_0)`$, where low $`T`$ transport is largely due to QE’s on the nodal FS arcs, and the rest of the cuprates phase diagram (for $`x>x_0^{}`$), where their role in transport less important.
### III.4 Spin excitations and the resonance mode
Tunneling results Zasadzinski show a correlation between the width of the SC gap-edge peak and the neutron-scattering resonance-mode energy Bourges1 $`E_{\mathrm{res}}`$. Such spin excitations are determined by the imaginary part of the spin susceptibility $`\chi ^{\prime \prime }(𝐤,\omega )`$, and an expression for the contribution of the large-$`U`$ orbitals to it has been derived in Ref. Ashk03 . It turns out that large contributions to it are obtained from double-svivon excitations, when both svivon states are close to $`𝐤_0`$, and their energies, $`\omega _1`$ and $`\omega _2`$, either have the same sign, and contribute to $`\chi ^{\prime \prime }(𝐤,\omega )`$ at $`\omega =\pm (\omega _1+\omega _2)`$, or they have opposite signs, and contribute to it at $`\omega =\pm (\omega _1\omega _2)`$. This results (see Fig. 2) in two branches of spin excitations, a low-$`\omega `$ branch, having a maximum $`2\overline{ϵ}^\zeta (𝐤_0)`$ at $`𝐤=𝐐=2𝐤_0`$ (identified as the resonance mode Bourges1 ; Reznik ), and a high-$`\omega `$ wide branch with an extensive minimum, spreading over the first branch Ashk03 ; Ashk04 . These branches, and also their linewidths, below and above $`T_c`$, correspond to neutron-scattering results Reznik ; Tran2 ; Hayden ; Birgeneau . \[A more quantitative calculation of $`\chi ^{\prime \prime }(𝐤,\omega )`$ is in preparation\].
The width of the spin excitations is determined Ashk04 by the scattering between QE, stripon, and svivon states, which is strong when $`E^qE^p\pm \overline{ϵ}^\zeta `$, for svivon states close to $`𝐤_0`$ \[where the $`\mathrm{cosh}(\xi _𝐤)`$ and $`\mathrm{sinh}(\xi _𝐤)`$ factors are large – see Eqs. (1) and (2) and the discussion following Fig. 2\]. The existence of a pairing gap limits (especially below $`T_c`$) the scattering of the svivon states around $`𝐤_0`$, resulting in a decrease in their linewidth. Let $`𝐤_{\mathrm{min}}`$ be the points of small svivon linewidth, for which $`\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})`$ is the closest to the energy minimum $`\overline{ϵ}^\zeta (𝐤_0)`$ (see Fig. 2). Often one has $`𝐤_{\mathrm{min}}=𝐤_0`$, but there are cases, like that of LSCO Ashk04 , where the linewidth of $`\overline{ϵ}^\zeta `$ is small not at $`𝐤_0`$, but at close points $`𝐤_{\mathrm{min}}=𝐤_0\pm 𝐪`$. The resonance mode energy $`E_{\mathrm{res}}`$ is taken here as $`2\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})`$, accounting both for the often observed “commensurate mode” at $`Q`$, and for cases of an “incommensurate mode” Ashk04 , as observed in LSCO at $`𝐐\pm 2𝐪`$ Tran3 ; Wakimoto ; Christ .
The determination of $`𝐤_{\mathrm{min}}`$ is Ashk04 through the condition $`E_{\mathrm{res}}=2|\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})|2\stackrel{~}{\mathrm{\Delta }}^{\mathrm{SC}}`$, where $`2\stackrel{~}{\mathrm{\Delta }}^{\mathrm{SC}}`$ is somewhat smaller than the maximal SC gap $`2\mathrm{\Delta }_{\mathrm{max}}^{\mathrm{SC}}`$. Since $`\overline{ϵ}^\zeta (𝐤_0)`$ is zero for an AF (see Fig. 2), its value (and thus $`E_{\mathrm{res}}`$) is expected to increase with $`x`$, distancing from an AF state. However, since its linewidth cannot remain small if $`|\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})|`$ exceeds the value of $`\stackrel{~}{\mathrm{\Delta }}^{\mathrm{SC}}`$, which decreases with $`x`$, the energy $`E_{\mathrm{res}}`$ of a sharp resonance mode is expected Ashk04 to cross over from an increase to a decrease with $`x`$ when it approaches the value of $`2\stackrel{~}{\mathrm{\Delta }}^{\mathrm{SC}}`$, as has been observed Bourges1 . This crossover could be followed Ashk04 by a shift of the resonance wave vector $`2𝐤_{\mathrm{min}}`$ from the AF wave vector $`𝐐`$ to incommensurate wave vectors.
### III.5 The gap-edge peak
By the above scattering conditions Ashk04 , svivon energies $`\overline{ϵ}^\zeta `$ have a small linewidth within the range $`|\overline{ϵ}^\zeta ||\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})|`$. Consequently, the gap-edge peaks of the convoluted stripon–svivon states of energies $`E_\pm ^p\pm \overline{ϵ}^\zeta `$, centered at $`\pm E_{\mathrm{peak}}`$ \[see Eq. (8)\], have a “basic” width:
$$W_{\mathrm{peak}}=2\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})=E_{\mathrm{res}}.$$
(9)
Additional contributions to the width of this peak come from the svivon and stripon linewidths, and from the dispersion of $`E_\pm ^p(𝐤\pm 𝐤^{})`$ when $`\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})\mathrm{}<\overline{ϵ}^\zeta (𝐤^{})\mathrm{}<0`$ (see Fig. 2). This result explains Ashk04 the observed correlation (for different doping levels) between the peak’s width, and $`E_{\mathrm{res}}`$ Zasadzinski .
Studies of the pair-breaking excitations in the SC state by ARPES Gromko ; Feng1 ; Borisenko ; Janowitz ; Plate confirm also the $`𝐤`$ dependence predicted by the GTC Ashk04 . The convoluted stripon–svivon states contribute over a range of the BZ around the antinodal points a single weakly-dispersive gap-edge peak at $`E_{\mathrm{peak}}(𝐤)`$. Most of the measurements were performed on bilayer BSCCO, where in addition to this peak there are around the antinodal points bilayer-split QE bands, the bonding band (BB) and the antibonding band (AB), contributing two humps around the gap Gromko ; Feng1 ; Borisenko . The AB lies very close to $`E__\mathrm{F}`$ on the SC gap edge through the antinodal BZ range, and it almost overlaps with $`E_{\mathrm{peak}}(𝐤)`$ in the OD regime, where they both appear as narrow peaks Gromko ; Feng1 .
The BB, on the other hand, disperses considerably, crossing $`E__\mathrm{F}`$, and contributes a clearly distinguished hump Gromko ; Feng1 ; Borisenko . In the range where the QE BB approaches $`E_{\mathrm{peak}}(𝐤)`$, the fact that the electron band is formed by their hybridized contributions results in the appearance of the antinodal kink Gromko ; Sato ; Lanzara1 (discussed above), due to the narrowing of the peak, as $`T`$ is lowered below $`T_c`$. The width of the hump states depends on the rate of their scattering to stripon–svivon states. Since the svivons are dressed by phonons, an anomalous isotope effect is obtained for the width, and thus also for the position of the hump states Lanzara1 .
The peak-dip-hump structure has been observed also in tunneling Kugler and ARPES Janowitz measurements in single-layer BSCO and BSLCO, proving that it is not just the effect of bilayer splitting Ashk04 . Recent ARPES measurements in OD single-layer TBCO Plate show the existence of a gap-edge peak, which could be identified with $`E_{\mathrm{peak}}(𝐤)`$, over a range of the BZ around the antinodal points. It disperses over a range $`0.02`$eV, which is comparable with the SC gap, consistently with the GTC prediction Ashk04 for the OD regime.
Low-temperature ARPES results for the spectral weight within the sharp SC gap-edge peak, integrated over the antinodal BZ area Feng2 ; Feng3 , show a maximum for $`x0.19`$. This is expected by the GTC Ashk04 , assuming that the integrated spectral weight is counted within the stripon–svivon peak around $`E_{\mathrm{peak}}`$ in regions where this peak is sharp. If the number of svivon $`𝐤`$ states contributing to that peak (see Fig. 2) does not vary significantly with doping in the range of interest, than the measured integrated peak counts the number of hole-like pair-breaking excitations of stripons within the $`E_{}^p`$ band \[see Eq. (6)\], in regions where the peak is sharp.
For the intrinsically heterogenous $`x\mathrm{}<0.19`$ regime, discussed above, this number increases with $`x`$ because of the increase in the fraction of space covered by the perfectly SC regions (where the stripon band is approximately half full and the stripon–svivon states contribute a sharp gap-edge peak). For the $`x\mathrm{}>0.19`$ regime there are no PG-like regions, and the peak is sharp wherever it exists. Thus the measured integrated peak counts there the number of hole-like pair-breaking excitations of stripons within the $`E_{}^p`$ band, which is decreasing with the increase of $`x`$ below half filling of the stripon band Ashk04 . Note that the contribution of the QE AB to the ARPES peak had to be omitted Feng3 in order to get the decrease of the peak weight for $`x\mathrm{}>0.19`$, confirming the GTC prediction that this behavior is due to the stripon–svivon gap-edge peak.
### III.6 Asymmetry of the tunneling spectrum
One of the features of the tunneling spectrum in the cuprates Davis2 ; Renner is its asymmetry with higher DOS for hole- than particle-excitations. This asymmetry is extending beyond the limit of the presented spectrum, few tenths of an eV on both sides of $`E__\mathrm{F}`$. Within the GTC, high-$`T_c`$ in the cuprates is occurring in the regime of a Mott transition, where the spectral function $`A_e(\omega )`$ is spanned between the upper and lower Hubbard bands and the vicinity of $`E__\mathrm{F}`$. As was discussed above Ashk03 , $`A_e(\omega )`$ has a large incoherent part, and its magnitude is decreasing when $`\omega `$ is varied from the Hubbard bands to the energy space between them.
In p-type cuprates $`E__\mathrm{F}`$ is closer to the (Zhang-Rice-type) lower Hubbard band, and thus $`A_e(\omega )`$ is descending when $`\omega `$ is varied from below $`E__\mathrm{F}`$ to above it (Anderson and Ong Anderson1 demonstrated such a behavior within the $`t`$$`J`$ model), in agreement with the observed asymmetry in the tunneling spectrum Davis2 ; Renner . In n-type cuprates, $`E__\mathrm{F}`$ is closer to the upper Hubbard band, and thus an opposite asymmetry is expected in the tunneling spectrum there (with higher DOS for particle- than hole-excitations), in agreement with results in n-type NCCO Kashiwaya and “infinite-layer” SLCO Chen .
### III.7 Transport properties
#### III.7.1 Expressions within the one-band approximation
The normal-state transport properties have been a major mystery in the cuprates. The author has been involved from the start in attempts to understand the anomalous TEP, resistivity, and Hall constant Fisher ; Bar-Ad . Their correct treatment, within the GTC, though not including the effects of the PG, was first presented in Ref. Ashk01 . Even though, the dynamics of the stripe-like inhomogeneities is fast above $`T_{\mathrm{pair}}`$, their adiabatic treatment is still expected to be a reasonable approximation for the evaluation of the transport properties. Their derivation, within the $`ab`$ plane, is based on linear-response theory, A condition used is that the direct current (DC) $`𝐣`$ could be expressed as a sum of QE ($`𝐣^q`$) and stripon ($`𝐣^p`$) terms which are proportional to each other with an approximately $`T`$-independent proportionality factor $`\alpha `$. Thus:
$$𝐣=𝐣^q+𝐣^p\frac{𝐣^q}{1\alpha }\frac{𝐣^p}{\alpha }.$$
(10)
As is discussed further below, in relation to optical conductivity, this condition is a consequence of the assumption that the contribution of the CuO<sub>2</sub> planes to transport is derived, dominantly, from one electron band of the homogeneous planes. The formation of stripe-like inhomogeneities results in separate contributions of QE’s and stripons, but the stripes dynamics results in DC based on the band of the averaged homogeneous planes. The coefficient $`\alpha `$ depends on the inhomogeneous structure, which determines how $`𝐣`$ is split between $`𝐣^q`$ and $`𝐣^p`$.
Transport can be treated using the separate QE and stripon chemical potentials, $`\mu ^q`$ and $`\mu ^p`$ (due to the large-$`U`$ constraint), discussed above. When an electric field $`𝐄`$ is applied, Eq. (10) is satisfied by the formation of gradients $`\mathbf{}\mu ^q`$ and $`\mathbf{}\mu ^p`$ Ashk01 , where the homogeneity of charge neutrality imposes:
$$\frac{n_e^q}{\mu ^q}\mathbf{}\mu ^q+\frac{n_e^p}{\mu ^p}\mathbf{}\mu ^p=0.$$
(11)
Here $`n_e^q`$ and $`n_e^p`$ are the contributions of QE and stripon states to the electrons occupation. Because of the small stripon bandwidth, it is convenient to introduce Ashk01 :
$$N_e^q\frac{n_e^q}{\mu ^q},M_e^p(T)T\frac{n_e^p}{\mu ^p},$$
(12)
where $`N_e^q`$ is almost temperature independent, and $`M_e^p(T)`$ saturates at temperatures above the stripon bandwidth to a constant $`n^p(1n^p)`$ Ashk01 .
The chemical potential gradients introduce “chemical fields”, resulting in different effective fields for QE’s and stripons Ashk01 :
$$𝓔^q=𝐄+\mathbf{}\mu ^q/\mathrm{e},𝓔^p=𝐄+\mathbf{}\mu ^p/\mathrm{e},$$
(13)
and by Eqs. (11), (12), and (13), one gets:
$$𝐄=\frac{M_e^p(T)𝓔^p+N_e^qT𝓔^q}{M_e^p(T)+N_e^qT}.$$
(14)
Effective QE and stripon conductivities $`\sigma ^q`$ and $`\sigma ^p`$ are introduced through the expressions:
$$𝓔^q=𝐣/\sigma ^q(T),𝓔^p=𝐣/\sigma ^p(T),$$
(15)
and related expressions are used Ashk01 to introduce effective QE and stripon Hall numbers $`n__\mathrm{H}^q`$ and $`n__\mathrm{H}^p`$, and TEP’s $`S^q`$ and $`S^p`$. Using Eq. (14), the following expressions are then obtained Ashk01 for the electrical conductivity ($`\sigma `$), Hall number ($`n__\mathrm{H}`$), and TEP ($`S`$):
$`\sigma `$ $`=`$ $`{\displaystyle \frac{M_e^p(T)+N_e^qT}{M_e^p(T)/\sigma ^p(T)+N_e^qT/\sigma ^q(T)}},`$ (16)
$`n__\mathrm{H}`$ $`=`$ $`{\displaystyle \frac{M_e^p(T)+N_e^qT}{M_e^p(T)/n__\mathrm{H}^p+N_e^qT/n__\mathrm{H}^q}},`$ (17)
$`S`$ $`=`$ $`{\displaystyle \frac{M_e^p(T)S^p(T)+N_e^qTS^q(T)}{M_e^p(T)+N_e^qT}}.`$ (18)
At low $`T`$, for $`xx_0^{}`$, the value of $`M_e^p(T)`$ is generally about an order of magnitude greater than $`N_e^qT`$. Consequently, by Eqs. (11) and (12), $`\mathbf{}\mu ^p0`$ and $`\mathbf{}\mu ^q\mathrm{e}𝐄`$. Thus, by Eq. (13), $`𝓔^p𝐄`$, $`𝓔^q0`$, and by Eqs. (16), (17) and (18), the transport coefficients are dominantly stripon-like at low $`T`$, both above and below $`T_c`$, and a crossover to a QE-like behavior is occurring when $`T`$ is increased through the stability temperature regime. When pairing gaps exist, $`\mu ^q`$ and $`\mu ^p`$ are within these gaps, and their existence results in some decrease in the low-$`T`$ value of $`M_e^p(T)`$ in Eq. (12), resulting in some increase in the QE contribution to transport at low $`T`$. The assumption that $`M_e^p(T)N_e^qT`$, at low $`T`$, breaks down when $`x`$ approaches $`x_0^{}`$ (towards the heavily UD regime), where the PG and SC gap are large. As was discussed above, low-$`T`$ transport for $`x<x_0^{}`$ is largely due to QE’s, which are on the nodal FS arcs, for the unpaired carriers Basov5 .
#### III.7.2 Transport results
The above expressions were applied Ashk01 , using a minimal set of realistic parameters. The transport coefficients were calculated as a function of $`T`$, for five stoichiometries of p-type cuprates, ranging from $`n^p=0.8`$, corresponding to the UD regime, to $`n^p=0.4`$, corresponding to the OD regime. GTC auxiliary spectral functions, and zero-energy $`T`$-dependent scattering rates were used, and also impurity-scattering $`T`$-independent terms. The stripon band was modeled by a “rectangular” spectral function $`A^p`$ of width $`\omega ^p`$, for which was a value of $`0.02`$eV was found to be consistent with experiment.
The TEP results depend strongly on $`n^p`$, and reproduce very well Ashk01 the doping-dependent experimental behavior Fisher ; Tanaka , which has been used to determine stoichiometry. The low-$`T`$ stripon-like result saturates at $`T\mathrm{}>200`$K to a constant $`(k__\mathrm{B}/\mathrm{e})\mathrm{ln}[n^p/(1n^p)]`$ Ashk01 . But then a crossover starts towards a QE-like linear behavior, for which a negative slope is predicted Ashk01 , in agreement with experiment Tanaka . The stripon term $`S^p(T)`$ vanishes for $`n^p={\scriptscriptstyle \frac{1}{2}}`$, and then the TEP is determined by negative-slope QE term $`S^p(T)`$, corresponding experimentally to $`x=x_c0.19`$. Also the results for the Hall coefficients Ashk01 reproduce very well the experimental behavior Kubo ; Hwang . The approximately linear increase of $`n__\mathrm{H}`$ with $`T`$ is due to its crossover from a low-$`T`$ stripon-like $`n__\mathrm{H}^p`$ towards a QE-like $`n__\mathrm{H}^q`$. The signs of $`n__\mathrm{H}^p`$ and $`n__\mathrm{H}^q`$ are the same, both determined by the nature of the one-band current $`𝐣`$ of the averaged homogeneous planes. Thus $`n__\mathrm{H}`$ is not expected to change sign with $`T`$ within the assumed one-band approximation \[leading to Eq. (10)\].
The calculated $`T`$ dependence of $`\rho `$ Ashk01 is linear at high $`T`$, flattening to (stripon-like) “superlinearity” at low $`T`$ (for all stoichiometries). Experimentally Takagi , the low-$`T`$ behavior crosses over from “sublinearity” in the UD regime (and a non-metallic upturn at lower $`T`$ if $`T_c`$ is low enough), to superlinearity in the OD regime. As is discussed below, in relation to optical conductivity, the sublinear behavior is the effect of reduced scattering rate in the PG state. The low-$`T`$ upturn results from the localization in this state, discussed above. The crossover to superlinear behavior (predicted here) in the OD regime is a natural consequence of the disappearance of the PG with increasing $`x`$ (see Fig. 3). The linear $`T`$-dependence of $`\rho `$ persists to low-$`T`$ for $`x0.19`$ Oh , which is due to quantum criticality Marel6 close to the QCP (discussed above). In the critical region there is only one energy scale, which is the temperature, resulting in MFL-type behavior Varma .
The TEP in n-type cuprates is normally expected Ashk01 to behave similarly to the TEP in p-type cuprates, but with an opposite sign and slope. Results for NCCO Takeda show such behavior for low doping levels, but in SC doping levels the high-$`T`$ slope of $`S`$ is changing from positive to negative, and its behavior resembles that of OD p-type cuprates. This led Ashk01 to the suggestion that NCCO may be not a real n-type cuprate, its stripons being based on holon states (as in p-type cuprates).
More recent results on the n-type infinite-layer SLCO Williams do show TEP results for an SC cuprate which have the opposite sign and slope than those calculated for p-type cuprates Ashk01 , as is expected for real n-type cuprates (thus with stripons based on excession states). It is likely that the sign of the TEP slope in NCCO changes with doping because the one-band approximation, leading to Eq. (10), becomes invalid. This is indicated by ARPES Armitage2 , and by the change with temperature of the sign of the Hall constant of NCCO Takeda , at the stoichiometries where the sign of the slope of the TEP has changed (see above). Also, in YBCO, the contribution of an additional band of the chains’ carriers results in an almost zero high-$`T`$ slope of the TEP Fisher , rather than the negative slope predicted by the GTC Ashk01 .
### III.8 Superfluid density
The Uemura’s plots Uemura give information about the effective density of SC pairs $`n_s^{}`$ through Eq. (4). One of the mysteries of the cuprates has been the boomerang-type behavior Niedermayer of these plots for $`x\mathrm{}>0.19`$. The author connected this behavior, as early 1994 Ashk94 , with the fact that the (presently called) stripon band passes through half filling. As was mentioned above, the effect of pairing is the hybridization between stripon and QE pairs Ashk03 . The determination of $`\lambda `$ in the $`\mu SR`$ measurements Uemura is through DC in a magnetic field, and as was determined in Eqs. (16)–(18), and the discussion following them, DC below $`T_c`$ is stripon like for lightly UD, optimal (OPT), and OD stoichiometries.
Thus $`n_s^{}`$ determined from $`\lambda `$, through $`\mu SR`$ measurements Uemura , approximately corresponds, in this regime, to the density of stripon pairs. As was mentioned above, the stripon band is half full for $`x=x_c0.19`$, and consequently $`n_s^{}`$ is maximal around this stoichiometry, being determined (for p-type cuprates) by the density of hole-like stripon pairs for $`x<x_c`$, and of particle-like stripon pairs for $`x>x_c`$. This result is not changed by the intrinsic heterogeneity Ashk04 for $`x\mathrm{}<0.19`$, discussed above; $`T_c`$ is determined there, through Eq. (3), by the average value of $`n_s^{}`$, determined by measuring the penetration depth which Tajima is larger than the sizes of the heterogenous regions Davis2 .
In the crossover to the heavily UD ($`x<x_0^{}`$) regime, the behavior of low-$`T`$ DC crosses over, as was discussed above, from being stripon-like to being QE-like. Thus, the effective superfluid density $`n_{s,\mathrm{eff}}`$ (derived from the measured $`\lambda `$) becomes dominated, for $`x<x_0^{}`$, by the QE contribution to the superfluid. As was discussed above, the glassy structure in PG-like SC regions results in the formation of localization gaps for QE states around the antinodal points, thus preventing their contribution to the superfluid when their localization gaps are, approximately, greater than their pairing gaps. Since the fraction of the QE states with localization gaps, which are small enough to contribute to the superfluid, decreases as $`xx_0`$ (and $`n_s^{}0`$), the contribution of QE’s to the superfluid decreases then faster than its density $`n_s^{}`$. Thus, it is reasonable to assume in this regime an approximate expression of the form $`n_{s,\mathrm{eff}}(n_s^{})^\beta `$, where $`\beta >1`$. Consequently, Eqs. (3) and (4) yield in this regime $`T_cn_s^{}n_{s,\mathrm{eff}}^{1/\beta }`$. And indeed, penetration-depth measurements in the heavily UD regime Zuev reveal such a behavior with $`\beta =2.3\pm 0.4`$.
## IV Optical Conductivity within the $`𝒂𝒃`$ Plane
### IV.1 One-band formalism
It is assumed that the optical conductivity, within the $`ab`$ plane, is dominantly contributed (within the frequency range of interest) by one band Macridin of the homogeneous planes. Thus the relevant Hamiltonian is expressed as:
$$=\underset{𝐤,\sigma }{}ϵ_\mathrm{b}(𝐤)d_\sigma ^{}(𝐤)d_\sigma ^{}(𝐤)+_{\mathrm{int}}.$$
(19)
It includes a “bare band” \[$`ϵ_\mathrm{b}(𝐤)`$\] one-particle term, and a two-particle interaction term $`_{\mathrm{int}}`$. $``$ determines electron velocities $`𝐯(𝐤)`$, in terms of which the electrical current operator (due to it) is expressed as:
$$\widehat{𝐣}\mathrm{e}\underset{𝐤,\sigma }{}𝐯(𝐤)d_\sigma ^{}(𝐤)d_\sigma ^{}(𝐤).$$
(20)
The electron creation and annihilation operators $`d_\sigma ^{}(𝐤)`$, and $`d_\sigma ^{}(𝐤)`$, are expressed in terms of the auxiliary-space QE and convoluted stripon–svivon operators. Consequently the $`\omega `$-dependent electrical current, in the presence of an electric field $`𝐄(\omega )`$, can be expressed as a sum of a QE ($`𝐣^q`$), a stripon–svivon ($`𝐣^p`$), and a mixed ($`𝐣^{qp}`$)term:
$$𝐣(\omega )=\widehat{𝐣}(\omega )=𝐣^q(\omega )+𝐣^p(\omega )+𝐣^{qp}(\omega ).$$
(21)
Since the svivons carry no charge, their effect on $`𝐣^p`$ is through field-independent spin occupation factors. On the other hand, $`𝐣^{qp}`$ involves field-induced transitions between QE and stripon–svivon states, and thus svivon excitations. Since $`A^\zeta (\omega =0)=0`$, the contribution of $`𝐣^{qp}`$ to Eq. (21) vanishes for $`\omega 0`$, and Eq. (10) can be used for DC.
### IV.2 The f–sum rule
If the entire frequency spectrum were considered, the real part of the electrons’ contribution to the optical conductivity $`\sigma (\omega )=𝐣(\omega )/𝐄(\omega )`$ should obey Kubo1 ; Tinkham ; Norman ; Hanke ; Marel3 the f–sum rule $`_0^{\mathrm{}}\mathrm{}\sigma (\omega )𝑑\omega =\pi n\mathrm{e}^2/2m_e`$, where $`n`$ and $`m_e`$ are the electrons’ (total) density and (unrenormalized) mass. By considering only the one-band contribution of $``$ in Eq. (19), one obtains the partial f–sum rule Kubo1 ; Tinkham ; Norman ; Hanke ; Marel3 ($`V`$ is the volume):
$`{\displaystyle _0^{\mathrm{}}}\mathrm{}\sigma (\omega )𝑑\omega `$ $`=`$ $`{\displaystyle \frac{\pi \mathrm{e}^2}{2V}}{\displaystyle \underset{𝐤,\sigma }{}}{\displaystyle \frac{d_\sigma ^{}(𝐤)d_\sigma ^{}(𝐤)}{m_\mathrm{b}(𝐤)}},`$ (22)
$`{\displaystyle \frac{1}{m_\mathrm{b}(𝐤)}}`$ $``$ $`{\displaystyle \frac{1}{\mathrm{}^2}}{\displaystyle \frac{^2ϵ_\mathrm{b}(𝐤)}{𝐤^2}}.`$ (23)
The effective mass $`m_\mathrm{b}(𝐤)`$, appearing in Eq. (22), is within the bare band in $``$ \[Eq. (19)\]. The effect of the $`_{\mathrm{int}}`$ term there is through the occupation factors $`d_\sigma ^{}(𝐤)d_\sigma ^{}(𝐤)`$ in Eq. (22). Since a large-$`U`$-limit method is applied here for $`_{\mathrm{int}}`$, an optical determination of $`d_\sigma ^{}(𝐤)d_\sigma ^{}(𝐤)`$ requires the consideration of transitions including states in the lower, as well as the upper Hubbard band (both in p-type cuprates and in n-type cuprates), which means the inclusion of energies $`\omega \mathrm{}>2`$eV in the integral in Eq. (22). This conclusion is supported by the XAS results Sawatzky for the position of the chemical potential in p- and n-type cuprates.
The optical conductivity can be determined experimentally by reflectance measurements and a Kramers–Kronig analysis Tanner1 ; Marel3 , and the integral in Eq. (22) can be carried out up to an experimental cutoff frequency $`\omega _{\mathrm{co}}`$. This procedure has to be modified in the SC state Tanner1 ; Marel3 , where the superfluid spectral weight, contributing to the integral, comes from a $`\delta `$-function term in $`\mathrm{}\sigma (\omega )`$ Tinkham . This spectral weight equals $`\omega _{\mathrm{ps}}^2/8`$, where $`\omega _{\mathrm{ps}}`$ is the superfluid plasma frequency, and it can be determined Tanner1 ; Marel3 by ellipsometry measurements, including the evaluation of the real part of of the dielectric function $`ϵ(\omega )`$, and the application of the relation: $`\mathrm{}ϵ(\omega )=ϵ(\mathrm{})(\omega _{\mathrm{ps}}/\omega )^2`$.
This derivation of the superfluid plasma frequency provides an optical method to determine the value of the parameter $`\rho _s`$ \[see Eq. (4)\], through the relation $`\rho _s=\omega _{\mathrm{ps}}^2`$. There is a question concerning the relation between this optically-derived value and the value of $`\rho _s`$ derived through measurements of the penetration depth $`\lambda `$ in a magnetic field, by methods like $`\mu SR`$. The carrier masses, determined by the two methods, are not identical; while the optical measurements yield the bare-band mass $`m_\mathrm{b}(𝐤)`$, appearing in Eq. (22), carriers dynamics in a magnetic field depends on the effective mass $`m_s^{}`$ of the pairs. The considerable reduction of $`m_s^{}`$ compared to the stripon’s mass is the driving force of the GTC pairing mechanism Ashk03 .
And indeed, it turns out Tajima that the $`ab`$ plane penetration depth determined optically, is about twice the value determined by $`\mu SR`$. Also, as was discussed above, the stripons’ band is half full for $`x0.19`$, and $`m_s^{}`$ changes there from being hole like to being particle like, resulting in a boomerang-type behavior of $`\rho _s`$ in the OD regime (thus changing from increasing to decreasing with $`x`$), in agreement with $`\mu SR`$ results Niedermayer . On the other hand $`m_\mathrm{b}(𝐤)`$, which determines the optical $`\rho _s`$, corresponds to the bare mass of the averaged homogeneous CuO<sub>2</sub> planes, which is hole like, and does not pass through half filling within the SC doping range. Thus the optically determined $`\rho _s`$ is expected to rise with $`x`$, and have no boomerang-type behavior, as has been observed Timusk1 .
### IV.3 “Violations” of the f–sum rule
In ordinary SC’s Tinkham , the formation of an SC gap in $`\mathrm{}\sigma (\omega )`$ is followed by the transfer of spectral weight of magnitude $`\omega _{\mathrm{ps}}^2/8`$ from it to the $`\delta `$-function term. Thus, the value of the integral in Eq. (22) (including the $`\omega _{\mathrm{ps}}^2/8`$ contribution in the SC state) is not expected then to change between the normal and the SC state, when $`\omega _{\mathrm{co}}`$ is taken sufficiently above the gap energy. Ellipsometric measurements in p-type cuprates Basov3 ; Marel1 ; Bontemps ; Homes2 confirm such behaviour in the OD regime, while for UD and OPT stoichiometries, this behavior was found to be “violated” even when $`\omega _{\mathrm{co}}`$ values above $`2`$eV were used. This “violation” points to the transfer of spectral weight below $`T_c`$ from energies $`\mathrm{}>2`$eV to the vicinity of $`E__\mathrm{F}`$ (its validity is confirmed in a recent debate about it Boris ).
There have been theoretical suggestions trying to explain this transfer of spectral weight as being due to a mechanism of pairing from a non-FL normal state, to an “FL SC” state Norman , due to pairing phase fluctuations Hanke , or due to pairing via spin fluctuations within the nearly AF FL model Carbotte1 . But the high energy scale $`\mathrm{}>2`$eV involved is hard to understand unless it is assumed that the spectral weight is transferred from both the lower and the upper Hubbard bands Ashk03 ; Ashk04 (thus beyond the range of applicability of the $`t`$$`t^{}`$$`J`$ model). As was discussed above, the spectral weight in the vicinity of $`E__\mathrm{F}`$ is increasing with $`x`$ at the expense of the weight in the upper and lower Hubbard bands far from it, resulting in an increase in the itineracy of the carriers. The change in $`x`$ is provided by doping atoms out of the CuO<sub>2</sub> planes. These atoms contribute electronic states close to $`E__\mathrm{F}`$ \[and thus also contribute to $`\sigma (\omega )`$\], from which charge is transferred to the CuO<sub>2</sub> planes.
The UD and OPT stoichiometries, where the transfer of spectral weight at $`T_c`$ has been observed Basov3 ; Marel1 ; Bontemps ; Homes2 , are those where the SC transition is from the PG state (see Fig. 3). As was discussed above, the transition there is due the establishment of phase coherence of existing localized pairs, turning them into a superfluid. There are two possible mechanisms (or their combination) for an accompanying transfer of spectral weight from the Hubbard bands to the vicinity of $`E__\mathrm{F}`$, both driven by the free energy gain in the SC state. The first one is that this transfer of spectral weight is associated with the increased itineracy of the pairs in the superfluid. The second mechanism is that since the free energy gain due to SC is determined (like $`T_c`$) by the phase stiffness, which scales \[see Eqs. (3) and (4)\] with the density of pairs within the CuO<sub>2</sub> planes, it drives further charge transfer from the doped atoms to the planes below $`T_c`$. This results in transfer of spectral weight from the Hubbard bands to the vicinity of $`E__\mathrm{F}`$, as if $`x`$ were increased. On the other hand, in the OD regime the SC transition is more BCS-like, due to pairing of electrons in an FL state, and such a transfer of spectral weight is expected less, if at all, within both of the above mechanisms.
### IV.4 Optical carriers around optimal stoichiometry
#### IV.4.1 Contributions to the effective density
An experimental study of the optical conductivity, and of the partial f–sum rule, as a function $`\omega _{\mathrm{co}}`$, for a eight different cases of p-type cuprates near OPT doping, was carried out by Tanner et al. Tanner1 . A typical curve of $`\mathrm{}\sigma (\omega )`$, for different temperatures, is shown in Fig. 4. The electron mass $`m_e`$ was chosen for $`m_\mathrm{b}(𝐤)`$ in Eq. (22), which is unjustified beside being a working assumption (it is not clear at this point whether an effective mass derived in an LDA-based calculation would be an appropriate choice either).
A value $`\omega _{\mathrm{co}}12000\mathrm{cm}^1`$ was chosen Tanner1 in order to avoid the effect of the “charge transfer” band, and count mainly the carriers in the conduction band. As can be concluded, e.g., from the measured position of the chemical potential Sawatzky , such a choice of $`\omega _{\mathrm{co}}`$ in Eq. (22) does not count the contribution of most of the excession-based states of the upper Hubbard band (ignored in the $`t`$$`t^{}`$$`J`$ model), while it does count the contribution of most of the holon-based states of the (Zhang-Rice-type) lower Hubbard band (considered in the $`t`$$`t^{}`$$`J`$ model). Also the contribution of states of the doped atoms (out of the CuO<sub>2</sub> planes) is counted if they are close enough to $`E__\mathrm{F}`$. In YBCO this corresponds to the chains states, and they contribute considerably in measurements with polarization in the chains direction ($`b`$).
Optical carriers’ densities derived through this experimental analysis Tanner1 are presented in Table I. One derived quantity is the effective number $`N_{\mathrm{eff}}/\mathrm{Cu}`$ of carriers per Cu atom at $`T100`$K (above $`T_c`$). Most values are found to be around 0.4–0.5, except for a lower value of 0.15 for monolayer LCO, and a higher value of 0.59 (due to the contribution of the chains’ carriers) for YBCO with the polarization taken in the $`b`$ direction. It may be misleading to connect this $`N_{\mathrm{eff}}`$ with an actual number of carriers, because the contribution of the upper-Hubbard-band states is not integrated in Eq. (22). These states are essential to determine the number of carriers contributed by the QE’s, which form bands combining states of the lower and the upper Hubbard bands. On the other hand the number of carriers contributed by stripons could be well described just within the frame of the integrated lower-Hubbard-band states.
Two other types of carriers’ densities Tanner1 , presented in Table I, may have more physical significance. The first one is the number $`N_\mathrm{s}/\mathrm{Cu}`$ of the superfluid carriers per Cu atom, which was determined (as was explained above) from the $`\delta `$-function term in $`\sigma (\omega )`$ at $`T10`$K. The second one is the number $`N_\mathrm{D}/\mathrm{Cu}`$ of Drude carriers per Cu atom (at $`T100`$K). It was determined by fitting $`\sigma (\omega )`$ to the contributions of a number of oscillators, including a Drude oscillator at zero frequency, and Lorentzian oscillators at higher frequencies. $`N_\mathrm{D}`$ was obtained by integrating over the Drude contribution, which introduces (for $`T>T_c`$) the major low-$`\omega `$ contribution to $`\sigma (\omega )`$ Tanner1 in Fig. 4.
Both $`N_\mathrm{D}`$ and $`N_\mathrm{s}`$ are related to the number of carriers involved in DC conductivity, which as was shown in Eq. (16), and the discussion following it, is dominantly stripon-like at low $`T`$, for the stoichiometry studied. Thus $`N_\mathrm{D}`$ and $`N_\mathrm{s}`$ approximately correspond to the number of stripon “holes”, where $`N_\mathrm{s}`$ has somewhat larger QE contribution due to the effect of the gap on Eq. (12). The contribution to the current in Eq. (21) of $`𝐣^{qp}(\omega )`$ (due to transitions between QE and stripon–svivon states) could be neglected in the treatment of DC above, but it does result in much of the non-Drude (often called mid-IR) contribution to $`\sigma (\omega )`$ Tanner1 in Fig. 4. The dressing by phonons, discussed above, is reflected in signatures of their structure there. Also are included in this term QE contributions which (as was discussed above) are largely “blocked” at low $`\omega `$ and $`T`$ by the QE chemical potential gradient \[see Eq. (13)\], satisfying there $`\mathbf{}\mu ^q\mathrm{e}𝐄`$.
#### IV.4.2 Stripon-like carriers
Thus, it is not surprising that the ratios $`N_\mathrm{s}/N_\mathrm{D}`$ Tanner1 , presented in Table I, are close to one (ranging between 0.87 and 0.95, except for values around 0.80 for LCO and YBCO with polarization in the $`b`$ direction). If all the Drude carriers were paired, the somewhat larger QE contribution to $`N_\mathrm{s}`$ would have yielded for it a larger value than $`N_\mathrm{D}`$. However, since as was discussed above, there are some unpaired carriers for OPT stoichiometry, in the PG-like heterogenous regions \[as can be seen in $`\sigma (\omega )`$ below $`T_c`$ Tanner1 in Fig. 4\], one gets somewhat smaller $`N_\mathrm{s}`$ than $`N_\mathrm{D}`$. The lower $`N_\mathrm{s}/N_\mathrm{D}`$ ratio in YBCO is due to the contribution of the chain carriers, whose pairing could be approximately regarded as induced by proximity Kresin . Since the QE contribution to the carriers’ density above $`T_c`$ increases with $`T`$, the lower $`T_c`$ in LCO is consistent with a lower $`N_\mathrm{s}/N_\mathrm{D}`$ ratio in it.
As was discussed above, the TEP results Ashk01 indicate a half-filled stripon band for $`x=x_c0.19`$. Thus, for the striped structure shown in Fig. 1, the OPT stoichiometry (see Fig. 3), around which the measurements by Tanner et al. Tanner1 were carried out, corresponds to a number of $`N_{\mathrm{con}}/\mathrm{Cu}0.125\times 0.17/0.19=0.11`$ stripon hole carriers per Cu atom. The values of $`N_\mathrm{s}/\mathrm{Cu}`$ and $`N_\mathrm{D}/\mathrm{Cu}`$ for most of the cases Tanner1 in Table I are quite close to 0.10, except for smaller values for LCO, and larger values for YBCO with polarization in the $`b`$ direction.
This overall agreement is surprisingly good, considering the fact that $`m_e`$ was used for $`m_\mathrm{b}(𝐤)`$ in Eq. (22). The deviation in YBCO is understood due to the contribution of the chains’ carriers. The deviation in LCO could be because $`m_\mathrm{b}(𝐤)/m_e`$ is significantly larger than one there. Pavarini et al. Andersen have shown that when the parameters for a one-band approximation are derived from first-principles calculations, cuprates with larger maximal $`T_c`$ have larger $`t^{}`$ hopping parameters, and thus smaller $`m_\mathrm{b}(𝐤)`$. So having larger $`m_\mathrm{b}(𝐤)/m_e`$ for LCO than for the other cuprates studied (in agreement with Ref. Andersen ) is consistent with its considerably lower $`T_c`$. Also for TBCCO, which has higher $`T_c`$ than the other cuprates studied, its somewhat larger values of $`N_\mathrm{s}/\mathrm{Cu}`$ and $`N_\mathrm{D}/\mathrm{Cu}`$, than of the other cuprates presented in Table I, corresponds to its smaller $`m_\mathrm{b}(𝐤)/m_e`$.
#### IV.4.3 Tanner’s law and its resolution
The puzzling result of Tanner et al. Tanner1 (known as Tanner’s law) is that, for all the cases presented in Table I, $`N_\mathrm{s}/N_{\mathrm{eff}}`$ ranges between 0.19 and 0.23, which means that that $`N_{\mathrm{eff}}`$ equals 4-5 times the number of stripon carriers. If the integration through Eq. (22) were extended to include the contribution of the upper-Hubbard-band states, but omitting the contribution of bands not included in $``$ in Eq. (19), than the number of carriers per Cu atom would have been $`1x`$ (since the bare band is short by $`x/2`$ from being half full for each spin state). For OPT stoichiometry this corresponds to about 0.83 carriers per Cu atom, which is greater than the $`(N_{\mathrm{eff}}/\mathrm{Cu})`$ values measured Tanner1 on the basis of partial integration (omitting the upper-Hubbard-band states), and presented in Table I. It is, however, unrealistic to count just the contribution of the conduction-band carriers if the integration range in Eq. (22) is extended to include the upper-Hubbard-band states, since they overlap with other bands.
An analysis of $`\sigma (\omega )`$ on the basis of the lower-Hubbard-band states alone (counted in Ref. Tanner1 ) could be made, in analogy to the Kohn-Sham approach Kohn , by replacing $``$ in Eq. (19) by an effective Hamiltonian of small-$`U`$ carriers, of the same spin symmetry as the carriers in the real system. In this system an effective time- and spin-dependent single-particle potential is introduced in order to simulate (at least approximately) the many-body effects (causing the existence of the upper-Hubbard-band) occurring in the real system. As was shown in Eq. (16), and the discussion following it, such a many-body effect is that at low $`\omega `$ (thus DC) and $`T`$ the contribution of QE’s to transport is largely blocked around OPT stoichiometry by the QE chemical potential gradient.
In the effective system the large-$`U`$ constraint, which results in different chemical potentials for QE’s and stripons, doesn’t exist, and a single-particle potential has to replace it as a mechanism for blocking the contribution of QE’s, but not of stripons, to conductivity at low $`T`$ and $`\omega `$. This effective potential has to be dynamical to simulate the effect of the dynamical stripe-like inhomogeneities, and its time average should maintain the translational symmetry of the CuO<sub>2</sub> planes, as in the bare band in Eq. (19). In the effective system, as in the real one, the stripons correspond (see Fig. 1) to about a quarter of the relevant orbitals of the Cu atoms in the planes, while the QE’s to the remaining three quarters.
Such an effective dynamical potential exists, and it induces an effective (dynamical) SDW, where minigaps are created between the states corresponding to three quarters of the Cu atoms’ states, thus blocking their contribution to conductivity at low $`T`$ and $`\omega `$, while those corresponding to the remaining quarter of Cu atoms’ states continue contributing to conductivity. It also induces a dynamical charge transfer between these two types of Cu atoms. However, this effect is minor since the bare QE and stripon states are strongly renormalized, through the coupling between them Ashk01 ; Ashk03 ; Ashk04 \[see Eq. (2)\]. Thus the (dynamical) charge transfer between Cu atoms in the AF and the charged stripes is considerably smaller than what would be concluded from the occupation of stripon and QE states, if they were (wrongly) approximated by their bare states. Such a reduced charge transfer is estimated from the neutron-scattering Tran1 and STM Davis1 results.
In the effective system, interband transitions across the minigaps would recover in higher $`\omega `$ and $`T`$ the contribution to conductivity of the three quarters of the orbitals, blocked at low $`T`$ and $`\omega `$, and the f–sum rule would indicate a total number of carriers $`N_{\mathrm{eff}}`$ which is about 4 times the number of (stripon) carriers, contributing to transport at low $`T`$ and $`\omega `$. The contribution of carriers residing in the inter-planar layers, discussed above, alters the factor 4 to a somewhat higher factor, but this is partly compensated by the minor (dynamical) charge transfer between planar Cu atoms, mentioned above, and Tanner’s law Tanner1 is obtained.
### IV.5 Optical carriers in the heavily UD regime
#### IV.5.1 Low-energy excitations
An experimental study of the optical conductivity, together with DC transport, for LSCO and YBCO in the heavily UD regime, through $`x_0^{}`$, $`x_0`$, and the AF phase boundary, was carried out by Padilla et al. Basov5 . Results at low $`T`$ (7K or just above $`T_c`$) and $`x<x_0^{}`$ (thus in the regime where transport is dominated by QE’s on the nodal FS arcs) show Basov5 that the Drude term in $`\sigma (\omega )`$ is fairly separated there from the mid-IR contribution to it. In this case, $`\omega _{\mathrm{co}}`$ in the integration in Eq. (22) could be chosen low enough to include, approximately, only the Drude term in $`\sigma (\omega )`$.
Then, similarly to the above analysis of Tanner’s law, $``$ in Eq. (19) could be replaced (in analogy to the Kohn-Sham approach Kohn ) by an effective Hamiltonian, for which the Drude term would be the only contribution to $`\sigma (\omega )`$, and the effective periodic potential would yield carriers with the effective mass $`m__\mathrm{D}`$ of the Drude carriers. Thus, by choosing such $`\omega _{\mathrm{co}}`$, this effective system would yield: $`\rho __\mathrm{D}/8=_0^{\omega _{\mathrm{co}}}\mathrm{}\sigma (\omega )𝑑\omega =\pi n__\mathrm{D}\mathrm{e}^2/2m__\mathrm{D}`$, where $`n__\mathrm{D}`$ is the density of the Drude carriers.
Results obtained by Padilla et al. Basov5 are presented in Fig. 5. The density $`n__\mathrm{D}`$ was estimated by measuring the low-$`T`$ Hall constant (thus assuming $`n__\mathrm{D}=n__\mathrm{H}`$), which by Eq. (17), and the discussion following it, corresponds in this regime primarily to the QE’s of the nodal FS arcs, as well. The effective mass $`m__\mathrm{D}`$ of these carriers was estimated by combining the optical and the Hall results, presented, respectively, by red empty squares and dots in the bottom panels in Fig. 5 ($`N_{\mathrm{eff}}`$ there is $`\rho __\mathrm{D}`$).
The evaluated values of $`m__\mathrm{D}`$ are presented as red empty circles in the top panels of Fig. 5, and they turn out Basov5 to be $`4m_e`$ in LSCO, and $`2m_e`$ in YBCO, with almost no change over the range $`0.01\mathrm{}<x<x_0^{}0.09`$. The larger mass in LSCO is consistent with the results of Tanner et al. Tanner1 discussed above (see Table I). The constant value found for $`m__\mathrm{D}`$ in this regime confirms the GTC scenario that no divergence of the effective mass is occurring for $`x0`$ Basov5 , but that carriers are doped within the Hubbard gap, and their density is increasing with doping. These carriers become localized at low $`T`$ in the low-$`x`$ regime, and form an FL in the high-$`x`$ regime. The low-$`x`$ and high-$`x`$ limits are separated by the SC phase, or by a QCP in the case that SC is suppressed (see Fig. 3).
Padilla et al. Basov5 found that when $`x`$ is increased above $`x_0^{}`$, the overlap between the energy ranges of the Drude and the mid-IR contributions to $`\sigma (\omega )`$ is growing, and thus the above analysis in terms of a separate Drude term is becoming inappropriate. Nevertheless, a trend can still be observed Basov5 in Fig. 5 that the low-$`T`$ $`\rho __\mathrm{D}`$ continues its increase with $`x`$ at about the same rate as for $`x<x_0^{}`$, while the rate of increase of the low-$`T`$ $`n__\mathrm{H}`$ grows substantially for $`x>x_0^{}`$. This may indicate an increase in the effective mass of the carriers for $`x>x_0^{}`$, confirming the GTC prediction of a crossover in the type of carriers dominating low-$`T`$ transport, between QE’s on the nodal arcs, and the higher-effective-mass stripons.
#### IV.5.2 High-energy excitations
The study by Padilla et al. Basov5 included also the determination of the effective mass $`m_{\mathrm{eff}}`$, on the basis of a high-$`\omega `$ $`\rho _{\mathrm{eff}}`$, and a high-$`T`$ $`n_{\mathrm{eff}}`$. The first was obtained by integrating on $`\mathrm{}\sigma (\omega )`$ in Eq. (22) up to $`\omega _{\mathrm{co}}12000\mathrm{cm}^1`$ (in order to avoid the effect of the “charge transfer” band, as was discussed above Tanner1 ). The effective density ($`n_{\mathrm{eff}}=n__\mathrm{H}`$) was obtained by measuring $`n__\mathrm{H}`$ in the high-$`T`$ Hall-constant plateau, reached at about $`800`$$`900`$K Basov5 (thus above $`T^{}`$). These optical and the Hall results are presented, respectively, by blue empty squares and dots in the bottom panels in Fig. 5 ($`N_{\mathrm{eff}}`$ there is $`\rho _{\mathrm{eff}}`$). The obtained values Basov5 of $`m_{\mathrm{eff}}`$ for $`x<x_0^{}`$ in LSCO, are about the same as the low-$`T`$ $`m__\mathrm{D}`$ in this regime, discussed above (presented as red empty circles in the top panels of Fig. 5).
Both the high-$`T`$ plateau in $`n__\mathrm{H}`$ \[see Eq. (17)\], and the high-$`\omega `$ $`\rho _{\mathrm{eff}}`$, in the heavily UD ($`x<x_0^{}`$) regime, correspond within the GTC mainly to the contribution of QE’s. But while at low $`T`$ and $`\omega `$, the QE’s determining transport are those on the nodal FS arcs, at high $`T`$ (above $`T^{}`$) and $`\omega `$ all the QE’s in the conduction band become available for transport, due to the closing of the PG. Thus, the the result: $`m_{\mathrm{eff}}m__\mathrm{D}`$ Basov5 indicates that about the same effective mass is relevant for both cases. Temperatures in the $`800`$$`900`$K range still correspond to a lower energy than that determined by the AF exchange coupling in the cuprates Anisimov1 ; Munoz , and thus both $`m__\mathrm{D}`$ and $`m_{\mathrm{eff}}`$ are expected to be determined mainly by $`t^{}`$ processes Ashk03 , which do not compete with the effect of AF exchange.
An analysis of the relation between $`\rho _{\mathrm{eff}}`$ and $`\rho __\mathrm{D}`$, in the $`x<x_0^{}`$ regime, can be carried out similarly to the analysis of Tanner’s law above. But here the contribution of the carriers to transport, except for the QE’s on the nodal FS arcs, is blocked, at low $`T`$ and $`\omega `$, by (real) gaps, and becomes available when the wide frequencies range is considered. Thus, a different factor is expected for $`\rho _{\mathrm{eff}}/\rho __\mathrm{D}`$ in this regime, than the 4-5 factor Tanner1 obtained above in Tanner’s law around the OPT regime. And indeed, Padilla et al. Basov5 got $`\rho _{\mathrm{eff}}/\rho __\mathrm{D}6`$-7, for $`x<x_0^{}`$, and a crossover to Tanner’s 4-5 factor, as $`x`$ is increased above $`x_0^{}`$, as is expected from the crossover to a stripon-dominated low-$`T`$ transport in the rest of the doping regime.
Padilla et al. Basov5 also applied the “extended Drude model” Puchkov ; Basov :
$`{\displaystyle \frac{m^{}(\omega )}{m_\mathrm{b}}}`$ $`=`$ $`{\displaystyle \frac{\omega _{\mathrm{pn}}^2}{4\pi \omega }}\mathrm{}\left[{\displaystyle \frac{1}{\sigma (\omega )}}\right],`$ (24)
$`\omega _{\mathrm{pn}}^2`$ $`=`$ $`\rho __\mathrm{D}={\displaystyle \frac{4\pi \mathrm{e}^2n__\mathrm{D}}{m__\mathrm{D}}},`$ (25)
($`\omega _{\mathrm{pn}}`$ is the normal-state plasma frequency) to estimate the mass renormalization $`m^{}(\omega =0)/m_\mathrm{b}`$, in LSCO and YBCO, about room temperature, for $`x`$ below and above $`x_0^{}`$. The obtained values are presented as black empty circles in the top panels of Fig. 5. The results for $`m^{}/m_\mathrm{b}`$ are close to the values (presented in red-circles) obtained for $`m__\mathrm{D}/m_e`$ for $`x<x_0^{}`$, at low $`T`$ and $`\omega `$.
As is demonstrated in the bottom row of Fig. 6, for results of Puchkov et al. Puchkov for UD cuprates, such an estimate of $`m^{}/m_\mathrm{b}`$ involves an extrapolation to $`\omega =0`$ from the mid-IR range, reflecting again a dominant contribution of QE’s, both for $`x<x_0^{}`$ and $`x>x_0^{}`$. Thus the agreement with the effective mass of the QE’s on the nodal FS arcs, which are dominant for $`x<x_0^{}`$ at low $`T`$, is expected. This mass renormalization reflects the omission of the contribution of $`t`$ processes Ashk03 (which, unlike $`t^{}`$ processes Ashk03 do compete with the effect of AF exchange). Thus the $`m^{}(\omega )/m_\mathrm{b}`$ ratio is increased for $`\omega `$ smaller than the energy effect of AF exchange Anisimov1 ; Munoz , as is seen in Fig. 6. This ratio is greater in LSCO than in YBCO (see Fig. 5) because of the smaller $`t^{}/t`$ ratio for LSCO Andersen .
### IV.6 In–gap states
#### IV.6.1 Low-$`T`$ unpaired carriers and their nature
As was discussed after Eq. (7), in the UD regime, where the $`\overline{ϵ}^p(𝐤)^2`$ term in Eq. (6) is considerably smaller than the $`\mathrm{\Delta }^p(𝐤)^2`$ term at $`T0`$, the Bogoliubov transformation dictates Ashk04 an approximate half filling of the stripon band. Such half filling (namely $`n^p={\scriptscriptstyle \frac{1}{2}}`$), corresponds by the TEP results Ashk01 to a lightly OD stoichiometry of $`x=x_c0.19`$. Thus, it was concluded Ashk04 that only a part of the stripons are paired at $`T0`$, in the UD regime.
The occupation of stripon states can be regarded as consisting of $`n^p`$ stripon “particles”, and $`(1n^p)`$ stripon “holes”, per stripon state. The $`T0`$ stripon-pairing scenario, within the UD regime, can be approximately described as a situation where each paired stripon “hole” is “coupled” with a stripon “particle” (thus yielding lower and upper Bogoliubov bands, each of approximately equal contributions of stripon “holes” or “particles”). Within this scenario, there are unpaired states of stripon “holes”, which are not “coupled” to stripon “particles”, and vice versa.
In p-type cuprates, a linear approximation, expressing the dependence of the number/state of stripon “holes” on $`x`$, yields: $`1n^px/2x_c`$. Since the $`T0`$ number/state of paired stripon “holes” drops from $`1n^p`$, at $`x=x_c`$, to zero at $`x=x_0`$, a linear approximation in $`x`$ would yield for it: $`(1n^p)(xx_0)/(x_cx_0)`$. Thus, the $`T0`$ number/state of unpaired stripon “holes” would then be approximated as: $`(1n^p)[1(xx_0)/(x_cx_0)]=(1n^p)(x_cx)/(x_cx_0)`$.
As was shown in Eq. (16), and the discussion following it, the major contribution to $`\sigma _{ab}(\omega )`$, for low $`T`$ and $`\omega `$, in the OPT and lightly UD regime (corresponding to $`x_c>x\mathrm{}>x_0^{\prime \prime }0.13`$ Basov5 ), is due to stripon “holes”, both above and below $`T_c`$. Thus (at low-$`\omega `$ in p-type cuprates in this regime), a density of approximately $`n_{\mathrm{con}}(1n^p)`$ carriers exists in $`\sigma _{ab}(\omega )`$ just above $`T_{\mathrm{pair}}`$. At $`T0`$, a part $`n_{\mathrm{prd}}`$ of this density is of paired carriers, while the other part, $`n_{\mathrm{unp}}`$, is of carriers which remain unpaired, and become localized below a temperature $`T_{\mathrm{loc}}`$. Using the above linear interpolation in $`x`$, one can express:
$`n_{\mathrm{con}}`$ $``$ $`1n^p{\displaystyle \frac{x}{2x_c}},`$
$`n_{\mathrm{prd}}`$ $``$ $`n_{\mathrm{con}}\left({\displaystyle \frac{xx_0}{x_cx_0}}\right),`$
$`n_{\mathrm{unp}}`$ $``$ $`n_{\mathrm{con}}\left({\displaystyle \frac{x_cx}{x_cx_0}}\right),`$ (26)
$`0.13`$ $``$ $`x_0^{\prime \prime }\mathrm{}<x<x_c0.19.`$
As was discussed above, the QE contribution to the Drude term in $`\sigma _{ab}(\omega )`$ (as carriers’ density $`n_{\mathrm{con}}`$ above $`T_{\mathrm{pair}}`$, and $`n_{\mathrm{unp}}`$ below it), and to the $`\delta `$-function superfluid term in it (as carriers’ density $`n_{\mathrm{prd}}`$), is growing when $`x`$ is decreased. The relative QE contribution to the Drude term above $`T_{\mathrm{pair}}`$ is growing when $`T`$ is increased. The low-$`T`$ unpaired carriers (of density $`n_{\mathrm{unp}}`$) become dominated by QE’s on the nodal FS arcs, for $`x<x_0^{}0.09`$. The crossover between the regimes of stripon- and QE-dominated low-$`T`$ transport occurs for $`x_0^{}<x<x_0^{\prime \prime }`$ Basov5 .
The observed behavior Tanner1 ; Lupi1 ; Lupi2 ; Kim of the low-$`\omega `$ $`\sigma _{ab}(\omega )`$ below $`T_{\mathrm{pair}}`$ (both in the SC and PG states) confirms the above predictions. This is viewed in the effect of the $`n_{\mathrm{unp}}`$ unpaired carriers, appearing as a low frequency Drude-like term Tanner1 ; Lupi1 ; Lupi2 ; Kim (see Figs. 4, 7, and 8). The width of this Drude term is smaller than that of the normal-state Drude term Tanner1 , and it is further decreased Tanner1 ; Lupi2 when the temperature is decreased, which will be shown below to be a consequence of the GTC. This Drude term turns into a low-$`\omega `$ peak in $`\sigma _{ab}(\omega )`$ for $`T<T_{\mathrm{loc}}`$ Tanner1 ; Lupi2 , as is expected for localized carriers. Also the spectral weight within this Drude-like term agrees with the trend predicted in Eq. (26) for $`n_{\mathrm{unp}}`$, being small for the OPT stoichiometry Tanner1 , where $`x`$ is only a little below $`x_c`$, and increasing Lupi2 as $`x`$ is decreased within the UD regime (see below).
#### IV.6.2 The contribution of $`𝐣^{qp}(\omega )`$
Another contribution to $`\sigma _{ab}(\omega )`$, which is in agreement with the GTC predictions, is the one due to transitions between QE and stripon–svivon states Lupi1 ; Lupi2 ; Kim ; Carbotte2 ; Timusk , and its evolution with $`T`$ for $`0<T<T_{\mathrm{pair}}`$. Since the stripe-like inhomogeneities are generated by the Bose condensation of the svivons, and their structure is determined by the details of the of the svivon spectrum, around its energy minimum (see Fig. 2), attempts to interpret the structure contributed to $`\sigma _{ab}(\omega )`$ as a signature of stripes Lupi1 ; Lupi2 are consistent with the present approach. Since the interval of svivon energies involved is between $`+\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})`$ and $`\overline{ϵ}^\zeta (𝐤_{\mathrm{min}})`$, and thus its width is equal by Eq. (9) to the resonance-mode energy, attempts to interpret the contributed optical structure in terms of spin excitations, and specifically the resonance mode Carbotte2 ; Timusk , are also consistent with the present approach.
The energies involved in transitions between QE and stripon–svivon states include the energy difference between an occupied QE, and an unoccupied stripon state, or vice versa, plus or minus a svivon energy. As was discussed above, the stripon bandwidth is $`0.02`$eV, in the normal state, and smaller in a pairing state due to the nature of the expression for Bogoliubov quasiparticle energies, in Eq. (6). The QE’s, on the other hand, have a wide energy spectrum around $`E__\mathrm{F}`$. The svivon spectrum is sketched in Fig. 2; its bottom is $`0.02`$eV below zero, and its top is few tenths of an eV above zero. So these transitions contribute to $`\sigma _{ab}(\omega )`$ a wide, almost featureless, spectrum, forming a major part of its mid-IR background Tanner1 (see Fig. 4). An important role in these transitions is played by svivons around their energy minimum at $`𝐤_0`$ where, by Eq. (1), the $`\mathrm{cosh}(\xi _𝐤)`$ and $`\mathrm{sinh}(\xi _𝐤)`$ factors appearing in the scattering Hamiltonian in Eq. (2) are large. As was discussed above, they are excited during transitions between stripons and QE’s around the antinodal points.
When the stripon and QE pairing gaps open below $`T_{\mathrm{pair}}`$, the width of svivon states around their minimum at $`𝐤_0`$ decreases, as was discussed before Eq. (9). In an analogous manner to the spectroscopic results in Eqs. (8) and (9), one gets that transitions between QE states on one side of the pairing gap, and stripon–svivon states on its other side, result in a contribution to $`\sigma _{ab}(\omega )`$ around the gap-edge energy. This contribution narrows down below $`T_c`$ to a peak of the width $`\mathrm{}>W_{\mathrm{peak}}`$, given in Eq. (9), which equals the resonance-mode energy. Note, however, that in the heavily OD regime, the dispersion in the stripon Bogoliubov band becomes larger than $`W_{\mathrm{peak}}`$, resulting in the smearing of the peak due to the $`𝐤`$-integration taking place when the optical spectrum is derived. Such a peak of the predicted width has been observed by Hwang et al. Timusk , over a wide range of doping, and was found to “disappear” in the heavily OD regime, which, as predicted above, is due to smearing.
#### IV.6.3 Evolution of the states with doping
The existence of unpaired stripon states within the pairing gap, for $`x<x_c`$, causes (as was discussed above) an increase in the width of this peak above $`W_{\mathrm{peak}}`$ (as is noticed in Ref. Timusk in the UD regime). The inclusion of transitions between QE and unpaired stripon–svivon states results in a wider peak in $`\sigma _{ab}(\omega )`$, centered within the energy range of the optical “gap” \[which, as will be discussed below, does not appear as a gap in $`\sigma _{ab}(\omega )`$\]. Furthermore, as was discussed above, the unpaired stripons \[see Eq. (26)\] introduce also a Drude-like term, turning for $`T<T_{\mathrm{loc}}`$ into a peak in $`\sigma _{ab}(\omega )`$, which merges with the one due to transitions between QE and stripon–svivon states into one peak, within the range of the optical “gap”. The appearance of this merged peak could be regarded as a signature of the glassy (checkerboard) structure intertwined with the localization of the stripon carriers, and the creation of an associated gap.
Measurements of the evolution of this structure in $`\sigma _{ab}(\omega )`$, as $`T`$ is lowered through $`T^{}`$ and $`T_c`$, were presented in Refs. Lupi1 ; Lupi2 ; Kim . The results of Lupi et al. Lupi1 on single-layer BSCO are presented in Fig. 7. They show a Drude term at very low energies, due to the density $`n_{\mathrm{unp}}`$ of unpaired carriers below $`T_{\mathrm{pair}}`$ given in Eq. (26). They also show a term due to transitions between QE and stripon–svivon states, which is moved from the mid-IR range, and narrows down, as $`T`$ is lowered, into a peak overlapping energies within the range of the optical gap, as is suggested above.
The results of Lucarelli et al. Lupi2 on LSCO for eight doping levels in the range $`0<x<0.26`$ are presented in Fig. 8. They show a Drude normal behavior for $`x=0.26`$, while for $`x=0.19`$, 0.15, they show the evolution of a peak due to transitions between QE and stripon–svivon states, similarly to the one shown in Fig. 7 for BSCO Lupi1 . A low-energy Drude term, due unpaired carriers below $`T_{\mathrm{pair}}`$ (as in Fig. 7) is observed in Fig. 8 for $`x=0.15`$. Sharp peaks due to phonon modes are observed too. For $`x=0.12`$, this Drude term turns (due to localization in the glassy structure) into a peak, merging (as was suggested above) with the other peak (due to transitions between QE and stripon–svivon states), and the center of the combined peak shifts to lower energy. A similar low-$`T`$ merged peak, within the SC optical gap, is observed also in Fig. 4 Tanner1 .
For $`x=0.07`$ and $`x=0.05`$ Lupi2 this peak in Fig. 8 shifts to a higher energy, with the decrease of $`x`$, which is consistent with the increase of the localization gap. This peak turns into a wider background for $`x=0.03`$ Lupi2 , consistently with the change of the structure into that of diagonal stripes, and it almost disappears for $`x=0`$. Padilla et al. Basov6 studied the evolution of this structure, and the phonon modes in it, in the range $`0x0.08`$, where $`x`$ passes through $`x_00.05`$. Lee et al. Basov4 found a similar behavior in the very low doping regime in YBCO.
Kim et al. Kim studied LSCO for $`x=0.11`$, 0.09, 0,07, 0.063, and also obtained in this regime a combined peak, as in Fig. 8 (interpreted here to be related to the glassy structure), which is shifting to a higher energy, with the decrease of $`x`$. At $`T=300`$K Kim they extrapolated a low-energy Drude term, due to the contribution of QE’s on the nodal FS arcs, discussed above Basov5 . This Drude term is of a different nature than that due to the stripon contribution to $`n_{\mathrm{unp}}`$ \[given in Eq. (26)\], which is observed at higher values of $`x`$.
## V Optical Conductivity in the $`𝒄`$ Direction
Optical conductivity out of the CuO<sub>2</sub> planes, and specifically in the $`c`$ direction, cannot be discussed within the one-band Hamiltonian of Eq. (19), and inter-planar orbitals are involved. The normal-state DC conductivity in the $`c`$ direction is coherent in cases like OD YBCO Homes4 , due to the chains orbitals, and incoherent in other cases. SC in the $`c`$ direction has been often attributed Shibauchi ; Basov1 ; Basov2 to Josephson tunneling between the CuO<sub>2</sub> planes Ambegaokar ; Lawrence . This scenario is not distinct from the description of the $`c`$-direction SC as bulk SC within the dirty limit. Both scenarios give the same temperature dependence of the $`c`$-direction penetration depth Shibauchi ; Basov1 .
The inter-planar layers may be too thin to be described as “metallic” or “insulating”. Since the orbitals of the SC electrons of the CuO<sub>2</sub> planes are hybridized with inter-planar orbitals, the Cooper pairs of the planes do penetrate the inter-planar layers, and in cases that these layers contribute many carriers around $`E__\mathrm{F}`$, proximity-induced pairing is expected to exist there Kresin . Upon crossing to the OD regime Homes4 , coherence is $`c`$ direction can be established in the normal state, and the dirty-limit scenario looks more natural than the Josephson scenario, but still the $`c`$-direction SC features are not significantly altered.
The SC transition is followed by the establishment of a $`c`$-direction SC plasmon mode Marel3 ; Dulic2 ; Dordevic1 , which is referred to, within the Josephson tunneling scenario, as a Josephson plasma resonance. But it is also expected within the GTC pairing mechanism, discussed above Ashk03 , due to the hybridization of orbitals of the inter-planar layers with the CuO<sub>2</sub>-plane QE states, resulting in some inter-plane pair hopping. Effects of the SC transition on the $`c`$-direction optical conductivity include a similar “violation” of the f–sum rule Basov3 ; Marel4 , as in the $`ab`$-plane optical conductivity, due to transfer of spectral weight from energies $`\mathrm{}>2`$eV, which was explained above within the GTC.
Another $`c`$-direction optical effect is the transfer of spectral weight from the SC gap to the mid- or far-IR range, which has been identified Marel3 ; Marel2 ; Dulic1 ; Marel4 ; Marel5 ; Dordevic2 as the signature of a $`c`$-axis collective mode. In bilayer structures it was attributed to a transverse out-of-phase bilayer plasmon, related to “excitons” first considered by Legget Legget , due to relative phase fluctuations of the condensates formed in two different bands. The observation of a similar effect in monolayer LSCO Marel4 suggests that the LaSrO layers, where the doped atoms reside, may introduce enough states close to $`E__\mathrm{F}`$ to provide the second band needed for such transverse plasmons to be formed. This collective mode was found Marel4 ; Dordevic2 to be well defined already in the PG state. This is consistent with the GTC under which the pairs already exist in this state, and even though they lack long-range phase coherence, short-range effects over the distance between close planes/layers, required for this mode to exist, are expected to be present.
## VI Optical Scattering Rates
Using the extended Drude model Puchkov ; Basov , the optical scattering rate can be expressed as \[see Eq. (25)\]:
$$\frac{1}{\tau (\omega )}=\frac{\omega _{\mathrm{pn}}^2}{4\pi }\mathrm{}\left[\frac{1}{\sigma (\omega )}\right]$$
(27)
(thus the carriers density dependence of the conductivity is eliminated to obtain the scattering rate). The $`ab`$-plane and $`c`$-direction scattering rates, $`\tau _{ab}(\omega )^1`$ and $`\tau _c(\omega )^1`$, behave differently in the cuprates. Quantum criticality Marel6 close to the QCP in the normal state (see Fig. 3) results in the existence of a critical region where only one energy scale, which is the temperature, exists. Consequently an MFL-type behavior Varma is obtained there (thus linearity in $`\omega `$) for $`\tau _{ab}(\omega )^1`$ and related quantities Timusk2 , and further optical quantities behave critically Marel6 .
As was shown in Eq. (16), and the discussion following it, the conductivity in the $`ab`$ plane, $`\sigma _{ab}(\omega )`$, is dominantly stripon-like for low $`T`$ and $`\omega `$ for $`x>x_0^{\prime \prime }`$ \[see Eq. (26)\]. Consequently, $`\tau _{ab}(\omega )^1`$ is determined in this regime by the scattering of stripons, through $`^{}`$ \[see Eq. (2)\] into QE and svivon states Ashk01 ; Ashk03 . As was discussed before Eq. (9), the main contribution to such scattering comes Ashk04 from QE states around the antinodal points, and svivon states around their energy minimum at $`𝐤_0`$ (see Fig. 2), where, by Eq. (1), the $`\mathrm{cosh}(\xi _𝐤)`$ and $`\mathrm{sinh}(\xi _𝐤)`$ factors are large. Below $`T_{\mathrm{pair}}`$ a QE pairing gap opens around the antinodal points, and as was discussed above, all the QE states there either become paired, or have a localization gap, for $`T0`$. This results in a reduction, below $`T_{\mathrm{pair}}`$, of $`\tau _{ab}(\omega )^1`$, for $`\omega `$ within the pairing gap. This reduction becomes drastic at low $`T`$, when the width of the gap-edge peak becomes small (because of the exclusion of such scattering), and the QE gap approaches its $`T0`$ value.
Since the width of the Drude term in $`\sigma _{ab}(\omega )`$ is $`\tau _{ab}^1`$, its decrease below $`T_{\mathrm{pair}}`$ results in a smaller Drude width for the unpaired carriers there \[see Eq. (26) and the above discussion\], than for the carriers above $`T_{\mathrm{pair}}`$ (see Fig. 4). Furthermore, this Drude width is expected to decrease with decreasing temperature, in agreement with experiment Tanner1 ; Lupi2 . As was discussed above, the relative contribution of QE’s to the Drude term is growing when $`T`$ is increasing above $`T_{\mathrm{pair}}`$. This results in an increase in the width $`\tau _{ab}^1`$ of the Drude term with $`T`$ above $`T_{\mathrm{pair}}`$, in agreement with experiment Tanner1 (see Fig. 4).
For $`x<x_0^{\prime \prime }`$, the nature of the unpaired carriers below $`T_{\mathrm{pair}}`$ crosses over from being dominantly stripon like, to being dominated by QE’s on the nodal FS arcs. However, this does not result in an increase in their scattering rate $`\tau _{ab}^1`$ (and thus Drude width), because (unlike the antinodal QE’s, which contribute significantly above $`T_{\mathrm{pair}}`$) these QE’s are not scattered to stripon–svivon states, around the svivons’ energy minimum at $`𝐤_0`$ (see Fig. 2), and the contribution of other svivons to their scattering is much smaller (as was discussed above). And indeed, the width of the Drude term in the heavily UD regime was found Basov5 to be pretty small and even fairly separated from the mid-IR term, as was discussed above.
A consequence of the existence of unpaired carriers within the gap is that the drop in $`\tau _{ab}^1`$ below $`T_{\mathrm{pair}}`$ is not followed by an equal relative drop in the effective density of carriers. Since these carriers contribute to DC conduction for $`T>T_{\mathrm{loc}}`$, the expected effect is an increase in the $`ab`$-plane DC conductivity above $`T_{\mathrm{loc}}`$ in the PG state. This explains the observed Takagi sublinear $`T`$-dependence of the DC resistivity for $`T_{\mathrm{loc}}<T<T^{}`$ in the UD regime. Also is explained the increase in this effect on the resistivity Takagi when $`x`$ is decreased, and thus $`T^{}`$ and the PG size are increased.
The $`c`$-direction conductivity $`\sigma _c(\omega )`$ is, on the other hand, determined by the QE states (with hybridization taking place between planar states and inter-planar orbitals). Thus $`\tau _c(\omega )^1`$ is determined by scattering of QE’s to stripon–svivon states, persisting below $`T_c`$ (and thus resulting in the wide hump, as was discussed above), and by processes within the inter-planar layers, which are unrelated to the pairing process, and it is not expected to vary significantly below $`T_{\mathrm{pair}}`$. Since and the number of unpaired QE’s drops below $`T_{\mathrm{pair}}`$ (while $`\tau _c^1`$ does not vary significantly), a drop in $`\sigma _c(\omega )`$ within the gap is occurring below $`T_{\mathrm{pair}}`$. The contribution to it from pairs is expected to be transferred to the $`\delta `$-function term in the SC state, and to higher frequencies in the PG state, where the pairs are localized.
The contribution of the QE localization gaps to the PG (discussed above), especially in the heavily UD regime, also contributes to the drop in $`\sigma _c(\omega )`$ within the gap. The spectral weight corresponding to these localized states is not transferred to the $`\delta `$-function term, when the localization gaps (approximately) exceed the pairing gaps, and this applies to most of the QE’s when $`xx_0`$.
These predictions are confirmed by experiment. A sharp drop has been observed Bonn ; Puchkov in $`\tau _{ab}(\omega )^1`$ of the quasiparticles below $`T_c`$, as can be seen in the top row of Fig. 6. Consequently a clean-limit treatment applies in the SC state within the $`ab`$ plane, while $`\tau _{ab}(\omega )^1`$ above $`T_c`$ corresponds to the intermediate scattering regime (see below) Homes5 . Also, as can be seen in Fig. 6, the reduction in $`\tau _{ab}(\omega )^1`$, for $`\omega `$ within the gap, starts in the PG state (thus below $`T^{}=T_{\mathrm{pair}}`$) Puchkov ; Basov .
On the other hand, as can be seen in Fig. 9, a gap-like depression has been observed below $`T^{}`$ Homes3 ; Puchkov in $`\sigma _c(\omega )`$ within the PG energy range, with the spectral weight from it transferred to higher energies (above $`T_c`$). The similarity, shown in the inset in Fig. 9, between the gap-like behavior in $`\sigma _c(\omega 0)`$, and the gap observed in Knight-shift results, is consistent with the GTC prediction that this gap is in the spectrum of spin-carrying QE-like carriers.
## VII Homes’ Law
Homes et al. Homes1 have demonstrated that, in many SC’s, the optical quantity $`\rho _s=\omega _{\mathrm{ps}}^2`$ \[based on a similar expression as in Eq. (4)\], obeys approximately (thus when presented on a log–log scale) the relation
$$\rho _s35\sigma _{_{\mathrm{DC}}}(T_c)T_c$$
(28)
(both sides in the equation possess the same units). This relation, known as Homes’ law, is obeyed in the cuprates in almost the entire doping regime, both in the $`ab`$-plane, and in the $`c`$-direction, and in other SC’s which behave according to dirty-limit approximations. In a subsequent paper, Homes et al. Homes5 demonstrated that this law should be valid in the dirty-limit, and in the intermediate-scattering regime, when it applies to transport in the normal state. The reason is that the $`\rho _s/8`$ spectral weight, which condenses into the superfluid $`\delta `$-funcion term in $`\sigma (\omega )`$, approximately scales then as the product of $`\sigma _{_{\mathrm{DC}}}(T_c)`$ and the SC gap, which approximately scales as $`T_c`$. Thus, in the regime of validity of Homes’ law, it could be expressed as:
$$k__\mathrm{B}T_c\mathrm{}<\mathrm{}\tau ^1(T_c),$$
(29)
As was discussed above, $`\sigma _{ab}(\omega )`$ in the cuprates corresponds to the intermediate-scattering regime in the normal state, with most of the Drude carriers turning into a superfluid. But the drop in $`\tau _{ab}(\omega )`$ below $`T_c`$ (seen in Fig. 6), predicted by the GTC, turns them into clean-limit SC’s. The $`c`$-direction conductivity was discussed above as being described by the dirty-limit both above and below $`T_c`$, and also through the Josephson tunneling scenario, which was shown Homes1 ; Homes4 to result in Eq. (28) too. Thus, the GTC is consistent with the applicability of Homes’ law in the cuprates, both in the $`ab`$-plane, and the $`c`$-direction, and it could be expressed there as:
$$k__\mathrm{B}T_c(\mathrm{cuprates})\mathrm{}\tau ^1(T_c).$$
(30)
This expression is the consequence of the MFL behavior Varma , due to quantum criticality, and thus Marel6 the existence of one energy scale, which is the temperature. The coexistence of Homes’ law with Uemura’s law \[see Eq. (4)\] in the UD regime implies, by Eq. (28), that $`\sigma _{_{\mathrm{DC}}}(T_c)`$ does not vary much with $`x`$ there (in agreement with experiment Takagi , except when localization starts above $`T_c`$), which was shown above to be the consequence of the decrease in $`\tau _{ab}^1`$ \[and by Eq. (30), also in $`T_c`$\] in the PG state.
Eq. (30) connects high $`T_c`$ in the cuprates with a high scattering rate, which lead also to an MIT. The existence of SC close to the boundary of MIT’s has been pointed out by Osofsky et al. Osofsky . Zaanen Zaanen suggested that Homes’ law implies that $`T_c`$ in the cuprates is determined, approximately by the condition that the scattering rate in the normal state is becoming at $`T_c`$ as high as is permitted by the laws of quantum physics. Here it was shown that the high scattering rate is strongly decreased below $`T_c`$, since it is determined by the same interactions which determine pairing.
## VIII Conclusions
The success of the GTC in providing the understanding of a variety of puzzling optical properties of the cuprates, in addition to its earlier success in explaining many other anomalous properties of these materials, strengthens the point of view that the occurrence of high-$`T_c`$ SC in them requires the proximity of a Mott transition.
It suggests the occurrence of spin-charge separation in the cuprates Anderson , but only due to the existence of dynamical inhomogeneities, which provide quasi-one-dimensional structures. It also predicts an intrinsic origin to the static nanoscale heterogeneity observed in the UD regime Davis2 .
Furthermore, the GTC supports the opinion that the same interactions which play a major role in the determination of the electronic structure of the cuprates, also primarily determine pairing, transport, and other anomalous properties in them.
Even though a complete rigorous first-principles proof on the validity of the GTC is still beyond reach, due to the complexity of the global scheme, and many of its results are still obtained on a qualitative level, the mounting evidence on the global scope of its applicability points very strongly to its validity for the cuprates.
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# Competing orderings in an extended Falicov-Kimball model
## Abstract
We present a Hartree-Fock study of the Falicov-Kimball model extended by both on-site and non-local hybridization. We examine the interplay between excitonic effects and the charge-density wave (CDW) instability known to exist at zero hybridization. It is found that the CDW state remains stable in the presence of finite hybridization; for on-site hybridization the Coulomb interaction nevertheless strongly enhances the excitonic average above its value in the non-interacting system. In contrast, for non-local hybridization, we observe no such enhancement of the excitonic average or a spontaneous on-site hybridization potential. Instead, we find only a significant suppression of the excitonic correlations in the CDW state. A phenomenological Ginzburg-Landau analysis is also provided to understand the interplay.
The Falicov-Kimball model (FKM) describes a tight-binding system of itinerant $`d`$-electrons interacting via on-site Coulomb repulsion $`U`$ with localized $`f`$-electrons of energy $`ϵ_f`$. The FKM was originally introduced as a minimal model of valence transitions in systems such as $`\text{SmB}_6`$ and Ce: by varying the inter-orbtial Coulomb repulsion $`U`$ or the $`f`$-level $`ϵ_f`$, both discontinuous and continuous changes in the distribution of the electrons across the localized and itinerant states were found. FKMoriginal It was soon realised, however, that some overlap between the $`d`$\- and $`f`$-wavefunctions was an essential feature of most systems displaying valence instabilities. LRP81 This “mixing” of the electron wavefunctions may be explicitly introduced by the inclusion of a hybridization potential $`V`$. A variety of methods, including Hartree-Fock, L78 real-space renormalization group HH82 and alloy-analog approximation, BC82 revealed that the hybridization removed the previously observed discontinuous valence transitions. Work on the FKM ceased in the mid-1980s as it became apparent that the Periodic Anderson Model offered a more realistic description of valence transition physics. C86
As interest in the FKM as a model of valence transitions waned, it was adopted as a model of a simple binary alloy. KL/BS86 In the limit of vanishing hybridization the $`f`$-electron occupation at each site is a good quantum number: fixing the $`f`$\- and $`d`$-populations, the ground state is identified as the configuration adopted by the $`f`$-electrons that minimizes the energy of the conduction electrons. In particular, for a bipartite lattice at half-filling and equal concentration of $`d`$\- and $`f`$-electrons, the $`f`$-electrons occupy the sites of one sublattice only, the so-called chequerboard phase. For dimension $`d2`$, this chequerboard charge-density wave (CDW) state obtains for temperatures below a critical temperature $`T_{CDW}`$; above this temperature a disordered phase is realized. For $`d=1`$ the critical temperature is zero. We note that the FKM as a binary alloy has been extensively studied in the case of infinite dimension $`d\mathrm{}`$: the dynamical mean-field theory (DMFT) gives an exact solution in this limit. FZ03
The FKM with hybridization has lately attracted renewed attention due to the investigation of optical properties in this model by Portengen *et al.* POS96 Following closely Leder’s Hartree-Fock (HF) work, L78 they found that the Coulomb repulsion induced an effective on-site hybridization; this effect was sufficiently strong that it persists in the limit of negligible hybridization. In fact, their calculations were performed exclusively in this limit: their solution with non-zero polarization or excitonic average $`d^{}f`$ is indistinguishable from the well-known excitonic insulator (EI) state. EI The “spontaneous” excitonic average was interpreted as evidence of electronic ferroelectricity. Their HF solution, however, assumed a homogeneous ground state for the system; the possibility of a CDW ground state was not considered.
The problem of reconciling the results of Portengen *et al.* with the known CDW instability has only been partially addressed. Since the DMFT equations are no longer exactly solvable for non-zero hybridization potential, Czycholl C99 performed a HF analysis for the $`d\mathrm{}`$ model. It was found that for $`V=0`$ there was no spontaneous excitonic average and that the CDW phase was stable against sufficiently small on-site hybridization. For a given $`U`$ there was a critical hybridization $`V_c(U)`$ under which the CDW phase prevails. Czycholl nevertheless concluded that the the inter-orbtial $`U`$ could strongly renormalize the hybridization, and so could be important in the description of the optical properties of strongly correlated electron systems. Also working in the limit of large spatial dimensions, Zlatić et al. noted that for $`V=0`$ the hybridization susceptibility diverges as $`T0`$, although they concluded that a generalization of the FKM would be required for $`d^{}f0`$ at finite temperatures. ZFLC01
Comparatively little work has been done on this problem in finite dimensions. Farkašovský has used exact-diagonalization and the density matrix renormalization group methods on small one-dimensional systems to rule out the possibility of a spontaneous excitonic average at zero temperature. FNEF By the same methods, Farkašovský has also analyzed the effect of local F97 and non-local F04 hybridization; these works are more concerned with the effect of the hybridization on valence transitions, ignoring the possibility of an excitonic renormalization of the hybridization potentials. Batista and co-workers have claimed that a non-local hybridization stabilizes ferroelectricity in a FKM extended by $`f`$-hopping; Batista Sarasua and Continentino have investigated a similar system. SC04
The FKM extended by hybridization cannot be solved exactly and so it is necessary to use approximate methods to understand the properties of the model. In this paper, we present a HF study of the effect of the hybridization upon the CDW state on a two-dimensional square lattice. The HF approximation is reliable for small temperatures. It tends, however, to overestimate the stability of ordered phases: in particular, the HF result for critical temperature $`T_{CDW}`$ is very likely to be larger than the exact value. Nevertheless, we can reasonably expect that the HF approximation will give at least a qualitatively correct account of the relative stability of ordered phases, even in 2D. The HF is therefore an appropriate tool to study the competition between the EI and CDW phases in the FKM. We consider only $`ϵ_f=0`$ and half-filling (the particle-hole symmetry point) as the CDW state here adopts the simple chequerboard form; for these parameters also the excitonic average takes its maximum as shown in the analysis of Portengen et al. POS96
The FKM Hamiltonian for spinless Fermions is written
$`_{}`$ $`=`$ $`t{\displaystyle \underset{i,j}{}}d_i^{}d_j+ϵ_f{\displaystyle \underset{j}{}}n_j^f+{\displaystyle \underset{i,j}{}}\{V_{ij}d_i^{}f_j+\text{H.c.}\}`$ (1)
$`+U{\displaystyle \underset{j}{}}n_j^dn_j^f.`$
Some overlap between the $`d`$\- and $`f`$-electron wavefunctions is assumed, hence the hybridization term $`V_{ij}`$. The concentration of electrons is fixed at $`1=\frac{1}{N}_j\left\{n_j^f+n_j^d\right\}`$ where $`N`$ is the number of sites. We measure all energies in terms of the $`d`$-electron hopping integral $`t`$.
In our HF decoupling of the Coulomb interaction, we include the possibility of the CDW state by allowing for a periodic modulation of the order parameters:
$`n_j^f`$ $`=`$ $`n^f+\delta _f\mathrm{cos}(𝐐𝐫_j),`$ (2)
$`n_j^d`$ $`=`$ $`n^d+\delta _d\mathrm{cos}(𝐐𝐫_j),`$ (3)
$`f_j^{}d_j`$ $`=`$ $`\mathrm{\Delta }+\mathrm{\Delta }_Q\mathrm{cos}(𝐐𝐫_j).`$ (4)
The nesting vector $`𝐐=(\frac{\pi }{a},\frac{\pi }{a})`$ where $`a`$ is the lattice constant. The order parameter of the CDW state is $`\delta _d`$ and $`\delta _f`$ for the $`d`$\- and $`f`$-electrons respectively. Note that we require $`\text{sgn}(\delta _f)=\text{sgn}(\delta _d)`$. $`\mathrm{\Delta }`$ is the excitonic average; in the absence of an on-site hybridization potential $`V`$, $`\mathrm{\Delta }0`$ indicates the EI phase. When $`V0`$, the EI-normal phase transition is lifted from criticality, in analogy to the ferromagnet-paramagnet transition in an external magnetic field. In this case, we cannot speak of an EI phase, but rather an excitonic enhancement of the hybridization. This will be apparent if $`\mathrm{\Delta }`$ exceeds its value in the $`U=0`$ system. The modulation factor $`\mathrm{\Delta }_Q`$ is included in Eq. (4) for completeness. In the usual HF treatment L78 ; POS96 a homogeneous solution is assumed and so $`\delta _d=\delta _f=\mathrm{\Delta }_Q=0`$ for all values of the Coulomb interaction.
We thus obtain for the HF Hamiltonian
$`_{\text{HF}}`$ $`=`$ $`t{\displaystyle \underset{i,j}{}}d_i^{}d_j+U{\displaystyle \underset{j}{}}(n^f+\delta _f\mathrm{cos}(𝐐𝐫_j))n_j^d`$ (5)
$`+U{\displaystyle \underset{j}{}}(n^d+\delta _d\mathrm{cos}(𝐐𝐫_j))n_j^f`$
$`+{\displaystyle \underset{ij}{}}\{(V_{ij}U[\mathrm{\Delta }+\mathrm{\Delta }_Q\mathrm{cos}(𝐐𝐫_j)]\delta _{ij})d_i^{}f_j`$
$`+\text{H.c.}\}.`$
An important feature of this Hamiltonian is the mean-field renormalization of the $`d`$-$`f`$ hybridization potential by the inter-orbital Coulomb interaction, $`V_{ij}V_{ij}U\left[\mathrm{\Delta }+\mathrm{\Delta }_Q\mathrm{cos}(𝐐𝐫_j)\right]\delta _{ij}`$. The effective on-site hybridization potential introduced by the decoupling of the interaction is responsible for the spontaneous polarization in Portengen et al.’s work. $`_{\text{HF}}`$ is diagonalized by the canonical transform
$$\gamma _𝐤^m=u_𝐤^md_𝐤+v_𝐤^md_{𝐤+𝐐}+\xi _𝐤^mf_𝐤+\zeta _𝐤^mf_{𝐤+𝐐},$$
(6)
where $`m=1,2,3,4`$. The coefficients in Eq. (6) are obtained by solving the associated Bogoliubov-de Gennes (BdG) eigenequations:
$$H_𝐤\mathrm{\Psi }_𝐤^m=E_𝐤^m\mathrm{\Psi }_𝐤^m,$$
(7)
where
$$H_𝐤=\left(\begin{array}{cccc}ϵ_𝐤+Un^f& U\delta _f& V_𝐤U\mathrm{\Delta }& U\mathrm{\Delta }_Q\\ U\delta _f& ϵ_{𝐤+𝐐}+Un^f& U\mathrm{\Delta }_Q& V_{𝐤+𝐐}U\mathrm{\Delta }\\ V_𝐤^{}U\mathrm{\Delta }^{}& U\mathrm{\Delta }_Q^{}& Un^d& U\delta _d\\ U\mathrm{\Delta }_Q^{}& V_{𝐤+𝐐}^{}U\mathrm{\Delta }^{}& U\delta _d& Un^d\end{array}\right)$$
(8)
and
$$\mathrm{\Psi }_𝐤^m=(u_𝐤^m,v_𝐤^m,\xi _𝐤^m,\zeta _𝐤^m)^{Transpose}$$
(9)
Here $`ϵ_𝐤=2t(\mathrm{cos}(k_xa)+\mathrm{cos}(k_ya))`$ is the $`d`$-electron energy dispersion. The self-consistency equations for the HF parameters may be written in terms of the BdG eigenvectors:
$`n^d`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{}_{}{}^{}\{d_𝐤^{}d_𝐤+d_{𝐤+𝐐}^{}d_{𝐤+𝐐}\}`$ (10)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle {}_{}{}^{}\underset{m}{}}\left\{u_𝐤^mu_𝐤^m+v_𝐤^mv_𝐤^m\right\}f(E_𝐤^m).`$
$`\delta _d`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{}_{}{}^{}\{d_{𝐤+𝐐}^{}d_𝐤+d_𝐤^{}d_{𝐤+𝐐}\}`$ (11)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle {}_{}{}^{}\underset{m}{}}\left\{v_𝐤^mu_𝐤^m+u_𝐤^mv_𝐤^m\right\}f(E_𝐤^m).`$
$`n^f`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{}_{}{}^{}\{f_𝐤^{}f_𝐤+f_{𝐤+𝐐}^{}f_{𝐤+𝐐}\}`$ (12)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle {}_{}{}^{}\underset{m}{}}\left\{\xi _𝐤^m\xi _𝐤^m+\zeta _𝐤^m\zeta _𝐤^m\right\}f(E_𝐤^m).`$
$`\delta _f`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{}_{}{}^{}\{f_{𝐤+𝐐}^{}f_𝐤+f_𝐤^{}f_{𝐤+𝐐}\}`$ (13)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle {}_{}{}^{}\underset{m}{}}\left\{\zeta _𝐤^m\xi _𝐤^m+\xi _𝐤^m\zeta _𝐤^m\right\}f(E_𝐤^m).`$
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{}_{}{}^{}\{f_𝐤^{}d_𝐤+f_{𝐤+𝐐}^{}d_{𝐤+𝐐}\}`$ (14)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle {}_{}{}^{}\underset{m}{}}\left\{\xi _𝐤^mu_𝐤^m+\zeta _𝐤^mv_𝐤^m\right\}f(E_𝐤^m).`$
$`\mathrm{\Delta }_Q`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{}_{}{}^{}\{f_{𝐤+𝐐}^{}d_𝐤+f_𝐤^{}d_{𝐤+𝐐}\}`$ (15)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle {}_{}{}^{}\underset{m}{}}\left\{\zeta _𝐤^mu_𝐤^m+\xi _𝐤^mv_𝐤^m\right\}f(E_𝐤^m).`$
The prime denotes summation over half the Brillouin zone; $`f(E)=1/\{1+\mathrm{exp}[\beta (E\mu )]\}`$ is the Fermi distribution function. The chemical potential $`\mu `$ is determined by the condition $`1=N^1_𝐤{}_{}{}^{}_{m}^{}f(E_𝐤^m)`$. We use an exact diagonalization method to solve the BdG equation (7) self-consistently. We start with an initial set of order parameters. By solving Eq. (7), the new order parameters are computed via Eqs. (10) to (15) and are substituted back into Eq. (7). The iteration is repeated until a desired accuracy is achieved.
We first consider the case of an on-site hybridization, $`V_{ij}=V\delta _{ij}`$. In agreement with previous work C99 ; ZFLC01 ; FNEF we find that for vanishing hybridization the CDW phase is always stable against the EI phase and there is no spontaneous excitonic average. The CDW order displayed by the $`V=0`$ ground state will persist in the presence of sufficiently small hybridization potentials, although the transition temperature $`T_{CDW}`$ will be considerably suppressed, see Fig. 1. We find, however, that for finite hybridization the Coulomb interaction will strongly enhance the magnitude of $`\mathrm{\Delta }`$. We plot the variation of $`|\mathrm{\Delta }|`$ with temperature in Fig. 2. Note that since $`\mathrm{\Delta }=|\mathrm{\Delta }|`$ there is a large renormalization of the hybridization potential due to the mean-field decoupling of the Coulomb interaction Eq. (5). Comparing the homogeneous solution without the CDW ordering (the solid line) with the solution with a coexisting CDW ordering (the dotted line), we find a significant suppression of $`|\mathrm{\Delta }|`$ at the onset of the CDW order at $`T0.1t`$. Even within the CDW phase, however, the excitonic enhancement of the on-site hybridization is still apparent as $`|\mathrm{\Delta }|`$ exceeds its value within the non-interacting system (the dashed line in Fig. 2). We do not find any evidence of non-zero $`\mathrm{\Delta }_Q`$.
This competition can be understood from a phenomenological Ginzburg-Landau (GL) theory. The GL free energy density, in terms of both the CDW ($`\delta _d`$) and EI ($`\mathrm{\Delta }`$) order parameters, can be constructed from a symmetry analysis:
$`f`$ $`=`$ $`\alpha _{EI}|\mathrm{\Delta }|^2+\alpha _{CDW}|\delta _d|^2+\beta _1|\mathrm{\Delta }|^4+\beta _2|\delta _d|^2`$ (16)
$`+\beta _3|\mathrm{\Delta }|^2|\delta _d|^2\beta _4(\mathrm{\Delta }^2\delta _d^2+\mathrm{\Delta }^2\delta _d^2),`$
where we assume $`\alpha _{EI}=\alpha _{EI}^{}(TT_{EI}^0)`$ and $`\alpha _{CDW}=\alpha _{CDW}^{}(TT_{CDW}^0)`$. We assume $`\beta _i`$ ($`i=1,2,3,4`$) are all positive. In the region where $`T_{EI}^0>T_{CDW}^0`$Note the second phase transition temperature for the CDW ordering is renormalized by the pre-existing EI order parameter:
$$T_{CDW}=T_{CDW}^0\frac{(\beta _32\beta _4)(T_{EI}^0T_{CDW}^0)}{2\beta _1\alpha _{CDW}^{}/\alpha _{EI}^{}(\beta _32\beta _4)}.$$
(17)
It means that when the EI order parameter pre-exists, the second phase transition temperature for the appearance of the CDW order parameter can be strongly suppressed by the dominant EI order parameter. This explains why the transition temperature $`T_{CDW}`$ decreases with increased hybridization potential $`V`$, as shown in Fig. 1. Below the second phase transition temperature $`T_{CDW}`$, a little algebra yields
$`\mathrm{\Delta }`$ $`=`$ $`\left[{\displaystyle \frac{2\beta _2\alpha _{EI}+\alpha _{CDW}(\beta _32\beta _4)}{4\beta _1\beta _2(\beta _32\beta _4)^2}}\right]^{1/2}`$ (18)
$`\delta _d`$ $`=`$ $`\left[{\displaystyle \frac{2\beta _1\alpha _{CDW}+\alpha _{EI}(\beta _32\beta _4)}{4\beta _1\beta _2(\beta _32\beta _4)^2}}\right]^{1/2}.`$ (19)
Under the condition that the temperature derivative $`\alpha _{CDW}^{}`$ is larger than $`\alpha _{EI}^{}`$, which is indeed confirmed by our numerical results near $`T_{CDW}`$ (see Figs. 1-2), $`\alpha _{CDW}`$ changes more rapidly than $`\alpha _{EI}`$ when the temperature is lowered. Consequently, the CDW order $`\delta _d`$ increases while the EI order $`\mathrm{\Delta }`$ decreases with the lowered temperature.
Despite the popularity of the on-site hybridization potential, this is actually forbidden in real $`d`$-$`f`$ systems by parity considerations. C86 We are instead required to consider a non-local hybridization with inversion symmetry: the simplest such potential is
$`V_{ij}`$ $`=`$ $`t_{df}[\delta _{i_xj_x}(\delta _{i_yj_y+1}`$ (20)
$`\delta _{i_yj_y1})+\delta _{i_yj_y}(\delta _{i_xj_x+1}\delta _{i_xj_x1})],`$
where any site on the lattice is given by $`𝐫_i=i_xa\widehat{𝐱}+i_ya\widehat{𝐲}`$. This is a particularly interesting case as the Coulomb-induced hybridization has a different ($`s`$-wave) symmetry. In the non-interacting system, the (on-site) excitonic average $`\mathrm{\Delta }`$ vanishes; the non-local hybridization potential instead gives rise to an anisotropic excitonic average
$$\mathrm{\Xi }=\mathrm{}\left\{\frac{1}{N}\underset{𝐤}{}\left(\mathrm{sin}(k_xa)+\mathrm{sin}(k_ya)\right)f_𝐤^{}d_𝐤\right\}.$$
(21)
The study of this quantity allows us to assess the effect of the inter-orbital Coulomb repulsion upon the $`d`$-$`f`$ hybridization.
As with the on-site hybridization, we find that for given $`U`$ the CDW phase is suppressed by the presence of the non-local hybridization (see Fig. 3). We do not, however, find any evidence for a Coulomb-induced on-site hybridization when the CDW instability is allowed: for all non-zero $`t_{df}`$ we have $`\mathrm{\Delta }=\mathrm{\Delta }_Q=0`$. Czycholl considered the appearance of an on-site average $`\mathrm{\Delta }`$ to be likely due to the substantial excitonic enhancement of the on-site hybridization potential by the Coulomb interaction. C99 Our results clearly demonstrate that this EI-like scenario, and the consequent formation of an electronic ferroelectric state, is severely compromised by the presence of non-local hybridization.
In Fig. 4 we plot $`\mathrm{\Xi }`$ as function of temperature in both the interacting and non-interacting systems. The onset of CDW order for the given Coulomb values occurs at the point of intersection of the broken lines with the non-interacting (solid) line. Remarkably, for the standard homogeneous solution there is no effect on $`\mathrm{\Xi }`$ due to the Coulomb interaction: the variation of $`\mathrm{\Xi }`$ with temperature exactly follows the curve for the non-interacting system. Within the CDW phase, however, $`\mathrm{\Xi }`$ is suppressed below its value in the non-interacting system. We offer the following explanation for this anomaly: the hybridization potential Eq. (20) connects the A-sublattice $`d`$-orbitals with B-sublattice $`f`$-orbitals and vice versa. Assume that in the CDW state the A-sublattice $`d`$-orbitals have $`n^d>0.5`$ and so the B-sublattice $`f`$-orbitals have $`n^f>0.5`$; clearly A-B sublattice $`d`$-$`f`$ hopping will be suppressed, hence also the reduction in $`\mathrm{\Xi }`$.
In conclusion, we have examined the competition between excitonic and CDW instabilities in the FKM extended by both on-site and non-local hybridization. In both cases, we find that the CDW phase remains stable at low temperatures even in the presence of a finite hybridization. For the local hybridization we find that the Coulomb interaction nevertheless strongly renormalizes the hybridization potential in agreement with previous work. C99 The situation is qualitatively different for the more realistic non-local hybridization: there is no enhancement of the non-local hybridization and the Coulomb interaction does not induce a spontaneous on-site hybridization. Within the CDW phase, the non-local hybridization is suppressed in line with the increasing localization of the $`d`$\- and $`f`$-electrons. The failure of the Coulomb interaction to induce an effective on-site hybridization except when such a term is already present casts significant doubt over the usefulness of Eq. (1) as a minimal model for electronic ferroelectricity. Inter-orbital Coulomb repulsion may nevertheless still be important for understanding optical properties of strongly-correlated electron systems: for example, a recent extension of the FKM by $`f`$-electron hopping offers a plausible scenario where the formation of an exciton BEC gives a spontaneous excitonic average. Batista
One of us (PMRB) acknowledges the hospitality of the Los Alamos National Laboratory where part of this work was carried out. This work was supported by the US DOE (JXZ and ARB).
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# Sector of the 2⁺⁺ mesons: observation of the tensor glueball
## 1 Introduction
A broad isoscalar-tensor resonance in the region of 2000 MeV is seen in various reactions . Recent measurements give:
$`M=2010\pm 25`$MeV, $`\mathrm{\Gamma }=495\pm 35`$MeV in $`p\overline{p}\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ ,
$`M=1980\pm 20`$MeV, $`\mathrm{\Gamma }=520\pm 50`$MeV in $`pppp\pi \pi \pi \pi `$ ,
$`M=2050\pm 30`$MeV, $`\mathrm{\Gamma }=570\pm 70`$MeV in $`\pi ^{}p\varphi \varphi n`$ ;
following them, we denote the broad resonance as $`f_2(2000)`$.
The large width of $`f_2(2000)`$ arouses the suspicion that this state is a tensor glueball. Such an opinion was expressed lately in different publications.
In , Chapter 5.4, it is said that the very broad isoscalar $`2^{++}`$ state observed in the region $`2000`$ MeV with a width of the order of $`400500`$ MeV could well be the trace of a tensor glueball lying on the Pomeron trajectory.
Another argument comes from the analysis of the mass shifts of the $`q\overline{q}`$ tensor mesons (, Section 12). It is stated here that the mass shift between $`f_2(1560)`$ and $`a_2(1700)`$ can not be explained by the mixing of non-strange and strange components in the isoscalar sector: in such a mixing the average mass squared does not change and we should find $`f_2(1750)`$ at a much higher mass. Instead, we observe a shift down in masses of both isoscalar states. Such a phenomenon can be an indication for the presence of a tensor glueball in the mass region 1800-2000 MeV.
In , the following argument is presented: a significant violation of the OZI-rule in the production of tensor mesons with dominant $`s\overline{s}`$ components (reactions $`\pi ^{}pf_2(2120)n`$, $`f_2(2340)n`$, $`f_2(2410)n\varphi \varphi n`$ ) is due to the presence of a broad glueball state $`f_2(2000)`$ in this region, resulting in a noticeable admixture of the glueball component in $`f_2(2120)`$, $`f_2(2340)`$, $`f_2(2410)`$.
The possibility that the broad resonance $`f_2(2000)`$ could be a glueball is discussed also in , Section 10. Here, however, a problem in the identification of $`f_2(2000)`$ as the tensor glueball is stressed. As it is written in Subsection 10.6, of the prediction for the branching fraction of the $`2^+`$ glueball is large if the width is taken to be the 500 MeV fitted to $`f_2(1950)`$. Observed decays to $`\sigma \sigma `$ and $`f_2(1270)\sigma `$ account for $`(10\pm 0.7\pm 3.6)10^4`$ of $`J/\mathrm{\Psi }`$ radiative decays and for a further $`(7\pm 1\pm 2)10^4`$ in $`K^{}\overline{K}^{}`$ decays. If one assumes flavour-blindness for vector-vector final states, the vector-vector contribution increases to $`(16\pm 2\pm 4.5)10^4`$. The total $`2.610^3`$ is still less by a factor 9 than predicted for a glueball; in this is considered as a problem in identifying $`f_2(1950)`$ with the $`2^+`$ glueball.
In it was emphasised that the $`f_2(2000)`$ being superfluous for $`q\overline{q}`$ systematics can be considered as the lowest tensor glueball. A recent re-analysis of the $`\varphi \varphi `$ spectra in the reaction $`\pi ^{}p\varphi \varphi n`$ , the study of the processes $`\gamma \gamma \pi ^+\pi ^{}\pi ^0`$ , $`\gamma \gamma K_SK_S`$ and the analysis of the $`p\overline{p}`$ annihilation in flight $`p\overline{p}\pi \pi ,\eta \eta ,\eta \eta ^{}`$ clarified essentially the status of the ($`J^{PC}=2^{++}`$)-mesons. This allows us to place the $`f_2`$ mesons reliably on the $`(n,M^2)`$-trajectories , where $`n`$ is the radial quantum number of the $`q\overline{q}`$-state. In the present review we discuss the data for $`\gamma \gamma \pi ^+\pi ^{}\pi ^0`$ , $`\gamma \gamma K_SK_S`$ and $`p\overline{p}\pi \pi ,\eta \eta ,\eta \eta ^{}`$ in Section 2.
In (see also ), the known $`q\overline{q}`$-mesons consisting of light quarks ($`q=u,d,s`$) are placed on the $`(n,M^2)`$ trajectories. Trajectories for mesons with various quantum numbers turn out to be linear with a good accuracy. In Section 3 we give a systematisation of tensor mesons, $`f_2`$ and $`a_2`$, on the $`(n,M^2)`$ planes.
The quark states with ($`I=0`$, $`J^{PC}=2^{++}`$) are determined by two flavour components $`n\overline{n}=(u\overline{u}+d\overline{d})/\sqrt{2}`$ and $`s\overline{s}`$ for which two states $`{}_{}{}^{2S+1}L_{J}^{}=^3P_2,^3F_2`$ are possible. Consequently, we have four trajectories on the $`(n,M^2)`$ plane. Generally speaking, the $`f_2`$-states are mixtures of both the flavour components and the $`L=1,3`$ waves. The real situation is, however, such that the lowest trajectory \[$`f_2(1275)`$, $`f_2(1580)`$, $`f_2(1920)`$, $`f_2(2240)`$\] consists of mesons with dominant $`{}_{}{}^{3}P_{2}^{}n\overline{n}`$ components, while the trajectory $`[f_2(1525)`$, $`f_2(1755)`$, $`f_2(2120)`$, $`f_2(2410)]`$ contains mesons with predominantly $`{}_{}{}^{3}P_{2}^{}s\overline{s}`$ components. The $`F`$-trajectories are presently represented by three resonances \[$`f_2(2020)`$, $`f_2(2300)`$\] and \[$`f_2(2340)`$\] with the corresponding dominant $`{}_{}{}^{3}F_{2}^{}n\overline{n}`$ and $`{}_{}{}^{3}F_{2}^{}s\overline{s}`$ states. In , it is shown that the broad resonance $`f_2(2000)`$ is not part of those states placed on the $`(n,M^2)`$ trajectories. In the region of 2000 MeV three $`n\overline{n}`$-dominant resonances, $`f_2(1920)`$, $`f_2(2000)`$ and $`f_2(2020)`$, are seen, while on the $`(n,M^2)`$-trajectories there are only two vacant places. This means that one state is obviously superfluous from the point of view of the $`q\overline{q}`$-systematics, i.e. it has to be considered as exotics.
There exist various arguments in favour of the assumption that $`f_2(2000)`$ is generated by a glueball. Still, it can not be a pure gluonic $`f_2(2000)`$ state: it follows from the $`1/N`$ expansion rules that the gluonic state $`(q\overline{q})`$ mixes with quarkonium systems $`(gg)`$ without suppression. The problem of the mixing of $`(gg)`$ and $`(q\overline{q})`$ systems is discussed in Section 4, where we present also the relations between decay constants of a glueball into two pseudoscalar mesons $`glueballPP`$ and into two vector mesons $`glueballVV`$.
In Section 5 we demonstrate that just $`f_2(2000)`$ is the glueball.
In it was pointed out that an exotic state has to be broad. Indeed, if an exotic resonance occurs among the standard $`q\overline{q}`$-states, they overlap, and their mixing becomes possible due to large distance transition: the $`resonance(1)realmesonsresonance(2)`$. Owing to these transitions, an exotic meson accumulates the widths of its neighbouring resonances. The phenomenon of the accumulation of widths was studied in the scalar sector near 1500 MeV . In , a model of mixing of the gluonium $`gg`$ with the neighbouring quarkonium states was considered. It was demonstrated that, as a result of mixing, it is precisely the gluonium state which transforms into a broad resonance. The reason is that $`q\overline{q}`$ states being orthogonal to each other mix weakly, while the gluonium mixes with neighbouring $`q\overline{q}`$ states without suppression. Therefore, the gluonium ”dives” more rapidly into the complex $`M`$-plane. The mixing of states is always accompanied by a repulsion of the corresponding poles: when poles are in the complex $`M`$-plane at approximately the same Re $`M`$, this repulsion results in ”sinking” one of them into the region of large and negative Im $`M`$ and ”pulling” others to the real $`M`$-axis.
Hence, the large width of $`f_2(2000)`$ can indicate that this state is an exotic one. Strictly speaking, this fact is not sufficient to prove that $`f_2(2000)`$ is a glueball. At the moment a variety of versions for exotic mesons is discussed; these are $`q\overline{q}g`$ hybrids as well as multiquark states (see, e.g. an references therein). Thus, in order to fix $`f_2(2000)`$ as a glueball, it is of great importance to investigate the decay couplings and prove that they satisfy relations characterising the glueball. The coupling constants for the transitions
$$f_2(1920),f_2(2000),f_2(2020),f_2(2240),f_2(2300)\pi \pi ,\eta \eta ,\eta \eta ^{}$$
are separated in on the basis of a partial wave analysis carried out earlier. The coupling constants obtained in indicate that only the decays $`f_2(2000)\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ obey the relations corresponding to a glueball, while the decay constants for other resonances do not fulfil such conditions. Note that the glueball decay couplings are close to those for the $`SU(3)`$-flavour singlet, but, because of the flavour symmetry violation caused by the strange quark, do not coincide with them exactly.
Let us remind that there are two more arguments in favour of the glueball nature of $`f_2(2000)`$:
(i) the Pomeron trajectory, determined on the basis of data on high-energy hadron decays (see, e.g., ), indicates that a tensor glueball has to have a mass of the order of $`1.72.3GeV`$;
(ii) lattice calculations lead to a similar value, $`M_{2^+glueball}2.32.5GeV`$.
We have one, sufficiently general, argument against the interpretation of $`f_2(2000)`$ as an exotic $`q\overline{q}g`$ or $`qq\overline{q}\overline{q}`$ state: the absence of any serious facts confirming their existence. Indeed, if such states existed, we would see a large number of them in the mass region above $`1500`$ MeV. Moreover, we could observe not only exotic states; the number of resonances with ”normal” quantum numbers would also be seriously increased. However, the systematics of quarks on the $`(n,M^2)`$-plane does not reveal such an increase: almost all observed resonances can be interpreted as $`q\overline{q}`$ states (see , Chapter 5). Apparently, Nature does not like coloured multiparticle objects. The same conclusion follows from the systematisation of baryons: experimental data give a much smaller amount of excited states, than calculations in the framework of a three-quark model do. One gets the impression that excited baryons are rather quark–diquark systems (see discussions in .
Owing to the $`1/N`$ expansion rules, the gluonium component is relatively small in the quark state $`f_2`$: its probability is suppressed as $`1/N_c`$. In Section 6 we determine the mixing angle of the $`n\overline{n}`$ and $`s\overline{s}`$ components in the quark $`f_2`$-mesons, making use of the relations between the decay constants $`f_2\pi \pi ,\eta \eta ,\eta \eta ^{}`$. Also, we estimate the possible changes in the mixing angle as a consequence of a gluonium component in the quark state $`f_2`$.
In the Conclusion, we discuss the situation in the glueball sector.
Up to now, two glueball states, the scalar meson $`f_0(12001600)`$ an the tensor $`f_2(2000)`$ are observed. Both states are broad ones, and the coupling constants corresponding to their decays into pseudoscalar mesons (channels $`glueballPP`$) satisfy just the relations characterising the glueball. The next states which are of interest are radial excitations of the scalar and tensor gluonia, and the pseudoscalar glueball. Taking the Pomeron trajectories on the $`(J,M^2)`$-plane as a basis, we predict the masses of excited scalar and tensor glueball states.
## 2 Analysis of the data for tensor mesons
We demonstrate here the results obtained from the data analysis used in performing the systematisation of the $`f_2`$ resonances and extracting the decay couplings $`f_2\pi \pi ,\eta \eta ,\eta \eta ^{}`$.
### 2.1 L3 data on the $`\gamma \gamma \pi ^+\pi ^{}\pi ^0`$ reaction
In this reaction the $`\gamma \gamma `$ channel couples only to states with C=+1 parity; $`3\pi `$ has a negative G-parity. For a $`q\overline{q}`$ system one has $`G=Ce^{i\pi I}`$, so the $`I=1`$ quark-antiquark states are produced only in the $`\gamma \gamma `$ channel. Due to $`C`$-parity conservation in neutral decay modes, only $`f`$-states with ($`J^{PC}=0^{++}`$, $`2^{++}`$, $`4^{++}\mathrm{}`$) are produced in the $`\pi ^+\pi ^{}`$ channel. In the $`\pi ^\pm \pi ^0`$ channel only isovector mesons with $`J^{PC}=1^{}`$, $`3^{}\mathrm{}`$ are produced.
The $`\gamma \gamma `$ mass distribution is dominated by the production of the $`a_2(1320)`$ resonance, see Fig. 1a. One can see a prominent structure in the mass region 1.6-1.8 GeV as well as a possible contribution at the $`a_2(1320)`$ signal.
The $`\pi ^+\pi ^{}`$ mass distribution is shown in Fig. 1b. There are no clear signals in the data coming from well known narrow scalar-isoscalar states $`f_0(980)`$ and $`f_0(1500)`$. Indeed, the partial wave analysis shows very small contributions of these mesons; such decay modes were omitted in the final fit.
A signal coming from $`f_2(1275)\pi ^0`$ is observed at high $`\gamma \gamma `$ energies; this is important to describe the two-pion mass spectrum and angular distributions.
The $`\pi \pi \pi \pi `$ S-wave amplitude has a broad component which covers the mass region from the $`\pi \pi `$ threshold up to 2 GeV. Such a component is introduced in the present analysis and is parametrized in two different ways.
The first parametrisation is taken from . It was introduced to describe the CERN-Munich data on ($`\pi ^{}p\pi ^+\pi ^{}n`$) and the Crystal Barrel data on proton-antiproton annihilation into $`3\pi ^0`$ and $`2\eta \pi ^0`$ channels simultaneously. To simulate a possible s-dependence of the vertex (which can be important for this very broad state), we use the method suggested in .
The second parametrisation was used in . It covers the mass region up to 1.9 GeV and describes, in the framework of the P-vector/K-matrix approach, a much larger number of two- and three-body reactions. To avoid an over-parametrisation of the fit, we vary only the production couplings of the two lowest K-matrix poles.
The main signal in the $`\pi ^\pm \pi ^0`$ mass spectrum is due to the production of $`\rho (770)`$. There is very little structure in the region higher than 1 GeV (see Fig. 1c) and the signal is almost zero at masses above 1.5 GeV. Because of this, neither $`\rho _3(1690)`$, nor $`\rho (1770)`$ has to be introduced in the analysis. A contribution from $`\rho (1450)`$ is found to be useful to describe the data: however, this state is quite broad and possibly simulates a non-resonant two pion production in this channel.
In the $`\pi ^+\pi ^{}\pi ^0`$ spectrum one can see a strong signal coming from $`a_2(1320)`$; the characteristics of this resonance were defined with high precision by the VES collaboration . It is not surprising that the $`\gamma \gamma 3\pi `$ data are dominated by the production of the $`a_2(1320)`$ state, since this resonance has the highest spin in the mass region below 1.6 GeV (the $`\gamma \gamma `$ cross section is proportional to $`(2J+1)`$) and it is a ground $`q\overline{q}`$ state with the radial quantum number $`n=1`$. Indeed, the $`\gamma \gamma resonance`$ production amplitude is a convolution of the photon and the quark-antiquark resonance wave functions . This provides a suppression of nearly an order for the production of the radially excited states $`(n2)`$ . Nevertheless, there is a manifest contribution of the higher tensor state. While $`a_2(1320)`$ decays practically only into the $`\rho (770)\pi `$ channel, the second tensor state decays almost equally to $`\rho (770)\pi `$ and $`f_2(1275)\pi `$. This fact allows us to identify this state with a good accuracy.
The solution reveals quite a large contribution coming from the $`0^+`$ partial wave decaying into $`f_0\pi `$. There is, however, a problem in distinguishing it from the experimental background: such a decay, giving S-wave amplitudes in all decay channels, provides very smooth structures in mass distributions and, moreover, these amplitudes do not interfere with the tensor amplitude.
The contribution of the $`2^+`$ state is found to be quite small when fitted to the $`\pi _2(1670)`$ state (see Table 1). If it is fitted with free Breit-Wigner parameters, it is always optimised at higher masses ($`1870`$ MeV).
The $`\pi ^+\pi ^{}\pi ^0`$ spectrum and the contributions of resonances with different $`J^{PC}`$ in the final solution are shown in Fig. 2. Masses, total widths, the $`\mathrm{\Gamma }_{\gamma \gamma }`$ partial width and the branching ratio into $`3\pi `$ are listed in Table 1 for the considered resonances.
### 2.2 L3 data on the reaction $`\gamma \gamma K_SK_S`$
Important data for the identification of the tensor mesons were obtained by the L3 collaboration on the reaction $`\gamma \gamma K_SK_S`$. Only states with even spin $`J`$ and positive parities $`P=C=+`$ contribute to two neutral pseudoscalar particles, what reduces the possible partial waves drastically. The $`(2J+1)`$ factor in the cross section favours the dominant production of the tensor states. The $`4^{++}`$ states are produced at high energies $`(M1.9GeV)`$.
The meson states consisting of light quarks $`u`$, $`d`$ $`s`$ form meson nonets: three isospin 1, four isospin $`1/2`$ and two isoscalar states. The isoscalar states can be a mixture of $`n\overline{n}=(u\overline{u}+d\overline{d})/\sqrt{2}`$ and $`s\overline{s}`$ components. The decay of $`(I=0)`$ and $`(I=1,I_3=0)`$-states into two kaons is defined by the production of a new $`s\overline{s}`$ pair (an $`s`$-quark exchange process) and has the following structure for different isospins: $`(u\overline{u}+d\overline{d})/\sqrt{2}K^+K^{}+K^0\overline{K}^0`$ for I=0, and $`(u\overline{u}d\overline{d})/\sqrt{2}K^+K^{}K^0\overline{K}^0`$ for I=1. As a result, a strong destructive interference occurs between the $`f_2(1275)`$ and $`a_2(1320)`$ mesons which is very sensitive to the mixing angle of the isoscalar state. The flavour content of isoscalar-scalar resonances belonging to the same nonet can be written in the form
$`f_2(q\overline{q})=`$ $`n\overline{n}\mathrm{cos}\phi `$ $`+s\overline{s}\mathrm{sin}\phi ,`$ (1)
$`f_2^{}(q\overline{q})=`$ $`n\overline{n}\mathrm{sin}\phi `$ $`+s\overline{s}\mathrm{cos}\phi .`$
Although the production of states with dominant $`s\overline{s}`$ components is suppressed by the smaller $`\gamma \gamma `$ coupling, these states usually have a larger branching ratio for the decay into the $`K_SK_S`$ channel, and can contribute appreciably into the total cross section.
There is no doubt about the nature of tensor resonances below 1600 MeV. The partial wave analysis showed that tensor resonances are produced dominantly in the $`{}_{}{}^{5}S_{2}^{}`$ $`\gamma \gamma `$ state, which was predicted by model calculations . The model gives directly the ratio between $`{}_{}{}^{5}S_{2}^{}`$ and $`{}_{}{}^{1}D_{2}^{}`$ waves and the $`\gamma \gamma `$ couplings. These values can be introduced in the analysis without changing the quality of the description. The data with the lowest tensor states are shown in Fig. 3a.
We have found that the $`0^{++}`$ partial wave can be fitted well in the framework of the P-vector/K-matrix approach. The form of the $`0^{++}`$ contribution in $`\gamma \gamma K\overline{K}`$ follows closely the $`\pi \pi K\overline{K}`$ cross section, which is not surprising: the production coupling to the $`s\overline{s}`$ system is strongly suppressed in both reactions. The values of the coupling constants of the reactions $`\gamma \gamma f_0,a_0`$ agree well with those given in the calculations : in the same way as in the tensor sector, the calculated $`\gamma \gamma `$ couplings can be used directly, not damaging the quality of the description.
A contribution of $`4^{++}`$ states is observed in high energy angular distributions. Some broad and some rather narrow components are seen in this wave. The broad state can be associated with $`n\overline{n}`$ and the narrow one with a $`4^{++}`$ $`s\overline{s}`$ state. The description of the data with tensor, scalar and $`4^{++}`$ states is shown in Fig. 3b.
There is a clear resonance structure in the $`1750`$ MeV region. The angular distribution in this region follows very closely the $`(1\mathrm{cos}^2\mathrm{\Theta })^2`$ shape which corresponds to the $`{}_{}{}^{5}S_{2}^{}\gamma \gamma `$ production of the tensor meson. However, the acceptance decreases rapidly in the forward and backward directions providing a very similar dependence. Due to such a behaviour and to the not too high statistics, the partial wave analysis produces almost the same angular distribution for a tensor state and for a scalar state. Still, the fit using a scalar state fails to reproduce the structure in the 1700-1800 MeV mass region. The description of the data with the best $`\chi ^2`$ is shown in Fig. 3c. The mass was optimised to $`1805\pm 30`$ MeV and the width to $`260\pm 30`$ MeV. With such a mass and width, the $`f_0`$ state can describe the slope in mass distribution above 1800 MeV. If the mass and width of the scalar state are fixed at the BES result $`M=1740`$ MeV, we obtain a description shown in Fig. 3d.
The main problem in a fit with a $`f_0`$ state is that there is no way to reproduce the dip in the 1700 MeV region and the slope above 1800 MeV using any (even very sophisticated) parametrisation of the resonance width. A $`f_0`$ state can interfere only with the $`{}_{}{}^{1}D_{2}^{}`$ component of a $`2^{++}`$ state, and this partial wave is very small in the data. Consequently, $`f_0`$ and $`f_2`$ contributions practically do not interfere and it is impossible to create a dip in the data. Contrary to this, a tensor resonance interfering with the tails of other tensor states naturally produces a dip and a good description of this mass region.
All isoscalar and isotensor states can contribute to the $`\gamma \gamma K_SK_S`$ cross section; this situation offers a very good possibility to study the reaction on the basis of the nonet classification. With $`SU(3)`$ relations imposed, the only parameters to fit the data for the first three states are masses, widths, the mixing angle and $`SU(3)`$ violation factors. We found all masses and widths for the members of the first tensor nonet to be in a very good agreement with PDG.
To describe the second nonet, we fixed parameters for $`f_2(1560)`$ from the Crystal Barrel results and for $`a_2(1700)`$ from the latest L3 analysis of the $`\pi ^+\pi ^{}\pi ^0`$ channel . At a given mixing angle the nonet coupling was calculated to reproduce the $`\pi \pi `$ width of 20-25 MeV for the $`f_2(1560)`$ state. We found a very good description of the data with $`SU(3)`$ relations imposed, similar to the fit described in the previous section. The masses, widths, radii, $`K\overline{K}`$ couplings, mixing angles and partial widths of the states are given in Table 2; the description of the data is shown in Fig. 4.
The $`f_2(1750)`$ state has a mass $`1755\pm 10`$ MeV and a total width $`67\pm 12`$ MeV. Its $`K\overline{K}`$ width is $`23\pm 7`$ MeV; the rest of the width is likely to be defined by the $`K^{}\overline{K}`$ channel. This resonance destructively interferes with the tail of the $`f_2^{}(1525)`$ state, creating a dip in the mass region 1700 MeV. When the sign of the real part of the $`f_2(1750)`$ amplitude changes, this interference becomes positive, producing a clear peak in the data.
A serious problem appears in the description of the data in the framework of the nonet approach, if the peak at 1750 MeV is assumed to be owing to a scalar state. If this state is a nonet partner of one of the known states, e.g. $`f_0(1370)`$ or $`f_0(1500)`$, then the calculated signal is too weak to fit the data. If the $`K\overline{K}`$ coupling of this scalar state is fitted freely, we find that it must be about four times larger than the total width of the resonance. This is due to the $`2J+1`$ suppression factor and to the absence of a positive interference with the tail of $`f_2^{}(1525)`$ which boosts the peak in the case of a tensor state. These are problems additional to those connected with the description of the dip in the 1700 MeV region.
### 2.3 Data for proton-antiproton annihilation in flight $`p\overline{p}\pi \pi ,\eta \eta ,\eta \eta ^{}`$
The $`p\overline{p}`$ annihilation in flight gives information about resonances with $`M>1900`$ MeV. High statistical data taken at antiproton momenta 600, 900, 1150, 1200, 1350, 1525, 1640, 1800 and 1940 MeV/c were used for the analysis $`p\overline{p}\pi ^0\pi ^0`$, $`\eta \eta ,`$ $`\eta \eta ^{}`$ . The combined analysis was performed together with $`\overline{p}p\pi ^+\pi ^{}`$ data obtained with a polarised target . Five tensor states are required to describe the data, $`f_2(1920)`$, $`f_2(2000)`$, $`f_2(2020)`$, $`f_2(2240)`$, $`f_2(2300)`$:
$`\mathrm{Resonance}`$ $`\mathrm{Mass}(\mathrm{MeV})`$ $`\mathrm{Width}(\mathrm{MeV})`$ (2)
$`f_2(1920)`$ $`1920\pm 30`$ $`230\pm 40`$
$`f_2(2000)`$ $`2010\pm 30`$ $`495\pm 35`$
$`f_2(2020)`$ $`2020\pm 30`$ $`275\pm 35`$
$`f_2(2240)`$ $`2240\pm 40`$ $`245\pm 45`$
$`f_2(2300)`$ $`2300\pm 35`$ $`290\pm 50.`$
The description of the data is illustrated by Figs. 5,6,7 and 8. In Fig. 9 we show the cross sections for $`p\overline{p}\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ in $`{}_{}{}^{3}P_{2}^{}\overline{p}p`$ and $`{}_{}{}^{3}F_{2}^{}\overline{p}p`$ waves (dashed and dotted lines) and the total $`(J=2)`$ cross section (solid line) as well as Argand-plots for the $`{}_{}{}^{3}P_{2}^{}`$ and $`{}_{}{}^{3}F_{2}^{}`$ wave amplitudes at invariant masses $`M=1.962`$, $`2.050`$, $`2.100`$, $`2.150`$, $`2.200`$, $`2.260`$, $`2.304`$, $`2.360`$, $`2.410`$ GeV.
The $`\overline{p}p\pi ^0\pi ^0`$, $`\eta \eta `$, $`\eta \eta ^{}`$ amplitudes provide the following ratios for the $`f_2`$ resonance couplings $`g_{\pi ^0\pi ^0}:g_{\eta \eta }:g_{\eta \eta ^{}}`$:
$`f_2(1920)`$ $`1:0.56\pm 0.08:0.41\pm 0.07`$
$`f_2(2000)`$ $`1:0.82\pm 0.09:0.37\pm 0.22`$
$`f_2(2020)`$ $`1:0.70\pm 0.08:0.54\pm 0.18`$
$`f_2(2240)`$ $`1:0.66\pm 0.09:0.40\pm 0.14`$
$`f_2(2300)`$ $`1:0.59\pm 0.09:0.56\pm 0.17.`$
These coupling ratios allow one to estimate the quarkonium-gluonium c content of the $`f_2`$ \- resonances.
## 3 Systematisation of tensor mesons on the $`(n,M^2)`$ trajectories
In (see also ), the known $`q\overline{q}`$-mesons consisting of light quarks ($`q=u,d,s`$) were put on the $`(n,M^2)`$ trajectories, where $`n`$ is the radial quantum number of the $`q\overline{q}`$ system with mass $`M`$. The trajectories for mesons with various quantum numbers turn out to be linear with a good accuracy:
$$M^2=M_0^2+(n1)\mu ^2$$
(4)
where $`\mu ^2=1.2\pm 0.1\mathrm{GeV}^2`$ is a universal slope, and $`M_0`$ is the mass of the lowest state with $`n=1`$.
In Fig. 10a we demonstrate the present status of the $`(n,M^2)`$ trajectories for the $`f_2`$ mesons (i.e. we use the results given by ). To avoid confusion, we list here the experimentally observable masses. First, it concerns the resonances seen in the $`\varphi \varphi `$ spectrum . In the $`\varphi \varphi `$ spectra were re-analysed, taking into account the existence of the broad $`f_2(2000)`$ resonance. As a result, the masses of three relatively narrow resonances are shifted compared to those given in the compilation PDG :
$$f_2(2010)|_{PDG}f_2(2120)\text{[4]},f_2(2300)|_{PDG}f_2(2340)\text{[4]},f_2(2340)|_{PDG}f_2(2410)\text{[4]}.$$
The trajectory for the $`a_2`$-mesons, Fig. 10b, is drawn on the basis of the recent data .
The quark states with ($`I=0`$, $`J^{PC}=2^{++}`$) are determined by two flavour components $`n\overline{n}`$ and $`s\overline{s}`$ for which two states $`{}_{}{}^{2S+1}L_{J}^{}=^3P_2,^3F_2`$ are possible. Consequently, we have four trajectories on the $`(n,M^2)`$ plane. Generally speaking, the $`f_2`$-states are mixtures of both the flavour components and the $`L=1,3`$ waves. The real situation is, however, such that the lowest trajectory \[$`f_2(1275)`$, $`f_2(1580)`$, $`f_2(1920)`$, $`f_2(2240)`$\] consists of mesons with dominant $`{}_{}{}^{3}P_{2}^{}n\overline{n}`$ components (we denote $`n\overline{n}=(u\overline{u}+d\overline{d})/\sqrt{2}`$), while the trajectory $`[f_2(1525)`$, $`f_2(1755)`$, $`f_2(2120)`$, $`f_2(2410)]`$ contains mesons with predominantly $`{}_{}{}^{3}P_{2}^{}s\overline{s}`$ components, and the $`F`$-trajectories are represented by three resonances \[$`f_2(2020)`$,$`f_2(2300)`$\] and \[$`f_2(2340)`$\] with the corresponding dominant $`{}_{}{}^{3}F_{2}^{}n\overline{n}`$ and $`{}_{}{}^{3}F_{2}^{}s\overline{s}`$ states. Following , we can state that the broad resonance $`f_2(2000)`$ is not part of those states placed on the $`(n,M^2)`$ trajectories. In the region of 2000 MeV three $`n\overline{n}`$-dominant resonances, $`f_2(1920)`$, $`f_2(2000)`$ and $`f_2(2020)`$, were seen, while on the $`(n,M^2)`$-trajectories there are only two vacant places. This means that one state is obviously ”superfluous” from the point of view of the $`q\overline{q}`$-systematics, i.e. it has to be considered as exotics. The large value of the width of the $`f_2(2000)`$ strengthen the suspicion that, indeed, this state is an exotic one.
## 4 Quarkonium and qluonium states: mixing and decay
On the basis of the $`1/N`$-expansion rules, we estimate here effects of mixing of quarkonium and qluonium states. Then, making use of the rules of quark combinatorics, we give the relations for decay constants of these states.
### 4.1 Mixing of $`q\overline{q}`$ and $`gg`$ states
The rules of the $`1/N`$-expansion , where $`N=N_c=N_f`$ are numbers of colours and light flavours, provide a possibility to estimate the mixing of the gluonium ($`gg`$) with the neighbouring quarkonium states ($`q\overline{q}`$).
The admixture of the $`gg`$ component in a $`q\overline{q}`$-meson is small, of the order of $`1/N_c`$ :
$`f_2(q\overline{q}\mathrm{meson})`$ $`=`$ $`q\overline{q}\mathrm{cos}\alpha +gg\mathrm{sin}\alpha `$ (5)
$`\mathrm{sin}^2\alpha `$ $``$ $`1/N_c.`$
The quarkonium component in the glueball should be larger, it is of the order of $`N_f/N_c`$ :
$`f_2(\mathrm{glueball})`$ $`=`$ $`gg\mathrm{cos}\gamma +(q\overline{q})_{glueball}\mathrm{sin}\gamma ,`$ (6)
$`\mathrm{sin}^2\gamma `$ $``$ $`N_f/N_c,`$
where $`(q\overline{q})_{glueball}`$ is a mixture of $`n\overline{n}=(u\overline{u}+d\overline{d})/\sqrt{2}`$ and $`s\overline{s}`$ components:
$$(q\overline{q})_{glueball}=n\overline{n}\mathrm{cos}\phi _{glueball}+s\overline{s}\mathrm{sin}\phi _{glueball},$$
(7)
with $`\mathrm{sin}\phi _{glueball}=\sqrt{\lambda /(2+\lambda )}`$. If the flavour SU(3) symmetry were satisfied, the quarkonium component $`(q\overline{q})_{glueball}`$ would be a flavour singlet, $`\phi _{glueball}\phi _{singlet}37^o`$. In reality, the probability of strange quark production in a gluon field is suppressed: $`u\overline{u}:d\overline{d}:s\overline{s}=1:1:\lambda `$, where $`\lambda 0.50.85`$. Hence, $`(q\overline{q})_{glueball}`$ differs slightly from the flavour singlet, it is determined by the parameter $`\lambda `$ as follows :
$$(q\overline{q})_{glueball}=(u\overline{u}+d\overline{d}+\sqrt{\lambda }s\overline{s})/\sqrt{2+\lambda }.$$
(8)
The suppression parameter $`\lambda `$ was estimated both in multiple hadron production processes , and in hadronic decay processes . In hadronic decays of mesons with different $`J^{PC}`$ the value of $`\lambda `$ can be, in principle, different. Still, the analyses of the decays of the $`2^{++}`$-states and $`0^{++}`$-states show that the suppression parameters are of the same order, 0.5–0.85, leading to
$$\phi _{glueball}26^033^o.$$
(9)
Let us explain now Eqs. (5)-(8) in detail.
First, let us evaluate the transition couplings using the rules of $`1/N`$-expansion; this evaluation will be done for the decay transitions gluonium $``$ two $`q\overline{q}`$-mesons and quarkonium$``$ two $`q\overline{q}`$-mesons. For this purpose, we consider the gluon loop diagram which corresponds to the two–gluon self–energy part: gluonium $``$ two gluons $``$ gluonium (see Fig. 11a). This loop diagram $`B(gluoniumgggluonium)`$ is of the order of unity, provided the gluonium is a two–gluon composite system: $`B(gluoniumgggluonium)g_{gluoniumgg}^2N_c^21`$, where $`g_{gluoniumgg}`$ is a coupling constant for the transition of a gluonium to two gluons. Therefore,
$$g_{gluoniumgg}1/N_c.$$
(10)
The coupling constant for the $`gluoniumq\overline{q}`$ transition is determined by the diagrams of Fig. 11b type. A similar evaluation gives:
$$g_{gluoniumq\overline{q}}g_{gluoniumgg}g^2N_c1/N_c.$$
(11)
Here $`g`$ is the quark–gluon coupling constant, which is of the order of $`1/\sqrt{N_c}`$ . The coupling constant for the gluonium $``$ two $`q\overline{q}`$-mesons transition in the leading $`1/N_c`$ terms is governed by diagrams of Fig. 11c type:
$$g_{gluoniumtwomesons}^Lg_{gluoniumq\overline{q}}g_{mesonq\overline{q}}^2N_c1/N_c.$$
(12)
In (12) the following evaluation of the coupling for transition $`q\overline{q}mesonq\overline{q}`$ has been used:
$$g_{mesonq\overline{q}}1/\sqrt{N_c},$$
(13)
which follows from the fact that the loop diagram of the $`q\overline{q}`$-meson propagator (see Fig. 12a) is of the order of unity: $`B(q\overline{q}mesonq\overline{q}meson)g_{mesonq\overline{q}}^2N_c1.`$
The diagram of the type of Fig. 11d governs the couplings for the transition gluonium $``$ two $`q\overline{q}`$-mesons in the next-to-leading terms of the $`1/N_c`$-expansion:
$$g_{gluoniumtwomesons}^{NL}g_{gluoniumgg}g_{mesongg}^2N_c^21/N_c^2,$$
(14)
where the coupling $`g_{mesongg}`$ has been estimated following the diagram in Fig. 12b:
$$g_{mesongg}g_{mesonq\overline{q}}g^21/N_c^{3/2}.$$
(15)
Decay couplings of $`q\overline{q}`$-meson into two mesons in leading and next-to-leading terms of $`1/N_c`$ expansion are determined by diagrams of the type of Figs. 12c and 12d, respectively. This gives
$`g_{mesontwomesons}^L`$ $``$ $`g_{mesonq\overline{q}}^3N_c1/\sqrt{N_c},`$ (16)
$`g_{mesontwomesons}^{NL}`$ $``$ $`g_{mesonq\overline{q}}^2g_{mesongg}g^2N_c^21/N_c^{3/2}.`$
Now we can estimate the order of the value of $`\mathrm{sin}^2\gamma `$ which defines the probability $`(q\overline{q})_{glueball}`$, see Eq. (6). This probability is determined by the self-energy part of the gluon propagator (diagram in Fig. 11e)— it is of the order of $`N_f/N_c`$, the factor $`N_f`$ being the light flavour number in the quark loop. Let us emphasise that the diagram in Fig. 11e stands for only one of the contributions of that type; indeed, contributions of the same order are also given by diagrams with all possible (but planar) gluon exchanges in the quark loop.
One can also evaluate $`\mathrm{sin}^2\gamma `$ using the transition amplitude $`gluoniumquarkonium`$ (see Fig. 12e), which is of the order of $`1/\sqrt{N_c}`$. The value $`\mathrm{sin}^2\gamma `$ is determined by the transition amplitude squared, summed over the flavours of all quarkonia, thus resulting in Eq. (6).
The probability of the gluonium component in the quarkonium, $`\mathrm{sin}^2\alpha `$, is of the order of the diagram in Fig. 12f, $`1/N_c`$, giving us the estimate (5). Here, as in the self-energy gluonium block, planar-type gluon exchanges are possible. Because of this, in the intermediate $`q\overline{q}`$ state all the interactions are taken into account.
The diagram in Fig. 11e defines also the flavour content of $`(q\overline{q})_{glueball}`$ — we see that the gluon field produces light quark pairs with probabilities $`u\overline{u}:d\overline{d}:s\overline{s}=1:1:\lambda `$, so $`(q\overline{q})_{glueball}`$ is determined by Eq. (8) not being a flavour singlet.
### 4.2 Quark combinatorial relations for decay constants
The rules of quark combinatorics lead to relations between decay couplings for mesons which belong to the same SU(3) nonet. The violation of the flavour symmetry is taken into account by introducing a suppression parameter $`\lambda `$ for the production of the strange quarks by gluons.
In the leading terms of the $`1/N`$ expansion, the main contribution to the decay coupling constant comes from planar diagrams. Examples of the production of new $`q\overline{q}`$-pairs by intermediate gluons are shown in Figs. 13a and 12b. When an isoscalar $`q\overline{q}`$-meson disintegrates, the coupling constants can be determined up to a common factor, by two characteristics of a meson. The first is the quark content of the $`q\overline{q}`$-meson, $`q\overline{q}=n\overline{n}\mathrm{cos}\phi +s\overline{s}\mathrm{sin}\phi `$, the second is the parameter $`\lambda `$. Experimental data provide the following values for this parameter: $`\lambda 0.5`$ in central hadron production in high–energy hadron–hadron collisions, $`\lambda =0.8\pm 0.2`$ for the decays of tensor mesons and $`\lambda =0.50.9`$ for the decays of $`0^{++}`$ mesons.
Let us consider in more detail the production of two pseudoscalar mesons $`P_1P_2`$ by $`f_2`$-quarkonium and $`f_2`$-gluonium:
$`f_2(\mathrm{quarkonium})`$ $``$ $`\pi \pi ,K\overline{K},\eta \eta ,\eta \eta ^{},\eta ^{}\eta ^{}`$ (17)
$`f_2(\mathrm{gluonium})`$ $``$ $`\pi \pi ,K\overline{K},\eta \eta ,\eta \eta ^{},\eta ^{}\eta ^{}.`$
The coupling constants for the decay into channels (17), which in the leading terms of the $`1/N`$ expansion are determined by diagrams of the type shown in Fig. 13, may be presented as
$`g^L(q\overline{q}P_1P_2)`$ $`=`$ $`C_{P_1P_2}^{q\overline{q}}(\phi ,\lambda )g_P^L,`$ (18)
$`g^L(ggP_1P_2)`$ $`=`$ $`C_{P_1P_2}^{gg}(\lambda )G_P^L,`$
where $`C_{P_1P_2}^{q\overline{q}}(\phi ,\lambda )`$ and $`C_{P_1P_2}^{gg}(\lambda )`$ are wholly calculable coefficients depending on the mixing angle $`\phi `$ and parameter $`\lambda `$; $`g_P^L`$ and $`G_P^L`$ are common factors describing the unknown dynamics of the processes.
Dealing with processes of the Fig. 13b type, one should bear in mind that they do not contain $`(q\overline{q})_{quarkonium}`$ components in the intermediate state but $`(q\overline{q})_{continuousspectrum}`$ only. The states $`(q\overline{q})_{quarkonium}`$ in this diagram would lead to processes of Fig. 13c, namely, to a diagram with the quarkonium decay vertex and the mixing block of $`gg`$ and $`q\overline{q}`$ components. All these sub-processes are taken into account separately.
The contributions of the diagrams of the type of Fig. 11d and 12d, which give the next-to-leading terms, $`g^{NL}(q\overline{q}P_1P_2)`$ and $`g^{NL}(ggP_1P_2)`$, may be presented in a form analogous to (18). The decay constant to the channel $`P_1P_2`$ is a sum of both contributions:
$`g^L(q\overline{q}P_1P_2)`$ $`+`$ $`g^{NL}(q\overline{q}P_1P_2),`$ (19)
$`g^L(ggP_1P_2)`$ $`+`$ $`g^{NL}(ggP_1P_2).`$
The second terms are suppressed compared to the first ones by a factor $`N_c`$; the experience in the calculation of quark diagrams teaches us that this suppression is of the order of 1/10.
Coupling constants for gluonium decays, $`g^L(ggP_1P_2)`$ and $`g^{NL}(ggP_1P_2)`$, are presented in Table 3 while those for quarkonium decays, $`g^L(q\overline{q}P_1P_2)`$ and $`g^{NL}(q\overline{q}P_1P_2)`$, are given in Table 4.
In Table 5 we give the couplings for decays of the gluonium state into channels of the vector mesons: $`ggV_1V_2`$.
Table 3
Coupling constants of the $`f_2`$-gluonium decaying to two pseudoscalar mesons, in the leading
and next-to leading terms of $`1/N`$ expansion. $`\mathrm{\Theta }`$ is here the mixing angle for $`\eta \eta ^{}`$ mesons: $`\eta =n\overline{n}\mathrm{cos}\mathrm{\Theta }s\overline{s}\mathrm{sin}\mathrm{\Theta }`$ and $`\eta ^{}=n\overline{n}\mathrm{sin}\mathrm{\Theta }+s\overline{s}\mathrm{cos}\mathrm{\Theta }`$.
Gluonium decay Gluonium decay Iden- couplings in the couplings in the tity Channel leading term of next-to-leading term factor $`1/N`$ expansion. of $`1/N`$ expansion. $`\pi ^0\pi ^0`$ $`G^L`$ 0 1/2 $`\pi ^+\pi ^{}`$ $`G^L`$ 0 1 $`K^+K^{}`$ $`\sqrt{\lambda }G^L`$ 0 1 $`K^0K^0`$ $`\sqrt{\lambda }G^L`$ 0 1 $`\eta \eta `$ $`G^L\left(\mathrm{cos}^2\mathrm{\Theta }+\lambda \mathrm{sin}^2\mathrm{\Theta }\right)`$ $`2G^{NL}(\mathrm{cos}\mathrm{\Theta }\sqrt{\frac{\lambda }{2}}\mathrm{sin}\mathrm{\Theta })^2`$ 1/2 $`\eta \eta ^{}`$ $`G^L(1\lambda )\mathrm{sin}\mathrm{\Theta }\mathrm{cos}\mathrm{\Theta }`$ $`2G^{NL}(\mathrm{cos}\mathrm{\Theta }\sqrt{\frac{\lambda }{2}}\mathrm{sin}\mathrm{\Theta })\times `$ 1 $`(\mathrm{sin}\mathrm{\Theta }+\sqrt{\frac{\lambda }{2}}\mathrm{cos}\mathrm{\Theta })`$ $`\eta ^{}\eta ^{}`$ $`G^L\left(\mathrm{sin}^2\mathrm{\Theta }+\lambda \mathrm{cos}^2\mathrm{\Theta }\right)`$ $`2G^{NL}\left(\mathrm{sin}\mathrm{\Theta }+\sqrt{\frac{\lambda }{2}}\mathrm{cos}\mathrm{\Theta }\right)^2`$ 1/2
Table 4
Coupling constants of the $`f_2`$-quarkonium decaying to two pseudoscalar mesons in the leading and next-to-leading terms of the $`1/N`$ expansion. The flavour content of the $`f_2`$-quarkonium is determined by the mixing angle $`\phi `$ as follows: $`f_2(q\overline{q})=n\overline{n}\mathrm{cos}\phi +s\overline{s}\mathrm{sin}\phi `$ where $`n\overline{n}=(u\overline{u}+d\overline{d})/\sqrt{2}`$.
| | Decay couplings of | Decay couplings of |
| --- | --- | --- |
| | quarkonium | quarkonium |
| Channel | in leading term | in next-to-leading term |
| | of $`1/N`$ expansion. | of $`1/N`$ expansion. |
| $`\pi ^0\pi ^0`$ | $`g^L\mathrm{cos}\phi /\sqrt{2}`$ | 0 |
| $`\pi ^+\pi ^{}`$ | $`g^L\mathrm{cos}\phi /\sqrt{2}`$ | 0 |
| $`K^+K^{}`$ | $`g^L(\sqrt{2}\mathrm{sin}\phi +\sqrt{\lambda }\mathrm{cos}\phi )/\sqrt{8}`$ | 0 |
| $`K^0K^0`$ | $`g^L(\sqrt{2}\mathrm{sin}\phi +\sqrt{\lambda }\mathrm{cos}\phi )/\sqrt{8}`$ | 0 |
| $`\eta \eta `$ | $`g^L(\mathrm{cos}^2\mathrm{\Theta }\mathrm{cos}\phi /\sqrt{2}+`$ | $`\sqrt{2}g^{NL}(\mathrm{cos}\mathrm{\Theta }\sqrt{\frac{\lambda }{2}}\mathrm{sin}\mathrm{\Theta })\times `$ |
| | $`\sqrt{\lambda }\mathrm{sin}\phi \mathrm{sin}^2\mathrm{\Theta })`$ | $`(\mathrm{cos}\phi \mathrm{cos}\mathrm{\Theta }\mathrm{sin}\phi \mathrm{sin}\mathrm{\Theta })`$ |
| $`\eta \eta ^{}`$ | $`g^L\mathrm{sin}\mathrm{\Theta }\mathrm{cos}\mathrm{\Theta }(\mathrm{cos}\phi /\sqrt{2}`$ | $`\sqrt{\frac{1}{2}}g^{NL}[(\mathrm{cos}\mathrm{\Theta }\sqrt{\frac{\lambda }{2}}\mathrm{sin}\mathrm{\Theta })\times `$ |
| | $`\sqrt{\lambda }\mathrm{sin}\phi )`$ | $`(\mathrm{cos}\phi \mathrm{sin}\mathrm{\Theta }+\mathrm{sin}\phi \mathrm{cos}\mathrm{\Theta })`$ |
| | | $`+(\mathrm{sin}\mathrm{\Theta }+\sqrt{\frac{\lambda }{2}}\mathrm{cos}\mathrm{\Theta })\times `$ |
| | | $`(\mathrm{cos}\phi \mathrm{sin}\mathrm{\Theta }\mathrm{sin}\phi \mathrm{cos}\mathrm{\Theta })]`$ |
| $`\eta ^{}\eta ^{}`$ | $`g^L(\mathrm{sin}^2\mathrm{\Theta }\mathrm{cos}\phi /\sqrt{2}+`$ | $`\sqrt{2}g^{NL}(\mathrm{sin}\mathrm{\Theta }+\sqrt{\frac{\lambda }{2}}\mathrm{cos}\mathrm{\Theta })\times `$ |
| | $`\sqrt{\lambda }\mathrm{sin}\phi \mathrm{cos}^2\mathrm{\Theta })`$ | $`(\mathrm{cos}\phi \mathrm{cos}\mathrm{\Theta }+\mathrm{sin}\phi \mathrm{sin}\mathrm{\Theta })`$ |
Table 5
The constants of the tensor glueball decay into two vector mesons in the leading (planar diagrams) and next-to-leading (non-planar diagrams) terms of $`1/N`$-expansion. The mixing angle for $`\omega \varphi `$ mesons is defined as: $`\omega =n\overline{n}\mathrm{cos}\phi _Vs\overline{s}\mathrm{sin}\phi _V`$, $`\varphi =n\overline{n}\mathrm{sin}\phi _V+s\overline{s}\mathrm{cos}\phi _V`$. Because of the small value of $`\phi _V`$, we keep in the Table only terms of the order of $`\phi _V`$.
| | Constants for | Constants for | Identity factor |
| --- | --- | --- | --- |
| | glueball decays in | glueball decays in | for decay |
| Channel | the leading order | next-to-leading order | products |
| | of $`1/N`$ expansion | of $`1/N`$ expansion | |
| $`\rho ^0\rho ^0`$ | $`G_V^L`$ | 0 | 1/2 |
| $`\rho ^+\rho ^{}`$ | $`G_V^L`$ | 0 | 1 |
| $`K^+K^{}`$ | $`\sqrt{\lambda }G_V^L`$ | 0 | 1 |
| $`K^0\overline{K}^0`$ | $`\sqrt{\lambda }G_V^L`$ | 0 | 1 |
| $`\omega \omega `$ | $`G_V^L`$ | $`2G_V^{NL}`$ | 1/2 |
| $`\omega \varphi `$ | $`G_V^L(1\lambda )\phi _V`$ | $`2G_V^{NL}\left(\sqrt{\frac{\lambda }{2}}+\phi _V\left(1\frac{\lambda }{2}\right)\right)`$ | 1 |
| $`\varphi \varphi `$ | $`\lambda G_V^L`$ | $`2G_V^{NL}\left(\frac{\lambda }{2}+\sqrt{2\lambda }\phi _V\right)`$ | 1/2 |
## 5 The broad state $`f_2(2000)`$: the tensor glueball
In the leading terms of $`1/N_c`$-expansion we have definite ratios for the glueball decay couplings. The next-to-leading terms in the decay couplings give corrections of the order of $`1/N_c`$. Let us remind that, as we see in the numerical calculations of the diagrams, the $`1/N_c`$ factor leads to a smallness of the order of $`1/10`$, and we neglect them in the analysis of the decays $`f_2\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$.
Considering a glueball state which is also a mixture of the gluonium and quarkonium components, we have $`\phi \phi _{glueball}=\mathrm{sin}^1\sqrt{\lambda /(2+\lambda )}`$ for the latter. So we can write
$$\frac{g^L((q\overline{q})_{glueball}P_1P_2)}{g^L((q\overline{q})_{glueball}P_1^{}P_2^{})}=\frac{g^L(ggP_1P_2)}{g^L(ggP_1^{}P_2^{})}$$
(20)
Then the relations for decay couplings of the glueball in the leading terms of the $`1/N`$-expansion read:
$`g_{\pi ^0\pi ^0}^{glueball}={\displaystyle \frac{G_{glueball}^L}{\sqrt{2+\lambda }}},`$
$`g_{\eta \eta }^{glueball}={\displaystyle \frac{G_{glueball}^L}{\sqrt{2+\lambda }}}(\mathrm{cos}^2\mathrm{\Theta }+\lambda \mathrm{sin}^2\mathrm{\Theta })`$
$`g_{\eta \eta ^{}}^{glueball}={\displaystyle \frac{G_{glueball}^L}{\sqrt{2+\lambda }}}(1\lambda )\mathrm{sin}\mathrm{\Theta }\mathrm{cos}\mathrm{\Theta }.`$ (21)
Hence, in spite of the unknown quarkonium components in the glueball, there are definite relations between the couplings of the glueball state with the channels $`\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ which can serve as signatures to define it.
### 5.1 Ratios between coupling constants of $`𝒇_\mathrm{𝟐}\mathbf{(}\mathrm{𝟐𝟎𝟎𝟎}\mathbf{)}\mathbf{}𝝅^\mathrm{𝟎}𝝅^\mathrm{𝟎}\mathbf{,}𝜼𝜼\mathbf{,}𝜼𝜼^{\mathbf{}}`$ as indication of a glueball nature of this state
Eq. (21) tells us that for the glueball state the relations between the coupling constants are $`1:(\mathrm{cos}^2\mathrm{\Theta }+\lambda \mathrm{sin}^2\mathrm{\Theta }):(1\lambda )\mathrm{cos}\mathrm{\Theta }\mathrm{sin}\mathrm{\Theta }`$. For $`(\lambda =0.5`$, $`\mathrm{\Theta }=37^{})`$ we have $`1:0.82:0.24`$, and for $`(\lambda =0.85`$, $`\mathrm{\Theta }=37^{})`$, respectively, $`1:0.95:0.07`$. Consequently, the relations between the coupling constants $`g_{\pi ^0\pi ^0}:g_{\eta \eta }:g_{\eta \eta ^{}}`$ for the glueball have to be
$$2^{++}glueballg_{\pi ^0\pi ^0}:g_{\eta \eta }:g_{\eta \eta ^{}}=1:(0.820.95):(0.240.07).$$
(22)
It follows from (2.3) that only the coupling constants of the broad $`f_2(2000)`$ resonance are inside the intervals: $`0.82g_{\eta \eta }/g_{\pi ^0\pi ^0}0.95`$ and $`0.24g_{\eta \eta ^{}}/g_{\pi ^0\pi ^0}0.07`$. Hence, it is just this resonance which can be considered as a candidate for a tensor glueball, while $`\lambda `$ is fixed in the interval $`0.5\lambda 0.7`$. Taking into account that there is no place for $`f_2(2000)`$ on the $`(n,M^2)`$-trajectories (see Fig. 10), it becomes evident that indeed, this resonance is the lowest tensor glueball.
### 5.2 Mixing of the glueball with neighbouring $`q\overline{q}`$-resonances
The position of the $`f_2`$-poles on the complex $`M`$-plane is shown in Fig. 14. We see that the glueball state $`f_2(2000)`$ overlaps with a large group of $`q\overline{q}`$-resonances. This means that there is a considerable mixing with the neighbouring resonances. The mixing can take place both at relatively small distances, on the quark-gluon level (processes of the type shown in Fig. 12e), and owing to decay processes
$$f_2(glueball)realmesonsf_2(q\overline{q}meson).$$
(23)
Processes of the type of (23) are presented in Fig. 15.
The estimates which were carried out in Section 4 demonstrated that even the mixing at the quark-gluon level (diagrams of the types in Fig. 12e) leads to a sufficiently large admixture of the quark-antiquark component in the glueball: $`f_2(glueball)=\mathrm{cos}\gamma gg+\mathrm{sin}\gamma gg`$ with $`\mathrm{sin}\gamma \sqrt{N_f/N_c}`$. A mixing due to processes (23), apparently, enhances the quark-antiquark component. The main effect of the processes (23) is, however, that in the case of overlapping resonances one of them accumulates the widths of the neighbouring resonances. The position of the $`f_2`$-poles in Fig. 14 makes it obvious that such a state is the tensor glueball.
A similar situation was detected also in the sector of scalar mesons in the region $`10001700`$ MeV: the scalar glueball, being in the neighbourhood of $`q\overline{q}`$-resonances, accumulated a relevant fraction of their widths and transformed into a broad $`f_0(12001600)`$ state. Such a transformation of a scalar glueball into a broad state was observed in ; further investigations verified this observation. The possibility that a scalar (and tensor) glueball may considerably mix with $`q\overline{q}`$-states was discussed already for quite a long time, see, e.g., . At the same time there is a number of papers, e.g. , in which the mixing due to transitions (23) is not taken into account. Hence, in these papers relatively narrow resonances like $`f_0(1500)`$ and $`f_0(1710)`$ are suggested as possible scalar glueballs.
We see that both glueballs, the scalar $`f_0(12001600)`$ and the tensor $`f_2(2000)`$ one, reveal themselves as broad resonances. We can suppose that this is not accidental. In the transition of the lowest scalar glueball into a broad resonance was investigated in the framework of modelling decays by self-energy quark an gluon diagrams. As it was discovered, it was just the glueball which, appearing among the $`q\overline{q}`$-states, began to mix with them actively, accumulating their widths. Hence, the glueball turned out to be the broadest resonance.
#### 5.2.1 The mixing of two unstable states
In the case of two resonances, the propagator of the state 1 is determined by the diagrams of Fig. 16a. With all these processes taken into account, the propagator of the state 1 is equal to:
$$D_{11}(s)=\left(m_1^2sB_{11}(s)\frac{B_{12}(s)B_{21}(s)}{m_2^2sB_{22}(s)}\right)^1.$$
(24)
Here $`m_1`$ and $`m_2`$ are masses of the input states 1 and 2, and the loop diagrams $`B_{ij}(s)`$ are defined by the spectral integral
$$B_{ij}(s)=\underset{4m^2}{\overset{\mathrm{}}{}}\frac{d(s^{})}{\pi }\frac{g_i(s^{})g_j(s^{})\rho (s^{})}{s^{}si0},$$
(25)
where $`g_i(s^{})`$ and $`g_j(s^{})`$ are vertices and $`\rho (s^{})`$ is the phase space for the intermediate state. It is helpful to introduce the propagator matrix $`D_{ij}`$, where the non-diagonal elements $`D_{12}=D_{21}`$ correspond to the transitions $`12`$ and $`21`$ (see Fig. 16b). The matrix reads:
$$\widehat{D}=\left|\begin{array}{cc}D_{11}\hfill & D_{12}\hfill \\ D_{21}\hfill & D_{22}\hfill \end{array}\right|=\frac{1}{(M_1^2s)(M_2^2s)B_{12}B_{21}}\left|\begin{array}{cc}M_2^2s,& B_{12}\\ B_{21},& M_1^2s\end{array}\right|.$$
(26)
Here the following notation is used:
$$M_i^2=m_i^2B_{ii}(s)i=1,2.$$
(27)
Zeros in the denominator of the propagator matrix (26) define the complex resonance masses after the mixing:
$$\mathrm{\Pi }(s)=(M_1^2s)(M_2^2s)B_{12}B_{21}=0.$$
(28)
Let us denote the complex masses of mixed states as $`M_A`$ and $`M_B`$.
Consider a simple model, where the $`s`$-dependence of the function $`B_{ij}(s)`$ near the points $`sM_A^2`$ and $`sM_B^2`$ is assumed to be negligible. Let $`M_i^2`$ and $`B_{12}`$ be constants. Then one has:
$$M_{A,B}^2=\frac{1}{2}(M_1^2+M_2^2)\pm \sqrt{\frac{1}{4}(M_1^2M_2^2)^2+B_{12}B_{21}}.$$
(29)
In the case, when the widths of initial resonances 1 and 2 are small (hence the imaginary part of the transition diagram $`B_{12}`$ is also small), the equation (29) turns into the standard formula of quantum mechanics for the split of mixing levels, which become repulsive as a result of the mixing. If so,
$$\widehat{D}=\left|\begin{array}{cc}\mathrm{cos}^2\theta /(M_A^2s)+sin^2\theta /(M_B^2s)& \mathrm{cos}\theta \mathrm{sin}\theta /(M_A^2s)+\mathrm{sin}\theta \mathrm{cos}\theta /(M_B^2s)\\ \mathrm{cos}\theta \mathrm{sin}\theta /(M_A^2s)+\mathrm{sin}\theta \mathrm{cos}\theta /(M_B^2s)& \mathrm{sin}^2\theta /(M_A^2s)+\mathrm{cos}^2\theta /(M_B^2s)\end{array}\right|,$$
$$\mathrm{cos}^2\theta =\frac{1}{2}+\frac{1}{2}\frac{\frac{1}{2}(M_1^2M_2^2)}{\sqrt{\frac{1}{4}(M_1^2M_2^2)^2+B_{12}B_{21}}}.$$
(30)
The states $`|A>`$ and $`|B>`$ are superpositions of the initial levels, $`|1>`$ and $`|2>`$, as follows:
$$|A>=\mathrm{cos}\theta |1>\mathrm{sin}\theta |2>,|B>=\mathrm{sin}\theta |1>+\mathrm{cos}\theta |2>.$$
(31)
In general, the representation of states $`|A>`$ and $`|B>`$ as superpositions of initial states is valid, when the $`s`$-dependence of functions $`B_{ij}(s)`$ can not be neglected, and their imaginary parts are not small. Consider the propagator matrix near $`s=M_A^2`$:
$$\widehat{D}=\frac{1}{\mathrm{\Pi }(s)}\left|\begin{array}{cc}M_2^2(s)s& B_{12}(s)\\ B_{21}(s)& M_1^2(s)s\end{array}\right|\frac{1}{\mathrm{\Pi }^{}(M_A^2)(M_A^2s)}\left|\begin{array}{cc}M_2^2(M_A^2)M_A^2& B_{12}(M_A^2)\\ B_{21}(M_A^2)& M_1^2(M_A^2)M_A^2\end{array}\right|.$$
(32)
In the left-hand side of Eq. (32), only the singular (pole) terms survive. The matrix determinant in the right-hand side of (32) equals zero:
$$[M_2^2(M_A^2)M_A^2][M_1^2(M_A^2)M_A^2]B_{12}(M_A^2)B_{21}(M_A^2)=0,$$
(33)
This equality follows from Eq. (28), which fixes $`\mathrm{\Pi }(M_A^2)=0`$. It allows us to introduce the complex mixing angle:
$$|A>=\mathrm{cos}\theta _A|1>\mathrm{sin}\theta _A|2>.$$
(34)
The right-hand side of Eq. (30) can be rewritten by making use of the mixing angle $`\theta _A`$, as follows:
$$\left[\widehat{D}\right]_{sM_A^2}=\frac{N_A}{M_A^2s}\left|\begin{array}{cc}\mathrm{cos}^2\theta _A& \mathrm{cos}\theta _A\mathrm{sin}\theta _A\\ \mathrm{sin}\theta _A\mathrm{cos}\theta _A& \mathrm{sin}^2\theta _A\end{array}\right|,$$
(35)
where
$$N_A=\frac{1}{\mathrm{\Pi }^{}(M_A^2)}[2M_A^2M_1^2M_2^2],\mathrm{cos}^2\theta _A=\frac{M_A^2M_2^2}{2M_A^2M_1^2M_2^2},\mathrm{sin}^2\theta _A=\frac{M_A^2M_1^2}{2M_A^2M_1^2M_2^2}.$$
(36)
We remind that in the formula (36) the functions $`M_1^2(s)`$, $`M_2^2(s)`$ and $`B_{12}(s)`$ are fixed in the point $`s=M_A^2`$. In the case under consideration, when the angle $`\theta _A`$ is a complex quantity, the values $`\mathrm{cos}^2\theta _A`$ and $`\mathrm{sin}^2\theta _A`$ do not determine the probability of states $`|1>`$ and $`|2>`$ in $`|A>`$; indeed, the values $`\sqrt{N_A}\mathrm{cos}\theta _A`$ and $`\sqrt{N_A}\mathrm{sin}\theta _A`$ are the transition amplitudes $`|A>|1>`$ and $`|A>|2>`$. Therefore, the corresponding probabilities are equal to $`|\mathrm{cos}\theta _A|^2`$ and $`|\mathrm{sin}\theta _A|^2`$.
In order to analyse the content of the state $`|B>`$, an analogous expansion of the propagator matrix should be carried out near the point $`s=M_B^2`$. Introducing
$$|B>=\mathrm{sin}\theta _B|1>+\mathrm{cos}\theta _B|2>,$$
(37)
we have the following expression for $`\widehat{D}`$ in the vicinity of the second pole $`s=M_B^2`$:
$$\left[\widehat{D}\right]_{sM_B^2}=\frac{N_B}{M_B^2s}\left|\begin{array}{cc}\mathrm{sin}^2\theta _B\hfill & \mathrm{cos}\theta _B\mathrm{sin}\theta _B\hfill \\ \mathrm{sin}\theta _B\mathrm{cos}\theta _B\hfill & \mathrm{cos}^2\theta _B\hfill \end{array}\right|,$$
(38)
where
$$N_B=\frac{1}{\mathrm{\Pi }^{}(M_B^2)}\left[2M_B^2M_1^2M_2^2\right],\mathrm{cos}^2\theta _B=\frac{M_B^2M_1^2}{2M_B^2M_1^2M_2^2},\mathrm{sin}^2\theta _B=\frac{M_B^2M_2^2}{2M_B^2M_1^2M_2^2}.$$
(39)
In Eqs. (38), (39) the functions $`M_1^2(s)`$, $`M_2^2(s)`$ and $`B_{12}(s)`$ are fixed in the point $`s=M_B^2`$.
If $`B_{12}`$ depends on $`s`$ weakly and one can neglect this dependence, the angles $`\theta _A`$ and $`\theta _B`$ coincide. In general, however, they are different. So the formulae for the propagator matrix differ from the standard approach of quantum mechanics by this very point.
Another distinction is related to the type of the level shift afforded by mixing, namely, in quantum mechanics the levels ”repulse” each other from the mean value $`1/2(E_1+E_2)`$ (see also Eq. (29)). Generally speaking, the equation (28) can cause both a ”repulsion” of masses squared from the mean value, $`1/2(M_1^2+M_2^2)`$, and an ”attraction”.
Let us remind now, how to write the amplitudes in the one-channel and multi-channel cases. The scattering amplitude in the one-channel case is defined by the following expression:
$$A(s)=g_i(s)D_{ij}(s)g_j(s).$$
(40)
In the multi-channel case, $`B_{ij}(s)`$ is a sum of loop diagrams: $`B_{ij}(s)=_aB_{ij}^{(a)}(s)`$, where $`B_{ij}^{(a)}`$ is a loop diagram in the channel $`a`$ with vertex functions $`g_i^{(a)}`$, $`g_j^{(a)}`$ and a phase space factor $`\rho _a`$. The partial scattering amplitude in the channel $`ab`$ equals
$$A_{ab}(s)=g_i^{(a)}(s)D_{ij}(s)g_j^{(b)}(s).$$
(41)
#### 5.2.2 Construction of propagator matrix in a general case ($`N`$ resonances)
Consider the construction of the propagator matrix $`\widehat{D}`$ for an arbitrary number $`(N)`$ of resonances. The matrix elements, $`D_{ij}`$, describe the transition from the initial state $`i`$ (with the bare propagator $`(m_i^2s)^1`$) to the state $`j`$. They obey the system of linear equations as follows:
$$D_{ij}=D_{ik}B_{kj}(s)(m_j^2s)^1+\delta _{ij}(m_j^2s)^1,$$
(42)
where $`B_{ij}(s)`$ is the loop diagram for the transition $`ij`$ and $`\delta _{ij}`$ is the Kronecker symbol.
Let us introduce the diagonal propagator matrix $`\widehat{d}`$ for initial states :
$$\widehat{d}=diag((m_1^2s)^1,(m_2^2s)^1,(m_3^2s)^1\mathrm{}).$$
(43)
Then the system of linear equations (39) can be rewritten in the matrix form as
$$\widehat{D}=\widehat{D}\widehat{B}\widehat{d}+\widehat{d}.$$
(44)
One obtains
$$\widehat{D}=\frac{I}{(\widehat{d}^1\widehat{B})}.$$
(45)
The matrix $`\widehat{d}^1`$ is diagonal, hence $`\widehat{D}^1=(\widehat{d}^1\widehat{B})`$ is of the form
$$\widehat{D}^1=\left|\begin{array}{cccc}M_1^2s& B_{12}(s)& B_{13}(s)& \mathrm{}\\ B_{21}(s)& M_2^2s& B_{23}(s)& \mathrm{}\\ B_{31}(s)& B_{32}(s)& M_3^2s& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right|,$$
(46)
where $`M_i^2`$ is defined by Eq. (27). Inverting this matrix, we obtain a full set of elements $`D_{ij}(s)`$:
$$D_{ij}(s)=\frac{(1)^{i+j}\mathrm{\Pi }_{ji}^{(N1)}(s)}{\mathrm{\Pi }^{(N)}(s)}.$$
(47)
Here $`\mathrm{\Pi }^{(N)}(s)`$ is the determinant of the matrix $`\widehat{D}^1`$, and $`\mathrm{\Pi }_{ji}^{(N1)}(s)`$ is a matrix supplement to the element $`[\widehat{D}^1]_{ji}`$, i.e. the matrix $`\widehat{D}^1`$ with an excluded $`j`$-th line and $`i`$-th column.
The zeros of $`\mathrm{\Pi }^{(N)}(s)`$ define the poles of the propagator matrix which correspond to physical resonances formed by the mixing. We denote the complex resonance masses as:
$$s=M_A^2,M_B^2,M_C^2,\mathrm{}$$
(48)
Near the point $`s=M_A^2`$, one can leave in the propagator matrix the leading pole term only. This means that the free term in Eq. (44) can be neglected, so we get a system of homogeneous equations:
$$D_{ik}(s)\left(\widehat{d}^1\widehat{B}\right)_{kj}=0.$$
(49)
The solution of this system is defined up to the normalisation factor, and it does not depend on the initial index $`i`$. If so, the elements of the propagator matrix may be written in a factorised form as follows:
$$\left[\widehat{D}^{(N)}\right]_{sM_A^2}=\frac{N_A}{M_A^2s}\left|\begin{array}{cccc}\alpha _1^2,\hfill & \alpha _1\alpha _2,\hfill & \alpha _1\alpha _3,\hfill & \mathrm{}\hfill \\ \alpha _2\alpha _1,\hfill & \alpha _2^2,\hfill & \alpha _2\alpha _3,\hfill & \mathrm{}\hfill \\ \alpha _3\alpha _1,\hfill & \alpha _3\alpha _2,\hfill & \alpha _3^2,\hfill & \mathrm{}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \end{array}\right|,$$
(50)
where $`N_A`$ is the normalisation factor chosen to satisfy the condition
$$\alpha _1^2+\alpha _2^2+\alpha _3^2+\mathrm{}+\alpha _N^2=1.$$
(51)
The constants $`\alpha _i`$ are the normalised amplitudes for the resonance A $``$ state $`i`$ transitions. The probability to find the state $`i`$ in the physical resonance $`A`$ is equal to:
$$w_i=|\alpha _i|^2.$$
(52)
Analogous representations of the propagator matrix can be given also in the vicinity of other poles:
$$D_{ij}^{(N)}(sM_B^2)=N_B\frac{\beta _i\beta _j}{M_B^2s},D_{ij}^{(N)}(sM_C^2)=N_C\frac{\gamma _i\gamma _j}{M_C^2s}\mathrm{}.$$
(53)
The coupling constants satisfy normalisation conditions similar to that of Eq. (48):
$$\beta _1^2+\beta _2^2+\mathrm{}+\beta _N^2=1,\gamma _1^2+\gamma _2^2+\mathrm{}+\gamma _N^2=1,\mathrm{}.$$
(54)
In the general case, however, there is no completeness condition for the inverse expansion:
$$\alpha _i^2+\beta _i^2+\gamma _i^2+\mathrm{}1.$$
(55)
For two resonances this means that $`\mathrm{cos}^2\mathrm{\Theta }_A+\mathrm{sin}^2\mathrm{\Theta }_B1`$. Still, let us remind that the equality in the inverse expansion, which is relevant for the completeness condition, appears in models where the $`s`$-dependence of the loop diagrams is neglected.
#### 5.2.3 Full resonance overlapping: the accumulation of widths of neighbouring resonances by one of them
Let us consider two examples which describe the idealised situation of a full overlapping of two or three resonances. In these examples, the effect of accumulation of widths of neighbouring resonances by one of them can be seen in its original untouched form.
a) Full overlapping of two resonances.
For the sake of simplicity, let $`B_{ij}`$ be a weak $`s`$-dependent function; Eq. (29) may be used. We define:
$$M_1^2=M_R^2iM_R\mathrm{\Gamma }_1,M_2^2=M_R^2iM_R\mathrm{\Gamma }_2,$$
(56)
and put
$$\mathrm{Re}B_{12}(M_R^2)=P\underset{(\mu _1+\mu _2)^2}{\overset{\mathrm{}}{}}\frac{ds^{}}{\pi }\frac{g_1(s^{})g_2(s^{})\rho (s^{})}{s^{}M_R^2}0.$$
(57)
It is possible that Re$`B_{12}(M_R^2)`$ equals zero at positive $`g_1`$ and $`g_2`$, if the contribution from the integration region $`s^{}<M_R^2`$ cancels the contribution from the $`s^{}>M_R^2`$ region. In this case
$$B_{12}(M_R^2)ig_1(M_R^2)g_2(M_R^2)\rho (M_R^2)=iM_R\sqrt{\mathrm{\Gamma }_1\mathrm{\Gamma }_2}.$$
(58)
Substituting Eqs. (53)–(55) into Eq. (45), one has:
$$M_A^2M_R^2iM_R(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)M_B^2M_R^2.$$
(59)
Therefore, after mixing, one of the states accumulates the widths of primary resonances, $`\mathrm{\Gamma }_A\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$, and another state becomes a quasi-stable particle, with $`\mathrm{\Gamma }_B0`$.
b) Full overlapping of three resonances.
Consider the equation
$$\mathrm{\Pi }^{(3)}(s)=0$$
(60)
in the same approximation as in the above example. Correspondingly, we put:
$$\mathrm{Re}B_{ab}(M_R^2)0,(ab);M_i^2=M_R^2siM_R\mathrm{\Gamma }_i=xi\gamma _i.$$
(61)
A new variable, $`x=M_R^2s`$, is used, and we denote $`M_R\mathrm{\Gamma }_i=\gamma _i`$. Taking into account $`B_{ij}B_{ji}=\gamma _i\gamma _j`$ and $`B_{12}B_{23}B_{31}=i\gamma _1\gamma _2\gamma _3`$, we can rewrite the equation (60) as follows:
$$x^3+x^2(i\gamma _1+i\gamma _2+i\gamma _3)=0.$$
(62)
Therefore, at full resonance overlapping, one obtains:
$$M_A^2M_R^2iM_R(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2+\mathrm{\Gamma }_3),M_B^2M_R^2,M_C^2M_R^2.$$
(63)
Thus, the resonance $`A`$ has accumulated the widths of three primary resonances, and the states $`B`$ and $`C`$ became quasi-stable and degenerate.
## 6 The $`𝒒\overline{𝒒}`$-$`𝒈𝒈`$ content of $`𝒇_\mathrm{𝟐}`$-mesons, observed in the reactions $`𝒑\overline{𝒑}\mathbf{}𝝅^\mathrm{𝟎}𝝅^\mathrm{𝟎}\mathbf{,}𝜼𝜼\mathbf{,}𝜼𝜼^{\mathbf{}}`$
We determine here the $`q\overline{q}gg`$ content of $`f_2`$-mesons, observed in the reactions $`p\overline{p}\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ . This determination is based on experimentally observed relations (2.3) and the rules of quark combinatorics taken into account in the leading terms of the $`1/N`$-expansion.
For the $`f_2\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ transitions, when the $`q\overline{q}`$-meson is a mixture of quarkonium and gluonium components, the decay vertices read in the leading terms of the $`1/N`$-expansion (see Tables 3 and 4) as follows:
$`g_{\pi ^0\pi ^0}^{q\overline{q}meson}=g{\displaystyle \frac{\mathrm{cos}\phi }{\sqrt{2}}}+{\displaystyle \frac{G}{\sqrt{2+\lambda }}},`$
$`g_{\eta \eta }^{q\overline{q}meson}=g\left(\mathrm{cos}^2\theta {\displaystyle \frac{\mathrm{cos}\phi }{\sqrt{2}}}+\mathrm{sin}^2\mathrm{\Theta }\sqrt{\lambda }\mathrm{sin}\phi \right)+{\displaystyle \frac{G}{\sqrt{2+\lambda }}}(\mathrm{cos}^2\mathrm{\Theta }+\lambda \mathrm{sin}^2\mathrm{\Theta }),`$
$`g_{\eta \eta ^{}}^{q\overline{q}meson}=\mathrm{sin}\mathrm{\Theta }\mathrm{cos}\mathrm{\Theta }\left[g\left({\displaystyle \frac{\mathrm{cos}\phi }{\sqrt{2}}}\sqrt{\lambda }\mathrm{sin}\phi \right)+{\displaystyle \frac{G}{\sqrt{2+\lambda }}}(1\lambda )\right].`$ (64)
The terms proportional to $`g`$ stand for the $`q\overline{q}twomesons`$ transitions ($`g=g^L\mathrm{cos}\alpha `$), while the terms with $`G`$ represent the $`gluoniumtwomesons`$ transition ($`G=G^L\mathrm{sin}\alpha `$). Consequently, $`G^2`$ and $`g^2`$ are proportional to the probabilities for finding gluonium ($`W=\mathrm{sin}^2\alpha `$) and quarkonium $`(1W)`$ components in the considered $`f_2`$-meson. Let us remind that the mixing angle $`\mathrm{\Theta }`$ stands for the $`n\overline{n}`$ and $`s\overline{s}`$ components in the $`\eta `$ and $`\eta ^{}`$ mesons; we neglect the possible admixture of a gluonium component to $`\eta `$ and $`\eta ^{}`$ (according to , the gluonium admixture to $`\eta `$ is less than 5%, to $`\eta ^{}`$ — less than 20%). For the mixing angle $`\mathrm{\Theta }`$ we take $`\mathrm{\Theta }=37^{}`$.
### 6.1 The analysis of the quarkonium-gluonium contents of <br>the $`𝒇_\mathrm{𝟐}\mathbf{(}\mathrm{𝟏𝟗𝟐𝟎}\mathbf{)}`$, $`𝒇_\mathrm{𝟐}\mathbf{(}\mathrm{𝟐𝟎𝟐𝟎}\mathbf{)}`$, $`𝒇_\mathrm{𝟐}\mathbf{(}\mathrm{𝟐𝟐𝟒𝟎}\mathbf{)}`$, $`𝒇_\mathrm{𝟐}\mathbf{(}\mathrm{𝟐𝟑𝟎𝟎}\mathbf{)}`$
Making use of the data (2.3), the relations (64) allow us to to find $`\phi `$ as a function of the ratio $`G/g`$ of the coupling constants. The result for the resonances $`f_2(1920)`$, $`f_2(2020)`$, $`f_2(2240)`$, $`f_2(2300)`$ is shown in Fig. 17. Solid curves enclose the values of $`g_{\eta \eta }/g_{\pi ^0\pi ^0}`$ for $`\lambda =0.6`$ (this is the $`\eta \eta `$-zone in the $`(G/g,\phi )`$ plane) and dashed curves enclose $`g_{\eta \eta ^{}}/g_{\pi ^0\pi ^0}`$ for $`\lambda =0.6`$ (the $`\eta \eta ^{}`$-zone). The values of $`G/g`$ and $`\phi `$, lying in both zones, describe the experimental data (2.3): these regions are shadowed in Fig. 17.
The correlation curves in Fig. 17 enable us to give a qualitative estimate for the change of the angle $`\phi `$ (i.e. the relation of the $`n\overline{n}`$ and $`s\overline{s}`$ components in the $`f_2`$ meson) depending on the value of the gluonium admixture. The values $`g^2`$ and $`G^2`$ are proportional to the probabilities of having quarkonium and gluonium components in the $`f_2`$ meson, $`g^2=(g^L)^2(1W)`$ and $`G^2=(G^L)^2W`$. Here $`W`$ is the probability of a gluonium admixture in the considered $`q\overline{q}`$-meson; $`g^L`$ and $`G^L`$ are universal constants, see Tables 3 and 4. Since $`G^L/g^L1/\sqrt{N_c}`$ and $`W1/N_c`$, we can give a rough estimate:
$$\frac{G^2}{g^2}\frac{W}{N_c(1W)}\frac{W}{10}.$$
(65)
Let us remind that the numerical calculations of the diagrams indicate that $`1/N_c`$ leads to a smallness of the order of $`1/10`$ – this is taken into account in (65). Assuming that the gluonium components are less than 20% ($`W<0.2`$) in each of the $`q\overline{q}`$ resonances $`f_2(1920)`$, $`f_2(2020)`$, $`f_2(2240)`$, $`f_2(2300)`$, we put roughly $`W10G^2/g^2`$, and obtain for the angles $`\phi `$ the following intervals:
$`W_{gluonium}[f_2(1920)]<20\%:0.8^{}<\phi [f_2(1920)]<3.6^{},`$
$`W_{gluonium}[f_2(2020)]<20\%:7.5^{}<\phi [f_2(2020)]<13.2^{},`$
$`W_{gluonium}[f_2(2240)]<20\%:8.3^{}<\phi [f_2(2240)]<17.3^{},`$
$`W_{gluonium}[f_2(2300)]<20\%:25.6^{}<\phi [f_2(2300)]<9.3^{}`$ (66)
### 6.2 The $`𝒏\overline{𝒏}`$-$`𝒔\overline{𝒔}`$ content of the $`𝒒\overline{𝒒}`$-mesons
Let us summarise what we know about the status of the $`(I=0,J^{PC}=2^{++})`$ $`q\overline{q}`$-mesons. Estimating the $`n\overline{n}`$-$`s\overline{s}`$ content of the $`f_2`$-mesons, we ignore the $`gg`$ admixture (remembering that it is of the order of $`\mathrm{sin}^2\alpha 1/N_c`$).
1. The resonances $`f_2(1270)`$ and $`f_2^{}(1525)`$ are well-known partners of the basic nonet with $`n=1`$ and a dominant $`P`$-component, $`1^3P_2q\overline{q}`$. Their flavour content, obtained from the reaction $`\gamma \gamma K_SK_S`$, is
$`f_2(1270)`$ $`=`$ $`\mathrm{cos}\phi _{n=1}n\overline{n}+\mathrm{sin}\phi _{n=1}s\overline{s},`$
$`f_2(1525)`$ $`=`$ $`\mathrm{sin}\phi _{n=1}n\overline{n}+\mathrm{cos}\phi _{n=1}s\overline{s},`$ (67)
$`\phi _{n=1}=1\pm 3^{}.`$
2. The resonances $`f_2(1560)`$ and $`f_2(1750)`$ are partners in a nonet with $`n=2`$ and a dominant $`P`$-component, $`2^3P_2q\overline{q}`$. Their flavour content, obtained from the reaction $`\gamma \gamma K_SK_S`$, is
$`f_2(1560)`$ $`=`$ $`\mathrm{cos}\phi _{n=2}n\overline{n}+\mathrm{sin}\phi _{n=2}s\overline{s},`$
$`f_2(1750)`$ $`=`$ $`\mathrm{sin}\phi _{n=2}n\overline{n}+\mathrm{cos}\phi _{n=2}s\overline{s},`$ (68)
$`\phi _{n=1}=10_{10}^{+5}{}_{}{}^{}.`$
3. The resonances $`f_2(1920)`$ and $`f_2(2120)`$ (in they are denoted as $`f_2(1910)`$ and $`f_2(2010)`$) are partners in a nonet with $`n=3`$ and with a dominant $`P`$-component, $`3^3P_2q\overline{q}`$. Ignoring the contribution of the glueball component, their flavour content, obtained from the reactions $`p\overline{p}\pi ^0\pi ^0`$, $`\eta \eta `$, $`\eta \eta ^{}`$, is
$`f_2(1920)`$ $`=`$ $`\mathrm{cos}\phi _{n=3}n\overline{n}+\mathrm{sin}\phi _{n=3}s\overline{s},`$
$`f_2(2120)`$ $`=`$ $`\mathrm{sin}\phi _{n=3}n\overline{n}+\mathrm{cos}\phi _{n=3}s\overline{s},`$ (69)
$`\phi _{n=3}=0\pm 5^{}.`$
4. The next, predominantly $`{}_{}{}^{3}P_{2}^{}`$ states with $`n=4`$ are $`f_2(2240)`$ and $`f_2(2410)`$ . (By mistake, in the resonance $`f_2(2240)`$ is listed as $`f_2(2300)`$, while $`f_2(2410)`$ is denoted as $`f_2(2340)`$). Their flavour content at $`W=0`$ is determined as
$`f_2(2240)`$ $`=`$ $`\mathrm{cos}\phi _{n=4}n\overline{n}+\mathrm{sin}\phi _{n=4}s\overline{s},`$
$`f_2(2410)`$ $`=`$ $`\mathrm{sin}\phi _{n=4}n\overline{n}+\mathrm{cos}\phi _{n=4}s\overline{s},`$ (70)
$`\phi _{n=4}=5\pm 11^{}.`$
5. $`f_2(2020)`$ and $`f_2(2340)`$ belong to the basic $`F`$-wave nonet $`(n=1)`$ (in the $`f_2(2020)`$ is denoted as $`f_2(2000)`$ and is put in the section ”Other light mesons”, while $`f_2(2340)`$ is denoted as $`f_2(2300)`$). The flavour content of the $`1^3F_2`$ mesons is
$`f_2(2020)`$ $`=`$ $`\mathrm{cos}\phi _{n(F)=1}n\overline{n}+\mathrm{sin}\phi _{n(F)=1}s\overline{s},`$
$`f_2(2340)`$ $`=`$ $`\mathrm{sin}\phi _{n(F)=1}n\overline{n}+\mathrm{cos}\phi _{n(F)=1}s\overline{s},`$ (71)
$`\phi _{n(F)=1}=5\pm 8^{}.`$
6. The resonance $`f_2(2300)`$ has a dominant $`F`$-wave quark-antiquark component; its flavour content for $`W=0`$ is defined as
$$f_2(2300)=\mathrm{cos}\phi _{n(F)=2}n\overline{n}+\mathrm{sin}\phi _{n(F)=2}s\overline{s},\phi _{n(F)=2}=8^{}\pm 12^{}.$$
(72)
A partner of $`f_2(2300)`$ in the $`2^3F_2`$ nonet has to be a $`f_2`$-resonance with a mass $`M2570`$MeV.
## 7 Conclusion
The broad $`f_2(2000)`$ state is the descendant of the lowest tensor glueball. This statement is favoured by estimates of parameters of the Pomeron trajectory (e.g., see , Chapter 5.4, and references therein), according to which $`M_{2^{++}glueball}1.72.5`$ GeV. Lattice calculations result in a similar value, namely, 2.2–2.4 GeV . The corresponding coupling constants $`f_2(2000)\pi ^0\pi ^0,\eta \eta ,\eta \eta ^{}`$ satisfy the relations for the glueball, eq.(22), with $`\lambda 0.50.7`$. The admixture of the quarkonium component $`(q\overline{q})_{glueball}`$ in $`f_2(2000)`$ cannot be determined by the ratios of the coupling constants between the hadronic channels; to define it, $`f_2(2000)`$ has to be observed in $`\gamma \gamma `$-collisions. The value of $`(q\overline{q})_{glueball}`$ in $`f_2(2000)`$ may be rather large: the rules of $`1/N`$-expansion give a value of the order of $`N_f/N_c`$. It is, probably, just the largeness of the quark-antiquark component in $`f_2(2000)`$ which results in its suppressed production in the radiative $`J/\psi `$ decays (see discussion in ).
We have now two observed glueballs, a scalar $`f_0(12001600)`$ (see also ) and a tensor one, $`f_2(2000)`$. It is illustrative to present the situation with $`0^{++},2^{++}`$ glueballs on the $`(J,M^2)`$-plane, we demonstrate this in Fig. 18. According to various estimates, the leading Pomeron trajectory has an intercept at $`\alpha (0)1.101.30`$ (see, for example, ). Assuming that the Pomeron trajectory has a linear behaviour, which does not contradict experimental data, $`\alpha _P(M^2)=\alpha _P(0)+\alpha _P^{}(0)M^2`$, we have for the slope $`\alpha _P^{}(0)=0.20\pm 0.05`$. The scalar glueball $`f_0(12001600)`$ is located on the daughter trajectory which predicts the second tensor glueball at $`M3.45`$ GeV. If the Pomeron trajectories in $`(n,M^2)`$ plane are linear, similar to the $`q\overline{q}`$-trajectories, then the next scalar glueball (radial exitation of $`gg`$ gluonium) should be at $`M3.2GeV`$.
Observed glueball states have transformed into broad resonances owing to the accumulation of widths of their neighbours. The existence of a low-lying pseudoscalar glueball is also expected. It is natural to assume that it has also transformed into a broad resonance. Consequently, the question is, where to look for this broad $`0^+`$ state. There are two regions in which we can suspect the existence of a pseudoscalar glueball: in the region of 1700 MeV or much higher, at $`2300`$MeV, see the discussion in (Section 10.5). In it is suggested that the lowest scalar and pseudoscalar glueballs must have roughly equal masses. If so, a $`0^+`$ glueball has to occur in the 1700 MeV region.
The authors are grateful to D.V. Bugg, L.D. Faddeev and S.S. Gershtein for stimulating discussions. The paper was supported by the grant No. 04-02-17091 of the RFFI.
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# 1 Introduction
## 1 Introduction
History of the question. One of the main problems of constructive theory of functions<sup>1</sup><sup>1</sup>1The concept became well known through S.N. Bersntein’s paper \[9, V. 2, p.295-300; p. 349-360\]. is in finding a relationship between differential properties of functions and its structural or constructive characteristics. This topic started to develop more than a century ago and in many cases the research was conducted as follows: authors considered a given functional class and by investigating the properties of its elements obtained embedding theorems with other functional classes. We recommend the recent articles , , for a historical survey.
Below we write the three results, which influenced substantially further research and actually gave rise to the development of new areas within the approximation theory.
$`(\mathrm{A})`$ $`f^{(r)}\text{Lip}\alpha `$ $`E_n(f)=O\left({\displaystyle \frac{1}{n^{r+\alpha }}}\right)`$ $`(0<\alpha <1,r𝐙_+),`$
$`(\mathrm{B})`$ $`f^{(r)}\text{Lip}\alpha `$ $`\omega _{r+1}(f,{\displaystyle \frac{1}{n}})=O\left({\displaystyle \frac{1}{n^{r+\alpha }}}\right)`$ $`(0<\alpha <1,r𝐙_+),`$
$`(\mathrm{C})`$ $`f\text{Lip}\alpha `$ $`\stackrel{~}{f}\text{Lip}\alpha `$ $`(0<\alpha <1).`$
Criterion (A) was proved by <sup>2</sup><sup>2</sup>2 See for a detailed review of the question before the 30ths of the 20th century. D. Jackson (1911, ) and S.B. Stechkin () in the necessity part, and by S. N. Bernstein (1912, , ) and Ch. de la Vallée-Poussin (1919, ) in the sufficiency part.
The theorems of this type are called direct and inverse theorems of approximation theory. Direct theorems for $`L_p,1p\mathrm{}`$ (see e.g. , , ) are written as follows:
$`E_n(f)_p`$ $``$ $`C(k)\omega _k(f,{\displaystyle \frac{1}{n}})_p,k,n𝐍,`$ (1)
$`E_n(f)_p`$ $``$ $`{\displaystyle \frac{C(k)}{n^r}}\omega _k(f^{(r)},{\displaystyle \frac{1}{n}})_p,k,n,r𝐍.`$ (2)
Inverse theorems for $`L_p,1p\mathrm{}`$ (see e.g. , , ):
$$\omega _k(f,\frac{1}{n})_p\frac{C(k)}{n^k}\underset{\nu =0}{\overset{n}{}}(\nu +1)^{k1}E_\nu (f)_p,k,n𝐍,$$
(3)
$$\omega _k(f^{(r)},\frac{1}{n})_pC(k)\left(\frac{1}{n^k}\underset{\nu =0}{\overset{n}{}}(\nu +1)^{k+r1}E_\nu (f)_p+\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r1}E_\nu (f)_p\right),k,n𝐍.$$
(4)
Here and further, the best trigonometric approximation $`E_n(f)_p`$ and the modulus of smoothness $`\omega _k(f,\delta )_p`$ are defined as follows:
$$E_n(f)_p=\mathrm{min}\left(fT_p:T𝐓_n\right),𝐓_n=span\{\mathrm{cos}mx,\mathrm{sin}mx:|m|n\}$$
and
$$\omega _k(f,\delta )_p=\underset{|h|\delta }{sup}\mathrm{}_h^kf(x)_p,$$
(5)
$$\mathrm{}_h^kf(x)=\mathrm{}_h^{k1}\left(\mathrm{}_hf(x)\right)\text{ }\mathrm{}_hf(x)=f(x+h)f(x),$$
respectively.
In the case of $`1<p<\mathrm{}`$ one can write (\[15, p. 210\], , ) the following improvement of estimates (1) and (3):
$`{\displaystyle \frac{C(k)}{n^k}}\left({\displaystyle \underset{\nu =0}{\overset{n}{}}}(\nu +1)^{k\tau 1}E_\nu (f)_p^\tau \right)^{\frac{1}{\tau }}`$ $``$ $`\omega _k(f,{\displaystyle \frac{1}{n}})_p`$
$``$ $`{\displaystyle \frac{C(k)}{n^k}}\left({\displaystyle \underset{\nu =0}{\overset{n}{}}}(\nu +1)^{k\theta 1}E_\nu (f)_p^\theta \right)^{\frac{1}{\theta }}`$
where $`k,n𝐍`$, $`\theta =\mathrm{min}(2,p)`$, and $`\tau =\mathrm{max}(2,p)`$.
In view of estimates (2) and (4), we note that investigation of the question on existence of the $`r`$-th derivative of $`f`$ from a given function space has been initiated by Bernstein . He proved that the condition $`\underset{\nu =1}{\overset{\mathrm{}}{}}\nu ^{r1}E_\nu (f)_{\mathrm{}}<\mathrm{}`$ implies $`f^{(r)}C`$. Later on, for $`L_p(1p\mathrm{})`$, the following results were obtained (see the review and the paper by O.V. Besov ). For convenience, we write these embeddings in terms of the Besov space $`B_{p,\theta }^r`$ and the Sobolev space $`W_p^r`$:
$`B_{p,1}^rW_p^rB_{p,\mathrm{}}^rp=1,\mathrm{},`$
$`B_{p,p}^rW_p^rB_{p,2}^r1<p2,`$
$`B_{p,2}^rW_p^rB_{p,p}^r2p<\mathrm{}.`$
Criterion (B) was proved by A. Zygmund (1945, ). He was one of the first to use the modulus of smoothness concept of an integer order introduced by Bernstein in 1912 (). At present, the moduli of smoothness properties are well-studied (, ) and the result (B) follows from the following inequalities (see \[15, Chapters 2 and 6\], ): ($`1p\mathrm{}`$)
$$\omega _{k+r}(f,\frac{1}{n})_p\frac{C(k,r)}{n^r}\omega _k(f^{(r)},\frac{1}{n})_p,k,r,n𝐍$$
(7)
$$\omega _k(f^{(r)},\frac{1}{n})_pC(k,r)\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r1}\omega _{k+r}(f,\frac{1}{\nu })_p,k,r,n𝐍.$$
(8)
Comparing the last two inequalities and inequalities (2) and (4) we see that from (7) and (8), using (1) and (3), it is easy to get (2) and (4).
We also mention the paper by J. Marcinkiewicz (1938, ), where the following two inequalities were proved:
$$f^{}_pC(p)\left(\underset{0}{\overset{2\pi }{}}\frac{\omega _2(f,u)_p^p}{u^{p+1}}𝑑u\right)^{\frac{1}{p}},1<p2$$
(9)
and
$$\left(\underset{0}{\overset{2\pi }{}}\frac{\omega _2(f,u)_p^p}{u^{p+1}}𝑑u\right)^{\frac{1}{p}}C(p)f^{}_p,2p<\mathrm{}.$$
(10)
It is easy to show \[17, Rem 3.5\] that (8) and (9) are corollaries of the estimate \[17, Th 3.1\]
$$\omega _k(f^{(r)},\frac{1}{n})_pC(k,r)\left(\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r\theta 1}\omega _{k+r}^\theta (f,\frac{1}{\nu })_p\right)^{\frac{1}{\theta }},$$
(11)
where $`k,r,n𝐍`$, $`\theta =\mathrm{min}(2,p)`$, and $`1p<\mathrm{}`$.
Criterion (C) was proved by I.I. Privalov (1919, ). The following inequality, which implies embedding (C), was obtained by A. Zygmund () and N.K. Bary and S.B. Stechkin () ($`p=1,\mathrm{}`$):
$$\omega _k(\stackrel{~}{f}^{(r)},\frac{1}{n})_pC(k,r)\left(n^k\underset{\nu =1}{\overset{n}{}}\nu ^{k+r1}\omega _k(f,\frac{1}{\nu })_p+\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r1}\omega _k(f,\frac{1}{\nu })_p\right),k,n𝐍,r𝐙_+,$$
(12)
Here and further, $`\stackrel{~}{f}`$ denotes the conjugate function to $`f`$ \[59, V. 1, Ch. 2\].
Next, we note the paper by G.H. Hardy and J.S. Littlewood (1928, ) in which seemingly, for the first time, some problems of constructive approximation theory were formulated and solved in terms of embedding theorems. Some historical aspects of this approach were presented in the paper .
Finally, we mention some improvements of above written inequalities for the case of the generalized derivatives and moduli of smoothness, as these estimates are of a particular interest for us: (1) and (3) are proved for any $`k>0`$ in and ; analogues of inequality (2) for the $`f_\beta ^\psi `$-derivatives are shown in \[39, 6.3\], and analogues of (4) and (12) are proved in , -, .
Embedding theorems for functional classes. The results (A) - (C) as well as their generalizations mentioned above can be written as the embedding theorems of the following functional classes:
$`W_p^r`$ $`=\{fL_p:f^{(r)}L_p\},`$
$`\stackrel{~}{W}_p^r`$ $`=\{fL_p:\stackrel{~}{f}^{(r)}L_p\},`$
$`W_p^rH_\alpha [\phi ]`$ $`=\{fW_p^r:\omega _\alpha (f^{(r)},\delta )_p=O[\phi (\delta )]\},`$
$`\stackrel{~}{W}_p^rH_\alpha [\varphi ]`$ $`=\{f\stackrel{~}{W}_p^r:\omega _\alpha (\stackrel{~}{f}^{(r)},\delta )_p=O[\varphi (\delta )]\},`$
$`W_p^rE[\xi ]`$ $`=\{fW_p^r:E_n\left(f^{(r)}\right)_p=O[\xi (1/n)]\}.`$
We will study more general classes such that $`W_p^r`$, $`\stackrel{~}{W}_p^r`$, $`W_p^rH_\alpha [\phi ]`$, $`\stackrel{~}{W}_p^rH_\alpha [\varphi ]`$, $`W^rE_p[\xi ]`$ are their particular cases.
Transformed Fourier series. Let $`L_p=L_p[0,2\pi ](1p<\mathrm{})`$ be a space of $`2\pi `$-periodic measurable functions such that $`|f|^p`$ is integrable, and $`L_{\mathrm{}}C[0,2\pi ]`$ be the space of $`2\pi `$-periodic continuous functions with the uniform norm, that is, $`f_{\mathrm{}}=\mathrm{max}\{|f(x)|,0x2\pi \}.`$
Let the Fourier series of a summable function $`f(x)`$ be written as
$$f(x)\sigma (f):=\frac{a_0(f)}{2}+\underset{\nu =1}{\overset{\mathrm{}}{}}\left(a_\nu (f)\mathrm{cos}\nu x+b_\nu (f)\mathrm{sin}\nu x\right)\underset{\nu =0}{\overset{\mathrm{}}{}}A_\nu (f,x).$$
(13)
The transformed Fourier series for series (13) is defined as follows:
$$\sigma (f,\lambda ,\beta ):=\underset{\nu =1}{\overset{\mathrm{}}{}}\lambda _\nu \left[a_\nu \mathrm{cos}\left(\nu x+\frac{\pi \beta }{2}\right)+b_\nu \mathrm{sin}\left(\nu x+\frac{\pi \beta }{2}\right)\right],$$
where $`\beta 𝐑`$ and $`\lambda =\left\{\lambda _n\right\}`$ is a given sequence of positive numbers.
This definition is well-known in the literature (see, for example, \[59, ch. 12 §8-9\], , , , , , , and remark 3 in this paper). We also note that $`\sigma (f,\lambda ,\beta )`$ coincides up to notations with the Fourier series of so called $`f_\beta ^\psi `$-derivatives, using the terminology of \[39, p. 132\].
Studies of the transformed Fourier series are naturally related to the problems of Fourier multipliers theory (see , \[59, V. 1, Ch. III\], , \[54, Chapter 7\]), summability methods (see \[59, V. 1, Ch. III\], \[12, Chapter 1.2\], \[54, Chapter 8\]) and<sup>3</sup><sup>3</sup>3See also references to \[4, §13, Ch. II\]. the so-called fractional Sobolev classes or the Weyl classes \[39, V. 1, Ch. III\].
The function class
$$W_p^{\lambda ,\beta }=\{fL_p:gL_p,\sigma (g)=\sigma (f,\lambda ,\beta )\}$$
is called the Weyl class (see for example , , ). It is named so, because for $`\lambda _n=n^r,r>0`$ and $`\beta =r`$ the class $`W_p^{\lambda ,\beta }`$ coincides with the class $`W_p^r`$, which is defined in terms of fractional derivatives $`f^{(r)}`$ in the Weyl sense (\[59, V. 2, Ch. XII\]). In the case of $`\lambda _n=n^r,r>0`$ and $`\beta =r+1`$ the class $`W_p^{\lambda ,\beta }`$ coincides with the class $`\stackrel{~}{W}_p^r`$. A function $`g(x)\sigma (f,\lambda ,\beta )`$ is called the $`(\lambda ,\beta )`$-derivative of a function $`f(x)`$ and is denoted by $`f^{(\lambda ,\beta )}(x)`$. Using the terminology of , we have $`f_\beta ^\psi =f^{(\lambda ,\beta )}`$ for $`\psi ^1(k)=\lambda _k`$.
The generalized Weyl-Nikolskii class. In the definition of this functional class we use the modulus of smoothness concept $`\omega _\alpha (f,\delta )_p`$ of fractional <sup>4</sup><sup>4</sup>4 The term ”fractional” can be found in earlier papers ( and ) which used this definition. As in the case of fractional derivatives, the positive number $`\alpha `$ that defines the modulus order is not necessarily rational. order of a function $`f(x)L_p`$, i.e.,
$$\omega _\alpha (f,\delta )_p=\underset{|h|\delta }{sup}\mathrm{}_h^\alpha f(x)_p,$$
where
$$\mathrm{}_h^\alpha f(x)=\underset{\nu =0}{\overset{\mathrm{}}{}}(1)^\nu \left(\genfrac{}{}{0pt}{}{\alpha }{\nu }\right)f(x+(\alpha \nu )h),\alpha >0$$
is the $`\alpha `$-th difference<sup>5</sup><sup>5</sup>5As usual, $`\left(\genfrac{}{}{0pt}{}{\beta }{\nu }\right)=\frac{\beta (\beta 1)\mathrm{}(\beta \nu +1)}{\nu !}`$ for $`\nu >1`$, $`\left(\genfrac{}{}{0pt}{}{\beta }{\nu }\right)=\beta `$ for $`\nu =1`$, and $`\left(\genfrac{}{}{0pt}{}{\beta }{\nu }\right)=1`$ for $`\nu =0`$. of a function $`f`$ with step $`h`$ at the point $`x`$. It is clear that for $`\alpha 𝐍`$ this definition is the same as (5).
Let $`\mathrm{\Phi }_\alpha (\alpha >0)`$ be the class of functions $`\phi (\delta )`$, defined and non-negative on $`(0,\pi ]`$ such that
1. $`\phi (\delta )0(\delta 0)`$,
2. $`\phi (\delta )`$ is non-decreasing,
3. $`\delta ^\alpha \phi (\delta )`$ is non-increasing.
For functions $`\phi \mathrm{\Phi }_\alpha ,\alpha >0`$ and for $`\lambda =\left\{\lambda _n\right\}`$ we define the generalized Weyl-Nikolskii class similarly to the classes $`W_p^rH^\alpha [\phi ]`$ and $`\stackrel{~}{W}_p^rH^\alpha [\varphi ]`$ (see, for example, ):
$$W_p^{\lambda ,\beta }H_\alpha [\phi ]=\{fW_p^{\lambda ,\beta }:\omega _\alpha (f^{(\lambda ,\beta )},\delta )_p=O\left[\phi (\delta )\right],\delta +0\}.$$
It is clear that if $`\lambda _n=n^r,r>0`$ and $`\beta =r`$, then $`W_p^{\lambda ,\beta }H_\alpha [\phi ]W_p^rH_\alpha [\phi ]`$; and if $`\lambda _n=n^r,r>0`$ and $`\beta =r+1`$, then $`W_p^{\lambda ,\beta }H_\alpha [\phi ]\stackrel{~}{W}_p^rH_\alpha [\phi ]`$.
In case $`\lambda _n1`$ and $`\beta =0`$ the class $`W_p^{\lambda ,\beta }H_\alpha [\phi ]`$ coincides with the generalized Lipschitz class $`H_\alpha ^\phi `$, i.e.,
$$H_\alpha ^p[\phi ]=\{fL_p:\omega _\alpha (f,\delta )_p=O\left[\phi (\delta )\right]\delta +0\}.$$
In particular, for $`0<\gamma 1`$,
$$\text{Lip}(\gamma ,L_p)H_1^p[\delta ^\gamma ]=\{fL_p:\omega _1(f,\delta )_p=O\left[\delta ^\gamma \right]\delta +0\}.$$
The problem setting and the structure of the paper. In this paper, we obtain embedding theorems for the Weyl classes $`W_p^{\lambda ,\beta }`$, for the generalized Weyl-Nikolskii classes $`W_p^{\lambda ,\beta }H_\alpha [\phi ]`$ and for the generalized Lipschitz classes $`H_\gamma ^p[\omega ]`$. We show how the parameters $`\alpha `$ and $`\gamma `$ are related to each other depending on the behavior of the sequence $`\{\lambda _n\}`$ and on the choice of the metric $`L_p`$.
This paper is organized as follows. In section 2 we formulate the main theorem. Sections 3 and 4 contain the proofs of the sufficiency and necessity parts of the main theorem respectively. In section 5 we prove several corollaries. In particular, we describe the difference in results for metrics $`L_p,1<p<\mathrm{}`$ and $`L_p,p=1,\mathrm{}`$. Also, the estimates of $`\omega _\gamma (f^{(r)},\delta )_p`$ and $`\omega _\gamma (\stackrel{~}{f}^{(r)},\delta )_p`$ are written in terms of $`\omega _\beta (f,\delta )_p`$ for different values of $`r,\gamma `$, and $`\beta `$. The concluding remarks are given in section 6.
## 2 Embedding theorems for the generalized Lipschitz and <br>Weyl-Nikolskii classes
For $`\lambda =\left\{\lambda _n\right\}_{n𝐍}`$ we define $`\mathrm{}\lambda _n:=\lambda _n\lambda _{n+1}`$; $`\mathrm{}^2\lambda _n:=\mathrm{}\left(\mathrm{}\lambda _n\right)`$.
###### Theorem 1
Let $`\theta =\mathrm{min}(2,p)`$, $`\alpha 𝐑_+,`$ $`\beta 𝐑`$, $`\rho 𝐑_+\{0\}`$ and $`\lambda =\left\{\lambda _n\right\}`$ be a non-decreasing sequence of positive numbers such that $`\left\{n^\rho \lambda _n\right\}`$ is non-increasing.
I. If $`1<p<\mathrm{}`$, then
$`H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$ $``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}^\theta \lambda _n^\theta \right)\omega ^\theta \left({\displaystyle \frac{1}{n}}\right)<\mathrm{},`$ (14)
$`H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$ $``$ $`\{n^{\alpha \theta }{\displaystyle \underset{\nu =1}{\overset{n}{}}}\nu ^{(\rho +\alpha )\theta }(\nu ^{\rho \theta }\lambda _\nu ^\theta (\nu +1)^{\rho \theta }\lambda _{\nu +1}^\theta )\omega ^\theta \left({\displaystyle \frac{1}{\nu }}\right)`$ (15)
$`+`$ $`{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}(\lambda _{\nu +1}^\theta \lambda _\nu ^\theta )\omega ^\theta \left({\displaystyle \frac{1}{\nu }}\right)+\lambda _{n+1}^\theta \omega ^\theta \left({\displaystyle \frac{1}{n+1}}\right)\}^{\frac{1}{\theta }}`$
$`=`$ $`O\left[\phi \left({\displaystyle \frac{1}{n+1}}\right)\right],`$
$`W_p^{\lambda ,\beta }H_{\alpha +\rho }^p[\omega ]`$ $``$ $`{\displaystyle \frac{1}{\lambda _n}}=O\left[\omega \left({\displaystyle \frac{1}{n}}\right)\right],`$ (16)
$`W_p^{\lambda ,\beta }H_\alpha [\phi ]H_{\alpha +\rho }^p[\omega ]`$ $``$ $`{\displaystyle \frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}}=O\left[\omega \left({\displaystyle \frac{1}{n}}\right)\right].`$ (17)
II. Let $`p=1`$ or $`p=\mathrm{}`$.
(a) If $`\mathrm{}^2\lambda _n0`$ or $`\mathrm{}^2\lambda _n0`$, then
$`H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$ $``$ $`|\mathrm{cos}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}\lambda _n\right)\omega \left({\displaystyle \frac{1}{n}}\right)`$ (18)
$`+`$ $`|\mathrm{sin}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\lambda _n{\displaystyle \frac{\omega \left(\frac{1}{n}\right)}{n}}<\mathrm{};`$
and if, additionally, for some $`\tau >0`$ the following inequality holds,
$$\mathrm{}^2\left(\frac{\lambda _n}{n^r}\right)0\text{for}r=\rho +\tau sign\left|\mathrm{sin}\frac{(\beta \rho )\pi }{2}\right|,$$
then
$`H_{\alpha +r}^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$ $``$ $`n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\nu ^{r+\alpha }\left(\nu ^r\lambda _\nu (\nu +1)^r\lambda _{\nu +1}\right)\omega \left({\displaystyle \frac{1}{\nu }}\right)`$ (19)
$`+`$ $`|\mathrm{cos}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\omega \left({\displaystyle \frac{1}{\nu }}\right)`$
$`+`$ $`|\mathrm{sin}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{\omega \left(\frac{1}{\nu }\right)}{\nu }}+\lambda _{n+1}\omega \left({\displaystyle \frac{1}{n+1}}\right)`$
$`=`$ $`O\left[\phi \left({\displaystyle \frac{1}{n+1}}\right)\right].`$
(b) If for $`\beta =2k,k𝐙`$, the condition $`\mathrm{}^2\left(1/\lambda _n\right)0`$ holds, and for $`\beta 2k,k𝐙`$ conditions $`\mathrm{}^2\left(1/\lambda _n\right)0`$ and $`\underset{\nu =n+1}{\overset{\mathrm{}}{}}\frac{1}{\nu \lambda _\nu }\frac{C}{\lambda _n}`$ are fulfilled, then
$$W_p^{\lambda ,\beta }H_{\alpha +\rho }^p[\omega ]\frac{1}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right];$$
(20)
and if, additionally, for some $`\tau >0`$ the following inequality holds,
$$\mathrm{}^2\left(\frac{n^r}{\lambda _n}\right)0\text{or}\mathrm{}^2\left(\frac{n^r}{\lambda _n}\right)0\text{for}r=\rho +\tau sign\left|\mathrm{sin}\frac{(\beta \rho )\pi }{2}\right|,$$
then
$$W_p^{\lambda ,\beta }H_\alpha [\phi ]H_{\alpha +r}^p[\omega ]\frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right].$$
(21)
## 3 Proof of sufficiency in Theorem 1.
We will use the following notations.
Let series (13) be the Fourier series of a function $`f(x)L`$. Then $`S_n(f)`$ denotes the $`n`$-th partial sum of series (13), $`V_n(f)`$ denotes the de la Vallée-Poussin sum and $`K_n(x)`$ is the Fejér kernel, i.e.,
$$S_n(f)=\underset{\nu =0}{\overset{n}{}}A_\nu (x),V_n(f)=\frac{1}{n}\underset{\nu =n}{\overset{2n1}{}}S_\nu (f),K_n(x)=\frac{1}{n+1}\underset{\nu =0}{\overset{n}{}}\left(\frac{1}{2}+\underset{m=1}{\overset{\nu }{}}\mathrm{cos}mx\right).$$
The following lemmas will play an important role in the proof of the main theorem.
###### Lemma 3.1
If $`f(x)L_p,`$ $`1p\mathrm{}`$ and $`\alpha >0`$, then
$$C_1(p,\alpha )\omega _\alpha (f,\frac{1}{n})_p(n^\alpha V_n^{(\alpha )}(f,x))_p+f(x)V_n(f,x)_p)C_2(p,\alpha )\omega _\alpha (f,\frac{1}{n})_p.$$
(22)
If $`f(x)L_p,`$ $`1<p<\mathrm{}`$, then
$$C_1(p,\alpha )\omega _\alpha (f,\frac{1}{n})_p(n^\alpha S_n^{(\alpha )}(f,x))_p+f(x)S_n(f,x)_p)C_2(p,\alpha )\omega _\alpha (f,\frac{1}{n})_p.$$
(23)
Proof of Lemma 3.1. The estimate of $`\omega _\alpha (f,\frac{1}{n})_p`$ from above follows from the inequality (see ) $`\omega _\alpha (T_n,\frac{1}{n})_pC(p,\alpha )n^\alpha T_n^{(\alpha )}_p`$, where $`T_n`$ is a trigonometric polynomial of order $`n`$. Indeed,
$`\omega _\alpha (f,{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(p,\alpha )\left(\omega _\alpha (T_n,{\displaystyle \frac{1}{n}})_p+fT_n_p\right)`$
$``$ $`C(p,\alpha )\left(n^\alpha T_n^{(\alpha )}_p+fT_n_p\right).`$
To estimate $`\omega _\alpha (f,\frac{1}{n})_p`$ from below, we will use the generalized Nikol’skii-Stechkin inequality (see ) $`n^\alpha T_n^{(\alpha )}_pC(p,\alpha )\omega _\alpha (T_n,\frac{1}{n})_p`$ and the generalized Jackson inequality (see for $`\alpha >0`$)
$$E_n(f)_pC(\alpha )\omega _\alpha (f,\frac{1}{n+1})_p.$$
(24)
Also, it is well known that the Vallée-Poussin mean is the near best approximant, i.e.,
$$fV_n(f)_pCE_n(f)_p.$$
(25)
Then
$`n^\alpha V_n^{(\alpha )}(f,x)_p+f(x)V_n(f,x)_p`$ $``$ $`C(p,\alpha )\left(\omega _\alpha (V_n,{\displaystyle \frac{1}{n}})_p+E_n(f)_p\right)`$
$``$ $`C(p,\alpha )\left(\omega _\alpha (f,{\displaystyle \frac{1}{n}})_p+\omega _\alpha (fV_n,{\displaystyle \frac{1}{n}})_p\right)`$
$``$ $`C(p,\alpha )\omega _\alpha (f,{\displaystyle \frac{1}{n}})_p,`$
i.e., (22) is proved. Using
$$fS_n(f)_pC(p)E_n(f)_p$$
(26)
for $`1<p<\mathrm{}`$, we obtain (23) analogously. Lemma 3.1 is proved.
We note that (22) and (23) are the realization results for modulus of smoothness (see the original paper by Z. Ditzian, V. Hristov, K. Ivanov).
###### Lemma 3.2
(). Let $`f(x)L_p,`$ $`p=1,\mathrm{}`$ and let the condition $`\underset{n=1}{\overset{\mathrm{}}{}}n^1E_n(f)_p<\mathrm{}`$ hold. Then $`\stackrel{~}{f}(x)L_p`$ and
$$E_n(\stackrel{~}{f})_pC\left(E_n(f)_p+\underset{k=n+1}{\overset{\mathrm{}}{}}k^1E_k(f)_p\right),n𝐍.$$
###### Lemma 3.3
Let $`p=1,\mathrm{}`$ and let $`\{\lambda _n\}`$ be monotone concave (or convex) sequence. Let
$`T_n(x)`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{n}{}}}a_\nu \mathrm{cos}\nu x+b_\nu \mathrm{sin}\nu x,`$
$`T_n(\lambda ,x)`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{n}{}}}\lambda _\nu \left(a_\nu \mathrm{cos}\nu x+b_\nu \mathrm{sin}\nu x\right).`$
Then for any integer $`M>N+2`$ we have
$$T_M(\lambda ,x)T_N(\lambda ,x)_p\mu (M,N)T_M(x)T_N(x)_p,$$
where
$$\mu (M,N)=\{\begin{array}{cc}2M(\lambda _M\lambda _{M1})+\lambda _{N+1}(N+1)(\lambda _{N+2}\lambda _{N+1}),\hfill & \text{if }\lambda _n\text{ }\text{(}n\text{)}\text{}\mathrm{}^2\lambda _n0;\hfill \\ 2\lambda _M+(N+1)(\lambda _{N+2}\lambda _{N+1})\lambda _{N+1},\hfill & \text{if }\lambda _n\text{ }\text{(}n\text{)}\text{}\mathrm{}^2\lambda _n0;\hfill \\ (N+1)(\lambda _{N+1}\lambda _{N+2})+\lambda _{N+1},\hfill & \text{if }\lambda _n\text{ }\text{(}n\text{)}\text{}\mathrm{}^2\lambda _n0,\hfill \end{array}$$
for $`M=N+1`$, $`\mu (M,N)=\lambda _M`$, and for $`M=N+2`$, $`\mu (M,N)=2|\lambda _{N+1}\lambda _{N+2}|+\lambda _{N+2}.`$
Proof of Lemma 3.3. First we consider the case when $`M>N+2`$. Applying twice Abel’s transformation, we write
$`T_M(\lambda ,x)T_N(\lambda ,x)_p={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{\pi }{\overset{\pi }{}}}\left(T_MT_N\right)(x+u){\displaystyle \underset{\nu =N+1}{\overset{M}{}}}\lambda _\nu \mathrm{cos}\nu udu_p`$
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{\pi }{\overset{\pi }{}}}(T_MT_N)(x+u)\{{\displaystyle \underset{\nu =N+1}{\overset{M2}{}}}(\lambda _\nu 2\lambda _{\nu +1}+\lambda _{\nu +2})(\nu +1)K_\nu (u)`$
$`+`$ $`(\lambda _{N+2}\lambda _{N+1})(N+1)K_N(u)+(\lambda _{M1}\lambda _M)MK_{M1}(u)\}du+\lambda _M(T_MT_N)_p`$
$``$ $`T_M(x)T_N(x)_p\left\{{\displaystyle \underset{\nu =N+1}{\overset{M2}{}}}\right|\lambda _\nu 2\lambda _{\nu +1}+\lambda _{\nu +2}\left|(\nu +1)+\right|\lambda _{M1}\lambda _M|M+\lambda _M\}`$
$`=:`$ $`T_M(x)T_N(x)_pI(M,N).`$
We estimate $`I(M,N)`$ in the case of $`\lambda _n`$ ($`n`$), $`\mathrm{}^2\lambda _n0`$. Then
$`I(M,N)`$ $`=`$ $`{\displaystyle \underset{\nu =N+1}{\overset{M2}{}}}\left(\lambda _\nu 2\lambda _{\nu +1}+\lambda _{\nu +2}\right)(\nu +1)+\left(\lambda _M\lambda _{M1}\right)M+\lambda _M`$
$`=`$ $`(N+1)\left(\lambda _{N+2}\lambda _{N+1}\right)+\left(\lambda _{N+1}\lambda _{M1}\right)+\lambda _M+(2M1)\left(\lambda _M\lambda _{M1}\right)`$
$`=`$ $`(N+1)\left(\lambda _{N+2}\lambda _{N+1}\right)+\lambda _{N+1}+2M\left(\lambda _M\lambda _{M1}\right).`$
If $`\lambda _n`$ ($`n`$), $`\mathrm{}^2\lambda _n0`$, then
$$I(M,N)=(N+1)\left(\lambda _{N+1}\lambda _{N+2}\right)+\left(\lambda _{M1}\lambda _{N+1}\right)+\lambda _M+\left(\lambda _M\lambda _{M1}\right).$$
Finally, if $`\lambda _n`$ ($`n`$), $`\mathrm{}^2\lambda _n0`$, then
$$I(M,N)=(N+1)(\lambda _{N+2}\lambda _{N+1})+\lambda _{N+1}.$$
The estimate for the case of $`M=N+1`$ is trivial and for the case of $`M=N+2`$ immediately follows from the equation
$`T_M(\lambda ,x)T_N(\lambda ,x)`$ $`=`$ $`{\displaystyle \frac{\lambda _M}{\pi }}{\displaystyle _\pi ^\pi }(T_MT_N)(x+u)\left({\displaystyle \underset{m=N+1}{\overset{N+2}{}}}\mathrm{cos}mu\right)𝑑u`$
$`+`$ $`{\displaystyle \frac{\lambda _{N+1}\lambda _{N+2}}{\pi }}{\displaystyle _\pi ^\pi }(T_MT_N)(x+u)\mathrm{cos}(N+1)u𝑑u.`$
The proof of Lemma 3.3 is complete.
###### Lemma 3.4
Let $`p=1,\mathrm{}`$. If $`T_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)=\underset{\nu =2^n}{\overset{2^{n+1}}{}}\left(c_\nu \mathrm{cos}\nu x+d_\nu \mathrm{sin}\nu x\right),`$ then
$$C_1\stackrel{~}{T}_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)_pT_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)_pC_2\stackrel{~}{T}_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)_p.$$
(27)
Proof of Lemma 3.4. We rewrite $`T_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)`$ in the following way
$$T_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)=\underset{\nu =2^n}{\overset{2^{n+1}}{}}\frac{1}{\nu }\left(\nu c_\nu \mathrm{cos}\nu x+\nu d_\nu \mathrm{sin}\nu x\right).$$
Applying Lemma 3.3 and the Bernstein inequality, we have
$`T_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)_p`$ $``$ $`C{\displaystyle \frac{1}{2^n}}{\displaystyle \underset{\nu =2^n}{\overset{2^{n+1}}{}}}\left(\nu c_\nu \mathrm{cos}\nu x+\nu d_\nu \mathrm{sin}\nu x\right)_p`$
$`=`$ $`C{\displaystyle \frac{1}{2^n}}\left({\displaystyle \underset{\nu =2^n}{\overset{2^{n+1}}{}}}d_\nu \mathrm{cos}\nu x+c_\nu \mathrm{sin}\nu x\right)^{}_p`$
$``$ $`C\stackrel{~}{T}_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)_p.`$
Similar reasoning for $`\stackrel{~}{T}_{2^n,\mathrm{\hspace{0.17em}2}^{n+1}}(x)`$ allows us to obtain the left-hand side inequality in (27). The proof of lemma 3.4 is now complete.
Sufficiency in (14) - (21).
I. $`1<p<\mathrm{}`$. In this case, for $`\lambda _n1`$, the Riesz inequality (\[59, V. 1, p. 253\]) $`\stackrel{~}{f}_pC(p)f_p`$ implies
$$f^{(\lambda ,\beta )}_pC(p,\beta )f_p$$
(28)
(Here and henceforth, by $`C(s,t,\mathrm{})`$ we understand positive constants that depend only on $`s,t,\mathrm{}`$ and in general, may be different in different inequalities).
Let the series in the right part of (14) be convergent and $`fH_{\alpha +\rho }^p[\omega ].`$ We will use the following representation
$$\lambda _{2^n}^\theta =\{\begin{array}{cc}\lambda _1^\theta +\underset{\nu =2}{\overset{n+1}{}}\left(\lambda _{2^{\nu 1}}^\theta \lambda _{2^{\nu 2}}^\theta \right),\hfill & \text{if}n1;\hfill \\ \lambda _1^\theta ,\hfill & \text{if}n=0.\hfill \end{array}$$
Applying Minkowski’s inequality, we get $`(`$ here and further $`\mathrm{}_1:=A_1(f,x),\mathrm{}_{n+2}:=\underset{\nu =2^n+1}{\overset{2^{n+1}}{}}A_\nu (f,x)`$, where $`A_\nu (f,x)`$ is from (13) $`)`$
$`I_1`$ $`:=`$ $`\left\{{\displaystyle \underset{0}{\overset{2\pi }{}}}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\lambda _{2^{n1}}^2\mathrm{}_n^2\right]^{\frac{p}{2}}𝑑x\right\}^{\frac{\theta }{p}}`$ (29)
$``$ $`C(p)\left(\lambda _1^\theta \left\{{\displaystyle \underset{0}{\overset{2\pi }{}}}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{}_n^2\right]^{\frac{p}{2}}𝑑x\right\}^{\frac{\theta }{p}}+{\displaystyle \underset{s=2}{\overset{\mathrm{}}{}}}\left(\lambda _{2^{s1}}^\theta \lambda _{2^{s2}}^\theta \right)\left\{{\displaystyle \underset{0}{\overset{2\pi }{}}}\left[{\displaystyle \underset{n=s}{\overset{\mathrm{}}{}}}\mathrm{}_n^2\right]^{\frac{p}{2}}𝑑x\right\}^{\frac{\theta }{p}}\right)^{\frac{1}{\theta }}.`$
By the Littlewood-Paley theorem (see, for example, , p. 349) and inequality (26), we obtain
$$I_1C(p)\{\lambda _1^\theta f_p^\theta +\underset{s=1}{\overset{\mathrm{}}{}}\left(\lambda _{2^s}^\theta \lambda _{2^{s1}}^\theta \right)E_{2^{s1}}^\theta (f)_p\}^{\frac{1}{\theta }}.$$
(30)
Then, both the generalized Jackson inequality (24) and the condition $`fH_{\alpha +\rho }^p[\omega ]`$ imply $`I_1<\mathrm{}.`$ Thus, there exists a function $`gL_p`$ with the Fourier series
$$\underset{n=1}{\overset{\mathrm{}}{}}\lambda _{2^{n1}}\mathrm{}_n,$$
(31)
and $`g_pC(p)I_1`$. We write series (31) in the form of $`\underset{n=1}{\overset{\mathrm{}}{}}\gamma _nA_n(f,x)`$, where $`\gamma _i:=\lambda _i,i=1,2`$ and $`\gamma _\nu :=\lambda _{2^n}`$ for $`2^{n1}+1\nu 2^n(n=2,3,\mathrm{}).`$ Further, we consider the series
$$\underset{n=1}{\overset{\mathrm{}}{}}\lambda _nA_n(f,x)=\underset{n=1}{\overset{\mathrm{}}{}}\gamma _n\mathrm{\Lambda }_nA_n(f,x),$$
(32)
where $`\mathrm{\Lambda }_1=\mathrm{\Lambda }_2=1`$, $`\mathrm{\Lambda }_\nu :=\lambda _\nu /\gamma _n=\lambda _\nu /\lambda _{2^n}`$ for $`2^{n1}+1\nu 2^n(n=2,3,\mathrm{}).`$
Since the sequence $`\left\{\mathrm{\Lambda }_n\right\}`$ satisfies the conditions of the Marcinkiewicz multiplier theorem (, p.346), series (32) is the Fourier series of a function $`f^{(\lambda ,0)}L_p`$ and $`f^{(\lambda ,0)}_pC(p)g_p.`$ Then from inequalities (24), (28) and (30) we get
$`f^{(\lambda ,\beta )}_p`$ $``$ $`C(p,\beta )\{\lambda _1^\theta f_p^\theta +{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}E_{2^{s1}}^\theta (f)_p{\displaystyle \underset{n=2^{s1}}{\overset{2^s1}{}}}\left(\lambda _{n+1}^\theta \lambda _n^\theta \right)\}^{\frac{1}{\theta }}`$
$``$ $`C(p,\beta ,\alpha ,\rho )\{\lambda _1^\theta f_p^\theta +{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}^\theta \lambda _n^\theta \right)\omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{n}})_p\}^{\frac{1}{\theta }},`$
i.e., the sufficiency in (14) is proved.
Let the relation in the right-hand side in (15) hold, and $`fH_{\alpha +\rho }^p[\omega ]`$. Let us prove $`fW_p^{\lambda ,\beta }H_\alpha [\phi ].`$ First, we estimate $`\omega _\alpha (f^{(\lambda ,\beta )},\frac{1}{n})_p.`$ By Lemma 3.1,
$$\omega _\alpha (f^{(\lambda ,\beta )},\frac{1}{n})_pC(p,\alpha )\left(f^{(\lambda ,\beta )}S_n(f^{(\lambda ,\beta )})_p+n^\alpha S_n^{(\alpha )}(f^{(\lambda ,\beta )})_p\right).$$
(34)
Using (3) for the function $`\left(fS_n\right)`$, we have ($`[a]`$ is the integer part of $`a`$)
$`f^{(\lambda ,\beta )}S_n(f^{(\lambda ,\beta )})_pC(p,\beta ,\alpha ,\rho )\{\lambda _{n+1}^\theta fS_n_p^\theta `$ $`+E_{\left[\frac{n}{2}\right]}^\theta (f)_p{\displaystyle \underset{s=1}{\overset{2n}{}}}\left(\lambda _{s+1}^\theta \lambda _s^\theta \right)`$
$`+{\displaystyle \underset{s=n+1}{\overset{\mathrm{}}{}}}(\lambda _{s+1}^\theta \lambda _s^\theta )E_{\left[\frac{s}{2}\right]}^\theta (f)_p\}^{\frac{1}{\theta }}`$
$`C(p,\beta ,\alpha ,\rho )\{\lambda _{n+1}^\theta \omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{n}})_p`$ $`+{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}(\lambda _{\nu +1}^\theta \lambda _\nu ^\theta )\omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{\nu }})_p\}^{\frac{1}{\theta }}.`$ (35)
Further, we estimate the second term of (34). Let $`m`$ be an integer such that $`2^mn+1<2^{m+1}.`$ We will use the identity
$$2^{s\rho \theta }\lambda _{2^s}^\theta =2^{(m+1)\rho \theta }\lambda _{2^{m+1}}^\theta +\underset{\nu =s}{\overset{m}{}}\left(2^{\nu \rho \theta }\lambda _{2^\nu }^\theta 2^{(\nu +1)\rho \theta }\lambda _{2^{\nu +1}}^\theta \right).$$
Then using Lemmas 3.1 and 3.3, we follow the proof of typical estimates (29)-(3). Then we get
$`n^\alpha S_n^{(\alpha )}\left(f^{(\lambda ,\beta )}\right)_p`$ $``$ $`C(p,\beta ,\alpha ,\rho )\{\lambda _{n+1}^\theta \omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{n}})_p`$ (36)
$`+`$ $`n^{\alpha \theta }{\displaystyle \underset{\nu =1}{\overset{n}{}}}(\nu ^{\rho \theta }\lambda _\nu ^\theta (\nu +1)^{\rho \theta }\lambda _{\nu +1}^\theta )\nu ^{(\rho +\alpha )\theta }\omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{\nu }})_p\}^{\frac{1}{\theta }}.`$
Collecting estimates (3), (36) and the inequality in the right-hand side of (15), we get $`fW_p^{\lambda ,\beta }H_\alpha [\phi ].`$
Now we prove that conditions $`\frac{1}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]`$ and $`\frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]`$ are sufficient for embeddings $`W_p^{\lambda ,\beta }H_{\alpha +\rho }^p[\omega ]`$ and $`W_p^{\lambda ,\beta }H_\alpha [\phi ]H_{\alpha +\rho }^p[\omega ]`$, respectively.
Using the Littlewood-Paley and the Marcinkiewicz multiplier theorems and the properties of the sequence $`\left\{\lambda _n\right\}`$ (following the proof of (29)-(3)), we get
$`\omega _{\alpha +\rho }(f,{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(p,\alpha ,\rho )\left(fS_n(f)_p+n^{(\alpha +\rho )}S_n^{(\alpha +\rho )}(f)_p\right)`$
$``$ $`C(p,\beta ,\alpha ,\rho )\left(\lambda _n^1f^{(\lambda ,\beta )}S_n\left(f^{(\lambda ,\beta )}\right)_p+\lambda _n^1n^\alpha S_n^{(\alpha )}(f^{(\lambda ,\beta )})_p\right).`$
Then, by Lemma 3.1, we have the following inequalities
$$\omega _{\alpha +\rho }(f,\frac{1}{n})_pC(p,\beta ,\alpha ,\rho )\lambda _n^1\omega _\alpha (f^{(\lambda ,\beta )},\frac{1}{n})_pC(p,\beta ,\alpha ,\rho )\lambda _n^1f^{(\lambda ,\beta )}_p,$$
where the first one implies sufficiency in (17) and the second one shows sufficiency in (16).
II. $`p=1`$ or $`p=\mathrm{}`$. Let the series in the right-hand side of (18) converge, and let $`fH_{\alpha +\rho }^p[\omega ]`$. Consider the series
$`\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}V_1(\lambda ,f)`$ $``$ $`\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}\stackrel{~}{V_1}(\lambda ,f)`$
$`+`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left\{\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}\left(V_{2^n}(\lambda ,f)V_{2^{n1}}(\lambda ,f)\right)\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}\left(\stackrel{~}{V_{2^n}}(\lambda ,f)\stackrel{~}{V_{2^{n1}}}(\lambda ,f)\right)\right\},`$
where $`V_1(\lambda ,f):=\lambda _1A_1(f,x),`$
$$V_n(\lambda ,f):=\sigma (\lambda ,V_n(f))=\underset{m=1}{\overset{n}{}}\lambda _mA_m(f,x)+\underset{m=n+1}{\overset{2n1}{}}\lambda _m\left(1\frac{mn}{n}\right)A_m(f,x)(n2).$$
Let $`M>N>0`$.
Using the inequality $`fV_n(f)_pCE_n(f)_p`$, the Jackson theorem (24), and the properties of $`\left\{\lambda _n\right\}`$, and following the proof of Lemma 3.3, we get
$`A`$ $`:=`$ $`{\displaystyle \underset{n=N}{\overset{M}{}}}\left[\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}\left(V_{2^{n+1}}(\lambda ,f)V_{2^n}(\lambda ,f)\right)\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}\left(\stackrel{~}{V_{2^{n+1}}}(\lambda ,f)\stackrel{~}{V_{2^n}}(\lambda ,f)\right)\right]_p`$ (38)
$``$ $`{\displaystyle \underset{n=N}{\overset{M}{}}}[|\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}|V_{2^{n+1}}(f)V_{2^n}(f)_p({\displaystyle \underset{m=2^n}{\overset{2^{n+2}1}{}}}|\mathrm{}^2\lambda _m|(m+1)+2^{n+2}|\mathrm{}\lambda _{2^{n+2}}|)`$
$`+`$ $`|\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}|\stackrel{~}{V_{2^{n+1}}}(f)\stackrel{~}{V_{2^n}}(f)_p({\displaystyle \underset{m=2^n}{\overset{2^{n+2}1}{}}}|\mathrm{}^2\lambda _m|(m+1)+2^{n+2}|\mathrm{}\lambda _{2^{n+2}1}|)]`$
$`+`$ $`|\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}|{\displaystyle \underset{n=N}{\overset{M}{}}}\lambda _{2^{n+2}}\left(V_{2^{n+1}}V_{2^n}\right)(f)_p+|\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}|{\displaystyle \underset{n=N}{\overset{M}{}}}\lambda _{2^{n+2}}\left(\stackrel{~}{V_{2^{n+1}}}\stackrel{~}{V_{2^n}}\right)(f)_p`$
$``$ $`C\{\lambda _{2^N}(|\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}|E_{2^N}(f)_p+|\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}|E_{2^N}(\stackrel{~}{f})_p)`$
$`+`$ $`{\displaystyle \underset{n=2^N1}{\overset{\mathrm{}}{}}}(\lambda _{n+1}\lambda _n)(|\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}|\omega _{\alpha +\rho }(f,{\displaystyle \frac{1}{n}})_p+|\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}|\omega _{\alpha +\rho }(\stackrel{~}{f},{\displaystyle \frac{1}{n}})_p)\}.`$
To complete the proof of the sufficiency part in (18), we apply Lemma 3.2, inequality (3) (see for the case $`k>0`$), and inequality (24).
Then the convergence of the series in the right-hand side of (18) and the condition $`fH_{\alpha +\rho }^p[\omega ]`$ imply the fact that the sequence
$$\left\{V_{2^n}(\lambda ,\beta ,f):=\mathrm{cos}\frac{\pi \beta }{2}V_{2^n}(\lambda ,f)\mathrm{sin}\frac{\pi \beta }{2}\stackrel{~}{V_{2^n}}(\lambda ,f)\right\}$$
is fundamental in $`L_p`$. If $`p=1`$, since $`L_1`$ is complete, there exists a subsequence $`\{n_k\}`$ such that $`V_{2^{n_k}}(\lambda ,\beta ,f)`$ converges almost everywhere to a function $`\phi L_1`$. Then from the mean convergence we obtain that, say for cosine coefficients,
$$a_n(\phi )=\frac{1}{\pi }\underset{\pi }{\overset{\pi }{}}\phi (x)\mathrm{cos}nxdx=\underset{k\mathrm{}}{lim}\frac{1}{\pi }\underset{\pi }{\overset{\pi }{}}V_{2^{n_k}}(\lambda ,\beta ,f)\mathrm{cos}nxdx=a_n(f^{(\lambda ,\beta )}).$$
Therefore, $`\sigma (\phi )=\sigma (f^{(\lambda ,\beta )}).`$
For $`p=\mathrm{}`$ the proof is similar. This completes the proof of the sufficiency part of (18).
Let now the condition in the right-hand side of (19) hold and $`fH_{\alpha +r}^p[\omega ]`$. Let us estimate $`\omega _\alpha (f^{(\lambda ,\beta )},\frac{1}{n})_p`$ from above. By Lemma 3.1,
$$\omega _\alpha (f^{(\lambda ,\beta )},\frac{1}{n})_pC(\alpha )\left(f^{(\lambda ,\beta )}V_n(f^{(\lambda ,\beta )})_p+n^\alpha V_n^{(\alpha )}(f^{(\lambda ,\beta )})_p\right).$$
Let us show that
$`f^{(\lambda ,\beta )}V_n(f^{(\lambda ,\beta )})_p`$ $``$ $`C(\beta ,\alpha ,r)(\lambda _n\omega _{\alpha +r}(f,{\displaystyle \frac{1}{n}})_p+|\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}|{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}(\lambda _{\nu +1}\lambda _\nu )\omega _{\alpha +r}(f,{\displaystyle \frac{1}{\nu }})_p`$ (39)
$`+`$ $`|\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}|{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\omega _{\alpha +r}(f,{\displaystyle \frac{1}{\nu }})_p).`$
It has been proved above that
$$A=\underset{n=N}{\overset{M}{}}\left(V_{2^{n+1}}(f^{(\lambda ,\beta )})V_{2^n}(f^{(\lambda ,\beta )})\right)_p.$$
Then for $`2^mn<2^{m+1}`$
$`f^{(\lambda ,\beta )}V_n(f^{(\lambda ,\beta )})_p`$ $``$ $`V_n(f^{(\lambda ,\beta )})V_{2^{m+1}}(f^{(\lambda ,\beta )})_p+{\displaystyle \underset{\nu =m+1}{\overset{\mathrm{}}{}}}\left(V_{2^\nu }(f^{(\lambda ,\beta )})V_{2^{\nu +1}}(f^{(\lambda ,\beta )})\right)_p`$
$`=:`$ $`I_1+I_2.`$
By Lemma 3.4, we have
$`I_1`$ $`=`$ $`\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}\left(V_n(\lambda ,f)V_{2^{m+1}}(\lambda ,f)\right)\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}\left(\stackrel{~}{V_n}(\lambda ,f)\stackrel{~}{V_{2^{m+1}}}(\lambda ,f)\right)_p`$
$``$ $`V_n(\lambda ,f)V_{2^{m+1}}(\lambda ,f)_p+\stackrel{~}{V_n}(\lambda ,f)\stackrel{~}{V_{2^{m+1}}}(\lambda ,f)_p`$
$``$ $`CV_n(\lambda ,f)V_{2^{m+1}}(\lambda ,f)_p`$
and, by Lemma 3.3, we write $`I_1\lambda _n\omega _{\alpha +r}(f,\frac{1}{n})_p`$.
Further, we estimate
$`I_2`$ $``$ $`|\mathrm{cos}{\displaystyle \frac{\pi \beta }{2}}|{\displaystyle \underset{\nu =m+1}{\overset{\mathrm{}}{}}}\left(V_{2^\nu }(\lambda ,f)V_{2^{\nu +1}}(\lambda ,f)\right)_p+|\mathrm{sin}{\displaystyle \frac{\pi \beta }{2}}|{\displaystyle \underset{\nu =m+1}{\overset{\mathrm{}}{}}}\left(\stackrel{~}{V_{2^\nu }}(\lambda ,f)\stackrel{~}{V_{2^{\nu +1}}}(\lambda ,f)\right)_p.`$
As in (38), we write
$`{\displaystyle \underset{\nu =N}{\overset{M}{}}}\left(V_{2^{\nu +1}}V_{2^\nu }\right)(\lambda ,f)_pC\left(\lambda _{2^N}\omega _{\alpha +r}(f,{\displaystyle \frac{1}{2^N}})_p+{\displaystyle \underset{\nu =2^N+1}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\omega _{\alpha +r}(f,{\displaystyle \frac{1}{\nu }})_p\right);`$
$`{\displaystyle \underset{\nu =N}{\overset{M}{}}}\left(\stackrel{~}{V_{2^\nu }}(\lambda ,f)\stackrel{~}{V_{2^{\nu +1}}}(\lambda ,f)\right)_p{\displaystyle \underset{\nu =N}{\overset{M}{}}}\lambda _{2^{\nu +2}}\left(\stackrel{~}{V_{2^\nu }}(f)\stackrel{~}{V_{2^{\nu +1}}}(f)\right)_p`$
$`+{\displaystyle \underset{\nu =N}{\overset{M}{}}}\stackrel{~}{V_{2^\nu }}(f)\stackrel{~}{V_{2^{\nu +1}}}(f)_p({\displaystyle \underset{m=2^\nu }{\overset{2^{\nu +2}1}{}}}|\mathrm{}^2\lambda _m|(m+1)+2^{\nu +2}|\mathrm{}\lambda _{2^{\nu +2}}|)=:I_{21}+I_{22}.`$
Applying Lemma 3.2,
$`I_{21}`$ $``$ $`\lambda _{2^{N+1}}\stackrel{~}{V_{2^{M+1}}}(f)\stackrel{~}{V_{2^N}}(f)_p+{\displaystyle \underset{\nu =N}{\overset{M}{}}}\left(\lambda _{2^{\nu +2}}\lambda _{2^{\nu +1}}\right)\stackrel{~}{V_{2^{\nu +1}}}(f)\stackrel{~}{V_{2^\nu }}(f)_p`$
$``$ $`C\left(\lambda _{2^{N+1}}{\displaystyle \underset{\nu =N}{\overset{\mathrm{}}{}}}E_{2^\nu }(f)_p+{\displaystyle \underset{\nu =N}{\overset{\mathrm{}}{}}}\left(\lambda _{2^{\nu +2}}\lambda _{2^{\nu +1}}\right){\displaystyle \underset{s=\nu }{\overset{\mathrm{}}{}}}E_{2^s}(f)_p\right)`$
$``$ $`C{\displaystyle \underset{\nu =2^N}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\omega _{\alpha +r}(f,{\displaystyle \frac{1}{\nu }})_p,`$
$`I_{22}`$ $``$ $`C{\displaystyle \underset{\nu =N}{\overset{M}{}}}\left(\lambda _{2^{\nu +3}}\lambda _{2^{\nu 1}}\right)E_{2^\nu }(\stackrel{~}{f})_pC{\displaystyle \underset{\nu =2^N}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\omega _{\alpha +r}(f,{\displaystyle \frac{1}{\nu }})_p,`$
and (39) follows.
Repeating the arguments which were used in (38), we estimate $`n^\alpha V_n^{(\alpha )}(f^{(\lambda ,\beta )})_p`$. Using Lemma 3.3 and inequalities (24) and (25), we write
$$V_n^{(\alpha )}(f^{(\lambda ,\beta )})_pC(\beta ,\alpha ,r)\left(n^\alpha \lambda _n\omega _{\alpha +r}(f,\frac{1}{n})_p+\underset{\nu =1}{\overset{n}{}}\left(\frac{\lambda _\nu }{\nu ^r}\frac{\lambda _{\nu +1}}{(\nu +1)^r}\right)\nu ^{\alpha +r}\omega _{\alpha +r}(f,\frac{1}{\nu })_p\right).$$
(40)
Collecting (39) and (40) and using the condition in the right-hand side of (19), we get $`fW_p^{\lambda ,\beta }H_\alpha [\phi ].`$
Let us prove (20). Let $`fW_p^{\lambda ,\beta }.`$ To establish $`fH_{\alpha +\rho }^p(\omega )`$, we will first estimate $`E_n(f)_p`$ for the case $`\mathrm{sin}\frac{\pi \beta }{2}0`$. Then repeating the outline of the proof of estimate (38), and taking into account (25), we get ( $`2^mn<2^{m+1}`$ )
$`E_n(f)_p`$ $``$ $`C{\displaystyle \underset{\nu =m}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\lambda _{2^\nu }}}\left(V_{[2^{\nu 1}]}V_{2^\nu }\right)(f^{(\lambda ,\beta )})_p`$
$``$ $`CE_{[2^{m1}]}(f^{(\lambda ,\beta )})_p{\displaystyle \underset{\nu =m}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\lambda _{2^\nu }}}C(\rho ,p){\displaystyle \frac{f^{(\lambda ,\beta )}_p}{\lambda _n}}.`$
If $`\mathrm{sin}\frac{\pi \beta }{2}=0`$, then it is easy to see that
$$E_n(f)_p\frac{C}{\lambda _n}E_n(f^{(\lambda ,\beta )})_p\frac{C}{\lambda _n}f^{(\lambda ,\beta )}_p.$$
Hence, substituting the obtained bound for $`E_n(f)_p`$ into inequality (3) and using the fact that $`n^\rho \lambda _n^1`$ ($`n`$), we obtain
$$\omega _{\alpha +\rho }(f,\frac{1}{n})_pC(\alpha ,\rho )\frac{1}{n^{\alpha +\rho }}\underset{\nu =0}{\overset{n}{}}\nu ^{\alpha +\rho 1}E_{\nu 1}(f)_p\frac{C(\alpha ,\rho )}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right],$$
i.e., $`fH_{\alpha +\rho }^p(\omega ).`$ This completes the proof of the sufficiency part in (20).
Let the right-hand side part of (21) hold true and $`fH_{\alpha +\rho }^p[\omega ]`$. First let us prove that
$`V_n^{(\alpha +r)}(f)_pC(\alpha ,r){\displaystyle \frac{n^r}{\lambda _n}}V_{2n}^{(\alpha )}(f^{(\lambda ,\beta )})_p.`$ (41)
If $`\beta =\rho +2m`$, and therefore, $`r=\rho `$, then
$$V_n^{(\alpha +r)}(f)=V_n^{(\alpha )}(\{\frac{\nu ^r}{\lambda _\nu }\},f^{(\lambda ,\beta )})$$
and, by Lemma 3.3
$`V_n^{(\alpha +r)}(f)_pC(\alpha ,r){\displaystyle \frac{n^r}{\lambda _n}}V_n^{(\alpha )}(f^{(\lambda ,\beta )})_pC(\alpha ,r){\displaystyle \frac{n^r}{\lambda _n}}V_{2n}^{(\alpha )}(f^{(\lambda ,\beta )})_p.`$
If $`\beta \rho +2m`$, and therefore, $`r>\rho `$, then by Lemma 3.4, we write
$`V_{2^n}^{(\alpha +r)}(f)_p`$ $``$ $`{\displaystyle \underset{\nu =1}{\overset{n}{}}}V_{2^\nu }^{(\alpha +r)}(f)V_{2^{\nu 1}}^{(\alpha +r)}(f)_p+V_1^{(\alpha +r)}(f)_p`$
$``$ $`C{\displaystyle \underset{\nu =1}{\overset{n}{}}}{\displaystyle \frac{2^{\nu r}}{\lambda _{2^\nu }}}V_{2^\nu }^{(\alpha )}(f^{(\lambda ,\beta )})V_{2^{\nu 1}}^{(\alpha )}(f^{(\lambda ,\beta )})_p+{\displaystyle \frac{1}{\lambda _1}}V_1^{(\alpha )}(f^{(\lambda ,\beta )})_p`$
$``$ $`CV_{2^{n+1}}^{(\alpha )}(f^{(\lambda ,\beta )})_p\left({\displaystyle \underset{\nu =1}{\overset{n}{}}}{\displaystyle \frac{2^{\nu r}}{\lambda _{2^\nu }}}+{\displaystyle \frac{1}{\lambda _1}}\right)`$
$``$ $`C(\alpha ,r){\displaystyle \frac{2^{nr}}{\lambda _{2^n}}}V_{2^{n+1}}^{(\alpha )}(f^{(\lambda ,\beta )})_p.`$
Thus, by (41) and the estimate
$$E_n(f)_p\frac{C}{\lambda _n}E_{\left[\frac{n}{4}\right]}(f^{(\lambda ,\beta )})_p$$
we can write
$`\omega _{\alpha +r}(f,{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(\alpha ,r)\left(n^{(\alpha +r)}V_n^{(\alpha +r)}(f)_p+E_n(f)_p\right)`$
$``$ $`{\displaystyle \frac{C(\alpha ,r)}{\lambda _n}}\left(n^\alpha V_n^{(\alpha )}(f^{(\lambda ,\beta )})_p+E_{\left[\frac{n}{4}\right]}(f^{(\lambda ,\beta )})_p\right)`$
$``$ $`{\displaystyle \frac{C(\alpha ,r)}{\lambda _n}}\omega _\alpha (f^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_p=O\left[{\displaystyle \frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}}\right]=O\left[\omega \left({\displaystyle \frac{1}{n}}\right)\right].`$
Thus, the condition $`\frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]`$ is sufficient for the embedding $`W_p^{\lambda ,\beta }H_\alpha [\phi ]H_{\alpha +r}^p[\omega ]`$.
## 4 Proof of necessity in Theorem 1.
We define the trigonometric polynomials $`\tau _{n+1}(x)`$:
$$\tau _{n+1}(x)=\underset{j=1}{\overset{n+1}{}}\alpha _j^n\mathrm{sin}jx,\text{where}\alpha _j^n=\{\begin{array}{cc}\frac{j}{n+2},\hfill & \text{ }1j\frac{n+2}{2}\hfill \\ 1\frac{j}{n+2},\hfill & \text{ }\frac{n+2}{2}jn+1.\hfill \end{array}$$
We will use the following lemmas, as well as Lemmas 3.1-3.2.
###### Lemma 4.1
(). Let series (13) be the Fourier series of a function $`f(x)L_1`$. Then
$$E_n(f)_1C\left|\underset{\nu =n+1}{\overset{\mathrm{}}{}}\frac{b_\nu }{\nu }\right|.$$
###### Lemma 4.2
(\[59, V.1, p. 345; V.2, p. 198\]). Let $`1p<\mathrm{}`$.
(a) If the series $`\underset{\nu =1}{\overset{\mathrm{}}{}}(a_\nu \mathrm{cos}2^\nu x+b_\nu \mathrm{sin}2^\nu x)`$ is the Fourier series of a function $`f(x)L_p`$, then
$$\left\{\underset{\nu =1}{\overset{\mathrm{}}{}}\left(a_\nu ^2+b_\nu ^2\right)\right\}^{\frac{1}{2}}Cf_p.$$
(b) Let $`a_n,b_n(n𝐍)`$ be real numbers such that $`\underset{\nu =1}{\overset{\mathrm{}}{}}\left(a_\nu ^2+b_\nu ^2\right)<\mathrm{}.`$ Then the series $`\underset{\nu =1}{\overset{\mathrm{}}{}}(a_\nu \mathrm{cos}2^\nu x+b_\nu \mathrm{sin}2^\nu x)`$ is the Fourier series of a function $`f(x)L_p`$, and at the same time
$$f_pC\left\{\underset{\nu =1}{\overset{\mathrm{}}{}}\left(a_\nu ^2+b_\nu ^2\right)\right\}^{\frac{1}{2}}.$$
###### Lemma 4.3
(\[4, Ch.11,§12\]). If the series $`\underset{\nu =1}{\overset{\mathrm{}}{}}(a_\nu \mathrm{cos}2^\nu x+b_\nu \mathrm{sin}2^\nu x)`$, $`a_\nu ,b_\nu 0`$ is the Fourier series of a function $`f(x)L_{\mathrm{}}`$, then
$$C_1\underset{\xi =n}{\overset{\mathrm{}}{}}(a_\xi +b_\xi )E_{2^n1}(f)_{\mathrm{}}C_2\underset{\xi =n}{\overset{\mathrm{}}{}}(a_\xi +b_\xi ).$$
We will use the following definitions. Let $`\omega ()\mathrm{\Phi }_\alpha `$.
A sequence $`\psi `$ is called $`Q_{\alpha ,\theta }(\omega )`$-sequence if
$`0`$ $`<`$ $`\psi _nn^\alpha \omega \left({\displaystyle \frac{1}{n}}\right),\psi _n(n)`$ (42)
$`C_1\omega \left({\displaystyle \frac{1}{n}}\right)`$ $``$ $`\left\{{\displaystyle \underset{\nu =n}{\overset{\mathrm{}}{}}}\nu ^{\alpha \theta 1}\psi _\nu ^\theta \right\}^{\frac{1}{\theta }}C_2\omega \left({\displaystyle \frac{1}{n}}\right).`$ (43)
A sequence $`\epsilon `$ is called $`q_{\alpha ,\theta }(\omega )`$-sequence if
$`0`$ $`<`$ $`\epsilon _n\omega \left({\displaystyle \frac{1}{n+1}}\right),\epsilon _n(n)`$ (44)
$`C_1\omega \left({\displaystyle \frac{1}{n+1}}\right)`$ $``$ $`\left\{(n+1)^{\alpha \theta }{\displaystyle \underset{\nu =1}{\overset{n+1}{}}}\nu ^{\alpha \theta 1}\epsilon _\nu ^\theta \right\}^{\frac{1}{\theta }}C_2\omega \left({\displaystyle \frac{1}{n+1}}\right).`$ (45)
Necessity in (14) - (21).
We prove the necessity part by constructing corresponding examples. The proof consists of eight steps.
I. $`1<p<\mathrm{}`$. Step 1. Let us show the necessity part in (14).
Let $`\omega ()\mathrm{\Phi }_{\alpha +\rho }`$ and $`\theta =\mathrm{min}(2,p)`$. We will construct a $`Q_{\alpha +\rho ,\theta }(\omega )`$-sequence $`\psi `$.
Assume that integers $`1=n_1<n_2<\mathrm{}<n_s`$ are chosen. Then as $`n_{s+1}`$ we take the minimum number $`N>n_s`$ such that
$$\omega \left(\frac{1}{N}\right)<\frac{1}{2}\omega \left(\frac{1}{n_s}\right)\omega \left(\frac{1}{N1}\right).$$
We set
$$\psi _n=\{\begin{array}{cc}n_s^{\rho +\alpha }\omega \left(\frac{1}{n_s}\right),\hfill & \text{if}n_sn<n_{s+1},s=1,2,\mathrm{};\hfill \\ 0,\hfill & \text{if}n=0.\hfill \end{array}$$
It is easy to see that this sequence is what we need.
Let $`H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$ and let the series in (14) be divergent. By means of properties of sequence $`\left\{\psi _n\right\}`$, we have
$`\mathrm{}={\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}^\theta \lambda _\nu ^\theta \right)\omega ^\theta \left({\displaystyle \frac{1}{\nu }}\right)`$ $``$ $`C(\alpha ,\rho ,\theta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}^\theta \lambda _\nu ^\theta \right){\displaystyle \underset{m=\nu }{\overset{\mathrm{}}{}}}m^{(\alpha +\rho )\theta 1}\psi _m^\theta `$
$``$ $`C(\alpha ,\rho ,\theta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\lambda _\nu ^\theta \nu ^{(\alpha +\rho )\theta 1}\psi _\nu ^\theta .`$
Step 1(a): $`2p<\mathrm{}`$. We consider the series
$$\underset{\nu =1}{\overset{\mathrm{}}{}}2^{\nu (\alpha +\rho )}\left(\psi _{2^\nu }^2\psi _{2^{\nu 1}}^2\right)^{\frac{1}{2}}\mathrm{cos}2^\nu x.$$
(46)
Since
$`{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}2^{2\nu (\alpha +\rho )}\left(\psi _{2^\nu }^2\psi _{2^{\nu 1}}^2\right)`$ $``$ $`{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(\psi _{2^\nu }^2\psi _{2^{\nu 1}}^2\right){\displaystyle \underset{\xi =\nu }{\overset{\mathrm{}}{}}}2^{2\xi (\alpha +\rho )}`$ (47)
$``$ $`{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}2^{2\nu (\alpha +\rho )}\psi _{2^\nu }^2C\omega ^2(1),`$
then, by Zygmund’s Lemma 4.2, series (46) is the Fourier series of a function $`f_1(x)L_p`$. Applying Lemmas 3.1 and 4.2, we get
$$C(\alpha ,\rho )\omega _{\alpha +\rho }(f_1,\frac{1}{2^n})_p2^{n(\alpha +\rho )}\left(\underset{\nu =1}{\overset{n}{}}a_\nu ^22^{2(\alpha +\rho )\nu }\right)^{\frac{1}{2}}+\left(\underset{\nu =n+1}{\overset{\mathrm{}}{}}a_\nu ^2\right)^{\frac{1}{2}}=:I_1+I_2,$$
where $`a_\nu =2^{(\alpha +\rho )}\left(\psi _{2^\nu }^2\psi _{2^{\nu 1}}^2\right)^{\frac{1}{2}}.`$ Similarly to estimate (47), by means of (42), we get
$$I_12^{n(\alpha +\rho )}\psi _{2^n}\omega \left(\frac{1}{2^n}\right),$$
and by means of (43),
$$I_2C(\alpha ,\rho )\left(\underset{\nu =n+1}{\overset{\mathrm{}}{}}2^{2\nu (\alpha +\rho )}\psi _{2^\nu }^2\right)^{\frac{1}{2}}C(\alpha ,\rho )\omega \left(\frac{1}{2^n}\right).$$
Thus, $`f_1(x)H_{\alpha +\rho }^p[\omega ].`$ Then from our assumption, $`f_1(x)W_p^{\lambda ,\beta }`$. On the other hand,
$$f_1^{(\lambda ,\beta )}_pC(\alpha ,\rho ,\theta )\left(\underset{\nu =1}{\overset{\mathrm{}}{}}\lambda _\nu ^2\nu ^{2(\alpha +\rho )1}\psi _\nu ^2\right)^{\frac{1}{2}}=\mathrm{}.$$
This contradiction proves the convergence of series in (14).
Step 1(b): $`1<p2`$. Consider series<sup>6</sup><sup>6</sup>6Series of this type was considered in .
$$\psi _1\mathrm{cos}x+\underset{\nu =1}{\overset{\mathrm{}}{}}2^{\nu (\alpha +\rho )}2^{\nu (\frac{1}{p}1)}\left(\psi _{2^\nu }^p\psi _{2^{\nu 1}}^p\right)^{\frac{1}{p}}\underset{\mu =2^{\nu 1}+1}{\overset{2^\nu }{}}\mathrm{cos}\mu x.$$
(48)
Using the Jensen inequality $`\left(\underset{n=1}{\overset{\mathrm{}}{}}a_n^\alpha \right)^{1/\alpha }\left(\underset{n=1}{\overset{\mathrm{}}{}}a_n^\beta \right)^{1/\beta }`$ ($`a_n0`$ and $`\mathrm{\hspace{0.17em}0}<\beta \alpha <\mathrm{}`$), we write
$`{\displaystyle \underset{0}{\overset{2\pi }{}}}\left[{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(2^{\nu (\alpha +\rho )}2^{\nu (\frac{1}{p}1)}\left(\psi _{2^\nu }^p\psi _{2^{\nu 1}}^p\right)^{\frac{1}{p}}{\displaystyle \underset{\mu =2^{\nu 1}+1}{\overset{2^\nu }{}}}\mathrm{cos}\mu x\right)^2\right]^{\frac{p}{2}}𝑑x`$
$``$ $`{\displaystyle \underset{0}{\overset{2\pi }{}}}\left[{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}2^{\nu p(\alpha +\rho )}2^{\nu (1p)}\left(\psi _{2^\nu }^p\psi _{2^{\nu 1}}^p\right)\left|{\displaystyle \underset{\mu =2^{\nu 1}+1}{\overset{2^\nu }{}}}\mathrm{cos}\mu x\right|^p\right]𝑑x`$
$``$ $`C(p){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(\psi _{2^\nu }^p\psi _{2^{\nu 1}}^p\right)2^{\nu p(\alpha +\rho )}C(p)\omega ^p(1),`$
because of $`C_1(p)2^{\nu (p1)}\underset{\mu =2^{\nu 1}+1}{\overset{2^\nu }{}}\mathrm{cos}\mu x_p^pC_2(p)2^{\nu (p1)}`$.
By the Littlewood-Paley theorem (see \[59, Vol. 2, p. 349\]), there exists a function $`f_2L_p`$ with the Fourier series (48). One can easily check that $`f_2H_{\alpha +\rho }^p[\omega ].`$ By our assumption, $`f_2(x)W_p^{\lambda ,\beta }`$. On the other hand, Paley’s theorem on Fourier coefficients \[59, V.2, p. 182\] implies that for $`f_2L_p`$
$`f_2^{(\lambda ,\beta )}_p^p`$ $``$ $`C(p){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}2^{\nu p(\alpha +\rho )}2^{\nu (1p)}\left(\psi _{2^\nu }^p\psi _{2^{\nu 1}}^p\right){\displaystyle \underset{\mu =2^{\nu 1}+1}{\overset{2^\nu }{}}}\lambda _\mu ^p\mu ^{p2}`$
$``$ $`C(\alpha ,\rho ,p){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(\psi _{2^\nu }^p\psi _{2^{\nu 1}}^p\right){\displaystyle \underset{\xi =\nu }{\overset{\mathrm{}}{}}}\left(2^{\xi (\alpha +\rho )p}\lambda _{2^\xi }^p2^{(\xi +1)(\alpha +\rho )p}\lambda _{2^{\xi +1}}^p\right)`$
$``$ $`C_1(\alpha ,\rho ,p){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\psi _\nu ^p\lambda _\nu ^p\nu ^{p(\alpha +\rho )1}C_2(\alpha ,\rho ,p)\psi _1^p\lambda _2^p2^{p(\alpha +\rho )1}=\mathrm{}.`$
This contradiction shows that the series in the right-hand side of (14) converges. This completes the proof of the necessity part of (14).
Step 2. Let us prove the necessity in (15) for the case $`2p<\mathrm{}.`$
We notice that using Lemmas 3.1 and 4.2, we get for
$$f(x)\underset{\nu =1}{\overset{\mathrm{}}{}}(a_\nu \mathrm{cos}2^\nu x+b_\nu \mathrm{sin}2^\nu x)$$
the following relation
$$\omega _\alpha (f,\frac{1}{2^m})_p\left(2^{2m\alpha }\underset{\nu =1}{\overset{m}{}}(a_\nu ^2+b_\nu ^2)2^{2\nu \alpha }\right)^{\frac{1}{2}}+\left(\underset{\nu =m+1}{\overset{\mathrm{}}{}}(a_\nu ^2+b_\nu ^2)\right)^{\frac{1}{2}}.$$
(49)
Let $`\omega ()\mathrm{\Phi }_{\alpha +\rho }`$. Then one can construct<sup>7</sup><sup>7</sup>7See, for example, , . a sequence $`\epsilon `$ such that it is a $`q_{\alpha +\rho ,\theta }(\omega )`$-sequence. In this case, we consider
$$\epsilon _0+\left(\epsilon _1^2\epsilon _2^2\right)^{\frac{1}{2}}\mathrm{cos}x+\underset{\nu =1}{\overset{\mathrm{}}{}}\left(\epsilon _{2^\nu }^2\epsilon _{2^{\nu +1}}^2\right)^{\frac{1}{2}}\mathrm{cos}2^\nu x.$$
(50)
Repeating the argument used for series (46), we obtain that series (50) is the Fourier series of a function $`f_3L_p`$. Since $`E_{2^{n1}}(f_3)_pC(p)\epsilon _{2^n}`$, then by (3) and (45) we have $`f_3H_{\alpha +\rho }^p[\omega ]`$. We define $`f_{13}:=f_1+f_3.`$ Then $`f_{13}H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$. It is easy to see from (49) that
$$C(\alpha ,\beta )\omega _\alpha (f_{13}^{(\lambda ,\beta )},\frac{1}{n+1})_p\omega _\alpha (f_1^{(\lambda ,\beta )},\frac{1}{n+1})_p+\omega _\alpha (f_3^{(\lambda ,\beta )},\frac{1}{n+1})_p.$$
(51)
Let us estimate $`\omega _\alpha (f_1^{(\lambda ,\beta )},\frac{1}{n+1})_p`$. Applying (49) and using the properties of the sequence $`\{\psi _\nu \}`$, we write ($`2^mn+1<2^{m+1}`$)
$`\omega _\alpha ^2(f_1^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_p`$ $``$ $`C(\alpha ,\beta ){\displaystyle \underset{\nu =m}{\overset{\mathrm{}}{}}}\psi _{2^\nu }^22^{2\nu (r+\alpha )}\left({\displaystyle \underset{k=m}{\overset{\nu }{}}}(\lambda _{2^k}^2\lambda _{2^{k1}}^2)+\lambda _{2^{m1}}^2\right)`$ (52)
$``$ $`C(\alpha ,\beta )\left(\lambda _{2^m}^2\omega ^2\left({\displaystyle \frac{1}{2^m}}\right)+{\displaystyle \underset{k=m}{\overset{\mathrm{}}{}}}(\lambda _{2^k}^2\lambda _{2^{k1}}^2)\omega ^2\left({\displaystyle \frac{1}{2^k}}\right)\right)`$
$``$ $`C(\alpha ,\beta )\left(\lambda _{n+1}^2\omega ^2\left({\displaystyle \frac{1}{n+1}}\right)+{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}^2\lambda _\nu ^2\right)\omega ^2\left({\displaystyle \frac{1}{\nu }}\right)\right).`$
Let us proceed to the estimation of $`\omega _\alpha (f_3^{(\lambda ,\beta )},\frac{1}{n+1})_p.`$ Using (49), we have ($`2^mn+1<2^{m+1}`$):
$`\omega _\alpha ^2(f_3^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_p`$ $``$ $`C(\alpha ,\beta )2^{2m\alpha }{\displaystyle \underset{\nu =0}{\overset{m}{}}}2^{2\nu \alpha }\lambda _{2\nu }^2\left(\epsilon _{2^\nu }^2\epsilon _{2^{\nu +1}}^2\right)`$ (53)
$``$ $`C_1(\alpha ,\beta )2^{2m\alpha }{\displaystyle \underset{\nu =0}{\overset{m}{}}}2^{2\nu \alpha }\lambda _{2^\nu }^2\epsilon _{2^\nu }^2C_2(\alpha ,p)\lambda _{2^{m+1}}^2\epsilon _{2^{m+1}}^2.`$
Then Jackson’s inequality implies
$`\omega _\alpha ^2(f_3^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_p`$ $``$ $`C(\alpha )E_{2^m1}^2\left(f_3^{(\lambda ,\beta )}\right)_p`$ (54)
$``$ $`C(\alpha ,\beta ){\displaystyle \underset{\nu =m}{\overset{\mathrm{}}{}}}\lambda _{2\nu }^2\left(\epsilon _{2^\nu }^2\epsilon _{2^{\nu +1}}^2\right)`$
$``$ $`C(\alpha ,\beta )\lambda _{2^m}^2\epsilon _{2^m}^2.`$
Applying estimates (53) and (54), we get
$$\omega _\alpha (f_3^{(\lambda ,\beta )},\frac{1}{n+1})_pC(\alpha ,\beta )\left((n+1)^{2\alpha }\underset{\nu =1}{\overset{n+1}{}}\lambda _\nu ^2\nu ^{2\alpha 1}\epsilon _\nu ^2\right)^{\frac{1}{2}}.$$
(55)
Further, (45) and $`\nu ^\rho \lambda _\nu `$ allow us to write the estimates
$`\lambda _{n+1}^2\omega ^2\left({\displaystyle \frac{1}{n+1}}\right)`$ $`+`$ $`(n+1)^{2\alpha }{\displaystyle \underset{\nu =1}{\overset{n+1}{}}}\nu ^{2(\rho +\alpha )}\omega ^2\left({\displaystyle \frac{1}{\nu }}\right)\left(\lambda _\nu ^2\nu ^{2\rho }\lambda _{\nu +1}^2(\nu +1)^{2\rho }\right)`$ (56)
$``$ $`C(\alpha ,\rho )(n+1)^{2\alpha }{\displaystyle \underset{\nu =1}{\overset{n+1}{}}}\nu ^{2\alpha 1}\lambda _\nu ^2\epsilon _\nu ^2.`$
Combining estimates (51), (52), (55), (56), and $`\omega _\alpha (f_{13}^{(\lambda ,\beta )},\frac{1}{n+1})_p=O\left[\phi \left(\frac{1}{n}\right)\right]`$, we obtain the condition in the right-hand side part of (15).
Step 3. Here we show the necessity in (15) for the case of $`1<p<2.`$ In this case the proof is similar to the proof of the case $`2p<\mathrm{}`$. The only difference is that we use Paley’s theorem on Fourier coefficients instead of Zygmund’s Lemma 4.2. In this case we consider the sum of $`f_2(x)`$ and the following function
$$\epsilon _0+\left(\epsilon _1^p\epsilon _2^p\right)^{\frac{1}{p}}\mathrm{cos}x+\underset{\nu =0}{\overset{\mathrm{}}{}}2^{\nu \left(\frac{1}{p}1\right)}\left(\epsilon _{2^{\nu +1}}^p\epsilon _{2^{\nu +2}}^p\right)^{\frac{1}{p}}\underset{\mu =2^\nu +1}{\overset{2^{\nu +1}}{}}\mathrm{cos}\mu x.$$
(57)
Step 4. To prove the necessity in (16) and (20), we consider the general case of $`1p\mathrm{}.`$ Let $`\mathrm{\Phi }`$ be the class of all decreasing null-sequences. It is clear that
$$\frac{1}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]\gamma =\{\gamma _n\}\mathrm{\Phi }\frac{\gamma _n}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right].$$
Let us assume that $`\frac{\gamma _n}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]`$ does not hold for all $`\gamma \mathrm{\Phi }`$ and $`W_p^{\lambda ,\beta }H_{\alpha +\rho }^p[\omega ]`$. Then there exist $`\gamma =\{\gamma _n\}\mathrm{\Phi }`$ and $`\{C_n\mathrm{}\}`$ such that $`\frac{\gamma _{m_n}}{\lambda _{m_n}}C_n\omega \left(\frac{1}{m_n}\right).`$ Further, we choose a subsequence $`\left\{m_{n_k}\right\}`$ such that $`\frac{m_{n_{k+1}}}{m_{n_k}}2`$ and $`\gamma _{m_{n_k}}2^k.`$ Consider the series
$$\underset{k=0}{\overset{\mathrm{}}{}}\frac{\gamma _{m_{n_k}}}{\lambda _{m_{n_k}}}\mathrm{cos}(m_{n_k}+1)x.$$
(58)
Since $`\underset{k=0}{\overset{\mathrm{}}{}}\frac{\gamma _{m_{n_k}}}{\lambda _{m_{n_k}}}\frac{1}{\lambda _{m_{n_0}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{2^k}<\mathrm{},`$ there exists a function $`f_4L_p`$ with the Fourier series (58). Because of $`\underset{k=0}{\overset{\mathrm{}}{}}\gamma _{m_{n_k}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{2^k}<\mathrm{},`$ we have $`f_4^{(\lambda ,\beta )}L_p`$, i.e., $`f_4W_p^{\lambda ,\beta }`$.
On the other hand, using (24) and $`E_{n1}(f)_pC\left(|a_n|+|b_n|\right)`$,
$`\omega _{\alpha +\rho }(f_4,{\displaystyle \frac{1}{m_{n_k}}})_pC(\alpha ,\rho )E_{m_{n_k}}(f_4)_pC(\alpha ,\rho ){\displaystyle \frac{\gamma _{m_{n_k}}}{\lambda _{m_{n_k}}}}C(\alpha ,\rho )C_{n_k}\omega \left({\displaystyle \frac{1}{m_{n_k}}}\right),`$
i.e., $`f_4H_{\alpha +\rho }^p[\omega ].`$ This contradiction proves that the condition $`\frac{1}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]`$ is necessary for $`W_p^{\lambda ,\beta }H_{\alpha +\rho }^p[\omega ]`$.
Step 5. To prove the necessity in (17) and (21), we verify that for any $`\rho >0`$ and $`1p\mathrm{}`$,
$$W_p^{\lambda ,\beta }H_\alpha [\phi ]H_{\alpha +\rho }^p[\omega ]\frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right].$$
(59)
First, we remark that
$$\frac{\phi \left(\frac{1}{n}\right)}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right]\gamma =\{\gamma _n\}\mathrm{\Phi }\frac{\gamma _n\phi \left(\frac{1}{n}\right)}{\lambda _n}=O\left[\omega \left(\frac{1}{n}\right)\right].$$
We assume that the relation in the right-hand side of (59) does not hold. Then there exist $`\gamma =\{\gamma _n\}\mathrm{\Phi }`$ and $`\{C_n\mathrm{}\}`$ such that $`\frac{\gamma _{m_n}\phi \left(\frac{1}{m_n}\right)}{\lambda _{m_n}}C_n\omega \left(\frac{1}{m_n}\right).`$ We choose a subsequence $`\left\{m_{n_k}\right\}`$ such that $`\frac{m_{n_{k+1}}}{m_{n_k}}2`$ $`\gamma _{m_{n_k}}2^k.`$ Because of $`\underset{k=0}{\overset{\mathrm{}}{}}\gamma _{m_{n_k}}\phi \left(\frac{1}{m_n}\right)\phi \left(\frac{1}{m_0}\right)\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{2^k}<\mathrm{},`$ there exists a function $`f_5L_p`$ with the Fourier series
$$\underset{k=0}{\overset{\mathrm{}}{}}\gamma _{m_{n_k}}\phi \left(\frac{1}{m_{n_k}}\right)\mathrm{cos}(m_{n_k}+1)x.$$
(60)
For $`m_{n_k}n<m_{n_{k+1}}`$ using Lemmas 3.1 and 4.3, we have
$`\omega _\alpha (f_5,{\displaystyle \frac{1}{n}})_p`$ $``$ $`C\omega _\alpha (f_5,{\displaystyle \frac{1}{n}})_{\mathrm{}}C\left(n^\alpha {\displaystyle \underset{s=0}{\overset{k}{}}}\gamma _{m_{n_s}}\phi \left({\displaystyle \frac{1}{m_{n_s}}}\right)m_{n_s}^\alpha +{\displaystyle \underset{s=k+1}{\overset{\mathrm{}}{}}}\gamma _{m_{n_s}}\phi \left({\displaystyle \frac{1}{m_{n_s}}}\right)\right)`$
$``$ $`C\left(\phi \left({\displaystyle \frac{1}{n}}\right){\displaystyle \underset{s=0}{\overset{k}{}}}\gamma _{m_{n_s}}+\phi \left({\displaystyle \frac{1}{n}}\right){\displaystyle \underset{s=k+1}{\overset{\mathrm{}}{}}}\gamma _{m_{n_s}}\right)C\phi \left({\displaystyle \frac{1}{n}}\right).`$
Then $`f_5H_\alpha ^p[\phi ]`$, i.e., setting $`\frac{1}{\lambda }:=\{\frac{1}{\lambda _1},\frac{1}{\lambda _2},\frac{1}{\lambda _3},\mathrm{}\}`$, we have $`f_5^{(\frac{1}{\lambda },\beta )}W_p^{\lambda ,\beta }H_\alpha [\phi ].`$
On the other hand,
$`\omega _{\alpha +\rho }(f_5^{(\frac{1}{\lambda },\beta )},{\displaystyle \frac{1}{m_{n_k}}})_pCE_{m_{n_k}}(f_5^{(\frac{1}{\lambda },\beta )})_pC{\displaystyle \frac{\gamma _{m_{n_k}}\phi \left(\frac{1}{m_{n_k}}\right)}{\lambda _{m_{n_k}}}}CC_{n_k}\omega \left({\displaystyle \frac{1}{m_{n_k}}}\right),`$
i.e., $`f_5^{(\frac{1}{\lambda },\beta )}H_{\alpha +\rho }^p[\omega ].`$ This contradicts our assumption. The proof of the necessity part in (16)-(17) and (20)-(21) is now complete.
II. $`p=1`$ or $`p=\mathrm{}`$. Step 6. Let us prove the necessity in (18). Let $`H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$ and the series in (18) be divergent.
Step 6(a): $`\mathrm{sin}\frac{\beta \pi }{2}0`$. In this case a divergence of the series in (18) is equivalent to the divergence of the series $`\underset{n=1}{\overset{\mathrm{}}{}}\frac{\lambda _n}{n}\omega \left(\frac{1}{n}\right).`$
Let $`p=1`$. We take a $`q_{\alpha +\rho ,\mathrm{\hspace{0.17em}1}}(\omega )`$-sequence $`\epsilon `$ and consider the series
$$\underset{\nu =1}{\overset{\mathrm{}}{}}\left(\epsilon _\nu \epsilon _{\nu +1}\right)K_\nu (x).$$
(61)
This series converges in $`L_1`$ (see ) to a function $`f_6(x)`$ and $`E_n(f_6)_1=O\left(\epsilon _n\right)`$. Applying (3) and (45), we get $`f_6H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$. One can also rewrite (61) in the following form
$$\underset{\nu =1}{\overset{\mathrm{}}{}}a_\nu \mathrm{cos}\nu x,\text{where}a_\nu =\epsilon _\nu \nu \underset{j=\nu }{\overset{\mathrm{}}{}}\frac{\epsilon _j\epsilon _{j+1}}{j+1}.$$
By Lemma 4.1,
$`f_6^{(\lambda ,\beta )}_1`$ $``$ $`C(\beta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}a_\nu =C(\beta )\left({\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\epsilon _\nu {\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \underset{j=\nu }{\overset{\mathrm{}}{}}}{\displaystyle \frac{\epsilon _j\epsilon _{j+1}}{j+1}}\right)`$
$`=`$ $`C(\beta )\left({\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\epsilon _\nu {\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}(a_\nu a_{\nu +1})\lambda _\nu \right)`$
$``$ $`C_1(\beta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\epsilon _\nu C_2(\beta )\left(\lambda _1a_1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(\lambda _{n+1}\lambda _n)a_n\right).`$
Further, using (44) and (45), we get
$`{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\epsilon _\nu `$ $``$ $`C(\rho ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu }}\omega \left({\displaystyle \frac{1}{\nu }}\right)C(\rho ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\lambda _\nu \nu ^{(\alpha +\rho )1}{\displaystyle \underset{m=1}{\overset{\nu }{}}}m^{\alpha +\rho 1}\epsilon _m`$
$`=`$ $`C(\rho ){\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}m^{\alpha +\rho 1}\epsilon _m{\displaystyle \underset{\nu =m}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu }{\nu ^\rho }}\nu ^{\alpha 1}C(\rho ,\alpha ){\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _m}{m}}\epsilon _m.`$
Then
$$f_6^{(\lambda ,\beta )}_1C_1(\alpha ,\rho ,\beta )\underset{\nu =1}{\overset{\mathrm{}}{}}\lambda _\nu \nu ^1\omega \left(\frac{1}{\nu }\right)C_2(\alpha ,\rho ,\beta )\left(\lambda _1a_1+\underset{\nu =1}{\overset{\mathrm{}}{}}(\lambda _{\nu +1}\lambda _\nu )a_\nu \right).$$
(62)
On the other hand, using monotonicity of $`\{a_\nu \}`$ and Lemma 4.1, we have
$`C(\beta ,\rho )f_6^{(\lambda ,\beta )}_1`$ $``$ $`C(\rho )\left(\lambda _1a_1+{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _\nu a_\nu }{\nu }}\right)C(\rho )\left(\lambda _1a_1+{\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}\lambda _{2^\nu }a_{2^\nu }\right)`$ (63)
$``$ $`\lambda _1a_1+{\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}a_{2^\nu }(\lambda _{2^{\nu +1}}\lambda _{2^\nu })\lambda _1a_1+{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}(\lambda _{\nu +1}\lambda _\nu )a_\nu .`$
Using (62) and (63), we get
$$f_6^{(\lambda ,\beta )}_1C(\alpha ,\rho ,\beta )\underset{\nu =1}{\overset{\mathrm{}}{}}\lambda _\nu \nu ^1\omega \left(\frac{1}{\nu }\right)=\mathrm{}.$$
The obtained contradiction implies the convergence of series in (18).
Let now $`p=\mathrm{}`$. Define the function (see ) $`f_7(x)=\underset{\nu =1}{\overset{\mathrm{}}{}}\epsilon _\nu \nu ^1\mathrm{sin}\nu x,`$ where $`\epsilon `$ is a $`q_{\alpha +\rho ,\mathrm{\hspace{0.17em}1}}(\omega )`$-sequence. Then $`E_n(f_7)_{\mathrm{}}C\epsilon _{n+1}`$. Using (3) and (45), we get $`f_7H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$. On the other hand,
$`\mathrm{}=f_7^{(\lambda ,\beta )}_{\mathrm{}}`$ $``$ $`C(\beta ){\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}2^{\nu (\alpha +\rho )}\epsilon _{2^\nu }{\displaystyle \frac{\lambda _{2^\nu }}{2^{\nu (\alpha +\rho )}}}`$
$``$ $`C(\alpha ,\rho ,\beta ){\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}2^{\nu (\alpha +\rho )}\epsilon _{2^\nu }{\displaystyle \underset{m=\nu }{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda _{2^m}}{2^{m(\alpha +\rho )}}}`$
$``$ $`C(\alpha ,\rho ,\beta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\lambda _\nu \nu ^1\omega \left({\displaystyle \frac{1}{\nu }}\right).`$
The obtained contradiction proves the convergence of the series in (18).
Step 6(b): $`\mathrm{sin}\frac{\beta \pi }{2}=0`$. Let the series in (18) be divergent. We consider only the non-trivial case of $`\rho >0.`$ Let $`\epsilon `$ be a $`q_{\alpha +\rho ,\mathrm{\hspace{0.17em}1}}(\omega )`$-sequence. By means of the properties $`\{\lambda _n\}`$, we have
$`{\displaystyle \underset{\nu =2}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\omega \left({\displaystyle \frac{1}{\nu }}\right)`$ $``$ $`C(\alpha ,\rho ){\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}(\epsilon _{2^s}\epsilon _{2^{s+1}})[{\displaystyle \underset{m=0}{\overset{s}{}}}2^{m(\alpha +\rho )}{\displaystyle \underset{\nu =m}{\overset{s}{}}}2^{\nu (\alpha +\rho )}(\lambda _{2^\nu }\lambda _{2^{\nu 1}})`$ (64)
$`+`$ $`{\displaystyle \underset{m=0}{\overset{s}{}}}2^{m(\alpha +\rho )}{\displaystyle \underset{\nu =s+1}{\overset{\mathrm{}}{}}}2^{\nu (\alpha +\rho )}(\lambda _{2^\nu }\lambda _{2^{\nu 1}})]`$
$``$ $`C(\alpha ,\rho ){\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}\left(\epsilon _\nu \epsilon _{\nu +1}\right)\lambda _\nu .`$
Let $`p=1`$. Then the series
$$\underset{\nu =1}{\overset{\mathrm{}}{}}\left(\epsilon _\nu \epsilon _{\nu +1}\right)\tau _\nu (x)$$
(65)
converges () to $`f_8L_1`$ and $`E_n(f_8)_1C\epsilon _{n+1}`$. Hence, we have $`f_8H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$.
We write series (65) in the following way
$$\underset{\nu =1}{\overset{\mathrm{}}{}}b_\nu \mathrm{sin}\nu x,\text{where}b_\nu =\underset{j=\nu }{\overset{2\nu 2}{}}\left(1\frac{\nu }{j+1}\right)\left(\epsilon _\nu \epsilon _{\nu +1}\right)+\underset{j=2\nu 1}{\overset{\mathrm{}}{}}\frac{\nu }{j+1}\left(\epsilon _\nu \epsilon _{\nu +1}\right).$$
By Lemma 4.1, we write
$$f_8^{(\lambda ,\beta )}_1C(\beta )\underset{\nu =1}{\overset{\mathrm{}}{}}\lambda _\nu \frac{b_\nu }{\nu }C(\beta )\underset{\nu =0}{\overset{\mathrm{}}{}}\lambda _\nu \left(\epsilon _\nu \epsilon _{\nu +1}\right).$$
This contradicts (64). Thus, the series in (18) converges.
Let now $`p=\mathrm{}`$. Let $`\psi `$ be a $`Q_{\alpha +\rho ,1}`$-sequence. Then, we define
$$f_9(x)=\psi _1\mathrm{cos}x+\underset{\nu =1}{\overset{\mathrm{}}{}}2^{\nu (\rho +\alpha )}\left(\psi _{2^\nu }\psi _{2^{\nu 1}}\right)\mathrm{cos}2^\nu x.$$
Similarly, as for the function $`f_1`$ in the case of $`2p<\mathrm{},`$ it is easy to check that $`f_9H_{\alpha +\rho }^p[\omega ]W_p^{\lambda ,\beta }`$. Applying Lemma 4.3, we get
$`f_9^{(\lambda ,\beta )}_{\mathrm{}}`$ $``$ $`C(\beta )\left(\lambda _1\psi _1+{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\lambda _{2^\nu }2^{\nu (\rho +\alpha )}\left(\psi _{2^\nu }\psi _{2^{\nu 1}}\right)\right)`$
$``$ $`C(\beta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\lambda _\nu \nu ^{(\rho +\alpha )1}\psi _\nu `$
$``$ $`C(\beta ){\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\omega (1/\nu )=\mathrm{},`$
which contradicts $`f_9W_p^{\lambda ,\beta }`$. This completes the proof of the necessity in (18).
Step 7. We will show the necessity in (19) for the case of $`\mathrm{sin}\frac{\pi \beta }{2}0`$. Let $`H_{\alpha +r}^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$. We take a $`q_{\alpha +r,\mathrm{\hspace{0.17em}1}}(\omega )`$-sequence $`\epsilon `$. Then (45) holds for $`\alpha +r`$ instead of $`\alpha `$ and for $`1`$ instead of $`\theta `$. Since $`\mathrm{sin}\frac{\pi \beta }{2}0`$,
$`J`$ $`:=`$ $`\lambda _{n+1}\omega \left({\displaystyle \frac{1}{n+1}}\right)+n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\nu ^{r+\alpha }\left(\nu ^r\lambda _\nu (\nu +1)^r\lambda _{\nu +1}\right)\omega \left({\displaystyle \frac{1}{\nu }}\right)`$ (66)
$`+`$ $`|\mathrm{cos}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\omega \left({\displaystyle \frac{1}{\nu }}\right)+|\mathrm{sin}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{\omega \left(\frac{1}{\nu }\right)}{\nu }}`$
$``$ $`C(\alpha ,\beta ,r)({\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{\epsilon _\nu }{\nu }}+n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\lambda _\nu \epsilon _\nu \nu ^{\alpha 1})=:C(\alpha ,\beta ,r)(J_1+J_2).`$
We will use several times the following evident relations:
$$\omega _\alpha (f,\frac{1}{n})_p\omega _\alpha (f_+,\frac{1}{n})_p+\omega _\alpha (f_{},\frac{1}{n})_p,\text{where}f_\pm (x):=\frac{f(x)\pm f(x)}{2}.$$
(67)
Step 7( ): $`p=\mathrm{}`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}0`$. In this case by Lemma 3.1, we have $`(`$ $`f_{7+}^{(\lambda ,\beta )}:=\left(f_7^{(\lambda ,\beta )}\right)_+`$ $`)`$
$`\omega _\alpha (f_7^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_p`$ $``$ $`\omega _\alpha (f_{7+}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_pC(\alpha )n^\alpha V_n^{(\alpha )}(f_{7+}^{(\lambda ,\beta )})()_p`$ (68)
$``$ $`C(\alpha )n^\alpha \left|V_n^{(\alpha )}(f_{7+}^{(\lambda ,\beta )})(0)\right|_pC(\alpha ,\beta )J_2.`$
Using (24) and $`\underset{k=2n}{\overset{\mathrm{}}{}}a_k4E_n(f)_{\mathrm{}}`$ (see ), we write
$`\omega _\alpha (f_7^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_pC(\alpha )E_{[\frac{n}{2}]}(f_{7+}^{(\lambda ,\beta )})_pC(\alpha ,\beta )J_1.`$ (69)
It is proved above that $`f_7H_{\alpha +r}^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ].`$ Collecting inequalities (68), (69), and (66), we obtain the estimate in the right-hand side of (19).
Step 7(b): $`p=\mathrm{}`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}=0`$. If $`\mathrm{cos}\frac{\pi \beta }{2}0,`$ then we use (69) and
$`\omega _\alpha (f_7^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(\alpha )n^\alpha V_n^{(\alpha )}(f_7^{(\lambda ,\beta )})()_pC(\alpha ,\beta )J_2.`$
If $`\mathrm{cos}\frac{\pi \beta }{2}=0,`$ then $`f_7=\pm f_{7+}`$ and
$`\omega _\alpha (f_7^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_pC(\alpha )E_{[\frac{n}{2}]}(f_{7+}^{(\lambda ,\beta )})_pC(\alpha ,\beta )J_1.`$
To obtain the estimate of $`J_2`$, we define
$$f_{10}(x)=\frac{\epsilon _0}{2}+\left(\epsilon _1\epsilon _2\right)\mathrm{cos}x+\underset{\nu =1}{\overset{\mathrm{}}{}}\left(\epsilon _{2^\nu }\epsilon _{2^{\nu +1}}\right)\mathrm{cos}2^\nu x.$$
It is clear that $`E_n(f_{10})_p\epsilon _{n+1}`$. Then using (45), we get $`f_{10}H_{\alpha +r}^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$. Further, applying $`\omega _\alpha (f,\frac{1}{n})_pC(\alpha )n^\alpha V_n^{(\alpha )}(f)_p`$ and Lemma 4.3, we write ($`2^mn+1<2^{m+1}`$)
$`\omega _\alpha (f_{10}^{(\lambda ,\beta )},{\displaystyle \frac{1}{2^m}})_p`$ $``$ $`C(\alpha ,\beta )2^{m\alpha }{\displaystyle \underset{\nu =0}{\overset{m}{}}}\lambda _{2^\nu }2^{\nu \alpha }\left(\epsilon _{2^\nu }\epsilon _{2^{\nu +1}}\right)`$ (70)
$``$ $`C_1(\alpha ,\beta ,r)2^{m\alpha }{\displaystyle \underset{\nu =0}{\overset{m}{}}}\lambda _{2^\nu }2^{\nu \alpha }\epsilon _{2^\nu }C_2(\alpha ,\beta ,r)\lambda _{2^{m+1}}\epsilon _{2^{m+1}}.`$
At the same time, $`\omega _\alpha \left({}_{}{}^{(\lambda ,\beta )},\frac{1}{2^m}\right)_pC(\alpha ,\beta )\lambda _{2^{m+1}}\epsilon _{2^{m+1}}.`$ Then we have
$`\omega _\alpha (f_{10}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_pC(\alpha ,\beta ,r)J_2`$ (71)
and applying (67),
$`C(\alpha ,\beta ,r)\left[J_1+J_2\right]`$ $``$ $`\omega _\alpha (f_{10}^{(\lambda ,\beta )},{\displaystyle \frac{1}{2^m}})_p+\omega _\alpha (f_7^{(\lambda ,\beta )},{\displaystyle \frac{1}{2^m}})_p`$
$``$ $`\omega _\alpha (\left(f_7+f_{10}\right)^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_p=O\left[\phi \left({\displaystyle \frac{1}{n}}\right)\right].`$
The necessity in (19) follows.
Step 7(c): $`p=1`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}0`$. In this case we use the function $`f_6`$.
It is known that $`f_6H_{\alpha +r}^p[\omega ]`$ and by Lemma 4.1, we have
$`\omega _\alpha (f_6^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_1`$ $``$ $`C(\alpha )E_n(f_6^{(\lambda ,\beta )})_pC(\alpha ,\beta ){\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{a_\nu }{\nu }}`$
$``$ $`C(\alpha ,\beta )\left({\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{\epsilon _\nu }{\nu }}a_{n+1}\lambda _n{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\left(\lambda _\nu \lambda _{\nu 1}\right)a_\nu \right).`$
On the other hand,
$$\omega _\alpha (f_6^{(\lambda ,\beta )},\frac{1}{n})_1C(\alpha ,\beta ,r)\left(a_{n+1}\lambda _n+\underset{\nu =n+1}{\overset{\mathrm{}}{}}\left(\lambda _\nu \lambda _{\nu 1}\right)a_\nu \right).$$
Therefore, the last two inequalities imply $`\omega _\alpha (f_6^{(\lambda ,\beta )},\frac{1}{n})_1C(\alpha ,\beta ,r)J_1`$. Moreover,
$`\omega _\alpha (f_6^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(\alpha )n^\alpha V_n^{(\alpha )}(f_6^{(\lambda ,\beta )})()_pC(\alpha ,\beta )J_2.`$ (72)
Step 7(d): $`p=1`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}=0`$. If $`\mathrm{cos}\frac{\pi \beta }{2}0,`$ then we use $`\omega _\alpha (f_6^{(\lambda ,\beta )},\frac{1}{n})_1C(\alpha ,\beta ,r)J_1`$ and
$`\omega _\alpha (f_6^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(\alpha )n^\alpha V_n^{(\alpha )}(f_{6+}^{(\lambda ,\beta )})()_pC(\alpha ,\beta )J_2.`$
If $`\mathrm{cos}\frac{\pi \beta }{2}=0,`$ we consider $`f_6+f_8`$. Using Lemmas 3.1 and 4.1, we get
$`\omega _\alpha (f_8^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_1`$ $``$ $`C(\alpha ,\beta )n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\lambda _\nu \nu ^{\alpha 1}b_\nu `$
$``$ $`C_1(\alpha ,\beta ,r)n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\lambda _\nu \nu ^{\alpha 1}\epsilon _{\nu 1}C_2(\alpha ,\beta ,r)\lambda _n\epsilon _n.`$
Since $`\omega _\alpha (f_8^{(\lambda ,\beta )},\frac{1}{n})_1C(\alpha ,\beta ,r)\lambda _n\epsilon _n`$, then $`\omega _\alpha (f_8^{(\lambda ,\beta )},\frac{1}{n})_1C(\alpha ,\beta ,r)J_2`$. Thus,
$$C(\alpha ,\beta ,r)\left[J_1+J_2\right]\omega _\alpha (\left(f_6+f_8\right)^{(\lambda ,\beta )},\frac{1}{n})_p=O\left[\phi \left(\frac{1}{n}\right)\right],$$
i.e., the necessity in (19) follows.
Step 8. We prove the necessity in (19) in the case of $`\mathrm{sin}\frac{\pi \beta }{2}=0`$. Let $`H_{\alpha +r}^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$ and let $`\epsilon `$ be a $`q_{\alpha +r,\mathrm{\hspace{0.17em}1}}(\omega )`$-sequence. Since $`\mathrm{sin}\frac{\pi \beta }{2}=0`$, from (45) it follows that
$`\lambda _{n+1}\omega \left({\displaystyle \frac{1}{n+1}}\right)`$ $`+`$ $`n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\nu ^{r+\alpha }\left(\nu ^r\lambda _\nu (\nu +1)^r\lambda _{\nu +1}\right)\omega \left({\displaystyle \frac{1}{\nu }}\right)`$ (73)
$`+`$ $`|\mathrm{cos}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\omega \left({\displaystyle \frac{1}{\nu }}\right)`$
$``$ $`C(\alpha ,\beta ,r)\left({\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu \left(\epsilon _\nu \epsilon _{\nu +1}\right)+n^\alpha {\displaystyle \underset{\nu =1}{\overset{n}{}}}\lambda _\nu \epsilon _\nu \nu ^{\alpha 1}\right)`$
$`=:`$ $`C(\alpha ,\beta ,r)\left(J_3+J_4\right).`$
Step 8(a): $`p=\mathrm{}`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}0`$. Applying the Jackson inequality and Lemma 4.3, we get
$`\omega _\alpha (f_{10}^{(\lambda ,\beta )},{\displaystyle \frac{1}{2^m}})_pC(\alpha ,\beta ){\displaystyle \underset{\nu =m}{\overset{\mathrm{}}{}}}\lambda _{2^\nu }\left(\epsilon _{2^\nu }\epsilon _{2^{\nu +1}}\right).`$ (74)
We also note that by Lemma (4.3), (71) holds for all $`\alpha >0`$. This and (74) allow us to write $`\omega _\alpha (f_{10}^{(\lambda ,\beta )},\frac{1}{n+1})_pC(\alpha ,\beta )\left(J_3+J_4\right)`$. Using condition (73) and $`f_{10}H_{\alpha +r}^p[\omega ]W_p^{\lambda ,\beta }H_\alpha [\phi ]`$, we obtain the relation in the right-hand side of (19).
Step 8(b): $`p=\mathrm{}`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}=0`$. Then we consider $`f_{10}`$ and $`f_{11}:=\stackrel{~}{f_{10}}`$. It is clear that $`f_{11}L_p`$ and $`f_{10}+f_{11}H_{\alpha +r}^p[\omega ]`$. Besides,
$`\omega _\alpha (f_{10}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_pC(\alpha ,\beta )J_3,`$
$`\omega _\alpha (f_{11}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_pC(\alpha ,\beta )J_4`$
and
$`\omega _\alpha (f_{10}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_p+\omega _\alpha (f_{11}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_p\omega _\alpha (\left(f_{10}+f_{11}\right)^{(\lambda ,\beta )},{\displaystyle \frac{1}{n+1}})_p.`$
Step 8(c): $`p=1`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}0`$. Since $`f_8^{(\lambda ,\beta )}(x)\pm \underset{\nu =1}{\overset{\mathrm{}}{}}\lambda _\nu b_\nu \mathrm{sin}\nu x,`$ then by Lemmas 3.1 and 4.1, we write (see Step 7(d))
$`\omega _\alpha (f_8^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_pC(\alpha )n^\alpha V_n^{(\alpha )}(f_8^{(\lambda ,\beta )})()_pC(\alpha ,\beta )J_4.`$ (75)
Further, using Lemma 4.1 and the Jackson inequality (24), we have
$`\omega _\alpha (f_8^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_1`$ $``$ $`C(\alpha ,\beta ){\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{b_\nu }{\nu }}`$
$``$ $`C(\alpha ,\beta ){\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \underset{j=2\nu 1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\epsilon _j\epsilon _{j+1}}{j+1}}`$
$``$ $`C(\alpha ,\beta ,r){\displaystyle \underset{j=4n1}{\overset{\mathrm{}}{}}}\left(\epsilon _j\epsilon _{j+1}\right)\lambda _\nu .`$
Using the properties of the modulus of smoothness, we get (19).
Step 8(d): $`p=1`$ and $`\mathrm{cos}\frac{\pi \alpha }{2}=0`$. We use the fact that $`f_{12}:=\stackrel{~}{f_8}L_1`$ and $`E_n(f_{12})_pC\epsilon _{n+1}`$ (see ). Then $`f_8+f_{12}H_{\alpha +r}^p[\omega ]`$ and
$`\omega _\alpha (f_8^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_1`$ $``$ $`C(\alpha ,\beta )J_3,`$
$`\omega _\alpha (f_{12}^{(\lambda ,\beta )},{\displaystyle \frac{1}{n}})_1`$ $``$ $`C(\alpha ,\beta )J_4.`$
Theorem 1 is fully proved.
## 5 Corollaries. Estimates of transformed Fourier series.
Theorem 1 actually provides estimates of the norms and moduli of smoothness of the transformed Fourier series, i.e., the estimates of $`\phi _p`$ and $`\omega _\alpha (\phi ,\delta )_p`$, where $`\phi \sigma (f,\lambda )`$ in terms of $`\omega _\gamma (f,\delta )_p`$. Analyzing the obtained results, one can see that the following two conditions play a crucial role for these estimates. The first is the behavior of the transforming sequence $`\{\lambda _n\}`$ and the second is the choice between the considered space (as the Riesz inequality (28) holds for $`L_p,`$ $`1<p<\mathrm{}`$ and no such inequality exists for $`L_p,`$ $`p=1,\mathrm{}`$).
We will investigate in detail some important examples for $`L_p,`$ $`1<p<\mathrm{}`$ and for $`L_p,`$ $`p=1,\mathrm{}`$, separately.
1. The case of $`\mathrm{𝟏}<𝐩<\mathrm{}`$.
###### Theorem 2
Let $`1<p<\mathrm{}`$, $`\theta =\mathrm{min}(2,p)`$ $`\tau =\mathrm{max}(2,p)`$, $`\alpha 𝐑_+`$, $`\rho 𝐑_+\{0\}`$ and $`\lambda =\left\{\lambda _n\right\}`$ be a non-decreasing sequence of positive numbers such that $`\left\{n^\rho \lambda _n\right\}`$ is non-increasing.
I. If for $`fL_p^0`$ the series
$$\underset{n=1}{\overset{\mathrm{}}{}}\left(\lambda _{n+1}^\theta \lambda _n^\theta \right)\omega _{\alpha +\rho }^\theta (f,\frac{1}{n})_p$$
converges, then there exists a function $`\phi L_p^0`$ with the Fourier series $`\sigma (f,\lambda )`$, and
$`\phi _p`$ $``$ $`C(p,\lambda ,\alpha ,\rho )\{\lambda _1^\theta f_p^\theta +{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}^\theta \lambda _n^\theta \right)\omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{n}})_p\}^{\frac{1}{\theta }},`$ (76)
$`\omega _\alpha (\phi ,{\displaystyle \frac{1}{n+1}})_p`$ $``$ $`C(p,\lambda ,\alpha ,\rho )\{n^{\alpha \theta }{\displaystyle \underset{\nu =1}{\overset{n}{}}}\nu ^{(\rho +\alpha )\theta }(\nu ^{\rho \theta }\lambda _\nu ^\theta (\nu +1)^{\rho \theta }\lambda _{\nu +1}^\theta )\omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{\nu }})_p`$ (77)
$`+`$ $`{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}(\lambda _{\nu +1}^\theta \lambda _\nu ^\theta )\omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{\nu }})_p+\lambda _{n+1}^\theta \omega _{\alpha +\rho }^\theta (f,{\displaystyle \frac{1}{n}})_p\}^{\frac{1}{\theta }}.`$
II. If for $`fL_p^0`$ there exists a function $`\phi L_p`$ with the Fourier series $`\sigma (f,\lambda )`$, then
$`\{\lambda _1^\tau f_p^\tau +{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}^\tau \lambda _n^\tau \right)\omega _{\alpha +\rho }^\tau (f,{\displaystyle \frac{1}{n}})_p\}^{\frac{1}{\tau }}`$ $``$ $`C(p,\lambda ,\alpha ,\rho )\phi _p,`$ (78)
$`\{n^{\alpha \tau }{\displaystyle \underset{\nu =1}{\overset{n}{}}}\nu ^{(\rho +\alpha )\tau }(\nu ^{\rho \tau }\lambda _\nu ^\tau (\nu +1)^{\rho \tau }\lambda _{\nu +1}^\tau )\omega _{\alpha +\rho }^\tau (f,{\displaystyle \frac{1}{\nu }})_p`$ $`+`$ (79)
$`{\displaystyle \underset{\nu =n+2}{\overset{\mathrm{}}{}}}(\lambda _{\nu +1}^\tau \lambda _\nu ^\tau )\omega _{\alpha +\rho }^\tau (f,{\displaystyle \frac{1}{\nu }})_p+\lambda _{n+1}^\tau \omega _{\alpha +\rho }^\tau (f,{\displaystyle \frac{1}{n}})_p\}^{\frac{1}{\tau }}`$ $``$ $`C(p,\lambda ,\alpha ,\rho )\omega _\alpha (\phi ,{\displaystyle \frac{1}{n+1}})_p,`$
$`\omega _{\alpha +\rho }(f,{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(p,\lambda ,\alpha ,\rho ){\displaystyle \frac{\phi _p}{\lambda _n}},`$ (80)
$`\omega _{\alpha +\rho }(f,{\displaystyle \frac{1}{n}})_p`$ $``$ $`C(p,\lambda ,\alpha ,\rho ){\displaystyle \frac{\omega _\alpha (\phi ,\frac{1}{n})_p}{\lambda _n}}.`$ (81)
Inequalities (76)-(77) and (80)-(81) were actually proved in Theorem 1 (see the proof of sufficiency in the part I). The estimates (78)-(79) are proved analogously, using Lemma 3.1, the theorems by Littlewood-Paley, Marcinkiewicz, and the Minkowski’s inequality (see also and ).
An important corollary of Theorem 2 is the following.
###### Corollary 1
Let $`1<p<\mathrm{}`$, $`\theta =\mathrm{min}(2,p)`$, and $`\tau =\mathrm{max}(2,p)`$. Then for any $`k,r>0`$ we have
$$\begin{array}{c}C_1\left\{\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r\tau 1}\omega _{k+r}^\tau (f,\frac{1}{\nu })_p\right\}^{\frac{1}{\tau }}\omega _k(f^{(r)},\frac{1}{n})_pC_2\left\{\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r\theta 1}\omega _{k+r}^\theta (f,\frac{1}{\nu })_p\right\}^{\frac{1}{\theta }},\hfill \end{array}$$
(82)
where $`C_1=C_1(p,k,r),C_2=C_2(p,k,r),n𝐍.`$
See also , where the right-hand side estimate in (82) was shown for integers $`k`$ and $`r`$. The last two inequalities provide sharper bounds in the sense of order than (7) and (8). Indeed, using properties of the modulus of smoothness and the Jensen inequality, we have
$`n^r\omega _{k+r}(f,{\displaystyle \frac{1}{n}})_pC(k,r)\left\{{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\nu ^{r\tau 1}\omega _{k+r}^\tau (f,{\displaystyle \frac{1}{\nu }})_p\right\}^{\frac{1}{\tau }};`$
$`\left\{{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\nu ^{r\theta 1}\omega _{k+r}^\theta (f,{\displaystyle \frac{1}{\nu }})_p\right\}^{\frac{1}{\theta }}C(k,r){\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\nu ^{r1}\omega _{k+r}(f,{\displaystyle \frac{1}{\nu }})_p.`$
Example. Let $`\psi (t)=t^r\mathrm{ln}^A(1/t)`$ and $`2p<\mathrm{},\frac{1}{2}<A<1.`$ If $`\omega _{k+r}(f,t)_p\psi (t)`$, then inequalities (7) and (8) give only $`C\mathrm{ln}^A(1/t)\omega _k(f^{(r)},t)_p`$. At the same time, (82) implies $`C_1\mathrm{ln}^{A+1/p}(1/t)\omega _k(f^{(r)},t)_pC_2\mathrm{ln}^{A+1/2}(1/t)`$, which is sharper.
Proof of Corollary 1 follows from (77) and (79) with $`r=\rho `$, because if $`fL_p,`$ $`1<p<\mathrm{}`$, then one can assume that $`f^{(r)}\sigma (f,\lambda )`$ for $`\{\lambda _n=n^r\}`$.
The estimates (7), (8) and (82) show that it is natural to estimate $`\omega _\alpha (f^{(\gamma )},\frac{1}{n})_p`$ in terms of $`\omega _{\alpha +r}(f,\frac{1}{\nu })_p`$. Further analysis allowed us to distinguish three different types of such estimates. It will be convenient for us to write inequalities in the integral form:
1. $`\gamma =r`$ (see Corollary 1)<sup>8</sup><sup>8</sup>8Here and further, $`\tau =\mathrm{max}(2,p),\theta =\mathrm{min}(2,p)`$. If $`A_1CA_2`$, $`C1`$, we write $`A_1A_2`$. Also, if $`A_1A_2`$ and $`A_2A_1`$, then $`A_1A_2`$.
$$\left\{\underset{0}{\overset{\delta }{}}t^{r\tau 1}\omega _{r+\alpha }^\tau (f,t)_p𝑑t\right\}^{\frac{1}{\tau }}\omega _\alpha (f^{(r)},\delta )_p\left\{\underset{0}{\overset{\delta }{}}t^{r\theta 1}\omega _{r+\alpha }^\theta (f,t)_p𝑑t\right\}^{\frac{1}{\theta }};$$
(83)
2. $`\gamma =r\epsilon ,`$$`0<\epsilon <r`$ (see Theorem 2 for $`\rho =r`$ and $`\lambda _n=n^{r\epsilon }`$):
$$\left\{\underset{0}{\overset{\delta }{}}t^{(r\epsilon )\tau 1}\omega _{r+\alpha }^\tau (f,t)_p𝑑t+\delta ^{\alpha \tau }\underset{\delta }{\overset{1}{}}t^{(r\epsilon +\alpha )\tau 1}\omega _{r+\alpha }^\tau (f,t)_p𝑑t\right\}^{\frac{1}{\tau }}\omega _\alpha (f^{(r\epsilon )},\delta )_p,$$
(84)
$$\omega _\alpha (f^{(r\epsilon )},\delta )_p\left\{\underset{0}{\overset{\delta }{}}t^{(r\epsilon )\theta 1}\omega _{r+\alpha }^\theta (f,t)_p𝑑t+\delta ^{\alpha \theta }\underset{\delta }{\overset{1}{}}t^{(r\epsilon +\alpha )\theta 1}\omega _{r+\alpha }^\theta (f,t)_p𝑑t\right\}^{\frac{1}{\theta }};$$
(85)
3. $`\gamma =r+\epsilon ,`$$`0<\epsilon <\alpha ,`$ (see ):
$$\left\{\underset{0}{\overset{\delta }{}}t^{(r+\epsilon )\tau 1}\omega _{r+\alpha }^\tau (f,t)_p𝑑t\right\}^{\frac{1}{\tau }}\delta ^{\alpha \epsilon }\left\{\underset{\delta }{\overset{1}{}}t^{(\alpha \epsilon )\theta 1}\omega _\alpha ^\theta (f^{(r+\epsilon )},t)_p𝑑t\right\}^{\frac{1}{\theta }},$$
(86)
$$\delta ^{\alpha \epsilon }\left\{\underset{\delta }{\overset{1}{}}t^{(\alpha \epsilon )\tau 1}\omega _\alpha ^\tau (f^{(r+\epsilon )},t)_p𝑑t\right\}^{\frac{1}{\tau }}\left\{\underset{0}{\overset{\delta }{}}t^{(r+\epsilon )\theta 1}\omega _{r+\alpha }^\theta (f,t)_p𝑑t\right\}^{\frac{1}{\theta }}.$$
(87)
Some more general estimates of the type (84)-(87) for moduli of smoothness of the transformed Fourier series can be obtained (, , ) using the sequences of the type (see, for example, , ) $`\{\mathrm{\Lambda }_n(s):=\mathrm{\Lambda }(s,\frac{1}{n})\}`$, where
$$\mathrm{\Lambda }(s,t)=\mathrm{\Lambda }(s,r,t)=\left(\underset{t}{\overset{1}{}}\xi (u)𝑑u+t^{rs}\underset{0}{\overset{t}{}}u^{rs}\xi (u)𝑑u\right)^{\frac{1}{s}}$$
(88)
and a non-negative function $`\xi (u)`$ on $`[0,1]`$ is such that $`u^{rs}\xi (u)`$ is summable.
2. The case of $`𝐩=\mathrm{𝟏},\mathrm{}`$.
Estimates of $`\omega _\alpha (\phi ,t)_p`$ in terms of $`\omega _{r+\alpha }(f,t)_p`$ for this case follow from Theorem 1 (see item II). We will write only the commonly used estimates of $`\omega _\alpha (f^{(r)},t)_p`$ and $`\omega _\alpha (\stackrel{~}{f}^{(r)},t)_p`$ in terms of $`\omega _{r+\alpha }(f,t)_p`$ (see also , -).
###### Corollary 2
If $`p=1,\mathrm{}`$, then inequalities (7), (8) hold true for any $`k,r>0.`$
If $`\left\{\lambda _n=n^\rho \right\}`$, $`\rho 0`$ and $`\beta =\rho +1`$, Theorem 1 implies the following
###### Corollary 3
Let $`p=1,\mathrm{}`$. Then
$$H_{\alpha +\rho }^p[\omega ]\stackrel{~}{W}_p^\rho \underset{n=1}{\overset{\mathrm{}}{}}n^{\rho 1}\omega \left(\frac{1}{n}\right)<\mathrm{}.$$
Note that in the case of $`p=1,\rho =0`$ and $`\alpha =1`$, Corollary 3 gives the answer for the question by F. Móricz (, 1995) on necessary conditions for the embedding $`H_{\alpha +\rho }^p[\omega ]\stackrel{~}{W}_p^\rho `$. We also mention the papers , , and , where the embedding theorems were proved in the necessity part.
###### Corollary 4
Let $`p=1,\mathrm{}`$ and $`r,\alpha ,\epsilon >0`$.
I. If for $`fL_p`$ the series $`\underset{\nu =1}{\overset{\mathrm{}}{}}\nu ^{r1}\omega _{r+\alpha +\epsilon }(f,\frac{1}{\nu })_p`$ converges, then there exists $`\stackrel{~}{f}^{(r)}L_p`$ and
$$\omega _\alpha (\stackrel{~}{f}^{(r)},\frac{1}{n})_pC(r,\alpha ,\epsilon )\left(n^\alpha \underset{\nu =1}{\overset{n}{}}\nu ^{r+\alpha 1}\omega _{r+\alpha +\epsilon }(f,\frac{1}{\nu })_p+\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r1}\omega _{r+\alpha +\epsilon }(f,\frac{1}{\nu })_p\right),n𝐍.$$
II. If for $`fL_p`$ there exists $`\stackrel{~}{f}^{(r)}L_p`$, then
$$\omega _{r+\alpha +\epsilon }(f,\frac{1}{n})_p\frac{C(r,\alpha ,\epsilon )}{n^r}\omega _\alpha (\stackrel{~}{f}^{(r)},\frac{1}{n})_p,n𝐍.$$
(89)
Using the direct and inverse approximation theorems, we can write inequality from the item I of the previous corollary in the following equivalent form (see also , , \[39, Vol. 2, Ch. 6 and 7\]-):
###### Corollary 5
Let $`p=1,\mathrm{}`$ and $`r0,\alpha >0`$. If for $`fL_p`$ the series $`\underset{\nu =1}{\overset{\mathrm{}}{}}\nu ^{r1}E_\nu \left(f\right)_p`$ converges, then there exists $`\stackrel{~}{f}^{(r)}L_p`$ and
$$\omega _\alpha (\stackrel{~}{f}^{(r)},\frac{1}{n})_pC(r,\alpha ,\epsilon )\left(n^\alpha \underset{\nu =1}{\overset{n}{}}\nu ^{r+\alpha 1}E_\nu \left(f\right)_p+\underset{\nu =n+1}{\overset{\mathrm{}}{}}\nu ^{r1}E_\nu \left(f\right)_p\right),n𝐍.$$
## 6 Remarks
1. The embedding theorems for the classes $`H_l^p[\omega ],W_p^r`$, and $`W_p^rH_k[\phi ]`$ in the necessity part were investigated, for example, in the papers by N.K. Bary and S.B. Stechkin (see , $`H_1^p[\phi ]H_1^p[\omega ]`$), V. È Geĭt (see , $`H_l^p[\phi ]H_k^p[\omega ]`$, $`H_l^p[\phi ]\stackrel{~}{H}_k^p[\omega ]`$, $`p=1,\mathrm{}`$), N. A. Il’yasov (see , $`H_l^p[\phi ]W_p^r`$, $`H_l^p[\omega ]W_p^rH_k[\phi ]`$, $`r,k𝐍`$). Note also that all estimates in Theorem 1 are $`\mathrm{"}`$correct$`\mathrm{"}`$ (in the terminology of Stechkin), that is, the sharpness from the point of view of order, is realized with the help of individual functions. We thank N. A. Il’yasov for this remark.
Theorem 1 specifies the previous results both in the necessity and sufficiency parts. Sufficient conditions for the embeddings $`H_l^p[\omega ]W_p^{\lambda ,\beta }H_k[\phi ]`$ and $`H_l^p[\omega ]W_p^{\lambda ,\beta }`$ were studied in the papers -. From these articles, particularly, for an important model example $`f^{(\lambda ,\beta )}f^{(r)}`$, $`r>0`$, we have the following estimates
$`E_n(f^{(r)})_p`$ $``$ $`C(r)\left(n^rE_n(f)_p+\left\{{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}k^{r\theta 1}E_k^\theta (f)_p\right\}^{\frac{1}{\theta }}\right),`$ (90)
$`\omega _k(f^{(r)},{\displaystyle \frac{1}{n}})_p`$ $``$ $`{\displaystyle \frac{C(k,r)}{n^k}}\left({\displaystyle \underset{\nu =0}{\overset{n}{}}}(\nu +1)^{(k+r)\theta 1}E_\nu ^\theta (f)_p\right)^{\frac{1}{\theta }}+C(k,r)\left({\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\nu ^{r\theta 1}E_\nu ^\theta (f)_p\right)^{\frac{1}{\theta }},`$ (91)
where $`\theta =1`$. At the same time, for the case of $`1<p<\mathrm{}`$, Theorem 3 and Theorem 1<sup>9</sup><sup>9</sup>9 Theorem 1 implies inequality (85), which is equivalent (using the direct and inverse theorems) to inequality (91). respectively, imply that these estimates hold for $`\theta =\mathrm{min}(2,p)`$ and this exponent is sharp the best possible.
2. A generalization of the class $`W_p^rE[\xi ]`$ is the class
$$W_p^{\lambda ,\beta }E[\omega ]=\{fW_p^{\lambda ,\beta }:E_n\left(f^{(\lambda ,\beta )}\right)_p=O\left[\omega _n\right]\}.$$
In this paper we do not consider in detail the embedding theorems between the classes $`W_p^{\lambda ,\beta }E[\omega ]`$, $`W_p^{\lambda ,\beta }H_\alpha [\phi ]`$, and $`E_p[\epsilon ]W_p^{\{1\},\mathrm{\hspace{0.17em}0}}E[\epsilon ]`$. We only notice that some results of such types easily follow from direct and inverse theorems (1)-(4), (1) and some are given in , \[39, Vol. 2, Ch. 6 and 7\], and . For the case when $`\lambda _n`$ satisfies the $`\mathrm{}_2`$-condition, a complete solution of the problem on embedding between the classes $`W_p^{\lambda ,\beta }`$, $`W_p^{\lambda ,\beta }E[\xi ]`$, and $`E_p[\epsilon ]`$ is described in the following result.
###### Theorem 3
Let $`1<p<\mathrm{}`$, $`\theta =\mathrm{min}(2,p)`$, $`\beta 𝐑`$, and $`\lambda =\left\{\lambda _n\right\}`$ be a non-decreasing sequence of positive numbers satisfying the $`\mathrm{}_2`$-condition, i.e., $`\lambda _{2n}C\lambda _n`$. Let also $`\epsilon =\{\epsilon _n\}`$ and $`\omega =\{\omega _n\}`$ be non-increasing null-sequences.
I. If $`1<p<\mathrm{}`$, then
$`E_p[\epsilon ]W_p^{\lambda ,\beta }`$ $``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}^\theta \lambda _n^\theta \right)\epsilon _n^\theta <\mathrm{},`$
$`E_p[\epsilon ]W_p^{\lambda ,\beta }E[\omega ]`$ $``$ $`\left\{{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}^\theta \lambda _\nu ^\theta \right)\epsilon _\nu ^\theta \right\}^{\frac{1}{\theta }}+\lambda _n\epsilon _n=O\left[\omega _n\right],`$
$`W_p^{\lambda ,\beta }E_p[\epsilon ]`$ $``$ $`{\displaystyle \frac{1}{\lambda _n}}=O\left[\epsilon _n\right],`$
$`W_p^{\lambda ,\beta }E[\omega ]E_p[\epsilon ]`$ $``$ $`{\displaystyle \frac{\omega _n}{\lambda _n}}=O\left[\epsilon _n\right].`$
II. Let $`p=1`$ or $`p=\mathrm{}`$.
(a) If $`\mathrm{}\lambda _nC\mathrm{}\lambda _{2n}`$ and $`\mathrm{}^2\lambda _n0`$ (or $`0`$), then
$`E_p[\epsilon ]W_p^{\lambda ,\beta }`$ $``$ $`|\mathrm{cos}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(\lambda _{n+1}\lambda _n\right)\epsilon _n`$
$`+`$ $`|\mathrm{sin}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\lambda _n{\displaystyle \frac{\epsilon _n}{n}}<\mathrm{},`$
$`E_p[\epsilon ]W_p^{\lambda ,\beta }E[\omega ]`$ $``$ $`|\mathrm{cos}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\left(\lambda _{\nu +1}\lambda _\nu \right)\epsilon _\nu +\lambda _n\epsilon _n`$
$`+`$ $`|\mathrm{sin}{\displaystyle \frac{\beta \pi }{2}}|{\displaystyle \underset{\nu =n+1}{\overset{\mathrm{}}{}}}\lambda _\nu {\displaystyle \frac{\epsilon _\nu }{\nu }}=O\left[\omega _n\right].`$
(b) If for $`\beta =2k,k𝐙`$ the condition $`\mathrm{}^2\left(1/\lambda _n\right)0`$ holds, and for $`\beta 2k,k𝐙`$ the conditions $`\mathrm{}^2\left(1/\lambda _n\right)0`$ and $`\underset{\nu =n+1}{\overset{\mathrm{}}{}}\frac{1}{\nu \lambda _\nu }\frac{C}{\lambda _n}`$ hold, then
$`W_p^{\lambda ,\beta }E_p[\epsilon ]`$ $``$ $`{\displaystyle \frac{1}{\lambda _n}}=O\left[\epsilon _n\right],`$
$`W_p^{\lambda ,\beta }E[\omega ]E_p[\epsilon ]`$ $``$ $`{\displaystyle \frac{\omega _n}{\lambda _n}}=O\left[\epsilon _n\right].`$
3. The Weyl class $`W_p^{\lambda ,\beta }`$ coincides with the class of functions from $`L(0,2\pi )`$ such that their Fourier series can be presented in the following form
$$\frac{a_0(f)}{2}+\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{1}{\pi \lambda _\nu }\underset{0}{\overset{2\pi }{}}\psi (xt)\mathrm{cos}\left(\nu t\frac{\pi \beta }{2}\right)𝑑t,\psi (x)L^0.$$
Further, consider the case when
$$\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{1}{\lambda _\nu }\mathrm{cos}\left(\nu t\frac{\pi \beta }{2}\right)$$
is the Fourier series of a summable function $`D_{\lambda ,\beta }(t)`$.
For example, it is certainly so if $`\{\lambda _\nu \mathrm{}\}`$ ($`n`$) and $`\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{1}{\nu \lambda _\nu }<\mathrm{}`$ (see \[59, Ch. 5\]). Then elements of $`W_p^{\lambda ,\beta }`$ can differ only by the mean value from functions $`f`$, which have the following representation by convolution,
$$f(x)=\frac{1}{\pi }\underset{0}{\overset{2\pi }{}}\psi (xt)D_{\lambda ,\beta }(t)𝑑t,\psi (x)L^0.$$
Here, $`\psi `$ coincides almost everywhere with $`f^{(\lambda ,\beta )}`$. See, for example, \[12, Ch. 11\]. Also, similar questions are discussed in detail in the book \[39, Vol. 1, Ch. 3\]. Note that a representation by convolution was first considered in .
4. The results of sections 2 and 5 can be extended to a multidimensional case. We only write the following estimates for the mixed modulus of smoothness $`\omega _{\alpha _1,\alpha _2}(f,\delta _1,\delta _2)_p`$ of orders $`\alpha _1`$ and $`\alpha _2`$ ($`\alpha _1,\alpha _2>0`$) of a function $`f`$ (in the $`L_p`$ metric) with respect to the variables $`x_1`$ and $`x_2`$, respectively<sup>10</sup><sup>10</sup>10The definition of the mixed modulus of smoothness and the mixed derivative in the Weyl sense can be found in, for example, ..
###### Theorem 4
(see also ) Let $`f(x_1,x_2)L_p`$, $`1<p<\mathrm{}`$, $`\theta =\mathrm{min}(2,p),\tau =\mathrm{max}(2,p)`$ and let $`\alpha _1,\alpha _2`$, $`r_1,r_2>0`$,
I. If
$$J_1(\theta ):=\left(\underset{0}{\overset{1}{}}\underset{0}{\overset{1}{}}t_1^{r_1\theta 1}t_2^{r_2\theta 1}\omega _{r_1+\alpha _1,r_2+\alpha _2}^\theta (f,t_1,t_2)_p𝑑t_1𝑑t_2\right)^{\frac{1}{\theta }}<\mathrm{},$$
then $`f`$ has the mixed derivative in the Weyl sense $`f^{(r_1,r_2)}L_p^0`$. Moreover,
$$f^{(r_1,r_2)}_pC(p,r_1,r_2)J_1(\theta )$$
and
$$\begin{array}{c}\omega _{\alpha _1,\alpha _2}(f^{(r_1,r_2)},\delta _1,\delta _2)_pC(p,\alpha _1,\alpha _2,r_1,r_2)\left(\underset{0}{\overset{\delta _1}{}}\underset{0}{\overset{\delta _2}{}}t_1^{r_1\theta 1}t_2^{r_2\theta 1}\omega _{r_1+\alpha _1,r_2+\alpha _2}^\theta (f,t_1,t_2)_p𝑑t_1𝑑t_2\right)^{\frac{1}{\theta }}\hfill \\ \hfill =:CJ_2(\theta ).\end{array}$$
II. If $`f`$ has the mixed derivative in the Weyl sense $`f^{(r_1,r_2)}L_p`$, then
$$J_1(\tau )C(p,r_1,r_2)f^{(r_1,r_2)}_p$$
and
$$J_2(\tau )C(p,\alpha _1,\alpha _2,r_1,r_2)\omega _{\alpha _1,\alpha _2}(f^{(r_1,r_2)},\delta _1,\delta _2)_p.$$
###### Theorem 5
Let $`f(x_1,x_2)L_p`$, $`p=1,\mathrm{}`$ and let $`\alpha _1,\alpha _2,r_1,r_2>0`$.
I. If
$$J_3:=\underset{0}{\overset{1}{}}\underset{0}{\overset{1}{}}t_1^{r_11}t_2^{r_21}\omega _{r_1+\alpha _1,r_2+\alpha _2}(f,t_1,t_2)_p𝑑t_1𝑑t_2<\mathrm{},$$
then $`f`$ has the mixed derivative in the Weyl sense $`f^{(r_1,r_2)}L_p`$. Moreover,
$$f^{(r_1,r_2)}_pC(r_1,r_2)J_3$$
and
$$\omega _{\alpha _1,\alpha _2}(f^{(r_1,r_2)},\delta _1,\delta _2)_pC(\alpha _1,\alpha _2,r_1,r_2)\underset{0}{\overset{\delta _1}{}}\underset{0}{\overset{\delta _2}{}}t_1^{r_11}t_2^{r_21}\omega _{r_1+\alpha _1,r_2+\alpha _2}(f,t_1,t_2)_p𝑑t_1𝑑t_2.$$
II. If $`f`$ has the mixed derivative in the Weyl sense $`f^{(r_1,r_2)}L_p`$, then
$$\omega _{r_1+\alpha _1,r_2+\alpha _2}(f,\delta _1,\delta _2)_pC(\alpha _1,\alpha _2,r_1,r_2)\delta _1^{r_1}\delta _2^{r_2}f^{(r_1,r_2)}_p$$
and
$$\omega _{r_1+\alpha _1,r_2+\alpha _2}(f,\delta _1,\delta _2)_pC(\alpha _1,\alpha _2,r_1,r_2)\delta _1^{r_1}\delta _2^{r_2}\omega _{\alpha _1,\alpha _2}(f^{(r_1,r_2)},\delta _1,\delta _2)_p.$$
For more details on the estimates of transformed series in a multidimensional case, see the articles -.
5. In view of inequalities (83) and (84) - (87), the problem of finding the estimates of $`\omega _\alpha (\phi ,t)_p`$ in terms of $`\omega _{\alpha +r}(f,t)_p`$ arises, e.g., in the case $`\phi \sigma (f,\lambda )`$, where $`\lambda _n=n^r\mathrm{ln}^An`$. If $`A<0`$ (which is an analogue of the case $`\lambda _n=n^{r\epsilon }`$), then estimates $`\omega _\alpha (\phi ,t)_p`$ follow from .
For example, if $`p=2`$ and $`\phi \sigma (f,n^r\mathrm{ln}^An),A<0`$, then
$$\omega _\beta ^2(\phi ,\delta )_2\underset{0}{\overset{\delta }{}}\frac{t^{2r1}}{\mathrm{ln}^{2|A|}\left(\frac{2}{t}\right)}\omega _{r+\beta }^2(f,t)_2𝑑t+\delta ^{2\beta }\underset{\delta }{\overset{1}{}}\frac{t^{2(r+\beta )1}}{\mathrm{ln}^{1+2|A|}\left(\frac{2}{t}\right)}\omega _{r+\beta }^2(f,t)_2𝑑t.$$
(92)
Note that the differences between (92) and (84)-(85) are related only to the replacement of $`n^\epsilon `$ by $`\mathrm{ln}^An`$.
The case $`A>0`$ (which is an analogue of the case $`\lambda _n=n^{r+\epsilon }`$) is interesting. For $`p=2`$ and $`\phi \sigma (f,n^r\mathrm{ln}^An),A>0`$ we have
$$\delta ^{2\beta }\mathrm{ln}^{2A}\left(\frac{2}{\delta }\right)\underset{\delta }{\overset{1}{}}\frac{t^{2\beta 1}}{\mathrm{ln}^{1+2A}\left(\frac{2}{t}\right)}\omega _\beta ^2(\phi ,t)_2𝑑t+\omega _\beta ^2(\phi ,\delta )_2\underset{0}{\overset{\delta }{}}t^{2r1}\mathrm{ln}^{2A}\left(\frac{2}{t}\right)\omega _{r+\beta }^2(f,t)_2𝑑t.$$
(93)
Comparing these relations with estimates (86)-(87), one can remark that the new term $`\omega _\beta ^2(\phi ,\delta )_2`$ appears in (93). Thus, this case has essential distinctions. See for detail the papers , .
6. Defining the class $`W_p^{\lambda ,\beta }H_\alpha [\phi ]`$ we assume that $`\phi \mathrm{\Phi }_\alpha `$. This restriction is natural for a majorant of the modulus of smoothness of order $`\alpha `$ (see ).
7. In Theorem 1 (item II) we used the inequality $`\underset{\nu =n+1}{\overset{\mathrm{}}{}}\frac{1}{\nu \lambda _\nu }\frac{C}{\lambda _n}`$. It is equivalent to the following condition: there exists $`\epsilon >0`$ such that the sequence $`\left\{n^\epsilon \lambda _n\right\}`$ is almost increasing, i.e., $`n^\epsilon \lambda _nCm^\epsilon \lambda _m`$, $`C1`$, $`nm.`$ This and other conditions can be found in and .
The paper was supported by RFFI (project N 06-01-00268), the programm Leading Scientific Schools (project NSH-2787.2008.1), Centre de Recerca Matematica and Scuola Normale Superiore.
Boris V. Simonov
Volgograd Technical University
Volgograd, 400131
Russia
e-mail: htf@vstu.ru
Sergey Yu. Tikhonov
ICREA and
Centre de Recerca Matemàtica
Apartat 50 08193 Bellaterra
Barcelona Spain
e-mail: stikhonov@crm.es
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# Quantum Random Walks without Coin Toss1footnote 11footnote 1Invited lecture at the Workshop on Quantum Information, Computation and Communication (QICC-2005), IIT Kharagpur, India, February 2005, quant-ph/0506221.
## I Introduction
Random walks are a fundamental ingredient of non-deterministic algorithms motwani , and are used to tackle a wide variety of problems—from graph structures to Monte Carlo samplings. Such algorithms have many evolutionary branches, which are explored probabilistically, to estimate the correct result. A classical computer can explore only one branch at a time, so the algorithm is executed several times, and the estimate of the final result is extracted from the ensemble of individual executions by methods of probability theory. Such algorithms are typically represented using graphs, with vertices denoting the states and the edges denoting the evolutionary routes. A particular evolution corresponds to a specific walk on the graph, and the final result is obtained by combining the results for many different walks. To ensure that different evolutionary branches are explored in different executions, one needs non-deterministic instructions, and they are provided in the form of random numbers. A coin toss is the simplest example of a random number generator, and it is included in the instruction set for a probabilistic Turing machine.
A quantum computer can explore multiple branches of a non-deterministic algorithm in a single attempt, by using a clever superposition of states. The probabilistic result can then be arrived at by interference of amplitudes corresponding to different branches. Thus as long as the means to construct a variety of superposed states exist, there is no a priori reason to include a coin toss as an instruction for a (probabilistic) quantum Turing machine.
In what follows, we construct a quantum random walk on a hypercubic lattice in arbitrary dimensions without using a coin toss instruction, analyse its properties, and use it to find a marked vertex on the lattice. More details are available in Refs.qwalk1 ; qwalk2 .
## II Diffusion
A random walk is a diffusion process, commonly described using the Laplacian operator in the continuum. To construct a discrete quantum random walk, we must discretise the diffusion process using evolution operators that are both unitary and ultra-local (an ultra-local operator vanishes outside a finite range).
On a periodic lattice, the spatial modes are characterised by discrete wave vectors $`\stackrel{}{k}`$. Quantum diffusion then depends on the energy of these modes according to $`U(\stackrel{}{k},t)=\mathrm{exp}(iE(\stackrel{}{k})t)`$. The lowest energy mode, $`\stackrel{}{k}=0`$, corresponding to a uniform distribution, is an eigenstate of the diffusion operator and does not propagate. The slowest propagating modes are the ones with smallest nonzero $`|\stackrel{}{k}|`$. The classical Laplacian operator gives $`E(\stackrel{}{k})|\stackrel{}{k}|^2`$ massterm , which translates to the characteristic Brownian motion signature, spread $`n_{\mathrm{rms}}\sqrt{t}`$. There is an alternative in quantum theory—instead of the non-relativistic Schödinger equation based on the Laplacian operator $`^2`$, one can use the relativistic Dirac equation based on the operator $`/`$. The Dirac operator gives $`E(\stackrel{}{k})|\stackrel{}{k}|`$, with the associated signature, spread $`n_{\mathrm{rms}}t`$. Clearly the Dirac operator, with its faster diffusion of the slowest modes compared to the Laplacian, is the operator of choice for constructing faster diffusion based quantum algorithms.
An automatic consequence of the Dirac operator is the appearance of an additional internal degree of freedom corresponding to spin, whereby the quantum state is described by a multi-component spinor. These spinor components were identified with the states of a coin in Refs.gridsrch1 ; gridsrch2 , with the coin evolution rule guiding the quantum diffusion process. While this is the correct procedure in the continuum theory, another option is available for a lattice theory, i.e. staggered fermions staggered . In this approach, the spinor degrees of freedom are spread out over an elementary hypercube, location dependent signs appear in the evolution operator, and translational invariance exists in steps of 2 instead of 1. We follow this approach to construct, a quantum diffusion process on a hypercubic lattice, without a coin toss instruction,
The free particle Dirac Hamiltonian in $`d`$-space dimensions is
$$H_{\mathrm{free}}=i\stackrel{}{\alpha }\stackrel{}{}+\beta m.$$
(1)
On a hypercubic lattice, the simplest discretisation of the derivative operator is
$$_kf(\stackrel{}{x})=\frac{1}{2}[f(\stackrel{}{x}+\widehat{k})f(\stackrel{}{x}\widehat{k})].$$
(2)
Then the anticommuting matrices $`\stackrel{}{\alpha },\beta `$ can be spin-diagonalised to the location dependent signs
$$\alpha _k=\underset{j=1}{\overset{k1}{}}(1)^{x_j},\beta =\underset{j=1}{\overset{d}{}}(1)^{x_j}.$$
(3)
Even when the Hamiltonian $`H`$ is ultra-local (i.e. has a finite range), the evolution operator $`U=\mathrm{exp}(iHt)`$ is not. To make the evolution operator ultra-local, we break up $`H`$ in to block-diagonal Hermitian parts, and then exponentiate each part separately. Partitioning of $`H`$ in to two parts (which we label “odd” and “even”) is sufficient for this purpose spinorsize . This partition is illustrated in Fig.1 for $`d=1`$ and $`d=2`$. Each part contains all the vertices but only half of the links attached to each vertex. Consequently, each link appears in only one of the two parts, and can be associated with a term in $`H`$ providing propagation along it, i.e.
$$H_{\mathrm{free}}=H_o+H_e.$$
(4)
The Hamiltonian is thus divided in to a set of non-overlapping blocks that can be exponentiated exactly. Each block is an elementary hypercube on the lattice, and the block matrices are of size $`2^d\times 2^d`$ in $`d`$ dimensions. The ultra-local quantum random walk on the lattice then evolves the amplitude distribution according to
$$\psi (\stackrel{}{x};t)=W^t\psi (\stackrel{}{x};0),$$
(5)
$$W=U_eU_o=e^{iH_e\tau }e^{iH_o\tau }.$$
(6)
Each block of the unitary matrices $`U_{o(e)}`$ mixes the amplitudes of vertices belonging to a single elementary hypercube, and the amplitude distribution spreads because the two alternating matrices do not commute. The random walk operator $`W`$ is translationally invariant in steps of 2, along each coordinate direction.
## III Quantum Random Walk on a Line
### III.1 Construction
To explicitly illustrate the above described procedure, let us consider the random walk on a line, with the allowed positions labeled by integers.
The simplest translation invariant ultra-local discretisation of the Laplacian operator is
$$H|n\left[|n1+2|n|n+1\right].$$
(7)
One may search for ultra-local translationally invariant unitary evolution operators using the ansatz
$$U|n=a|n1+b|n+c|n+1,$$
(8)
but then the orthogonality constraints between different rows of the unitary matrix make two of $`\{a,b,c\}`$ vanish, and one obtains a directed walk instead of a random walk. This problem can be bypassed, and an ultra-local unitary random walk can be constructed, by enlarging the Hilbert space with a quantum coin, e.g.
$$U=\underset{n}{}[|||n+1n|+|||n1n|].$$
(9)
This route nayak brings its own set of caveats, due to quantum entanglement between the coin and the position degrees of freedom.
We follow an alternate route familiar to lattice field theorists staggered . It has also been used to simulate quantum scattering with ultra-local operators richardson , and to construct quantum cellular automata meyer . The starting point is the decomposition of the Laplacian operator in to its even and odd parts, $`H=H_e+H_o`$,
$$H\left(\begin{array}{cccccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& 1& 2& 1& 0& 0& 0& \mathrm{}\\ \mathrm{}& 0& 1& 2& 1& 0& 0& \mathrm{}\\ \mathrm{}& 0& 0& 1& 2& 1& 0& \mathrm{}\\ \mathrm{}& 0& 0& 0& 1& 2& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),$$
(10)
$$H_e\left(\begin{array}{cccccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& 1& 1& 0& 0& 0& 0& \mathrm{}\\ \mathrm{}& 0& 0& 1& 1& 0& 0& \mathrm{}\\ \mathrm{}& 0& 0& 1& 1& 0& 0& \mathrm{}\\ \mathrm{}& 0& 0& 0& 0& 1& 1& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right),$$
(11)
$$H_o\left(\begin{array}{cccccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& 0& 1& 1& 0& 0& 0& \mathrm{}\\ \mathrm{}& 0& 1& 1& 0& 0& 0& \mathrm{}\\ \mathrm{}& 0& 0& 0& 1& 1& 0& \mathrm{}\\ \mathrm{}& 0& 0& 0& 1& 1& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(12)
While $`H`$ has the structure of a second derivative, its two parts, $`H_e`$ and $`H_o`$, have the structure of a first derivative. The above decomposition is indeed reminiscent of the “square-root” one takes to go from the Laplacian to the Dirac operator.
The two parts, $`H_e`$ and $`H_o`$, are individually Hermitian. They are block-diagonal with a constant $`2\times 2`$ matrix, and so they can be exponentiated while maintaining ultra-locality. The total evolution operator can therefore be easily truncated, without giving up either unitarity or ultra-locality,
$`U(\mathrm{\Delta }t)=e^{i(H_e+H_o)\mathrm{\Delta }t}`$ $`=`$ $`e^{iH_e\mathrm{\Delta }t}e^{iH_o\mathrm{\Delta }t}+O((\mathrm{\Delta }t)^2)`$
$`=`$ $`U_e(\mathrm{\Delta }t)U_o(\mathrm{\Delta }t)+O((\mathrm{\Delta }t)^2).`$
The quantum random walk can now be generated using $`U_eU_o`$ as the evolution operator for the amplitude distribution $`\psi (n,t)`$,
$$\psi (n,t)=[U_eU_o]^t\psi (n,0),$$
(14)
The fact that $`U_e`$ and $`U_o`$ do not commute with each other is enough for the quantum random walk to explore all possible states. The price paid for the above manipulation is that the evolution operator is translationally invariant along the line in steps of 2, instead of 1.
The $`2\times 2`$ matrix appearing in $`H_e`$ and $`H_o`$ is proportional to $`(1\sigma _1)`$, and so its exponential will be of the form $`(c1+is\sigma _1)`$, $`|c|^2+|s|^2=1`$. A random walk should have at least two non-zero entries in each row of the evolution operator. Even though our random walk treats even and odd sites differently by construction, we can obtain an unbiased random walk, by choosing the $`2\times 2`$ blocks of $`U_e`$ and $`U_o`$ as $`\frac{1}{\sqrt{2}}\left(\genfrac{}{}{0pt}{}{1i}{i\mathrm{\hspace{0.17em}1}}\right)`$. Furthermore, it is computationally more convenient to choose a basis where the unitary operators are all real. Performing a global phase transformation, $`|ni^n|n`$ phaseshift , the $`2\times 2`$ blocks of $`U_e`$ and $`U_o`$ become $`\frac{1}{\sqrt{2}}(1\pm i\sigma _2)`$. The discrete quantum random walk then evolves the amplitude distribution according to
$`U_o|n`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[|n(1)^n|n+(1)^n\right],`$ (15)
$`U_e|n`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[|n+(1)^n|n(1)^n\right],`$ (16)
$$U_eU_o|n=\frac{1}{2}\left[|n1+|n|n+1+|n+2(1)^n\right].$$
(17)
It is instructive to realise that, with the above choice, the unbiased quantum random walk represents the path integral for a relativistic particle with $`|p|=m`$. Its speed (in units of speed of light) is then $`|v|=1/\sqrt{2}`$. The directed walk, with the $`2\times 2`$ block matrix $`U\sigma _2`$, corresponds to $`|v|=1`$, and the stationary limit $`U=1`$ corresponds to $`v=0`$.
### III.2 Analysis
It is straightforward to analyse the properties of the walk in Eq.(17) using the Fourier transform:
$$\stackrel{~}{\psi }(k,t)=\underset{n}{}e^{ikn}\psi (n,t),$$
(18)
$$\psi (n,t)=_\pi ^\pi \frac{dk}{2\pi }e^{ikn}\stackrel{~}{\psi }(k,t).$$
(19)
The evolution of the amplitude distribution in Fourier space is easily obtained by splitting it in to its even and odd parts:
$$\psi \left(\begin{array}{c}\psi _e\\ \psi _o\end{array}\right),\psi (k,t)=[M(k)]^t\psi (k,0),$$
(20)
$$M(k)=\left(\begin{array}{cc}e^{ik}\mathrm{cos}k& i\mathrm{sin}k\\ i\mathrm{sin}k& e^{ik}\mathrm{cos}k\end{array}\right).$$
(21)
The unitary matrix $`M`$ has the eigenvalues, $`\lambda _\pm e^{\pm i\omega _k}`$ (this $`\pm `$ sign label continues in all the results that follow),
$$\lambda _\pm =\mathrm{cos}^2k\pm i\mathrm{sin}k\sqrt{1+\mathrm{cos}^2k},\omega _k=\mathrm{cos}^1(\mathrm{cos}^2k),$$
(22)
with the (unnormalised) eigenvectors,
$`e_\pm `$ $``$ $`\left(\begin{array}{c}\mathrm{cos}k\sqrt{1+\mathrm{cos}^2k}\\ 1\end{array}\right),`$ (23)
$``$ $`\left(\begin{array}{c}1\\ \mathrm{cos}k\sqrt{1+\mathrm{cos}^2k}\end{array}\right).`$
The evolution of amplitude distribution then follows
$$\stackrel{~}{\psi }(k,t)=e^{iw_kt}\stackrel{~}{\psi }_+(k,0)+e^{iw_kt}\stackrel{~}{\psi }_{}(k,0),$$
(24)
where $`\stackrel{~}{\psi }_\pm (k,0)`$ are the projections of the initial amplitude distribution along $`e_\pm `$. The amplitude distribution in the position space is given by the inverse Fourier transform of $`\stackrel{~}{\psi }(k,t)`$. While we are unable to evaluate it exactly, many properties of the quantum random walk can be extracted numerically as well as by suitable approximations.
A walk starting at the origin satisfies $`\psi _\mathrm{o}(n,0)=\delta _{n,0}`$. This walk is asymmetric because our definitions treat even and odd sites differently. We can construct a symmetric walk, using the initial condition $`\psi _\mathrm{s}(n,0)=(\delta _{n,0}+i\delta _{n,1})/\sqrt{2}`$. The resultant probability distribution is then symmetric under $`n(1n)`$. (Real and imaginary components of the amplitude distribution evolve independently because we have chosen the evolution operator to be real.) For both these initial conditions, by construction, the quantum random walk remains within the interval $`[2t+1,2t]`$ after $`t`$ time steps.
The escape probability of the quantum random walk can be calculated by introducing a fully absorbing wall, say between $`n=0`$ and $`n=1`$. Mathematically, this absorbing wall amounts to a projection operator for $`n0`$. The unabsorbed part of the walk is given by
$`\psi (n,t+1)`$ $`=`$ $`P_{n0}U_eU_o\psi (n,t),`$
$`=`$ $`U_eU_o\psi (n,t){\displaystyle \frac{1}{2}}\delta _{n,1}(\psi (0,t)+\psi (1,t)),`$
with the absorption probability,
$$P_{\mathrm{abs}}(t)=1\underset{n0}{}|\psi (n,t)|^2.$$
(26)
All these variations in initial and boundary conditions are easy to implement numerically, and examples are shown in Fig.2. We have used such simulations to study various properties of the quantum random walk.
For large $`t`$, a good approximation to the probability distributions can be obtained by the stationary phase method nayak ; qwalk1 . The smoothed probability distribution for the symmetric walk, obtained by replacing the highly oscillatory terms by their mean values, is
$$|\psi _\mathrm{s}|_{\mathrm{smooth}}^2=\frac{4t^2}{\pi \sqrt{4t^22n^2}(4t^2n^2)}.$$
(27)
(Here, the $`n(1n)`$ symmetry can be restored by replacing $`n`$ by $`(n\frac{1}{2})`$.) As shown in the top part of Fig.2, it represents the average behavior of the distribution very well. Its low order moments are easily calculated to be,
$`{\displaystyle _{n=\sqrt{2}t}^{\sqrt{2}t}}|\psi _\mathrm{s}|_{\mathrm{smooth}}^2𝑑n`$ $`=`$ $`1,`$ (28)
$`{\displaystyle _{n=\sqrt{2}t}^{\sqrt{2}t}}|n||\psi _\mathrm{s}|_{\mathrm{smooth}}^2𝑑n`$ $`=`$ $`t,`$ (29)
$`{\displaystyle _{n=\sqrt{2}t}^{\sqrt{2}t}}n^2|\psi _\mathrm{s}|_{\mathrm{smooth}}^2𝑑n`$ $`=`$ $`2(2\sqrt{2})t^2.`$ (30)
### III.3 Results
The following properties of the quantum random walk are easily deduced qwalk1 :
$``$ The probability distribution is double-peaked with maxima approximately at $`\pm \sqrt{2}t`$. The distribution falls off steeply beyond the peaks, while it is rather flat in the region between the peaks. With increasing $`t`$, the peaks become more pronounced, because the height of the peaks decreases more slowly than that for the flat region. The location of the peaks is in accordance with the propagation speed, $`|v|=1/\sqrt{2}`$, once we take in to account the fact that a single step of our walk is a product of two nearest neighbor operators, $`U_e`$ and $`U_o`$.
$``$ The size of the tail of the amplitude distribution is limited by $`(ϵt)^1t^{1/3}`$, which gives $`\mathrm{\Delta }n_>=\mathrm{\Delta }(ϵt)=O(t^{1/3})`$. On the inner side, the width of the peaks is governed by $`|\omega _k^{^{\prime \prime }}t|^{1/2}t^{1/3}`$. For $`|n|=(\sqrt{2}\delta )t`$, this gives $`\mathrm{\Delta }n_<=\mathrm{\Delta }(\delta t)=O(t^{1/3})`$. The peaks therefore make a negligible contribution to the probability distribution, $`O(t^{1/3})`$.
$``$ Rapid oscillations contribute to the probability distribution (and hence to its moments) only at subleading order. They can be safely ignored in an asymptotic analysis, retaining only the smooth part of the probability distribution.
$``$ The quantum random walk spreads linearly in time, with a speed smaller by a factor of $`\sqrt{2}`$ compared to a directed walk. This speed is a measure of its mixing behavior and hitting probability. The probability distribution is qualitatively similar to a uniform distribution over the interval $`[\sqrt{2}t,\sqrt{2}t]`$. In particular, the $`m^{\mathrm{th}}`$ moment of the probability distribution is proportional to $`t^m`$. This behaviour is in sharp contrast to that of the classical random walk. The classical random walk produces a binomial probability distribution, which in the symmetric case has a single peak centered at the origin and variance proportional to $`t`$. The linear spread in time of our quantum random walk is achieved even when $`\psi `$ has 50% probability to stay put at the same location at every step, as can be seen from Eqs.(15,16). This means that our walk is more directed and less of a zigzag.
$``$ Above properties agree with those obtained in Refs.nayak ; watrous for a quantum random walk with a coin-toss instruction (extra factors of $`2`$ appear in our results because of difference in our conventions), demonstrating that the coin offers no advantage in this particular set up. Essentially, we have absorbed the two states of the coin in to the even/odd site label at no extra cost. By making the coin states part of the position space, we have eliminated quantum entanglement between the coin and the position degrees of freedom completely—only superposition representing the amplitude distribution survives entangle . Such a reorganisation would be a tremendous advantage in any practical implementation of the quantum random walk, because quantum entanglement is highly fragile against environmental disturbances while mere superposition is much more stable. The cost for gaining this advantage is the loss of short distance homogeneity—translational invariance holds in steps of $`2`$ instead of $`1`$.
$``$ Comparison of the numerically evaluated probability distributions in Fig.2, without and with the absorbing wall, shows that the absorbing wall disturbs the evolution of the walk only marginally. The probability distribution in the region close to $`n=0`$ is depleted as anticipated, while it is a bit of a surprise that the peak height near $`n=\sqrt{2}t`$ increases slightly. As a result, the escape speed from the wall is little higher than the spreading speed without the wall. Overall, the part of the quantum random walk going away from the absorbing wall just takes off at a constant speed, hardly ever returning to the starting point. Again, this behavior is in a sharp contrast to that of the classical random walk, which always returns to the starting point, sooner or later. We also find that the first two time steps dominate absorption, $`P_{\mathrm{s},\mathrm{abs}}(t=1)=0.25`$ and $`P_{\mathrm{s},\mathrm{abs}}(t=2)=0.375`$, with very little absorption later on. Asymptotically, the net absorption probability approaches $`P_{\mathrm{s},\mathrm{abs}}(\mathrm{})0.4098`$ for the symmetric walk. This value is smaller than the corresponding result $`P_{\mathrm{abs}}(\mathrm{})=2/\pi `$ for the symmetric quantum random walk with a coin-toss instruction watrous .
## IV Quantum Random Walkon a Hypercubic Lattice
### IV.1 $`2`$-dim Lattice
Next let us consider the situation for $`d=2`$. The partitioned free Hamiltonian is given by
$`H_o|x,y`$ $`=`$ $`{\displaystyle \frac{i}{2}}[(1)^x|x+(1)^x,y`$ (31)
$`+(1)^{x+y}|x,y+(1)^y],`$
$`H_e|x,y`$ $`=`$ $`{\displaystyle \frac{i}{2}}[(1)^x|x(1)^x,y`$ (32)
$`+(1)^{x+y}|x,y(1)^y],`$
$`H|x,y`$ $`=`$ $`(H_o+H_e)|x,y`$ (33)
$`=`$ $`{\displaystyle \frac{i}{2}}[|x+1,y|x1,y`$
$`+(1)^x(|x,y+1|x,y1)].`$
More explicitly, the $`4\times 4`$ blocks of the Hamiltonian are:
$`H_o^B`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\begin{array}{cccc}0& 1& 1& 0\\ 1& 0& 0& 1\\ 1& 0& 0& 1\\ 0& 1& 1& 0\end{array}\right)\begin{array}{c}00\\ 10\\ 01\\ 11\end{array}`$ (34)
$`=`$ $`{\displaystyle \frac{1}{2}}(I\sigma _2+\sigma _2\sigma _3),`$ (35)
where the column on the right denotes the vertices of the elementary square on which $`H_o^B`$ operates. Similarly, $`H_e^B=H_o^B`$, when operating on the square with vertices {00,-10,0-1,-1-1}. Noting that $`H_o^2=H_e^2=\frac{1}{2}I`$, the block-diagonal matrices are easily exponentiated to
$$U_{o(e)}=cIis\sqrt{2}H_{o(e)},|c|^2+|s|^2=1.$$
(36)
The parameter $`c`$ (or $`s`$) is to be tuned to achieve the fastest diffusion across the lattice.
The quantum random walk with the Dirac operator spreads on a two-dimensional grid as illustrated in Fig.3. The continuum Dirac Hamiltonian has exact rotational symmetry, and that survives to an extent even after discretisation on a hypercubic lattice. After a point start, the random walk spreads essentially isotropically at distances much larger than the lattice spacing, while the hypercubic symmetry governs the random walk pattern at shorter distances. Of course, the hypercubic symmetry would be exact for a $`d`$-dim random walk constructed as a tensor product of $`d`$ one-dimensional random walks.
### IV.2 $`3`$-dim Lattice
Next let us look at the situation for $`d=3`$. The partitioned free Hamiltonian is given by
$`H_o|x,y,z`$ $`=`$ $`{\displaystyle \frac{i}{2}}[(1)^x|x+(1)^x,y,z`$
$`+(1)^{x+y}|x,y+(1)^y,z`$
$`+(1)^{x+y+z}|x,y,z+(1)^z],`$
$`H_e|x,y,z`$ $`=`$ $`{\displaystyle \frac{i}{2}}[(1)^x|x(1)^x,y,z`$
$`+(1)^{x+y}|x,y(1)^y,z`$
$`+(1)^{x+y+z}|x,y,z(1)^z],`$
$`H|x,y,z`$ $`=`$ $`(H_o+H_e)|x,y,z`$
$`=`$ $`{\displaystyle \frac{i}{2}}[|x+1,y,z|x1,y,z`$
$`+(1)^x(|x,y+1,z|x,y1,z`$
$`+(1)^{x+y}(|x,y,z+1|x,y,z1)].`$
More explicitly, the $`8\times 8`$ blocks of the Hamiltonian are:
$`H_o^B`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left(\begin{array}{cccccccc}0& 1& 1& 0& 1& 0& 0& 0\\ 1& 0& 0& 1& 0& 1& 0& 0\\ 1& 0& 0& 1& 0& 0& 1& 0\\ 0& 1& 1& 0& 0& 0& 0& 1\\ 1& 0& 0& 0& 0& 1& 1& 0\\ 0& 1& 0& 0& 1& 0& 0& 1\\ 0& 0& 1& 0& 1& 0& 0& 1\\ 0& 0& 0& 1& 0& 1& 1& 0\end{array}\right)\begin{array}{c}000\\ 100\\ 010\\ 110\\ 001\\ 101\\ 011\\ 111\end{array}`$ (40)
$`=`$ $`{\displaystyle \frac{1}{2}}(II\sigma _2+I\sigma _2\sigma _3+\sigma _2\sigma _3\sigma _3),`$
with the column on the right indicating the vertices of the elementary cube on which $`H_o^B`$ operates. Likewise, $`H_e^B=H_o^B`$, when operating on the elementary cube with vertices {000,-100,0-10,-1-10,00-1,-10-1,0-1-1,-1-1-1}. With $`H_o^2=H_e^2=\frac{3}{4}I`$, the block-diagonal matrices exponentiate to
$$U_{o(e)}=cIis\frac{2}{\sqrt{3}}H_{o(e)},|c|^2+|s|^2=1.$$
(41)
Again $`c`$ (or $`s`$) is a parameter to be tuned to achieve the fastest diffusion across the lattice.
### IV.3 $`d`$-dim Lattice
We can now observe a pattern in the explicit results for $`d=1,2,3`$ above. The $`2^d\times 2^d`$ blocks of the Hamiltonian can be written as sums of tensor products of Pauli matrices. As suggested by Eqs.(35,40),
$$H_o^B=\frac{1}{2}\underset{j=1}{\overset{d}{}}I^{(dj)}\sigma _2\sigma _3^{(j1)},$$
(42)
and $`H_e^B=H_o^B`$ when operating on the hypercube with coordinates flipped in sign. The block-diagonal matrices satisfy $`H_o^2=H_e^2=\frac{d}{4}I`$, and exponentiate to
$$U_{o(e)}=cIis\frac{2}{\sqrt{d}}H_{o(e)},|c|^2+|s|^2=1.$$
(43)
## V Search on a Hypercubic Latticeusing the Dirac Operator
### V.1 Strategy
A clear advantage of quantum random walks is their linear spread in time, compared to square-root spread in time for classical random walks. So they are expected to be useful in problems requiring fast hitting times. Several examples of this nature have been explored in graph theoretical and sampling problems (see Refs.kempe ; ambainis for reviews). Here we consider the particular case of using the quantum random walk to find a marked vertex on a hypercubic lattice (see also Refs.gridsrch1 ; gridsrch2 ).
Consider a $`d`$-dim hypercubic lattice with $`N=L^d`$ vertices, one of which is marked. The quantum algorithmic strategy for the search process is to construct a Hamiltonian evolution, where the kinetic part of the Hamiltonian diffuses the amplitude distribution all over the lattice while the potential part of the Hamiltonian attracts the amplitude distribution towards the marked vertex grover\_strategy . The optimisation criterion is to concentrate the amplitude distribution towards the marked vertex as quickly as possible. In his algorithm, Grover constructed a global operator that allows diffusion from any vertex to any other vertex in just one step. Under different circumstances, when diffusion is restricted to be ultra-local (i.e. one can only go from a vertex to its neighbours in one step), one must find an appropriate diffusion operator that provides fast propagation of spatial modes. Obviously, the Dirac operator is better suited to this task than the Laplacian operator.
To search for a marked vertex, say the origin, we need to attract the quantum random walk towards it. This can be accomplished by adding a potential to the free Hamiltonian,
$$V=V_0\delta _{\stackrel{}{x},0}.$$
(44)
Exponentiation of this potential produces a phase change for the amplitude at the marked vertex. It is optimal to choose the magnitude of the potential to make the phase maximally different from $`1`$, i.e. $`e^{iV_0\tau }=1`$, whereby the phase becomes a reflection operator (binary oracle),
$$R=I2|\stackrel{}{0}\stackrel{}{0}|.$$
(45)
The search algorithm alternates between the diffusion and the reflection operators, yielding the evolution
$$\psi (\stackrel{}{x};t_1,t_2)=[W^{t_1}R]^{t_2}\psi (\stackrel{}{x};0,0).$$
(46)
Here $`t_2`$ is the number of oracle calls, and $`t_1`$ is the number of random walk steps between the oracle calls. Both have to be optimised, in addition to $`c`$ and depending on the size and dimensionality of the lattice, to find the quickest solution to the search problem.
Fastest search amounts to finding the shortest unitary evolution path between the initial state, typically chosen as the uniform superposition state $`|s=_x|\stackrel{}{x}/\sqrt{N}`$, and the marked state $`|\stackrel{}{0}`$. This path is a circular arc (geodesic) from $`|s`$ to $`|\stackrel{}{0}`$. With the random walk diffusion operator $`W`$, evolution of the state $`|\psi `$ does not remain restricted to the two-dimensional subspace formed $`|s`$ and $`|\stackrel{}{0}`$. Thus to optimise our algorithm, we need to tune the parameters so as to
(a) maximise the projection of the state $`|\psi `$ on to the two-dimensional $`|s|\stackrel{}{0}`$ subspace, and
(b) maximise the angle of rotation by the operator $`W^{t_1}R`$, for the projected component of $`|\psi `$ in the $`|s`$-$`|\stackrel{}{0}`$ subspace.
We have explored this optimisation numerically.
### V.2 Numerical Results
We carried out computer simulations of the quantum random walk search problem with a single marked vertex, for $`d=2,3`$. The algorithm was optimised by tuning the parameters $`c`$ and $`t_1`$, so as to minimise the number of oracle calls $`t_2`$ required to find the marked vertex. The following is a summary of our observations:
$``$ An unbiased search starts with a uniform probability distribution over the whole lattice. Thereafter, the probability at the marked vertex goes through periodic cycles of rise and fall as a function of time step. It is crucial to stop the algorithm at the right instance to find the marked vertex with a significant probability.
$``$ Our best results are obtained with $`c=1/\sqrt{2}`$ and $`t_1=3`$. In this case, the probability at the marked vertex reaches its largest value, and $`t_2`$ achieves its smallest value. With these parameters, the probability at the marked vertex shows a periodic sinusoidal behaviour, which persists for more than 30 cycles without any visible deviation. Also, apart from the uniform background, the probability distribution shows a sharp single-point delta function at the marked vertex. These features indicate that the walk evolves largely in the two-dimensional subspace formed by the uniform state and the marked state.
$``$ For $`c<1/\sqrt{2}`$, the walk diffuses more slowly, and $`t_2`$ increases. For $`c>1/\sqrt{2}`$, the probability at the marked vertex loses its periodic sinusoidal behaviour, suggesting that the walk no longer remains confined to the two-dimensional subspace. For the optimal choice $`c=1/\sqrt{2}`$, the probability of the walk remaining at the same vertex equals that for moving to a neighbouring vertex, which corresponds to the most efficient mixing between odd and even sublattices.
$``$ For $`t_1<3`$, the probability distribution spreads out instead of being a delta function at the marked vertex, as illustrated in Fig.3. This decreases the peak probability at the marked vertex. Moreover, $`t_2`$ increases. For $`t_1>3`$, the probability at the marked vertex loses its sinusoidal behaviour, again with a decrease in the peak probability. In both cases, the changes indicate that the walk is drifting out of the two-dimensional subspace. An appropriate choice of $`t_1`$ is thus crucial to keep the walk close to the two-dimensional subspace.
$``$ For the $`2`$-dim walk, the largest probability at the marked vertex is predicted to be $`O(1/\mathrm{log}N)`$, which occurs after $`O(\sqrt{N\mathrm{log}N})`$ time steps gridsrch1 ; gridsrch2 . To make the marked vertex probability $`O(1)`$, an amplitude amplification procedure is required brassard , and the overall search algorithm scales as $`O(\sqrt{N}\mathrm{log}N)`$. Our numerical results, shown in Fig.5, are consistent with these expectations. Simple fits provide the parametrisations:
$`O(1/\mathrm{log}N)`$ $``$ $`2.12/\mathrm{log}_2N,`$ (47)
$`O(\sqrt{N\mathrm{log}N})`$ $``$ $`0.137\sqrt{N\mathrm{log}_2N}.`$
$``$ For the walk in more than two dimensions, the largest probability at the marked vertex is predicted to be $`O(1)`$, which occurs after $`O(\sqrt{N})`$ time steps gridsrch1 ; gridsrch2 . Our numerical results for the $`3`$-dim walk, also displayed in Fig.5, agree with these scaling rules. Simple fits provide the parametrisations:
$$O(1)0.0969,O(\sqrt{N})0.313\sqrt{N}.$$
(48)
These results demonstrate that our quantum random walk algorithm achieves the optimal scaling behaviour for the problem of finding a marked vertex on a hypercubic lattice. Thus our quantum random walk, based on the Dirac operator and not containing a coin toss instruction, is no less effective in its diffusion properties than the earlier quantum random walks that use a coin toss instruction.
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# Simulation of the liquid pool for VT3-1 titanium alloy during vacuum arc remelting process.
## 1 Introduction.
Today a lot of attention is given to automation of VAR furnaces and the focus is on development of the arc clearance control flowcharts, feedback controllers, program management systems for furnace electric modes . Together with the aforementioned issues, the development of the VAR automatic control systems requires a reliable theoretical description of the ingot formation process that underlies the automatic control system algorithms. The significance of through and reliable theoretical models of the ingot process formation is determined by the fact that the immediate control of the ingot parameters during solidification of titanium alloy (in particular, the depth of the metal pool, width of the mushy zone, etc.) using the measuring instruments, appears to be a problem so far. Development of the theoretical ideas about formation of titanium alloy ingots is rather important today since there are cases when different ingot-related defects are observed on billets, bars and eventually on critical parts. Such defects include $`\beta `$-flecks , zonal segregation , tree-like structure and dark/light spots that may be observed in the transverse section under the macro examination of mill products in different titanium alloys. Very often such defects are related to the final remelting stage therefore the theoretical description of the processes that take place during solidification of the VAR melted ingots, as well as understanding of the impact of various melting parameters on ingot metallurgical quality, is undoubtedly the paramount objective for creation of the fully automated VAR control system.
Many researchers studied the heat processes in ingots during VAR , electroslag remelt and continuous casting . Today, in view of development of advanced computers it is getting important to write science-based VAR simulation programs for their use in industrial environment in order to develop new remelting modes for complex titanium alloys.
Taking into account aforementioned facts we reviewed in this article the mathematical model of the vacuum arc remelting and compared the theoretical calculation of the metal pool depth with the radiographical test results obtained for Russian titanium alloy VT3-1 (Ti-6.5Al-2.5Mo-1.5Cr-0.5Fe-0.3Si) . Taking into consideration the non-trivial complexity and essential non-linearity of the task, it should be noted that firstly there is no exact analytical solution of the task and secondly today there are a lot of different methods of approximation of the boundary conditions and building of the difference schemes for the heat conductivity equation. It is obvious that in such situation the experimental check of the computational solutions obtained through different methods of approximation is getting important. Therefore, in this article we reviewed both aspects of the problem: making the mathematical model and comparison of this model with the experimental results.
## 2 Vacuum arc remelting process.
We will describe the heat flows during the vacuum arc remelting through the heat conductivity equation that includes release of latent heat within the interval of the alloy solidification. In order not to involve the complex interface conditions on the phase boundary, we reviewed the so-called method of apparent heat capacity that allows to include crystallization rate (i.e. heat source) into the system heat capacity as a complementary additive term . Such problem statement allows to make calculation in one domain without ”distinguishing” the liquid, solid and mushy zone.
Also, from this point on we considered that density, heat capacity, heat conductivity and heat-transfer coefficient depend on the temperature in an arbitrary way .
Heat transfer equation for VAR ingot can be written in a form
$$\rho (T)C(T)\frac{T}{\tau }=\frac{1}{r}\frac{}{r}\left(r\lambda (T)\frac{T}{r}\right)+\frac{}{z}\left(\lambda (T)\frac{T}{z}\right),$$
(1)
where $`T`$ is temperature, $`\tau `$ is time, $`C(T)`$ is heat capacity, $`\lambda (T)`$ is heat conductivity, $`\rho (T)`$ is density, $`r`$ and $`z`$ are radial and axial coordinates.
The boundary may be divided into several areas shown in Figure 1. Such division is based on physical processes that take place on the ingot surface under VAR. We will consider those processes and the relevant boundary conditions.
#### Boundary $`AB`$.
On the ingot axis we have the simple symmetry condition:
$$\frac{T}{r}=0.$$
(2)
#### Boundary $`BC`$.
The border $`BC`$ corresponds to a bath mirror surface which are being under an end face of the consumable electrode. On this border we shall set temperature. In the elementary kind the temperature of a surface of a bath is defined by an overheat above alloy liquidus temperature :
$$T=T_L+\mathrm{\Delta }T,$$
(3)
where $`T_L`$ is liquidus temperature, $`\mathrm{\Delta }T`$ is overheat.
#### Boundary $`CD`$.
The boundary $`CD`$ corresponds to a ring gap. Here we also set a boundary condition of the I-st kind
$$T=T_L+\frac{D_{cr}2r}{D_{cr}D_{el}}\mathrm{\Delta }T,$$
(4)
where $`D_{cr}`$ is crucible diameter, $`D_{el}`$ is electrod diameter.
#### Boundary $`DE`$.
The given boundary corresponds to a zone of a contact belt for which experimental values of a specific thermal flux $`q`$ are known, therefore we write down a boundary condition of II-nd kind
$$\lambda (T)\frac{T}{r}=q.$$
(5)
#### Boundary $`EF`$.
The site $`EF`$ corresponds to a zone of forming ingot on which has already partially or completely occured separation of solid ingot from a crucible surface due to passage of process volumetric shrinkage.
Thus three possible mechanisms of heat removal can be observed:
* Heat removal by radiation through the formed gap.
* Heat removal by contact way in places of contact of a surface of an ingot with an internal wall of copper crucible.
* The heat transfer by convection through a gas phase in a gap.
Three specified mechanisms, it is possible to write down the boundary condition considering all in the form of
$`\lambda (T){\displaystyle \frac{T}{r}}=\eta \left(ϵ\sigma _0\left(T^4T_{cr}^4\right)+{\displaystyle \frac{k}{d}}\left(TT_{cr}\right)\right)+(1\eta )\alpha \left(TT_{cr}\right),`$ (6)
where $`\eta `$ – factor describing a share of contact of a surface of an ingot with a crucible, $`ϵ`$ – emissivity, $`T_{cr}`$ is temperature of inner surface of the crucible, $`\alpha `$ is heat transfer coefficient in the contact zone, $`k`$ is heat conductivity of gas in gap, $`d`$ is gap. The parameters $`\eta `$, $`ϵ`$, $`T_{cr}`$, $`\alpha `$, $`k`$ $`d`$ are depending on temperature. For simplification of a record we shall indicate a following designation
$$\beta _{cr}(T)=\eta _{cr}\left(ϵ\sigma _0\frac{T^4T_{cr}^4}{TT_{cr}}+\frac{k}{d}\right)+(1\eta _{cr})\alpha _{cr},$$
(7)
and (6) is rewritted in a more compact form
$$\lambda (T)\frac{T}{r}=\beta _{cr}(T)\left(TT_{cr}\right).$$
(8)
Thus, complex process of the heat transfer on border $`EF`$ is described with the nonlinear law (8).
#### Boundary $`FA`$.
On a site $`FA`$ boundary condition is similar to (8), therefore we shall write down:
$$\lambda (T)\frac{T}{r}=\beta _{bot}(T)\left(TT_{bot}\right),$$
(9)
where
$$\beta _{bot}(T)=\eta _{bot}\left(ϵ\sigma _0\frac{T^4T_{bot}^4}{TT_{bot}}+\frac{k}{d}\right)+(1\eta _{bot})\alpha _{bot}.$$
(10)
Besides it is necessary for us to set the initial condition for the equation (1). During vacuum arc remelting melt an ingot occurs, i.e. the border $`BCD`$ moves upwards with the speed defined by remelting conditions. Melting rate of an ingot is defined by the mass speed of fusion depending on force of a current, an electrode and a filled crystallizer diameters, and also on some other parameters. We shall accept in the beginning of remelting is already available ingot of some small height with homogeneous distribution of temperature on all volume is considered to be as the initial condition. Initial temperature it is accepted equal $`T_L+\mathrm{\Delta }T`$. For the subsequent moments of time the initial condition is defined by distribution of temperature during the previous moment of time with additional small layer of an ingot which temperature is $`T_L+\mathrm{\Delta }T`$.
In the following sections we shall consider discretization of the given problem on a rectangular grid, including linearization of a system, and decision method of the received the linear equations system.
## 3 Discretization of heat conduction equation.
We will consider the most common case when $`\lambda (T)`$ and $`\rho (T)C(T)`$ are piecewise continuous functions. In this case it is convenient to use the integral identity method to the finite-difference scheme construct . We introduce the heat fluxes along $`r`$ and $`z`$ directions:
$$Q(r,z)=\lambda (r,z)\frac{T}{r},$$
(11)
$$P(r,z)=\lambda (r,z)\frac{T}{z}.$$
(12)
Taking into account these notations, Eq. (1) is rewrited to
$$\rho (T)C(T)\frac{T}{\tau }=\frac{1}{r}\frac{}{r}\left(rQ(r,z)\right)+\frac{}{z}\left(P(r,z)\right).$$
(13)
Integrating the Eq. (13) over the cylindric volume for $`r[r_{i1/2},r_{i+1/2}]`$ and $`z[z_{j1/2},z_{j+1/2}]`$ (see Figure 2) we can write heat balance equation can be put as follows
$`{\displaystyle \underset{r_{i+1/2}}{\overset{r_{i1/2}}{}}}r𝑑r{\displaystyle \underset{z_{j+1/2}}{\overset{z_{j1/2}}{}}}𝑑z\rho (r,z)C(r,z){\displaystyle \frac{T(r,z)}{\tau }}`$
$`={\displaystyle \underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}}\left[r_{i+1/2}Q(r_{i+1/2},z)r_{i1/2}Q(r_{i1/2},z)\right]𝑑z`$
$`+{\displaystyle \underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}}\left[P(r,z_{j+1/2})P(r,z_{j1/2})\right]r𝑑r.`$ (14)
This equation is exact. Now we need to evaluate the integrals, which contains Eq. (14). Because we have (11) than
$$\underset{r_i}{\overset{r_{i+1}}{}}\frac{Q(r,z)}{\lambda (r,z)}𝑑r=T(r_{i+1})T(r_i).$$
(15)
On the other hand
$$\underset{r_i}{\overset{r_{i+1}}{}}\frac{Q(r,z)}{\lambda (r,z)}𝑑rQ(r_{i+1/2},z)\underset{r_i}{\overset{r_{i+1}}{}}\frac{dr}{\lambda (r,z)}.$$
(16)
Thus we can write following expressions within the fluxes for internal region of ingot
$$Q(r_{i+1/2},z)\frac{T(r_{i+1},z)T(r_i,z)}{\underset{r_i}{\overset{r_{i+1}}{}}\frac{dr}{\lambda (r,z)}},$$
(17)
$$Q(r_{i1/2},z)\frac{T(r_i,z)T(r_{i1},z)}{\underset{r_{i1}}{\overset{r_i}{}}\frac{dr}{\lambda (r,z)}},$$
(18)
$$P(r,z_{j+1/2})\frac{T(r,z_{j+1})T(r,z_j)}{\underset{z_j}{\overset{z_{j+1}}{}}\frac{dr}{\lambda (r,z)}},$$
(19)
$$P(r,z_{j1/2})\frac{T(r,z_j)T(r,z_{j1})}{\underset{z_{j1}}{\overset{z_j}{}}\frac{dr}{\lambda (r,z)}}.$$
(20)
Let these expressions be substituted in Eq. (14) then
$`{\displaystyle \underset{r_{i+1/2}}{\overset{r_{i1/2}}{}}}r𝑑r{\displaystyle \underset{z_{j+1/2}}{\overset{z_{j1/2}}{}}}𝑑z\rho (r,z)C(r,z){\displaystyle \frac{T(r,z)}{\tau }}`$
$`=r_{i+1/2}\left[T_{i+1,j}T_{i,j}\right]{\displaystyle \underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}}{\displaystyle \frac{dz}{\underset{r_i}{\overset{r_{i+1}}{}}\frac{dr}{\lambda (r,z)}}}r_{i1/2}\left[T_{i,j}T_{i1,j}\right]{\displaystyle \underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}}{\displaystyle \frac{dz}{\underset{r_{i1}}{\overset{r_i}{}}\frac{dr}{\lambda (r,z)}}}`$
$`+\left[T_{i,j+1}T_{i,j}\right]{\displaystyle \underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}}{\displaystyle \frac{rdr}{\underset{z_j}{\overset{z_{j+1}}{}}\frac{dz}{\lambda (r,z)}}}+\left[T_{i,j}T_{i,j1}\right]{\displaystyle \underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}}{\displaystyle \frac{rdr}{\underset{z_{j1}}{\overset{z_j}{}}\frac{dz}{\lambda (r,z)}}}.`$ (21)
Making use the simplest approximation for the integral in the left side of equation we obtain
$`\left(z_{j1/2}z_{j+1/2}\right){\displaystyle \frac{r_{i+1/2}^2r_{i1/2}^2}{2}}\rho _{i,j}C_{i,j}{\displaystyle \frac{\widehat{T}_{i,j}T_{i,j}}{\tau }}`$
$`=r_{i+1/2}\left[T_{i+1,j}T_{i,j}\right]{\displaystyle \underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}}{\displaystyle \frac{dz}{\underset{r_i}{\overset{r_{i+1}}{}}\frac{dr}{\lambda (r,z)}}}r_{i1/2}\left[T_{i,j}T_{i1,j}\right]{\displaystyle \underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}}{\displaystyle \frac{dz}{\underset{r_{i1}}{\overset{r_i}{}}\frac{dr}{\lambda (r,z)}}}`$
$`+\left[T_{i,j+1}T_{i,j}\right]{\displaystyle \underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}}{\displaystyle \frac{rdr}{\underset{z_j}{\overset{z_{j+1}}{}}\frac{dz}{\lambda (r,z)}}}+\left[T_{i,j}T_{i,j1}\right]{\displaystyle \underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}}{\displaystyle \frac{rdr}{\underset{z_{j1}}{\overset{z_j}{}}\frac{dz}{\lambda (r,z)}}}.`$ (22)
Moreover the integrals containing heat conductivity is simplified up to:
$$\underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}\frac{dz}{\underset{r_i}{\overset{r_{i+1}}{}}\frac{dr}{\lambda (r,z)}}\frac{z_{j+1/2}z_{j1/2}}{r_{i+1}r_i}\lambda _{i+1/2,j},$$
$$\underset{z_{j1/2}}{\overset{z_{j+1/2}}{}}\frac{dz}{\underset{r_{i1}}{\overset{r_i}{}}\frac{dr}{\lambda (r,z)}}\frac{z_{j+1/2}z_{j1/2}}{r_ir_{i1}}\lambda _{i1/2,j},$$
$$\underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}\frac{rdr}{\underset{z_j}{\overset{z_{j+1}}{}}\frac{dz}{\lambda (r,z)}}\frac{r_{i+1/2}^2r_{i1/2}^2}{2(z_{j+1}z_j)}\lambda _{i,j+1/2},$$
$$\underset{r_{i1/2}}{\overset{r_{i+1/2}}{}}\frac{rdr}{\underset{z_{j1}}{\overset{z_j}{}}\frac{dz}{\lambda (r,z)}}\frac{r_{i+1/2}^2r_{i1/2}^2}{2(z_jz_{j1})}\lambda _{i,j1/2}.$$
After that the Eq. (22) can be rewritten as
$`\rho _{i,j}C_{i,j}{\displaystyle \frac{\widehat{T}_{i,j}T_{i,j}}{\tau }}=`$
$`\left(T_{i+1,j}T_{i,j}\right)\lambda _{i+1/2,j}{\displaystyle \frac{r_{i+1/2}}{r_{i+1}r_i}}{\displaystyle \frac{2}{r_{i+1/2}^2r_{i1/2}^2}}`$
$`\left(T_{i,j}T_{i1,j}\right)\lambda _{i1/2,j}{\displaystyle \frac{r_{i1/2}}{r_ir_{i1}}}{\displaystyle \frac{2}{r_{i+1/2}^2r_{i1/2}^2}}`$
$`+\left(T_{i,j+1}T_{i,j}\right)\lambda _{i,j+1/2}{\displaystyle \frac{1}{\left(z_{j+1}z_j\right)\left(z_{j+1/2}z_{j1/2}\right)}}`$
$`\left(T_{i,j}T_{i,j1}\right)\lambda _{i,j1/2}{\displaystyle \frac{1}{\left(z_jz_{j1}\right)\left(z_{j+1/2}z_{j1/2}\right)}}.`$ (23)
We define (for the simplest way)
$$r_{i+1/2}=\frac{r_i+r_{i+1}}{2},r_{i1/2}=\frac{r_{i1}+r_i}{2},$$
$$z_{j+1/2}=\frac{z_j+z_{j+1}}{2},z_{j1/2}=\frac{z_{j1}+z_j}{2}.$$
Thus for the rectangular grid Eq. (23) can be written in the common form as
$`\rho _{i,j}C_{i,j}{\displaystyle \frac{\widehat{T}_{i,j}T_{i,j}}{\tau }}`$
$`={\displaystyle \frac{r_{i+1}+r_i}{r_{i+1}r_i}}{\displaystyle \frac{4\lambda _{i+1/2,j}}{\left(r_{i+1}+r_i\right)^2\left(r_ir_{i1}\right)^2}}\left[T_{i+1,j}T_{i,j}\right]`$
$`{\displaystyle \frac{r_i+r_{i1}}{r_ir_{i1}}}{\displaystyle \frac{4\lambda _{i1/2,j}}{\left(r_{i+1}+r_i\right)^2\left(r_ir_{i1}\right)^2}}\left[T_{i,j}T_{i1,j}\right]`$
$`+{\displaystyle \frac{2\lambda _{i,j+1/2}}{\left(z_{j+1}z_j\right)\left(z_{j+1}z_{j1}\right)}}\left[T_{i,j+1}T_{i,j}\right]`$
$`{\displaystyle \frac{2\lambda _{i,j1/2}}{\left(z_jz_{j1}\right)\left(z_{j+1}z_{j1}\right)}}\left[T_{i,j}T_{i,j1}\right].`$ (24)
For constant space steps ($`h_r`$ and $`h_z`$) we have
$`\rho _{i,j}C_{i,j}{\displaystyle \frac{\widehat{T}_{i,j}T_{i,j}}{\tau }}=`$
$`{\displaystyle \frac{\lambda _{i+1/2,j}}{h_r^2}}\left(1+{\displaystyle \frac{1}{2i}}\right)\left[T_{i+1,j}T_{i,j}\right]{\displaystyle \frac{\lambda _{i1/2,m}}{h_r^2}}\left(1{\displaystyle \frac{1}{2i}}\right)\left[T_{i,j}T_{i1,j}\right]+`$
$`+{\displaystyle \frac{\lambda _{i,j+1/2}}{h_z^2}}\left[T_{i,j+1}T_{i,j}\right]{\displaystyle \frac{\lambda _{i,j1/2}}{h_z^2}}\left[T_{i,j}T_{i,j1}\right].`$ (25)
To evaluate heat conductivity at the semi - integer domian points we use approximation
$`\lambda _{i+1/2,j}={\displaystyle \frac{2\lambda _{i,j}\lambda _{i+1,j}}{\lambda _{i,j}+\lambda _{i+1,j}}},\lambda _{i1/2,j}={\displaystyle \frac{2\lambda _{i1,j}\lambda _{i,j}}{\lambda _{i1,j}+\lambda _{i,j}}},`$
$`\lambda _{i,j+1/2}={\displaystyle \frac{2\lambda _{i,j}\lambda _{i,j+1}}{\lambda _{i,j}+\lambda _{i,j+1}}},\lambda _{i,j1/2}={\displaystyle \frac{2\lambda _{i,j1}\lambda _{i,j}}{\lambda _{i,j}+\lambda _{i,j1}}}.`$ (26)
Equation (25) joins the temperature $`T_{i,j}`$ with $`T_{i1,j}`$, $`T_{i+1,j}`$, $`T_{i,j1}`$ and $`T_{i,j+1}`$, than Eq. (25) has 5-point space scheme and explicit scheme in time. We will use 2-cycle scheme . Let we rewrite Eq. (25) as follows
$$\rho _{i,j}C_{i,j}\frac{\widehat{T}_{i,j}T_{i,j}}{\tau }=\mathrm{\Lambda }_rT_{i,j}+\mathrm{\Lambda }_zT_{i,j},$$
(27)
where operators $`\mathrm{\Lambda }_{r,z}`$ are defined by Eq. (25). The main idea of the 2-cycle factorization scheme is making of the following steps ($`k`$ is time step index)
$$\rho _{i,j}C_{i,j}\frac{T_{i,j}^{k1/2}T_{i,j}^{k1}}{\tau }=\mathrm{\Lambda }_r\frac{T_{i,j}^{k1/2}+T_{i,j}^{k1}}{2}$$
(28)
$$\rho _{i,j}C_{i,j}\frac{T_{i,j}^kT_{i,j}^{k1/2}}{\tau }=\mathrm{\Lambda }_z\frac{T_{i,j}^k+T_{i,j}^{k1/2}}{2}$$
(29)
$$\rho _{i,j}C_{i,j}\frac{T_{i,j}^{k+1/2}T_{i,j}^k}{\tau }=\mathrm{\Lambda }_z\frac{T_{i,j}^{k+1/2}+T_{i,j}^k}{2}$$
(30)
$$\rho _{i,j}C_{i,j}\frac{T_{i,j}^{k+1}T_{i,j}^{k+1/2}}{\tau }=\mathrm{\Lambda }_r\frac{T_{i,j}^{k+1}+T_{i,j}^{k+1/2}}{2}$$
(31)
The cycle of evaluation is namely to solve equations (28)-(31). Let we rewite these equations in a more useful form
$`\left[B_{i,j}\right]T_{i1,j}^{k1/2}\left[A_{i,j}+B_{i,j}+2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^{k1/2}+\left[A_{i,j}\right]T_{i+1,j}^{k1/2}`$
$`=\left(\left[B_{i,j}\right]T_{i1,j}^{k1}\left[A_{i,j}+B_{i,j}2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^{k1}+\left[A_{i,j}\right]T_{i+1,j}^{k1}\right)`$ (32)
$`\left[\mathrm{\Delta }_{i,j}\right]T_{i,j1}^k\left[\mathrm{\Delta }_{i,j}+\mathrm{\Gamma }_{i,j}+2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^k+\left[\mathrm{\Gamma }_{i,j}\right]T_{i,j+1}^k`$
$`=\left(\left[\mathrm{\Delta }_{i,j}\right]T_{i,j1}^{k1/2}\left[\mathrm{\Delta }_{i,j}+\mathrm{\Gamma }_{i,j}2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^{k1/2}+\left[\mathrm{\Gamma }_{i,j}\right]T_{i,j+1}^{k1/2}\right)`$ (33)
$`\left[B_{i,j}\right]T_{i1,j}^{k+1/2}\left[A_{i,j}+B_{i,j}+2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^{k+1/2}+\left[A_{i,j}\right]T_{i+1,j}^{k+1/2}`$
$`=\left(\left[B_{i,j}\right]T_{i1,j}^k\left[A_{i,j}+B_{i,j}2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^k+\left[A_{i,j}\right]T_{i+1,j}^k\right)`$ (34)
$`\left[\mathrm{\Delta }_{i,j}\right]T_{i,j1}^{k+1}\left[\mathrm{\Delta }_{i,j}+\mathrm{\Gamma }_{i,j}+2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^{k+1}+\left[\mathrm{\Gamma }_{i,j}\right]T_{i,j+1}^{k+1}`$
$`=\left(\left[\mathrm{\Delta }_{i,j}\right]T_{i,j1}^{k+1/2}\left[\mathrm{\Delta }_{i,j}+\mathrm{\Gamma }_{i,j}2\rho _{i,j}C_{i,j}/\tau \right]T_{i,j}^{k+1/2}+\left[\mathrm{\Gamma }_{i,j}\right]T_{i,j+1}^{k+1/2}\right)`$ (35)
where we denoted
$$B_{i,j}=\frac{\lambda _{i1/2,j}}{h_r^2}\left(1\frac{1}{2i}\right),A_{i,j}=\frac{\lambda _{i+1/2,j}}{h_r^2}\left(1+\frac{1}{2i}\right).$$
(36)
$$\mathrm{\Delta }_{i,j}=\frac{\lambda _{i,j1/2}}{h_z^2},\mathrm{\Gamma }_{i,j}=\frac{\lambda _{i,j+1/2}}{h_z^2}.$$
(37)
The solving of these equations is strightforward, because equation matrixes are three-diagonal, then these equations have 3-point space scheme. It is a very important note for our method, because we will have possibility to decrease CPU time.
It might to prove 2-cycle method is absolutly time step stability and it has second-order in $`\tau ^2`$ like Crank-Nicholson scheme and second-order in space. Moreover defferential operators $`\mathrm{\Lambda }_{r,z}`$ in the 2-cycle method can be noncommutating ones.
## 4 Discretization of boundary conditions.
To solve the equations (32)-(35) it is necessary to define values of temperature on border or to set some relationships connecting temperature on border of a body with temperature of an environment. Boundary conditions in a continuous limit are certain by equations (2) - (9). We shall construct their discrete finite difference analogues on the uniform grid, using a method of finite volume . For this purpose we shall integrate the equation (13) on some volume $`G`$, adjoining to border of an ingot and on time within the limits of $`[t,t+\tau ]`$. Such volumes are shown on Fig. 2.
Integral balance of heat can be written as
$$\underset{t}{\overset{t+\tau }{}}\underset{G}{}\rho (T)C(T)\frac{T}{t}𝑑V𝑑t=\underset{t}{\overset{t+\tau }{}}\underset{G}{}\left(\frac{1}{r}\frac{}{r}(rQ)+\frac{P}{z}\right)𝑑V𝑑t,$$
(38)
where heat fluxes are defined by Eqs. (11) and (12).
After integration we shall receive the relationship connecting a thermal flux on border of a body, set by a corresponding boundary condition, and thermal fluxes in nearby points for which some approximations are required. Integration on time is led as follows: the left part of the Eq. (38) we approximate as follows
$$\underset{t}{\overset{t+\tau }{}}\underset{G}{}\rho (T)C(T)\frac{T}{t}𝑑V𝑑t\widehat{T}T,$$
(39)
and the right part of Eq. (38) we shall calculate in previous the moment of time. Now it is easy to calculate finite difference counterparts of boundary conditions.
### 4.1 Boundary $`AB`$.
Let area $`G=[0,h_r/2]\times [(j1/2)h_z,(j+1/2)h_z]`$ then the general expression for a boundary condition on a site $`AB`$ after integration can be written down in the form of
$`{\displaystyle \frac{h_r^2h_z}{8}}\rho C_{0,j}\left[\widehat{T}_{0,j}T_{0,j}\right]={\displaystyle \frac{\tau h_z}{2}}\lambda _{1/2,j}\left[T_{1,j}T_{0,j}\right]`$
$`+{\displaystyle \frac{\tau h_r^2}{8h_z}}\left(\lambda _{0,j+1/2}\left[T_{0,j+1}T_{0,j}\right]\lambda _{0,j1/2}\left[T_{0,j}T_{0,j1}\right]\right).`$ (40)
### 4.2 Boundary $`BC`$.
On a site $`BC`$ there is no necessity to use a method of finite volume as in this case the field of temperatures is set directly, therefore
$$\widehat{T}_{i,M}=T_L+\mathrm{\Delta }T.$$
(41)
### 4.3 Boundary $`CD`$.
On a site $`CD`$ we set the temperature, as well as on $`BC`$
$$\widehat{T}_{i,M}=T_L+\frac{D_{cr}2ih_r}{D_{cr}D_{el}}\mathrm{\Delta }T.$$
(42)
### 4.4 Boundary $`DE`$.
Let $`G=[(N1)h_r,Nh_r]\times [(j1/2)h_z,(j+1/2)h_z]`$ then the general expression for a boundary condition on a site $`DE`$ after integration it is possible to write down in the form of
$`{\displaystyle \frac{h_r^2h_z}{8}}(4N1)\rho C_{N,j}\left[\widehat{T}_{N,j}T_{N,j}\right]=\tau h_zNh_rq\tau h_z(N{\displaystyle \frac{1}{2}})\lambda _{N1/2,j}\left[T_{N,j}T_{N1,j}\right]`$
$`+{\displaystyle \frac{\tau h_r^2}{8h_z}}(4N1)\left(\lambda _{N,j+1/2}\left[T_{N,j+1}T_{N,j}\right]\lambda _{N,j1/2}\left[T_{N,j}T_{N,j1}\right]\right).`$ (43)
### 4.5 Boundary $`EF`$.
Let $`G=[(N1)h_r,Nh_r]\times [(j1/2)h_z,(j+1/2)h_z]`$ then the general expression for a boundary condition on a site $`EF`$ after integration it is possible to write down area in the form of
$`{\displaystyle \frac{h_r^2h_z}{8}}(4N1)\rho C_{N,j}\left[\widehat{T}_{N,j}T_{N,j}\right]`$
$`=\tau h_zNh_r\beta _{N,j}^{cr}\left[T_{N,j}T_{cr}\right]+\tau h_z(N{\displaystyle \frac{1}{2}})\lambda _{N1/2,j}\left[T_{N,j}T_{N1,j}\right]`$
$`+{\displaystyle \frac{\tau h_r^2}{8h_z}}(4N1)\left(\lambda _{N,j+1/2}\left[T_{N,j+1}T_{N,j}\right]\lambda _{N,j1/2}\left[T_{N,j}T_{N,j1}\right]\right)`$ (44)
### 4.6 Boundary $`AF`$.
Let $`G=[(i1/2)h_r,(i+1/2)h_r]\times [0,h_z/2]`$ then the general expression for a boundary condition on a site $`FA`$ after integration it is possible to write down area in the form of
$`{\displaystyle \frac{ih_r^2h_z}{2}}\rho C_{i,0}\left[\widehat{T}_{i,0}T_{i,0}\right]`$
$`={\displaystyle \frac{i\tau h_r^2}{h_z}}\left(\lambda _{i,1/2}\left[T_{i,1}T_{i,0}\right]h_z\beta _{i,0}\left[T_{i,0}T\right]\right)`$
$`=\tau h_z\left(\left(i+{\displaystyle \frac{1}{2}}\right)\lambda _{i+1/2,0}\left[T_{i+1,0}T_{i,0}\right]\left(i{\displaystyle \frac{1}{2}}\right)\lambda _{i1/2,0}\left[T_{i,0}T_{i1,0}\right]\right)`$ (45)
Further we use approximations for heat conductivity factors in semi-integer points, similar to Eq. (26)
## 5 The model parameters for VT3-1 alloy.
As it is possible to see from expressions for boundary conditions on a mirror of a liquid bath ($`BC`$ and $`CD`$) we set distribution of temperature which is defined by alloy liquidus temperature and an overheat depending on parameters of VAR. The overheat of the melt $`\mathrm{\Delta }`$ can be described by formula :
$$\mathrm{\Delta }T(J,D_{in})=400e^{12\frac{D_{in}}{J}},$$
(46)
where $`J`$ is arc current, kA; $`D_{in}`$ – ingot diameter, m. In a zone $`DE`$ the thermal flux $`q`$ to which measurements for alloy VT3-1 work is devoted is set. The parameter $`\eta `$ defines the relative contribution of each of mechanisms of heat removal: at $`\eta =1`$ the heat-conducting path goes only by radiation, and at $`\eta =0`$ only by convection. For $`EF`$ we accepted $`\eta _{EF}=1.0`$, and for $`FA`$$`\eta _{FA}=0.5`$.
$`C(T)`$, including rate of release of latent heat in an interval of temperatures between liquidus $`T_L`$ and solidus $`T_S`$, it is possible to present an effective specific thermal capacity in the form of :
$$C(T)=\{\begin{array}{cc}C_L(T)\hfill & \text{for }T>T_L\hfill \\ g(T)C_S(T)+[1g(T)]C_L(T)L\frac{dg(T)}{dT}\hfill & \text{for }T_STT_L\hfill \\ C_S(T)\hfill & \text{for }T<T_S\hfill \end{array}$$
(47)
where $`C_L`$ and $`C_S`$ are specific heat capacities for liquid and solid phases; $`L`$ is latent heat of fusion, $`g(T)`$ is solid fraction. In case of binary alloy for $`g(T)`$ we can obtain simple expression (under lever rule suggestion)
$$g(T)=\frac{T_mT_S}{T_LT_S}\frac{T_LT}{T_mT}.$$
(48)
For temperature derivative of $`g(T)`$ we can obtain
$$\frac{dg(T)}{dT}=\frac{T_mT_L}{T_LT_S}\frac{T_mT_S}{(T_mT)^2},$$
(49)
where $`T_m`$ is fusion temperature of a solvent, $`T_L`$ is liquidus temperature, $`T_S`$ is solidus temperature. Parameters of model for alloy VT3-1 have been chosen by the following $`T_m=1668^oC`$, $`T_S=1550^oC`$, $`T_L=1620^oC`$, $`L=355000J/kg`$, $`T_{cr}=70^oC`$, $`T_{bot}=70^oC`$, $`\alpha _{bot}=300W/m^2K`$. Some data can be found in . Position of a point $`E`$ was defined as in work .
For the account of increase in heat conductivity due to fluid flow we have increased heat conductivity of an alloy in a liquid phase . Temperature dependence of the resulted ingot surface emissivity degree and an internal surface of crucible has been chosen in the form of square-law dependence on temperature .
## 6 Simulation results and radiographical experiments.
Using the mathematical model described above, we have led solidification modelling of ingots from the alloy VT3-1. Modes of remeltings are specified in the Table 1.
Some works describe the experimental research of the metal pool depth and profile through fixation with radioactive isotopes. Below we compare the simulation results with the experimental data on liquid metal pool depth at different points of time. In order to adequately describe the experimental data we have carried out ”adjusting” of the model. The adjusting provided for rather exact match of the calculated and experimental metal pool depth, as well as liquid metal pool profiles at different ingot height. Due to the fact that the most extensive experimental data are available for 750 mm dia ingot melted at 37 kA, we carried out ”adjusting” for the above-mentioned melting mode. The main variable parameters were – heat conductivity in the liquid phase (that was considered not dependent on temperature) and the surface emissivity factor depending on temperature.
### 6.1 ø750 mm ingot. Arc current $`J=37kA`$.
One of the most important parameters assessed during analysis of one or another VAR mode for titanium alloys (as well as for nickel-base, iron-base and zirconium base alloys ) is the depth of the liquid metal pool. Figure 3 shows the theoretical metal pool depth with experimentally obtained depth (through fixation with tungsten radioactive isotopes) depending on the height of the melted ingot. As is obviously, the theory describes the metal pool depth behavior satisfactorily. During melting the quasi-steady state was achieved at the metal pool depth having been steady.
The description of the liquid metal pool only may be misleading as the pool depth does not actually reflect the volume of the liquid metal at the particular melt point. In order to assess feasibility of satisfactory simulation of the liquid metal pool volume using the model, we calculated the liquid metal pool profiles at the points of time recorded on the experimental radiogram. This shows that the model adequately simulates the profiles of the liquid metal pool during the whole melt process. All this together suggests that on the basis of the simulation approach in question it is feasible to simulate some most important parameters of metal solidification under VAR (width of mushy zone, temperature gradient, isotherm travel rate).
Figure 4 shows the calculated profiles of the liquid metal pool during first 140 minutes of melting (with 20 minute interval) to evaluate the linear solid-melt interface travel rate in different ingot zones. The comparison of the simulated profiles with the experimentally obtained metal pool profiles shows that the simulated metal pool volumes do not deviate from the experimental ones by more than 15%.
### 6.2 ø570 mm ingot. Arc current $`J=25kA`$.
Figure 5 shows the theoretical curves and experimental points for $`570mm`$ ingot melted at 25 kA. In this case we compared the simulated data with the experimental data according to location of the pool open surface (i.e. melt mirror) and pool bottom coordinate versus the melting time. It can be seen from Figure 5 that the quasi-steady state was not achieved.
### 6.3 ø435 mm ingot. Arc current $`J=15kA`$.
Figure 6 shows the theoretical curves and experimental points for 435 mm ingot melted at 15 kA. In this case as above we compared the simulated data with the experimental data according to location of the pool open surface (i.e. melt mirror) and pool bottom coordinate versus the melting time. It can be seen from Figure 6 that the quasi-steady state was achieved.
Therefore, a conclusion can be made that under selected parameters for Vt3-1 alloy, the simulation model can describe the main behavior tendencies of such VAR parameters as liquid metal pool depth and profile at different melting points. It is very important to note that the model that was ”adjusted” once, gives satisfactory description of the experimental data beyond the ”adjusting values”.
## 7 Conclusion.
This article describes the mathematical simulation model of heat processes that take place under VAR. The discretization has been made of the non-linear heat conductivity equation through the Marchuck’s method of integral identities and of non-linear boundary conditions through the final volume method. The solution algorithm in question has unconditional numerical stability and is applicable for tasks with non-commuting differentiation operators dependent on time.
For the purposes of the simulation model testing we calculated the liquid metal pool depth under VAR melting of titanium alloy VT3-1. It was observed that the model simulates the liquid metal pool (depth and profiles) rather adequately during the whole melting process for different ingot diameters and different current strength. The relative error in determining the liquid metal pool profile does not exceed 15%. This would enable to use the model in future to calculate various melting modes in order to assess such parameters as: liquid metal pool depth, liquid metal pool volume, width of two-phase zone, temperature gradient, isotherms travel rate and local solidification time.
## References
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# The MC@NLO 3.1 Event Generator11footnote 1Work supported in part by the UK Particle Physics and Astronomy Research Council and by the EU Fourth Framework Programme ‘Training and Mobility of Researchers’, Network ‘Quantum Chromodynamics and the Deep Structure of Elementary Particles’, contract FMRX-CT98-0194 (DG 12 - MIHT).
## 1 Generalities
In this documentation file, we briefly describe how to run the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ package, implemented according to the formalism introduced in ref. . When using $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$, please cite refs. . The production processes now available are listed in table 1. The process codes IPROC will be explained below. $`H_{1,2}`$ represent hadrons (in practice, $`p`$ or $`\overline{p}`$). The treatment of (undecayed) vector boson pair production within $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ has been described in ref. , that of heavy quark pair production in ref. . The NLO matrix elements for these processes have been taken from refs. . The information given in refs. allows the implementation in MC@NLO of any production process, provided that the formalism of refs. is used for the computation of cross sections to NLO accuracy. The matrix elements for Standard Model Higgs, single vector boson, lepton pair, and associated Higgs production have been taken from refs. , ref. , ref. , and ref. respectively.
This documentation refers to $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ version 3.1 (previous versions 1.0, 2.0, 2.2, and 2.3 are described in refs. respectively). The new processes implemented since version 3.1 are Higgs production associated with a vector boson, and $`W^+W^{}`$ production with spin correlations for the decay into leptons. For precise details of version changes, see app. A.1-A.5.
### 1.1 Mode of operation
In the case of standard MC, a hard kinematic configuration is generated on a event-by-event basis, and it is subsequently showered and hadronized. In the case of $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$, all of the hard kinematic configurations are generated in advance, and stored in a file (which we call event file – see sect. 3.1); the event file is then read by HERWIG, which showers and hadronizes each hard configuration. Since version 2.0, the events are handled by the “Les Houches” generic user process interface (see ref. for more details). Therefore, in $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ the reading of a hard configuration from the event file is equivalent to the generation of such a configuration in the standard MC.
The signal to HERWIG that configurations should be read from an event file using the Les Houches interface is a negative value of the process code IPROC; this accounts for the negative values in table 1. In the case of heavy quark pair, Higgs, Higgs in association with a $`W`$ or $`Z`$, and lepton pair (through $`Z/\gamma ^{}`$ exchange) production, the codes are simply the negative of those for the corresponding standard HERWIG MC processes. Where possible, this convention will be adopted for additional $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ processes. Consistently with what happens in standard HERWIG, by subtracting 10000 from IPROC one generates the same processes as in table 1, but eliminates the underlying event<sup>2</sup><sup>2</sup>2The same effect can be achieved by setting the HERWIG parameter PRSOF $`=0`$..
Higgs decays are controlled in the same way as in HERWIG, that is by adding -ID to the process code. The conventions for ID are the same as in HERWIG, namely ID $`=1\mathrm{}6`$ for $`u\overline{u}\mathrm{}t\overline{t}`$; $`7\mathrm{}9`$ for $`e^+e^{}\mathrm{}\tau ^+\tau ^{}`$; $`10,11`$ for $`W^+W^{},ZZ`$; and 12 for $`\gamma \gamma `$. Furthermore, ID $`=0`$ gives quarks of all flavours, and ID $`=99`$ gives all decays.
Process codes IPROC=$`1360`$IL and $`1370`$IL do not have an analogue in HERWIG; they are the same as $`1350`$IL, except for the fact that only a $`Z`$ or a $`\gamma ^{}`$ respectively is exchanged. The value of IL determines the lepton identities, and the same convention as in HERWIG is adopted: IL=$`1,\mathrm{},6`$ for $`l_{\mathrm{IL}}=e,\nu _e,\mu ,\nu _\mu ,\tau ,\nu _\tau `$ respectively. At variance with HERWIG, IL cannot be set equal to zero. Process codes IPROC=$`1460`$IL and $`1470`$IL are the analogue of HERWIG $`1450+`$IL; in HERWIG either $`W^+`$ or $`W^{}`$ can be produced, whereas MC@NLO treats the two vector bosons separately. For these processes, as in HERWIG, IL=$`1,2,3`$ for $`l_{\mathrm{IL}}=e,\mu ,\tau `$, but again the choice $`\mathrm{𝙸𝙻}=0`$ is not allowed.
The lepton pair processes IPROC=$`1350`$IL, $`\mathrm{}`$, $`1470`$IL include spin correlations when generating the angular distributions of the produced leptons. However, if spin correlations are not an issue, the single vector boson production processes IPROC= $``$1396,$``$1397,$``$1497,$``$1498 can be used, in which case the vector boson decay products are distributed according to phase space.
There are a number of other differences between the lepton pair and single vector boson processes. The latter do not feature the $`\gamma `$$`Z`$ interference terms. Also, their cross sections are fully inclusive in the final-state fermions resulting from $`\gamma ^{}`$, $`Z`$ or $`W^\pm `$. The user can still select a definite decay mode using the variable MODBOS (see sect. 3.2), but the relevant branching ratio will not be included automatically by MC@NLO. In the case of $`\gamma ^{}`$ production, the branching ratios are $`C_iq_i^2/(20/3)`$, $`q_i`$ being the electric charge (in units of the positron charge) of the fermion $`i`$ selected through MODBOS, and $`C_i=1`$ for leptons and 3 for quarks. Notice that $`20/3=_iC_iq_i^2`$, the sum including all leptons and quarks except the top. Thus, the total rate predicted by MC@NLO in the case of lepton pair production can also be recovered by multiplying the corresponding single vector boson total rate by the relevant branching ratio.
In the case of vector boson pair production, the process codes are the negative of those adopted in $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ 1.0 (for which the Les Houches interface was not yet available), rather than those of standard HERWIG.
Furthermore, in the case of Higgs production in association with a $`W`$ or $`Z`$, as well as vector boson pair production, the value of IPROC alone is not sufficient to fully determine the process type, and new variables IV, IL<sub>1</sub>, and IL<sub>2</sub> have been introduced (see table 1). The variables IL<sub>1</sub> and IL<sub>2</sub> can take the same values as IL relevant to lepton pair production (notice, however, that in the latter case IL is not an independent variable, and its value is included via IPROC); in addition, IL<sub>α</sub>=7 implies that lepton spin correlations for the decay products of the corresponding vector boson are not taken into account, as indicated in table 1.
Apart from the above differences, $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ and HERWIG behave in exactly the same way. Thus, the available user’s analysis routines can be used in the case of $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$. One should recall, however, that $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ always generates some events with negative weights (see refs. ); therefore, the correct distributions are obtained by summing weights with their signs (i.e., the absolute values of the weights must NOT be used when filling the histograms).
With such a structure, it is natural to create two separate executables, which we improperly denote as NLO and MC. The former has the sole scope of creating the event file; the latter is just HERWIG, augmented by the capability of reading the event file.
### 1.2 Package files
The package consists of the following files:
* Shell utilities
MCatNLO.Script
MCatNLO.inputs
Makefile
* Utility codes
MEcoupl.inc
alpha.f
dummies.f
linux.f
mcatnlo\_date.f
mcatnlo\_hbook.f
mcatnlo\_helas2.f
mcatnlo\_hwdummy.f
mcatnlo\_int.f
mcatnlo\_libofpdf.f
mcatnlo\_mlmtopdf.f
mcatnlo\_pdftomlm.f
mcatnlo\_str.f
mcatnlo\_uti.f
mcatnlo\_uxdate.c
sun.f
trapfpe.c
* General HERWIG routines
mcatnlo\_hwdriver.f
mcatnlo\_hwlhin.f
* Process-specific codes
mcatnlo\_hwanbtm.f
mcatnlo\_hwanhgg.f
mcatnlo\_hwanllp.f
mcatnlo\_hwantop.f
mcatnlo\_hwansvb.f
mcatnlo\_hwanvbp.f
mcatnlo\_hwanvhg.f
mcatnlo\_hgmain.f
mcatnlo\_hgxsec.f
mcatnlo\_llmain.f
mcatnlo\_llxsec.f
mcatnlo\_qqmain.f
mcatnlo\_qqxsec.f
mcatnlo\_sbmain.f
mcatnlo\_sbxsec.f
mcatnlo\_vbmain.f
mcatnlo\_vbxsec.f
mcatnlo\_vhmain.f
mcatnlo\_vhxsec.f
hgscblks.h
hvqcblks.h
llpcblks.h
svbcblks.h
vhgcblks.h
These files can be downloaded from the web page:
http://www.hep.phy.cam.ac.uk/theory/webber/MCatNLO
The files mcatnlo\_hwanxxx.f, which appear in the list of the process-specific codes, are sample HERWIG analysis routines. They are provided here to give the user a ready-to-run package, but they should be replaced with appropriate codes according to the user’s needs.
In addition to the files listed above, the user will need a version of the HERWIG code . As stressed in ref. , for the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ we do not modify the existing (LL) shower algorithm. However, since $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ versions 2.0 and higher make use of the Les Houches interface, first implemented in HERWIG 6.5, the version must be 6.500 or higher. On most systems, users will need to delete the dummy subroutines UPEVNT, UPINIT, PDFSET and STRUCTM from the standard HERWIG package, to permit linkage of the corresponding routines from the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ package. As a general rule, the user is strongly advised to use the most recent version of HERWIG (currently 6.507 – with versions lower than 6.504 problems can be found in attempting to specify the decay modes of single vector bosons through the variable MODBOS. Also, crashes in the shower phase have been reported when using HERWIG 6.505, and we therefore recommend not to use that version).
### 1.3 Working environment
We have written a number of shell scripts and a Makefile (all listed under Shell utilities above) which will simplify the use of the package considerably. In order to use them, the computing system must support bash shell, and gmake<sup>3</sup><sup>3</sup>3For Macs running under OSX v10 or higher, make can be used instead of gmake.. Should they be unavailable on the user’s computing system, the compilation and running of $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ requires more detailed instructions; in this case, we refer the reader to app. B. This appendix will serve also as a reference for a more advanced use of the package.
### 1.4 Source and running directories
We assume that all the files of the package sit in the same directory, which we call the source directory. When creating the executable, our shell scripts determine the type of operating system, and create a subdirectory of the source directory, which we call the running directory, whose name is Alpha, Sun, Linux, or Darwin, depending on the operating system. If the operating system is not known by our scripts, the name of the working directory is Run. The running directory contains all the object files and executable files, and in general all the files produced by the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ while running. It must also contain the relevant grid files (see sect. 2.1), or links to them, if the library of parton densities provided with the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ package is used.
## 2 Prior to running
Before running the code, the user needs to edit the following files:
mcatnlo\_hwanxxx.f
mcatnlo\_hwdriver.f
mcatnlo\_hwlhin.f
We do not assume that the user will adopt the latest release of HERWIG (although, as explained above, it must be version 6.500 or higher). For this reason, the files mcatnlo\_hwdriver.f and mcatnlo\_hwlhin.f must be edited, in order to modify the INCLUDE HERWIGXX.INC command to correspond to the version of HERWIG the user is going to adopt. mcatnlo\_hwdriver.f contains a set of read statements, which are necessary for the MC to get the input parameters (see sect. 3 for the input procedure); these read statements must not be modified or eliminated. Also, mcatnlo\_hwdriver.f calls the HERWIG routines which perform showering, hadronization, decays (see sect. 3.2 for more details on this issue), and so forth; the user can freely modify this part, as customary in MC runs. Finally, the sample codes mcatnlo\_hwanxxx.f contain analysis-related routines: these files must be replaced by files which contain the user’s analysis routines. We point out that, since version 2.0, the Makefile need not be edited any longer, since the corresponding operations are now performed by setting script variables (see sect. 4).
### 2.1 Parton densities
Since the knowledge of the parton densities (PDF) is necessary in order to get the physical cross section, a PDF library must be linked. The possibility exists to link the CERNLIB PDF library (PDFLIB); however, we also provide a self-contained PDF library with this package, which is faster than PDFLIB, and contains PDF sets released after the last and final PDFLIB version (8.04). A complete list of the PDFs available in our PDF library can be downloaded from the MC@NLO web page. The user may link either PDF library; all that is necessary is to set the variable PDFLIBRARY (in the file MCatNLO.inputs) equal to THISLIB if one wants to link to our PDF library, and equal to PDFLIB if one wants to link to PDFLIB. Our PDF library collects the original codes, written by the authors of the PDF fits; as such, for most of the densities it needs to read the files which contain the grids that initialize the PDFs. These files, which can be also downloaded from the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ web page, must either be copied into the running directory, or defined in the running directory as logical links to the physical files (by using ln -sn). We stress that if the user runs MC@NLO with the shell scripts, the logical links will be created automatically at run time.
As stressed before, consistent inputs must be given to the NLO and MC codes. However, in ref. we found that the dependence upon the PDFs used by the MC is rather weak. So one may want to run the NLO and MC adopting a regular NLL-evolved set in the former case, and the default HERWIG set in the latter (the advantage is that this option reduces the amount of running time of the MC). In order to do so, the user must set the variable HERPDF equal to DEFAULT in the file MCatNLO.inputs; setting HERPDF=EXTPDF will force the MC to use the same PDF set as the NLO code.
Regardless of the PDFs used in the MC run, users must delete the dummy PDFLIB routines PDFSET and STRUCTM from HERWIG, as explained earlier.
In MC@NLO 3.1, the PDF library LHAPDF is not supported.
## 3 Running
It is straightforward to run the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$. First, edit
MCatNLO.inputs
and write there all the input parameters (for the complete list of the input parameters, see sect. 4). As the last line of the file MCatNLO.inputs, write
runMCatNLO
Finally, execute MCatNLO.inputs from the bash shell. This procedure will create the NLO and MC executables, and run them using the inputs given in MCatNLO.inputs, which guarantees that the parameters used in the NLO and MC runs are consistent. Should the user only need to create the executables without running them, or to run the NLO or the MC only, he/she should replace the call to runMCatNLO in the last line of MCatNLO.inputs by calls to
compileNLO
compileMC
runNLO
runMC
which have obvious meanings. We point out that the command runMC may be used with IPROC=1350+IL, 1450+IL, 1600+ID, 1699, 1705, 1706, 2600+ID, 2699, 2700+ID, 2799 to generate $`Z/\gamma ^{}`$, $`W^\pm `$, Higgs, $`b\overline{b}`$, $`t\overline{t}`$, $`H^0W`$ or $`H^0Z`$ events with standard HERWIG (see the HERWIG manual for more details).
We stress that the input parameters are not solely related to physics (masses, CM energy, and so on); there are a few of them which control other things, such as the number of events generated. These must also be set by the user, according to his/her needs: see sect. 4.
Two such variables are HERWIGVER and HWUTI, which were moved in version 2.0 from the Makefile to MCatNLO.inputs. The former variable must be set equal to the object file name of the version of HERWIG currently adopted (matching the one whose common blocks are included in the files mentioned in sect. 2). The variable HWUTI must be set equal to the list of object files that the user needs in the analysis routines.
If the shell scripts are not used to run the codes, the inputs are given to the NLO or MC codes during an interactive talk-to phase; the complete sets of inputs for our codes are reported in app. B.2 for vector boson pair production.
### 3.1 Event file
The NLO code creates the event file. In order to do so, it goes through two steps; first it integrates the cross sections (integration step), and then, using the information gathered in the integration step, produces a set of hard events (event generation step). Integration and event generation are performed with a modified version of the SPRING-BASES package .
We stress that the events stored in the event file just contain the partons involved in the hard suprocesses. Owing to the modified subtraction introduced in the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ formalism (see ref. ) they do not correspond to pure NLO configurations, and should not be used to plot physical observables. Parton-level observables must be reconstructed using the fully-showered events.
The event generation step necessarily follows the integration step; however, for each integration step one can have an arbitrary number of event generation steps, i.e., an arbitrary number of event files. This is useful in the case in which the statistics accumulated with a given event file is not sufficient.
Suppose the user wants to create an event file; editing MCatNLO.inputs, the user sets BASES=ON, to enable the integration step, sets the parameter NEVENTS equal to the number of events wanted on tape, and runs the code; the information on the integration step (unreadable to the user, but needed by the code in the event generation step) is written on files whose name begin with FPREFIX, a string the user sets in MCatNLO.inputs; these files (which we denotes as data files) have extensions .data. The name of the event file is EVPREFIX.events, where EVPREFIX is again a string set by the user.
Now suppose the user wants to create another event file, to increase the statistics. The user simply sets BASES=OFF, since the integration step is not necessary any longer (however, the data files must not be removed: the information stored there is still used by the NLO code); changes the string EVPREFIX (failure to do so overwrites the existing event file), while keeping FPREFIX at the same value as before; and changes the value of RNDEVSEED (the random number seed used in the event generation step; failure to do so results in an event file identical to the previous one); the number NEVENTS generated may or may not be equal to the one chosen in generating the former event file(s).
We point out that data and event files may be very large. If the user wants to store them in a scratch area, this can be done by setting the script variable SCRTCH equal to the physical address of the scratch area (see sect. 3.3).
### 3.2 Decays
$`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ is intended primarily for the study of NLO corrections to production cross sections and distributions; NLO corrections to the decays of produced particles are not included. As for spin correlations in decays, the situation in version 3.1 is summarized in table 1: they are included for all processes except $`t\overline{t}`$, $`ZZ`$, and $`WZ`$ production<sup>4</sup><sup>4</sup>4Non-factorizable spin correlations of virtual origin are not included in $`W^+W^{}`$ production.. For the latter processes, quantities sensitive to the polarisation of produced particles are not given correctly even to leading order. For such quantities, it may be preferable to use the standard HERWIG MC, which does include leading-order spin correlations.
Particular decay modes of vector bosons may be forced in $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ in the same way as in standard HERWIG, using the MODBOS variables – see sect. 3.4 of ref. . However, top decays cannot be forced in this way because the decay is treated as a three-body process: the $`W^\pm `$ boson entry in HEPEVT is for information only. Instead, the top branching ratios can be altered using the HWMODK subroutine – see sect. 7 of ref. . This is done separately for the $`t`$ and $`\overline{t}`$. For example, CALL HWMODK(6,1.D0,100,12,-11,5,0,0) forces the decay $`t\nu _ee^+b`$, while leaving $`\overline{t}`$ decays unaffected. Note that the order of the decay products is important for the decay matrix element (NME = 100) to be applied correctly. The relevant statements should be inserted in the HERWIG main program (corresponding to mcatnlo\_hwdriver.f in this package) after the statement CALL HWUINC and before the loop over events. A separate run with CALL HWMODK(-6,1.D0,100,-12,11,-5,0,0) should be performed if one wishes to symmetrize the forcing of $`t`$ and $`\overline{t}`$ decays, since calls to HWMODK from within the event loop do not produce the desired result.
### 3.3 Results
As in the case of standard HERWIG the form of the results will be determined by the user’s analysis routines. However, in addition to any files written by the user’s analysis routines, the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ writes the following files:
$`\mathrm{}`$ FPREFIXNLOinput: the input file for the NLO executable, created according to the set of input parameters defined in MCatNLO.inputs (where the user also sets the string FPREFIX). See table 2.
$`\mathrm{}`$ FPREFIXNLO.log: the log file relevant to the NLO run.
$`\mathrm{}`$ FPREFIXxxx.data: xxx can assume several different values. These are the data files created by the NLO code. They can be removed only if no further event generation step is foreseen with the current choice of parameters.
$`\mathrm{}`$ FPREFIXMCinput: analogous to FPREFIXNLOinput, but for the MC executable. See table 4.
$`\mathrm{}`$ FPREFIXMC.log: analogous to FPREFIXNLO.log, but for the MC run.
$`\mathrm{}`$ EVPREFIX.events: the event file, where EVPREFIX is the string set by the user in MCatNLO.inputs.
$`\mathrm{}`$ EVPREFIXxxx.events: xxx can assume several different values. These files are temporary event files, which are used by the NLO code, and eventually removed by the shell scripts. They MUST NOT be removed by the user during the run (the program will crash or give meaningless results).
By default, all the files produced by the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ are written in the running directory. However, if the variable SCRTCH (to be set in MCatNLO.inputs) is not blank, the data and event files will be written in the directory whose address is stored in SCRTCH (such a directory is not created by the scripts, and must already exist at run time).
## 4 Script variables
In the following, we list all the variables appearing in MCatNLO.inputs; these can be changed by the user to suit his/her needs. This must be done by editing MCatNLO.inputs. For fuller details see the comments in MCatNLO.inputs.
* The CM energy of the colliding particles.
* The ratio between the renormalization scale, and a reference mass scale.
* As FREN, for the factorization scale.
* The mass (in GeV) of the top quark, except when IPROC=–(1)1705, when it is the mass of the bottom quark. In this case, HVQMASS must coincide with BMASS.
* The mass (in GeV) of the particle x, with x=HGG,W,Z,U,D,S,C,B,G.
* The physical (Breit-Wigner) width (in GeV) of the particle x, with x=HGG,W,Z.
* Valid entries are 1 and 2. If set to 1, the exact top mass dependence is retained at the Born level in Higgs production. If set to 2, the $`m_t\mathrm{}`$ limit is used.
* If xGAMMAX $`>0`$, controls the width of the mass range for Higgs (x=H) and vector bosons (x=V1,V2): the range is $`\mathrm{𝙼𝙰𝚂𝚂}\pm (\mathrm{𝙶𝙰𝙼𝙼𝙰𝚇}\times \mathrm{𝚆𝙸𝙳𝚃𝙷})`$.
* Lower limit of the Higgs (x=H) or vector boson (x=V1,V2) mass range; used only when xGAMMAX $`<0`$.
* Upper limit of the Higgs (x=H) or vector boson (x=V1,V2) mass range; used only when xGAMMAX $`<0`$.
* Set it to YES to use running $`\alpha _{em}`$ in lepton pair and single vector boson production, set it to NO to use $`\alpha _{em}=1/137.0359895`$.
* Process number that identifies the hard subprocess: see table 1 for valid entries.
* Identifies the nature of the vector boson in associated Higgs production. It corresponds to variable IV of table 1.
* Identify the nature of the leptons emerging from vector boson decays
(x $`=1,2`$). They correspond to variables IL<sub>1</sub> and IL<sub>2</sub> of table 1.
* The type of the incoming particle #n, with n=1,2. HERWIG naming conventions are used (P, PBAR, N, NBAR).
* The name of the group fitting the parton densities used; the labeling conventions of PDFLIB are adopted.
* The number of the parton density set; according to PFDLIB, the pair (PDFGROUP, PDFSET) identifies the densities for a given particle type.
* The value of $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, for five flavours and in the $`\overline{\mathrm{MS}}`$ scheme, used in the computation of NLO cross sections.
* The value of $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ used in MC runs; this parameter has the same meaning as $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ in HERWIG.
* The subtraction scheme in which the parton densities are defined.
* Our integration routine creates files with name beginning by the string FPREFIX. These files are not directly accessed by the user; for more details, see sect. 3.1.
* The name of the event file begins with this string; for more details, see sect. 3.1.
* The names of the NLO and MC executables begin with this string; this is useful in the case of simultaneous runs.
* The number of events stored in the event file, eventually processed by HERWIG .
* Valid entries are 0 and 1. When set to 0, the weights in the event file are $`\pm 1`$. When set to 1, they are $`\pm w`$, with $`w`$ a constant such that the sum of the weights gives the total NLO cross section. N.B. These weights are redefined by HERWIG at MC run time according to its own convention (see HERWIG manual).
* The seed for the random number generation in the event generation step; must be changed in order to obtain statistically-equivalent but different event files.
* Controls the integration step; valid entries are ON and OFF. At least one run with BASES=ON must be performed (see sect. 3.1).
* Valid entries are PDFLIB and THISLIB. In the former case, PDFLIB is used to compute the parton densities, whereas in the latter case the densities are obtained from our self-contained faster package.
* If set to DEFAULT, HERWIG uses its internal PDF set (controlled by NSTRU), regardless of the densities adopted at the NLO level. If set to EXTPDF, HERWIG uses the same PDFs as the NLO code (see sect. 2.1).
* The physical address of the directory where the user’s preferred version of HERWIG is stored.
* The physical address of the directory where the user wants to store the data and event files. If left blank, these files are stored in the running directory.
* This variables must be set equal to a list of object files, needed by the analysis routines of the user (for example, HWUTI=obj1.o obj2.o obj3.o is a valid assignment).
* This variable must to be set equal to the name of the object file corresponding to the version of HERWIG linked to the package (for example, HERWIGVER=herwig65.o is a valid assignment).
* The physical address of the directory where the PDF grids are stored.
## Acknowledgement
Many thanks to Paolo Nason for contributions to the heavy quark code and valuable discussions on all aspects of the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ project. We also thank V. Drollinger and B. Quayle for testing a preliminary version of the $`W^+W^{}`$ code with spin correlations.
## Appendices
## Appendix A Version changes
### A.1 From MC@NLO version 1.0 to version 2.0
In this appendix we list the changes that occurred in the package from version 1.0 to version 2.0.
$``$ The Les Houches generic user process interface has been adopted.
$``$ As a result, the convention for process codes has been changed: MC@NLO process codes IPROC are negative.
$``$ The code mcatnlo\_hwhvvj.f, which was specific to vector boson pair production in version 1.0, has been replaced by mcatnlo\_hwlhin.f, which reads the event file according to the Les Houches prescription, and works for all the production processes implemented.
$``$ The Makefile need not be edited, since the variables HERWIGVER and HWUTI have been moved to MCatNLO.inputs (where they must be set by the user).
$``$ A code mcatnlo\_hbook.f has been added to the list of utility codes. It contains a simplified version (written by M. Mangano) of HBOOK, and it is only used by the sample analysis routines mcatnlo\_hwanxxx.f. As such, the user will not need it when linking to a self-contained analysis code.
We also remind the reader that the HERWIG version must be 6.5 or higher since the Les Houches interface is used.
### A.2 From MC@NLO version 2.0 to version 2.1
In this appendix we list the changes that occurred in the package from version 2.0 to version 2.1.
$``$ Higgs production has been added, which implies new process-specific files (mcatnlo\_hgmain.f, mcatnlo\_hgxsec.f, hgscblks.h, mcatnlo\_hwanhgg.f), and a modification to mcatnlo\_hwlhin.f.
$``$ Post-1999 PDF sets have been added to the MC@NLO PDF library.
$``$ Script variables have been added to MCatNLO.inputs. Most of them are only relevant to Higgs production, and don’t affect processes implemented in version 2.0. One of them (LAMBDAHERW) may affect all processes: in version 2.1, the variables LAMBDAFIVE and LAMBDAHERW are used to set the value of $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ in NLO and MC runs respectively, whereas in version 2.0 LAMBDAFIVE controlled both. The new setup is necessary since modern PDF sets have $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ values which are too large to be supported by HERWIG. (Recall that the effect of using LAMBDAHERW different from LAMBDAFIVE is beyond NLO.)
$``$ The new script variable PDFPATH should be set equal to the name of the directory where the PDF grid files (which can be downloaded from the MC@NLO web page) are stored. At run time, when executing runNLO, or runMC, or runMCatNLO, logical links to these files will be created in the running directory (in version 2.0, this operation had to be performed by the user manually).
$``$ Minor bugs corrected in mcatnlo\_hbook.f and sample analysis routines.
### A.3 From MC@NLO version 2.1 to version 2.2
In this appendix we list the changes that occurred in the package from version 2.1 to version 2.2.
$``$ Single vector boson production has been added, which implies new process-specific files (mcatnlo\_sbmain.f, mcatnlo\_sbxsec.f, svbcblks.h, mcatnlo\_hwansvb.f), and a modification to mcatnlo\_hwlhin.f.
$``$ The script variables WWIDTH and ZWIDTH have been added to MCatNLO.inputs. These denote the physical widths of the $`W`$ and $`Z^0`$ bosons, used to generate the mass distributions of the vector bosons according to the Breit–Wigner function, in the case of single vector boson production (vector boson pair production is still implemented only in the zero-width approximation).
### A.4 From MC@NLO version 2.2 to version 2.3
In this appendix we list the changes that occurred in the package from version 2.2 to version 2.3.
$``$ Lepton pair production has been added, which implies new process-specific files (mcatnlo\_llmain.f, mcatnlo\_llxsec.f, llpcblks.h, mcatnlo\_hwanllp.f), and modifications to mcatnlo\_hwlhin.f and mcatnlo\_hwdriver.f.
$``$ The script variable AEMRUN has been added, since the computation of single vector boson and lepton pair cross sections is performed in the $`\overline{\mathrm{MS}}`$ scheme (the on-shell scheme was previously used for single vector boson production).
$``$ The script variables FRENMC and FFACTMC have been eliminated.
$``$ The structure of pseudo-random number generation in heavy flavour production has been changed, to avoid a correlation that affected the azimuthal angle distribution for the products of the hard partonic subprocesses.
$``$ A few minor bugs have been corrected, which affected the rapidity of the vector bosons in single vector boson production (a 2–3% effect), and the assignment of $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ for the LO and NLO PDF sets of Alekhin.
### A.5 From MC@NLO version 2.3 to version 3.1
In this appendix we list the changes that occurred in the package from version 2.3 to version 3.1.
$``$ Associated Higgs production has been added, which implies new process-specific files (mcatnlo\_vhmain.f, mcatnlo\_vhxsec.f, vhgcblks.h, mcatnlo\_hwanvhg.f), and modifications to mcatnlo\_hwlhin.f and mcatnlo\_hwdriver.f.
$``$ Spin correlations in $`W^+W^{}`$ production and leptonic decay have been added; the relevant codes (mcatnlo\_vpmain.f, mcatnlo\_vhxsec.f) have been modified; the sample analysis routines (mcatnlo\_hwanvbp.f) have also been changed. Tree-level matrix elements have been computed with MadGraph/MadEvent , which uses HELAS ; the relevant routines and common blocks are included in mcatnlo\_helas2.f and MEcoupl.inc.
$``$ The format of the event file has changed in several respects, the most relevant of which is that the four-momenta are now given as $`(p_x,p_y,p_z,m)`$ (up to version 2.3 we had $`(p_x,p_y,p_z,E)`$). Event files generated with version 2.3 or lower must not be used with version 3.1 or higher (the code will prevent the user from doing so).
$``$ The script variables GAMMAX, MASSINF, and MASSSUP have been replaced with xGAMMAX, xMASSINF and xMASSSUP, with x=H,V1,V2.
$``$ New script variables IVCODE, IL1CODE, and IL2CODE have been introduced.
$``$ Minor changes have been made to the routines that put the partons on the HERWIG mass shell for lepton pair, heavy quark, and vector boson pair production; effects are beyond the fourth digit.
$``$ The default electroweak parameters have been changed for vector boson pair production, in order to make them consistent with those used in other processes. The cross sections are generally smaller in version 3.1 wrt previous versions, the dominant effect being the value of $`\mathrm{sin}\theta _\mathrm{W}`$: we have now $`\mathrm{sin}^2\theta _\mathrm{W}=0.2311`$, in lower versions $`\mathrm{sin}^2\theta _\mathrm{W}=1m_W^2/m_Z^2`$. The cross sections are inversely proportional to $`\mathrm{sin}^4\theta _\mathrm{W}`$.
## Appendix B Running the package without the shell scripts
In this appendix, we describe the actions that the user needs to take in order to run the package without using the shell scripts, and the Makefile. Examples are given for vector boson pair production, but only trivial modifications are necessary in order to treat other production processes.
### B.1 Creating the executables
An $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ run requires the creation of two executables, for the NLO and MC codes respectively. The files to link depend on whether one uses PDFLIB, or the PDF library provided with this package; we list them below:
* NLO without PDFLIB: mcatnlo\_vbmain.o mcatnlo\_vbxsec.o mcatnlo\_date.o mcatnlo\_int.o mcatnlo\_uxdate.o mcatnlo\_uti.o mcatnlo\_str.o
mcatnlo\_pdftomlm.o mcatnlo\_libofpdf.o dummies.o SYSFILE
* NLO with PDFLIB: mcatnlo\_vbmain.o mcatnlo\_vbxsec.o mcatnlo\_date.o mcatnlo\_int.o mcatnlo\_uxdate.o mcatnlo\_uti.o mcatnlo\_str.o
mcatnlo\_mlmtopdf.o dummies.o SYSFILE CERNLIB
* MC without PDFLIB: mcatnlo\_hwdriver.o mcatnlo\_hwlhin.o
mcatnlo\_hwanvbp.o mcatnlo\_hbook.o mcatnlo\_str.o mcatnlo\_pdftomlm.o
mcatnlo\_libofpdf.o dummies.o HWUTI HERWIGVER
* MC with PDFLIB: mcatnlo\_hwdriver.o mcatnlo\_hwlhin.o
mcatnlo\_hwanvbp.o mcatnlo\_hbook.o mcatnlo\_str.o mcatnlo\_mlmtopdf.o
dummies.o HWUTI HERWIGVER CERNLIB
The process-specific codes mcatnlo\_vbmain.o and mcatnlo\_vbxsec.o (for the NLO executable) and mcatnlo\_hwanvbp.o (the HERWIG analysis routines in the MC executable) need to be replaced by their analogues for other production processes, which can be easily read from the list given in sect. 1.2.
The variable SYSFILE must be set either equal to alpha.o, or to linux.o, or to sun.o, according to the architecture of the machine on which the run is performed. For any other architecture, the user should provide a file corresponding to alpha.f etc., which he/she will easily obtain by modifying alpha.f. The variables HWUTI and HERWIGVER have been described in sect. 4. Finally, CERNLIB must be set in order to link the local version of CERN PDFLIB. In order to create the object files eventually linked, static compilation is always recommended (for example, g77 -Wall -fno-automatic on Linux).
### B.2 The input files
In this appendix, we describe the inputs to be given to the NLO and MC executables in the case of vector boson pair production. The case of other production processes is completely analogous. When the shell scripts are used to run the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$, two files are created, FPREFIXNLOinput and FPREFIXMCinput, which are read by the NLO and MC executable respectively. We start by considering the inputs for the NLO executable, presented in table 2.
The variables whose name is in uppercase characters have been described in sect. 4. The other variables are assigned by the shell script. Their default values are given in table 3.
Users who run the package without the script should use the values given in table 3. The variable zi controls, to a certain extent, the number of negative-weight events generated by the $`\mathrm{MC}\mathrm{@}\mathrm{NLO}`$ (see ref. ). Therefore, the user may want to tune this parameter in order to reduce as much as possible the number of negative-weight events. We stress that the MC code will not change this number; thus, the tuning can (and must) be done only by running the NLO code. The variables nitn<sub>i</sub> control the integration step (see sect. 3.1), which can be skipped by setting nitn$`{}_{i}{}^{}=0`$. If one needs to perform the integration step, we suggest setting these variables as indicated in table 3.
We now turn to the inputs for the MC executable, presented in table 4. The variables whose names are in uppercase characters have been described in sect. 4. The other variables are assigned by the shell script. Their default values are given in table 5.
The user can freely change the values of esctype and pdftype; on the other hand, the value of beammom must always be equal to half of the hadronic CM energy.
In the case of $`\gamma /Z`$, $`W^\pm `$, Higgs or heavy quark production, the MC executable can be run with the corresponding positive input process codes IPROC = 1350, 1399, 1499, 1600+ID, 1705, 1706, 2600+ID or 2700+ID, to generate a standard HERWIG run for comparison purposes<sup>5</sup><sup>5</sup>5For vector boson pair production, for historical reasons, the different process codes 2800–2825 must be used.. Then the input event file will not be read: instead, parton configurations will be generated by HERWIG according to the LO matrix elements.
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# 1 Introduction
## 1 Introduction
Given a generic Lagrangian system on a fiber bundle, its Euler–Lagrange operator $``$ obeys the Noether identities. They need not be independent, but satisfy the first-stage Noether identities, which in turn are subject to the second-order Noether ones, and so on. The hierarchy of these Noether identities characterizes the degeneracy of a Lagrangian system in full. Noether’s second theorem states the relation between Noether identities and gauge symmetries of a Lagrangian , but Noether identities can be introduced without regarding $``$ as an Euler–Lagrange operator. Therefore, one can extend the notion of Noether identities to a generic differential operator on a fiber bundle. Our goal is the following.
Let $`YX`$ be a smooth fiber bundle and $`J^rY`$, $`r=1,\mathrm{},`$ the jet manifolds of its sections (the index $`r=0`$ further stands for $`Y`$). Let $`EX`$ be a vector bundle. A $`E`$-valued $`r`$-order differential operator $``$ on $`Y`$ is conventionally defined as a bundle morphism of $`J^rY`$ to $`E`$ over $`X`$ . We associate to $``$ the chain complex (10) whose boundaries vanish on $`\mathrm{Ker}`$ (Proposition 3). It is a complex of a certain ring $`𝒮_{\mathrm{}}[E;Y]`$ of Grassmann-graded functions and their jets on the infinite order jet manifold $`J^{\mathrm{}}Y`$ of $`Y`$. For our purpose, this complex can be replaced with the short zero-exact complex $`𝒮_{\mathrm{}}[E;Y]_2`$ (11). Recall that a chain complex is called $`r`$-exact if its homology of degree $`kr`$ is trivial.
Noether identities of a differential operator $``$ are defined as nontrivial elements of the first homology $`H_1(\delta )`$ of the complex (11) (Definition 3). If this homology is finitely generated, the complex (11) can be extended to the one exact-complex $`𝒫_{\mathrm{}}[E;Y;E_0]_3`$ (19) with the boundary operator $`\delta _0`$ (18) whose nilpotency conditions are equivalent to the above-mentioned Noether identities (Proposition 3). First-stage Noether identities are defined as two-cycles of this complex. They are trivial if two-cycles are boundaries, but the converse need not be true. Trivial first-stage Noether identities are boundaries iff a certain homology condition (called the two-homology regularity condition) holds (Proposition 4). In this case, the first-stage Noether identities are identified to nontrivial elements of the second homology of the complex (19). If this homology is finitely generated, the complex (19) is extended to the two-exact complex $`𝒫_{\mathrm{}}[E_1E;Y;E_0]_4`$ (30) with the boundary operator $`\delta _1`$ (29) whose nilpotency conditions are equivalent to the Noether and first-stage Noether identities (Proposition 4). If the third homology of this complex is not trivial, the second-stage Noether identities exist, and so on. Iterating the arguments, we come to the following.
Let we have the $`(N+1)`$-exact complex $`𝒫_{\mathrm{}}\{N\}_{N+3}`$ (33) such that: (i) the nilpotency conditions of its boundary operator $`\delta _N`$ (31) reproduce Noether and $`k`$-stage Noether identities for $`kN`$, (ii) the $`(N+1)`$-homology regularity condition (Definition 4) holds. Then the $`(N+1)`$-stage Noether identities are defined as $`(N+2)`$-cycles of this complex. They are trivial if cycles are boundaries, while the converse is true iff the $`(N+2)`$-homology regularity condition is satisfied. In this case, $`(N+1)`$-stage Noether identities are identified to nontrivial elements of the $`(N+2)`$-homology of the complex (33) (item (i) of Theorem 4). If this homology is finitely generated, this complex is extended to the $`(N+2)`$-exact complex $`𝒫_{\mathrm{}}\{N+1\}_{N+4}`$ (41) with the boundary operator $`\delta _{N+1}`$ (40) whose nilpotency restarts all the Noether identities up to stage $`(N+1)`$ (item (ii) of Theorem 4).
This iteration procedure results in the exact Koszul–Tate complex with the boundary operator whose nilpotency conditions reproduce all Noether and higher Noether identities characterizing the degeneracy of a differential operator $``$
## 2 The ring of Grassmann-graded functions and their jets
All chain complexes considered in the article are complexes of certain rings of Grassmann-graded function and their jets on the infinite order jet manifold $`J^{\mathrm{}}Y`$ of a fiber bundle $`YX`$. Let us describe such a ring.
Recall that $`J^{\mathrm{}}Y`$, is the projective limit $`(\pi _r^{\mathrm{}}:J^{\mathrm{}}YJ^rY)`$ of the inverse system of jet manifolds
$$X\stackrel{\pi }{}Y\stackrel{\pi _0^1}{}J^1Y\mathrm{}J^{r1}Y\stackrel{\pi _{r1}^r}{}J^rY\mathrm{},$$
(1)
of $`J^{\mathrm{}}Y`$, where $`\pi _{r1}^r`$ are affine bundles. It is a Fréchet manifold. A bundle atlas $`\{(U_Y;x^\lambda ,y^i)\}`$ of $`YX`$ induces the coordinate atlas
$$\{((\pi _0^{\mathrm{}})^1(U_Y);x^\lambda ,y_\mathrm{\Lambda }^i)\},y_{}^{}{}_{\lambda +\mathrm{\Lambda }}{}^{i}=\frac{x^\mu }{x^\lambda }d_\mu y_\mathrm{\Lambda }^i,0|\mathrm{\Lambda }|,$$
(2)
where $`\mathrm{\Lambda }=(\lambda _1\mathrm{}\lambda _k)`$, $`\lambda +\mathrm{\Lambda }=(\lambda \lambda _1\mathrm{}\lambda _k)`$ are symmetric multi-indices and
$`d_\lambda =_\lambda +{\displaystyle \underset{0|\mathrm{\Lambda }|}{}}y_{\lambda +\mathrm{\Lambda }}^i_i^\mathrm{\Lambda },d_\mathrm{\Lambda }=d_{\lambda _1}\mathrm{}d_{\lambda _k},`$
are total derivatives. We further assume that the cover $`\{\pi (U_Y)\}`$ of $`X`$ is also the cover of atlases of all vector bundles over $`X`$ in question. The inverse system (1) yields the direct system
$`C^{\mathrm{}}(X)\stackrel{\pi ^{}}{}C^{\mathrm{}}(Y)\stackrel{\pi _0^1^{}}{}C^{\mathrm{}}(J^1Y)\mathrm{}C^{\mathrm{}}(J^{r1}Y)\stackrel{\pi _{r1}^r^{}}{}\mathrm{}`$
of rings of smooth real functions on jet manifolds $`J^rY`$ with respect to the pull-back monomorphisms $`\pi _{r1}^r^{}`$. Its direct limit is the ring $`𝒪_{\mathrm{}}Y`$ of all smooth real functions on finite order jet manifolds modulo the pull-back identification.
Let us extend the ring $`𝒪_{\mathrm{}}Y`$ to a ring of graded functions on graded manifolds whose bodies are jet manifolds $`J^rY`$ of $`Y`$ . We restrict our consideration to graded manifolds $`(Z,𝔄)`$ with structure sheaves $`𝔄`$ of Grassmann algebras of finite rank, and refer to the following Serre–Swan theorem.
Theorem 1. Let $`Z`$ be a smooth manifold. A Grassmann algebra $`𝒜`$ over the ring $`C^{\mathrm{}}(Z)`$ of smooth real functions on $`Z`$ is isomorphic to the Grassmann algebra of graded functions on a graded manifold with a body $`Z`$ iff it is the exterior algebra of some projective $`C^{\mathrm{}}(Z)`$-module of finite rank.
Proof. The proof follows at once from the Batchelor theorem and the Serre-Swan theorem generalized to an arbitrary smooth manifold . The Batchelor theorem states that any graded manifold $`(Z,𝔄)`$ with a body $`Z`$ is isomorphic to the one $`(Z,𝔄_Q)`$ with the structure sheaf $`𝔄_Q`$ of germs of sections of the exterior bundle
$`Q^{}=R\underset{Z}{}Q^{}\underset{Z}{}\stackrel{2}{}Q^{}\underset{Z}{}\mathrm{},`$
where $`Q^{}`$ is the dual of some vector bundle $`QZ`$. We agree to call $`(Z,𝔄_Q)`$ the simple graded manifold with the structure vector bundle $`Q`$. Its structure ring $`𝒜_Q`$ of graded functions (sections of $`𝔄_Q`$) is the $`Z_2`$-graded exterior algebra of the $`C^{\mathrm{}}(Z)`$-module of sections of $`Q^{}Z`$. By virtue of the Serre–Swan theorem, a $`C^{\mathrm{}}(Z)`$-module is isomorphic to the module of sections of a smooth vector bundle over $`Z`$ iff it is a projective module of finite rank. $`\mathrm{}`$
Remark 1. With respect to bundle coordinates $`(z^A,q^a)`$ on $`Q`$ and the corresponding fiber basis $`\{c^a\}`$ for $`Q^{}X`$, graded functions read
$$f=\underset{k=0}{}\frac{1}{k!}f_{a_1\mathrm{}a_k}c^{a_1}\mathrm{}c^{a_k},fC^{\mathrm{}}(Z),$$
(3)
where we omit the symbol of the exterior product of elements $`c^a`$. Let $`u𝔡𝒜_Q`$ be a graded derivation of the $`R`$-ring $`𝒜_Q`$. Due to the canonical splitting $`VQ=Q\times Q`$, the fiber basis $`\{_a\}`$ for vertical tangent bundle $`VQQ`$ of $`QZ`$ is the dual of $`\{c^a\}`$ and, therefore, $`u`$ takes the local form $`u=u^A_A+u^a_a`$, where $`u^A,u^a`$ are local graded functions . It acts on graded functions (3) by the rule
$`u(f_{a\mathrm{}b}c^a\mathrm{}c^b)=u^A_A(f_{a\mathrm{}b})c^a\mathrm{}c^b+u^df_{a\mathrm{}b}_d(c^a\mathrm{}c^b).`$
Given a vector bundle $`EX`$, let us consider the simple graded manifold $`(J^rY,𝔄_{E_r})`$ whose body is $`J^rY`$ and the structure bundle is the pull-back
$`E_r=J^rY\underset{X}{\times }J^rE`$
onto $`J^rY`$ of the jet bundle $`J^rEX`$, which is a vector bundle. There is an epimorphism of graded manifolds
$`(J^{r+1}Y,𝔄_{E_{r+1}})(J^rY,𝔄_{E_r}).`$
It consists of the open surjection $`\pi _r^{r+1}`$ and the sheaf monomorphism $`\pi _r^{r+1}:𝔄_{E_r}𝔄_{E_{r+1}}`$, where $`\pi _r^{r+1}𝔄_{E_r}`$ is the pull-back onto $`J^{r+1}Y`$ of the topological fiber bundle $`𝔄_{E_r}J^rY`$. These sheaf monomorphisms yield monomorphisms of the corresponding canonical presheaves $`\overline{𝔄}_{E_r}\overline{𝔄}_{E_{r+1}}`$, which make up a direct system
$$\overline{𝔄}_{Y\times E}\overline{𝔄}_{E_1}\mathrm{}\overline{𝔄}_{E_r}\mathrm{},$$
(4)
and the monomorphisms of graded commutative rings $`𝒜_{E_r}𝒜_{E_{r+1}}`$ assembled into the direct system
$$𝒜_{Y\times E}𝒜_{E_1}\mathrm{}𝒜_{E_r}\mathrm{}.$$
(5)
A direct limit of this direct system is a graded commutative ring $`𝒮_{\mathrm{}}[E;Y]`$ of all graded functions $`f𝒜_{E_r}`$ on jet manifolds $`J^rY`$ modulo monomorphisms $`\pi _r^{r+1}`$. The monomorphisms $`C^{\mathrm{}}(J^rY)𝒜_{E_r}`$ provide the monomorphism $`𝒪_{\mathrm{}}Y𝒮_{\mathrm{}}[E;Y]`$, while the body epimorphisms $`𝒜_{E_r}C^{\mathrm{}}(J^rY)`$ yield the epimorphism $`𝒮_{\mathrm{}}[E;Y]𝒪_{\mathrm{}}Y`$.
One can think of elements of $`𝒮_{\mathrm{}}[E;Y]`$ as being graded functions on $`J^{\mathrm{}}Y`$ as follows. A direct limit of the direct system of presheaves (4) is a presheaf on the infinite order jet manifold $`J^{\mathrm{}}Y`$. Let $`𝔗_{\mathrm{}}[E;Y]`$ be the sheaf of germs of this presheaf. The structure module $`\mathrm{\Gamma }(𝔗_{\mathrm{}}[E;Y])`$ of sections of $`𝔗_{\mathrm{}}[E;Y]`$ is a ring such that, given an element $`f\mathrm{\Gamma }(𝔗_{\mathrm{}}[E;Y])`$ and a point $`zJ^{\mathrm{}}Y`$, there exist an open neighbourhood $`U`$ of $`z`$ and a graded function $`f^{(k)}`$ on some finite order jet manifold $`J^kY`$ so that $`f|_U=\pi _k^{\mathrm{}}f^{(k)}|_U`$. In particular, there is the monomorphism $`𝒮_{\mathrm{}}[E;Y]\mathrm{\Gamma }(𝔗_{\mathrm{}}[E;Y])`$. Due to this monomorphism, one can restrict $`𝒮_{\mathrm{}}[E;Y]`$ to the coordinate chart (2), and say that $`𝒮_{\mathrm{}}[E;Y]`$ as an $`𝒪_{\mathrm{}}Y`$-ring is locally generated by the elements $`c_\mathrm{\Lambda }^a`$, where $`\{c^a\}`$ is a a local fiber basis for $`E^{}`$ over $`\pi _{\mathrm{}}(U)`$. We agree to call $`(y_\mathrm{\Lambda }^i,c_\mathrm{\Lambda }^a)`$, $`0|\mathrm{\Lambda }|`$, the local basis for $`𝒮_{\mathrm{}}[E;Y]`$. Let the collective symbol $`s_\mathrm{\Lambda }^A`$ stand for its elements.
Remark 2. One can think of $`c_\mathrm{\Lambda }^a`$ as being jets of graded functions $`c^a`$ . This definition differs from the notion of jets of a graded commutative ring and that of jets of a graded fiber bundle , but reproduces the heuristic notion of jets of odd fields in Lagrangian field theory .
Remark 3. Let $`V`$, $`V^{}`$ and $`V^{\prime \prime }`$ be vector bundles over $`X`$. Let us consider the ring
$`𝒮_{\mathrm{}}[V\underset{X}{\times }E;Y\underset{X}{\times }V^{}\underset{X}{\times }V^{\prime \prime }],`$
and its subring of graded functions which are polynomial in fiber coordinates of the vector bundle $`Y\underset{X}{\times }V^{}\underset{X}{\times }V^{\prime \prime }Y`$. We denote the latter by $`𝒫_{\mathrm{}}[VE;Y;V^{}V^{\prime \prime }]`$. One can think of its elements as being graded functions on $`J^{\mathrm{}}Y`$, too.
Let $`\vartheta 𝔡𝒮_{\mathrm{}}[E;Y]`$ be a graded derivation of the $`R`$-ring $`𝒮_{\mathrm{}}[E;Y]`$ . With respect to the local basis $`(s_\mathrm{\Lambda }^A)`$, it takes the form
$`\vartheta =\vartheta ^\lambda d_\lambda +(\vartheta ^A_A+{\displaystyle \underset{|\mathrm{\Lambda }|>0}{}}d_\mathrm{\Lambda }\vartheta ^A_A^\mathrm{\Lambda }),`$
where the tuple of graded derivations $`\{_a^\mathrm{\Lambda }\}`$ is the dual of the tuple $`\{c_\mathrm{\Lambda }^a\}`$ of generating elements of the $`𝒪_{\mathrm{}}[F;Y]`$-algebra $`𝒮_{\mathrm{}}[E;Y]`$, and $`\vartheta ^\lambda `$, $`\vartheta ^A`$ are local graded functions. We further restrict our consideration to vertical contact graded derivations
$$\vartheta =\underset{0|\mathrm{\Lambda }|}{}d_\mathrm{\Lambda }\upsilon ^A_A^\mathrm{\Lambda }.$$
(6)
Such a derivation is completely determined by its first summand
$$\upsilon =\upsilon ^A(x^\lambda ,s_\mathrm{\Lambda }^A)_A,0|\mathrm{\Lambda }|k,$$
(7)
called a generalized graded vector field. For the sake of simplicity, the common symbol $`\upsilon `$ further stands both for a contact graded derivation (6) and a generalized vector field (7), which is also called a graded derivation. A graded derivation $`\upsilon `$ is said to be nilpotent if $`\upsilon (\upsilon (f))=0`$ for any graded function $`f𝒮_{\mathrm{}}[E;Y]`$. One can show that $`\upsilon `$ is nilpotent only if it is odd and iff all $`\upsilon ^A`$ obey the equality
$$\upsilon (\upsilon ^A)=\underset{0|\mathrm{\Sigma }|}{}\upsilon _\mathrm{\Sigma }^B_B^\mathrm{\Sigma }(\upsilon ^A)=0.$$
(8)
## 3 Noether identities
Let $``$ be a $`E`$-valued differential operator on a smooth fiber bundle $`YX`$.
Proposition 2. One can associate to $``$ a chain complex whose boundaries vanish on $`\mathrm{Ker}`$.
Proof. A differential operator $``$ is locally represented by a set of functions $`^a𝒪_{\mathrm{}}Y`$, possessing the corresponding coordinate transformation law. Let us provide the ring $`𝒮_{\mathrm{}}[E;Y]`$ with the nilpotent graded derivation
$$\delta =\stackrel{}{}{}_{a}{}^{}_{}^{a},$$
(9)
whose definition is independent of the choice of a local basis. It is convenient to deal with a graded derivation $`\delta `$ (9) acting on graded functions on the right by the rule
$`\delta (ff^{})=(1)^{[f^{}]}\delta (f)f^{}+f\delta (f^{}),`$
where $`[f^{}]`$ denotes the Grassmann parity. We call $`\delta `$ (9) the Koszul–Tate differential. With this differential, the ring $`𝒮_{\mathrm{}}[E;Y]`$ is split into the chain complex
$$0𝒪_{\mathrm{}}Y\stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_1\mathrm{}\stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_k\mathrm{}$$
(10)
graded in polynomials of the odd elements $`c_\mathrm{\Lambda }^a`$. Following the physical literature , we assign to $`c_\mathrm{\Lambda }^a`$ the antifield number Ant$`[c_\mathrm{\Lambda }^a]=1`$. It is readily observed that the boundaries of the complex (10) vanish on $`\mathrm{Ker}`$. $`\mathrm{}`$
Note that homology groups $`H_{}(\delta )`$ of the complex (10) are $`𝒪_{\mathrm{}}Y`$-modules, but these modules fail to be torsion-free. Indeed, given a cycle $`\mathrm{\Phi }𝒮_{\mathrm{}}[E;Y]_k`$ and an element $`f=\delta \sigma `$ of the ring $`𝒪_{\mathrm{}}Y𝒮_{\mathrm{}}[E;Y]`$, we obtain that $`f\mathrm{\Phi }=\delta (\sigma \mathrm{\Phi })`$ is a boundary. Therefore, one can not apply the Künneth formula to the homology of this complex, though any its term $`𝒮_{\mathrm{}}[E;Y]_k`$ is isomorphic to the graded commutative $`k`$-tensor product of the $`𝒪_{\mathrm{}}Y`$-module $`𝒮_{\mathrm{}}[E;Y]_1`$.
The homology $`H_0(\delta )`$ of the complex (10) is not trivial, but this homology and the higher ones $`H_{k2}(\delta )`$ are not essential for our consideration. Therefore, we replace the complex (10) with the finite one
$$0\mathrm{Im}\delta \stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_1\stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_2$$
(11)
of graded functions of antifield number $`k2`$. It is exact at $`\mathrm{Im}\delta `$, and its first homology coincides with that of the complex (10). Let us consider this homology.
A generic one-chain of the complex (11) takes the form
$$\mathrm{\Phi }=\underset{0|\mathrm{\Lambda }|}{}\mathrm{\Phi }_a^\mathrm{\Lambda }c_\mathrm{\Lambda }^a,\mathrm{\Phi }_a^\mathrm{\Lambda }𝒪_{\mathrm{}}Y,$$
(12)
and the cycle condition $`\delta \mathrm{\Phi }=0`$ reads
$$\underset{0|\mathrm{\Lambda }|}{}\mathrm{\Phi }_a^\mathrm{\Lambda }d_\mathrm{\Lambda }^a=0.$$
(13)
One can think of this equality as being a reduction condition on a differential operator $``$. Conversely, any reduction condition of form (13) comes from some cycle (12). The reduction condition (13) is trivial if a cycle is a boundary, i.e., it takes the form
$$\mathrm{\Phi }=\underset{0|\mathrm{\Lambda }|,|\mathrm{\Sigma }|}{}T_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}d_\mathrm{\Sigma }^bc_\mathrm{\Lambda }^a,T_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}=T_{ba}^{\mathrm{\Sigma }\mathrm{\Lambda }}.$$
(14)
If $``$ is an Euler–Lagrange operator of some Lagrangian system on a fiber bundle $`YX`$, the nontrivial reduction condition (13) is a Noether identity . Therefore, we come to the following definition.
Definition 3. A differential operator $``$ is called degenerate if the homology $`H_1(\delta )`$ of the complex (14) (or (12)) is not trivial. We agree to call a cycle condition (13) the Noether identity.
One can say something more if the $`𝒪_{\mathrm{}}Y`$-module $`H_1(\delta )`$ is finitely generated, i.e., it possesses the following particular structure. There are elements $`\mathrm{\Delta }H_1(\delta )`$ making up a projective $`C^{\mathrm{}}(X)`$-module $`𝒞_{(0)}`$ of finite rank which, by virtue of the Serre–Swan theorem, is isomorphic to the module of sections of some vector bundle $`E_0X`$. Let $`\{\mathrm{\Delta }^r\}`$ be local bases for this $`C^{\mathrm{}}(X)`$-module. Then every element $`\mathrm{\Phi }H_1(\delta )`$ factorizes
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Xi }|}{}}G_r^\mathrm{\Xi }d_\mathrm{\Xi }\mathrm{\Delta }^r,G_r^\mathrm{\Xi }𝒪_{\mathrm{}}Y,`$ (15)
$`\mathrm{\Delta }^r={\displaystyle \underset{0|\mathrm{\Lambda }|}{}}\mathrm{\Delta }_a^{\mathrm{\Lambda }r}c_\mathrm{\Lambda }^a,\mathrm{\Delta }_a^{\mathrm{\Lambda }r}𝒪_{\mathrm{}}Y,`$ (16)
via elements of $`𝒞_{(0)}`$, i.e., any Noether identity (13) is a corollary of the Noether identities
$$\underset{0|\mathrm{\Lambda }|}{}\mathrm{\Delta }_a^{\mathrm{\Lambda }r}d_\mathrm{\Lambda }^a=0.$$
(17)
Clearly, the factorization (15) is independent of specification of local bases $`\{\mathrm{\Delta }^r\}`$. We say that the Noether identities (17) are complete, and call $`\mathrm{\Delta }𝒞_{(0)}`$ the Noether operators. Note that, being representatives of $`H_1(\delta )`$, the graded functions $`\mathrm{\Delta }^r`$ (16) are not $`\delta `$-exact.
Proposition 4. If the homology $`H_1(\delta )`$ of the complex (11) is finitely generated, this complex can be extended to a one-exact complex with a boundary operator whose nilpotency conditions (8) are just complete Noether identities (see the complex (19) below).
Proof. Let us consider the ring $`𝒫_{\mathrm{}}[E;Y;E_0]`$ of graded functions on $`J^{\mathrm{}}Y`$ (see Remark 2). It possesses local bases $`\{y_\mathrm{\Lambda }^i,c_\mathrm{\Lambda }^a,c_\mathrm{\Lambda }^r\}`$, where $`[c_\mathrm{\Lambda }^r]=0`$ and Ant$`[c_\mathrm{\Lambda }^r]=2`$. This ring is provided with the nilpotent graded derivation
$$\delta _0=\delta +\stackrel{}{}{}_{r}{}^{}\mathrm{\Delta }_{}^{r},$$
(18)
called the extended Koszul–Tate differential. Its nilpotency conditions (8) are equivalent to the complete Noether identities (17). Then the module $`𝒫_{\mathrm{}}[E;Y;E_0]_3`$ of graded functions of antifield number $`k3`$ is split into the chain complex
$$0\mathrm{Im}\delta \stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_1\stackrel{\delta _0}{}𝒫_{\mathrm{}}[E;Y;E_0]_2\stackrel{\delta _0}{}𝒫_{\mathrm{}}[E;Y;E_0]_3.$$
(19)
Let $`H_{}(\delta _0)`$ denote its homology. We have $`H_0(\delta _0)=H_0(\delta )=0`$. Furthermore, any one-cycle $`\mathrm{\Phi }`$ up to a boundary takes the form (15) and, therefore, it is a $`\delta _0`$-boundary
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Sigma }|}{}}G_r^\mathrm{\Xi }d_\mathrm{\Xi }\mathrm{\Delta }^r=\delta _0({\displaystyle \underset{0|\mathrm{\Sigma }|}{}}G_r^\mathrm{\Xi }c_\mathrm{\Xi }^r).`$
Hence, $`H_1(\delta _0)=0`$, i.e., the complex (19) is one-exact. $`\mathrm{}`$
## 4 The Koszul–Tate complex
Turn now to the second homology $`H_2(\delta _0)`$ of the complex (19). A generic two-chain reads
$$\mathrm{\Phi }=G+H=\underset{0|\mathrm{\Lambda }|}{}G_r^\mathrm{\Lambda }c_\mathrm{\Lambda }^r+\underset{0|\mathrm{\Lambda }|,|\mathrm{\Sigma }|}{}H_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}c_\mathrm{\Lambda }^ac_\mathrm{\Sigma }^b,G_r^\mathrm{\Lambda },H_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}𝒪_{\mathrm{}}Y.$$
(20)
The cycle condition $`\delta _0\mathrm{\Phi }=0`$ takes the form
$$\underset{0|\mathrm{\Lambda }|}{}G_r^\mathrm{\Lambda }d_\mathrm{\Lambda }\mathrm{\Delta }^r+\delta H=0.$$
(21)
One can think of this equality as being the reduction condition on the Noether operators $`\mathrm{\Delta }^r`$ (16). Conversely, let
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Lambda }|}{}}G_r^\mathrm{\Lambda }c_\mathrm{\Lambda }^r𝒫_{\mathrm{}}[E;Y;E_0]_2`$
be a graded function such that the reduction condition (21) holds. Obviously, it is a cycle condition of the two-chain (20). The reduction condition (21) is trivial either if a two-cycle $`\mathrm{\Phi }`$ (20) is a boundary or its summand $`G`$ vanishes on $`\mathrm{Ker}`$.
Definition 5. A degenerate differential operator is said to be one-stage reducible if there exist non-trivial reduction conditions (21), called first-stage Noether identities.
Proposition 6. First-stage Noether identities can be identified to nontrivial elements of the homology $`H_2(\delta _0)`$ iff any $`\delta `$-cycle $`\mathrm{\Phi }𝒮_{\mathrm{}}[E;Y]_2`$ is a $`\delta _0`$-boundary.
Proof. It suffices to show that, if the summand $`G`$ of a two-cycle $`\mathrm{\Phi }`$ (20) is $`\delta `$-exact, then $`\mathrm{\Phi }`$ is a boundary. If $`G=\delta \mathrm{\Psi }`$, then
$$\mathrm{\Phi }=\delta _0\mathrm{\Psi }+(\delta \delta _0)\mathrm{\Psi }+H.$$
(22)
The cycle condition reads
$`\delta _0\mathrm{\Phi }=\delta ((\delta \delta _0)\mathrm{\Psi }+H)=0.`$
Then $`(\delta \delta _0)\mathrm{\Psi }+H`$ is $`\delta _0`$-exact since any $`\delta `$-cycle $`\varphi 𝒮_{\mathrm{}}[E;Y]_2`$ by assumption is a $`\delta _0`$-boundary. Consequently, $`\mathrm{\Phi }`$ (22) is $`\delta _0`$-exact. Conversely, let $`\mathrm{\Phi }𝒮_{\mathrm{}}[E;Y]_2`$ be an arbitrary $`\delta `$-cycle. The cycle condition reads
$$\delta \mathrm{\Phi }=2\mathrm{\Phi }_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}c_\mathrm{\Lambda }^a\delta c_\mathrm{\Sigma }^b=2\mathrm{\Phi }_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}c_\mathrm{\Lambda }^ad_\mathrm{\Sigma }^b=0.$$
(23)
It follows that $`\mathrm{\Phi }_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}\delta c_\mathrm{\Sigma }^b=0`$ for all induces $`(a,\mathrm{\Lambda })`$. Omitting a $`\delta `$-boundary term, we obtain
$`\mathrm{\Phi }_{ab}^{\mathrm{\Lambda }\mathrm{\Sigma }}c_\mathrm{\Sigma }^b=G_{ar}^{\mathrm{\Lambda }\mathrm{\Xi }}d_\mathrm{\Xi }\mathrm{\Delta }^r.`$
Hence, $`\mathrm{\Phi }`$ takes the form
$$\mathrm{\Phi }=G_{ar}^{\mathrm{\Lambda }\mathrm{\Xi }}d_\mathrm{\Xi }\mathrm{\Delta }^rc_\mathrm{\Lambda }^a.$$
(24)
We can associate to it the three-chain $`\mathrm{\Psi }=G_{ar}^{\mathrm{\Lambda }\mathrm{\Xi }}c_\mathrm{\Xi }^rc_\mathrm{\Lambda }^a`$ such that
$`\delta _0\mathrm{\Psi }=\mathrm{\Phi }+\sigma =\mathrm{\Phi }G_{ar}^{\mathrm{\Lambda }\mathrm{\Xi }}d_\mathrm{\Lambda }^ac_\mathrm{\Xi }^r.`$
Owing to the equality $`\delta \mathrm{\Phi }=0`$, we have $`\delta _0\sigma =0`$. Since $`\sigma `$ is $`\delta `$-exact, it by assumption is $`\delta _0`$-exact, i.e., $`\sigma =\delta _0\psi `$. Then we obtain that $`\mathrm{\Phi }=\delta _0\mathrm{\Psi }\delta _0\psi `$. $`\mathrm{}`$
Lemma 7. It is easily justified that a two-cycle $`\mathrm{\Phi }𝒮_{\mathrm{}}[E;Y]_2`$ is $`\delta _0`$-exact iff $`\mathrm{\Phi }`$ up to a $`\delta `$-boundary takes the form
$$\mathrm{\Phi }=\underset{0|\mathrm{\Lambda }|,|\mathrm{\Sigma }|}{}G_{rr^{}}^{\mathrm{\Sigma }\mathrm{\Lambda }}d_\mathrm{\Sigma }\mathrm{\Delta }^rd_\mathrm{\Lambda }\mathrm{\Delta }^r^{}.$$
(25)
If the condition of Proposition 4 (called the two-homology regularity condition) is satisfied, let us assume that first-stage Noether identities are finitely generated as follows. There are elements $`\mathrm{\Delta }_{(1)}H_2(\delta _0)`$ making up a projective $`C^{\mathrm{}}(X)`$-module $`𝒞_{(1)}`$ of finite rank which is isomorphic to the module of sections of some vector bundle $`E_1X`$. Let $`\{\mathrm{\Delta }^{r_1}\}`$ be local bases for this $`C^{\mathrm{}}(X)`$-module. Every element $`\mathrm{\Phi }H_2(\delta _0)`$ factorizes
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Xi }|}{}}\mathrm{\Phi }_{r_1}^\mathrm{\Xi }d_\mathrm{\Xi }\mathrm{\Delta }^{r_1},\mathrm{\Phi }_{r_1}^\mathrm{\Xi }𝒪_{\mathrm{}}Y,`$ (26)
$`\mathrm{\Delta }^{r_1}=G^{r_1}+h^{r_1}={\displaystyle \underset{0|\mathrm{\Lambda }|}{}}\mathrm{\Delta }_r^{\mathrm{\Lambda }r_1}c_\mathrm{\Lambda }^r+h^{r_1},h^{r_1}𝒮_{\mathrm{}}[E;Y]_2,`$ (27)
via elements of $`𝒞_{(1)}`$, i.e., any first-stage Noether identity (21) results from the equalities
$$\underset{0|\mathrm{\Lambda }|}{}\mathrm{\Delta }_r^{r_1\mathrm{\Lambda }}d_\mathrm{\Lambda }\mathrm{\Delta }^r+\delta h^{r_1}=0,$$
(28)
called the complete first-stage Noether identities. Elements of $`𝒞_{(1)}`$ are said to be the first-stage Noether operators. Note that the first summands $`G^{r_1}`$ of the operators $`\mathrm{\Delta }^{r_1}`$ (27) are not $`\delta `$-exact.
Proposition 8. Given a reducible degenerate differential operator $``$, let the associated one-exact complex (19) obey the two-homology regularity condition and let its homology $`H_2(\delta _0)`$ be finitely generated. Then this complex is extended to the two-exact one with a boundary operator whose nilpotency conditions are equivalent to complete Noether and first-stage Noether identities (see the complex (30) below).
Proof. Let us consider the ring $`𝒫_{\mathrm{}}^{}[E_1E;Y;E_0]`$ of graded functions on $`J^{\mathrm{}}Y`$ possessing local bases $`\{y_\mathrm{\Lambda }^i,c_\mathrm{\Lambda }^a,c_\mathrm{\Lambda }^r,c_\mathrm{\Lambda }^{r_1}\}`$, where $`[c_\mathrm{\Lambda }^{r_1}]=1`$ and Ant$`[c_\mathrm{\Lambda }^{r_1}]=3`$. It can be provided the first-stage Koszul–Tate differential defined as the nilpotent graded derivation
$$\delta _1=\delta _0+\stackrel{}{}{}_{r_1}{}^{}\mathrm{\Delta }_{}^{r_1}.$$
(29)
Its nilpotency conditions (8) are equivalent to complete Noether identities (17) and complete first-stage Noether identities (28). Then the module $`𝒫_{\mathrm{}}[E_1E;Y;E_0]_4`$ of graded functions of antifield number Ant$`[\varphi ]4`$ is split into the chain complex
$`0\mathrm{Im}\delta \stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_1\stackrel{\delta _0}{}𝒫_{\mathrm{}}[E;Y;E_0]_2\stackrel{\delta _1}{}𝒫_{\mathrm{}}[E_1E;Y;E_0]_3`$ (30)
$`\stackrel{\delta _1}{}𝒫_{\mathrm{}}[E_1E;Y;E_0]_4.`$
Let $`H_{}(\delta _1)`$ denote its homology. It is readily observed that
$`H_0(\delta _1)=H_0(\delta )=0,H_1(\delta _1)=H_1(\delta _0)=0.`$
By virtue of the expression (26), any two-cycle of the complex (30) is a boundary
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Xi }|}{}}\mathrm{\Phi }_{r_1}^\mathrm{\Xi }d_\mathrm{\Xi }\mathrm{\Delta }^{r_1}=\delta _1({\displaystyle \underset{0|\mathrm{\Xi }|}{}}\mathrm{\Phi }_{r_1}^\mathrm{\Xi }c_\mathrm{\Xi }^{r_1}).`$
It follows that $`H_2(\delta _1)=0`$, i.e., the complex (30) is two-exact. $`\mathrm{}`$
If the third homology $`H_3(\delta _1)`$ of the complex (30) is not trivial, there are reduction conditions on the first-stage Noether operators, and so on. Iterating the arguments, we come to the following.
Let $``$ be a degenerate differential operator whose Noether identities are finitely generated. In accordance with Proposition 3, we associates to it the one-exact chain complex (19). Given an integer $`N1`$, let $`E_1,\mathrm{},E_N`$ be some vector bundles over $`X`$. Let us consider the ring
$`𝒫_{\mathrm{}}\{N\}=𝒫_{\mathrm{}}[E_{N1}\mathrm{}E_1E;Y;E_0\mathrm{}E_N]`$
of graded functions on $`J^{\mathrm{}}Y`$ if $`N`$ is even or the ring
$`𝒫_{\mathrm{}}\{N\}=𝒫_{\mathrm{}}[E_N\mathrm{}E_1E;Y;E_0\mathrm{}E_{N1}]`$
if $`N`$ is odd. It possesses local bases $`\{y_\mathrm{\Lambda }^i,c_\mathrm{\Lambda }^a,c_\mathrm{\Lambda }^r,c_\mathrm{\Lambda }^{r_1},\mathrm{},c_\mathrm{\Lambda }^{r_N}\}`$ where $`[c_\mathrm{\Lambda }^{r_k}]=(k+1)`$mod2 and Ant$`[c_\mathrm{\Lambda }^{r_k}]=k+2`$. Let $`k=1,0`$ further stand for $`y^i`$ and $`c^r`$, respectively. We assume that:
(i) the ring $`𝒫_{\mathrm{}}\{N\}`$ is provided with the nilpotent graded derivation
$`\delta _N=\delta _0+{\displaystyle \underset{1kN}{}}\stackrel{}{}{}_{r_k}{}^{}\mathrm{\Delta }_{}^{r_k},`$ (31)
$`\mathrm{\Delta }^{r_k}=G^{r_k}+h^{r_k}={\displaystyle \underset{0|\mathrm{\Lambda }|}{}}\mathrm{\Delta }_{r_{k1}}^{\mathrm{\Lambda }r_k}c_\mathrm{\Lambda }^{r_{k1}}+{\displaystyle \underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}}(h_{ar_{k2}}^{\mathrm{\Xi }\mathrm{\Sigma }r_k}c_\mathrm{\Xi }^ac_\mathrm{\Sigma }^{r_{k2}}+\mathrm{}),`$ (32)
of antifield number -1;
(ii) the module $`𝒫_{\mathrm{}}\{N\}_{N+3}`$ of graded functions of antifield number $`kN+3`$ is split into the $`(N+1)`$-exact chain complex
$`0\mathrm{Im}\delta \stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_1\stackrel{\delta _0}{}𝒫_{\mathrm{}}\{0\}_2\stackrel{\delta _1}{}𝒫_{\mathrm{}}\{1\}_3\mathrm{}`$ (33)
$`\stackrel{\delta _{N1}}{}𝒫_{\mathrm{}}\{N1\}_{N+1}\stackrel{\delta _N}{}𝒫_{\mathrm{}}\{N\}_{N+2}\stackrel{\delta _N}{}𝒫_{\mathrm{}}\{N\}_{N+3},`$
which satisfies the $`(N+1)`$-homology regularity condition, introduced below.
Definition 9. One says that the complex (33) obeys the $`(N+1)`$-homology regularity condition if any $`\delta _{k<N1}`$-cycle $`\mathrm{\Phi }𝒫_{\mathrm{}}\{k\}_{k+3}𝒫_{\mathrm{}}\{k+1\}_{k+3}`$ is a $`\delta _{k+1}`$-boundary.
Remark 4. The $`(N+1)`$-exactness of the complex (33) implies that any $`\delta _{k<N1}`$-cycle $`\mathrm{\Phi }𝒫_{\mathrm{}}^{0,n}\{k\}_{k+3}`$, $`k<N`$, is a $`\delta _{k+2}`$-boundary, but not necessary a $`\delta _{k+1}`$-one.
If $`N=1`$, the complex $`𝒫_{\mathrm{}}\{1\}_4`$ (33) restarts the complex (30) associated to a first-stage reducible differential operator in accordance with Proposition 4. Therefore, we agree to call $`\delta _N`$ (31) the $`N`$-stage Koszul–Tate differential. Its nilpotency implies complete Noether identities (17), first-stage Noether identities (28) and the equalities
$$\underset{0|\mathrm{\Lambda }|}{}\mathrm{\Delta }_{r_{k1}}^{\mathrm{\Lambda }r_k}d_\mathrm{\Lambda }(\underset{0|\mathrm{\Sigma }|}{}\mathrm{\Delta }_{r_{k2}}^{\mathrm{\Sigma }r_{k1}}c_\mathrm{\Sigma }^{r_{k2}})+\delta (\underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}h_{ar_{k2}}^{\mathrm{\Xi }\mathrm{\Sigma }r_k}c_\mathrm{\Xi }^ac_\mathrm{\Sigma }^{r_{k2}})=0,$$
(34)
for $`k=2,\mathrm{},N`$. One can think of the equalities (34) as being complete $`k`$-stage Noether identities because of their properties which we will justify in the case of $`k=N+1`$. Accordingly, $`\mathrm{\Delta }^{r_k}`$ (32) are said to be the $`k`$-stage Noether operators.
A generic $`(N+2)`$-chain $`\mathrm{\Phi }𝒫_{\mathrm{}}\{N\}_{N+2}`$ takes the form
$$\mathrm{\Phi }=G+H=\underset{0|\mathrm{\Lambda }|}{}G_{r_N}^\mathrm{\Lambda }c_\mathrm{\Lambda }^{r_N}+\underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}(H_{ar_{N1}}^{\mathrm{\Xi }\mathrm{\Sigma }}c_\mathrm{\Xi }^ac_\mathrm{\Sigma }^{r_{N1}}+\mathrm{}).$$
(35)
Let it be a cycle. The cycle condition $`\delta _N\mathrm{\Phi }=0`$ implies the equality
$$\underset{0|\mathrm{\Lambda }|}{}G_{r_N}^\mathrm{\Lambda }d_\mathrm{\Lambda }(\underset{0|\mathrm{\Sigma }|}{}\mathrm{\Delta }_{r_{N1}}^{\mathrm{\Sigma }r_N}c_\mathrm{\Sigma }^{r_{N1}})+\delta (\underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}H_{ar_{N1}}^{\mathrm{\Xi }\mathrm{\Sigma }}c_\mathrm{\Xi }^ac_\mathrm{\Sigma }^{r_{N1}})=0.$$
(36)
One can think of this equality as being the reduction condition on the $`N`$-stage Noether operators (32). Conversely, let
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Lambda }|}{}}G_{r_N}^\mathrm{\Lambda }c_\mathrm{\Lambda }^{r_N}𝒫_{\mathrm{}}\{N\}_{N+2}`$
be a graded function such that the reduction condition (36) holds. Then this reduction condition can be extended to a cycle one as follows. It is brought into the form
$`\delta _N({\displaystyle \underset{0|\mathrm{\Lambda }|}{}}`$ $`G_{r_N}^\mathrm{\Lambda }c_\mathrm{\Lambda }^{r_N}+{\displaystyle \underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}}H_{ar_{N1}}^{\mathrm{\Xi }\mathrm{\Sigma }}c_\mathrm{\Xi }^ac_\mathrm{\Sigma }^{r_{N1}})=`$
$`{\displaystyle \underset{0|\mathrm{\Lambda }|}{}}G_{r_N}^\mathrm{\Lambda }d_\mathrm{\Lambda }h^{r_N}+{\displaystyle \underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}}H_{ar_{N1}}^{\mathrm{\Xi }\mathrm{\Sigma }}c_\mathrm{\Xi }^ad_\mathrm{\Sigma }\mathrm{\Delta }^{r_{N1}}.`$
A glance at the expression (32) shows that the term in the right-hand side of this equality belongs to $`𝒫_{\mathrm{}}\{N2\}_{N+1}`$. It is a $`\delta _{N2}`$-cycle and, consequently, a $`\delta _{N1}`$-boundary $`\delta _{N1}\mathrm{\Psi }`$ in accordance with the $`(N+1)`$-homology regularity condition. Then the reduction condition (36) is a $`c_\mathrm{\Sigma }^{r_{N1}}`$-dependent part of the cycle condition
$`\delta _N({\displaystyle \underset{0|\mathrm{\Lambda }|}{}}G_{r_N}^\mathrm{\Lambda }c_\mathrm{\Lambda }^{r_N}+{\displaystyle \underset{0\mathrm{\Sigma },0\mathrm{\Xi }}{}}H_{ar_{N1}}^{\mathrm{\Xi }\mathrm{\Sigma }}c_\mathrm{\Xi }^ac_\mathrm{\Sigma }^{r_{N1}}\mathrm{\Psi })=0,`$
but $`\delta _N\mathrm{\Psi }`$ does not make a contribution to this reduction condition.
Being a cycle condition, the reduction condition (36) is trivial either if a cycle $`\mathrm{\Phi }`$ (35) is a $`\delta _N`$-boundary or its summand $`G`$ is $`\delta `$-exact, i.e., it is a boundary, too, as we have stated above. Then Definition 4 can be generalized as follows.
Definition 10. A degenerate differential operator is said to be $`(N+1)`$-stage reducible if there exist non-trivial reduction conditions (36), called the $`(N+1)`$-stage Noether identities.
Theorem 11. (i) The $`(N+1)`$-stage Noether identities can be identified to nontrivial elements of the homology $`H_{N+2}(\delta _N)`$ of the complex (33) iff this homology obeys the $`(N+2)`$-homology regularity condition. (ii) If the homology $`H_{N+2}(\delta _N)`$ is finitely generated as defined below, the complex (33) admits an $`(N+2)`$-exact extension.
Proof. (i) The $`(N+2)`$-homology regularity condition implies that any $`\delta _{N1}`$-cycle $`\mathrm{\Phi }𝒫_{\mathrm{}}\{N1\}_{N+2}𝒫_{\mathrm{}}\{N\}_{N+2}`$ is a $`\delta _N`$-boundary. Therefore, if $`\mathrm{\Phi }`$ (35) is a representative of a nontrivial element of $`H_{N+2}(\delta _N)`$, its summand $`G`$ linear in $`c_\mathrm{\Lambda }^{r_N}`$ does not vanish. Moreover, it is not a $`\delta `$-boundary. Indeed, if $`\mathrm{\Phi }=\delta \mathrm{\Psi }`$, then
$$\mathrm{\Phi }=\delta _N\mathrm{\Psi }+(\delta \delta _N)\mathrm{\Psi }+H.$$
(37)
The cycle condition takes the form
$`\delta _N\mathrm{\Phi }=\delta _{N1}((\delta \delta _N)\mathrm{\Psi }+H)=0.`$
Hence, $`(\delta \delta _0)\mathrm{\Psi }+H`$ is $`\delta _N`$-exact since any $`\delta _{N1}`$-cycle $`\varphi 𝒫_{\mathrm{}}\{N1\}_{N+2}`$ is a $`\delta _N`$-boundary. Consequently, $`\mathrm{\Phi }`$ (37) is a boundary. If the $`(N+2)`$-homology regularity condition does not hold, trivial reduction conditions (36) also come from nontrivial elements of the homology $`H_{N+2}(\delta _N)`$.
(ii) Let the $`(N+1)`$-stage Noether identities be finitely generated. Namely, there exist elements $`\mathrm{\Delta }_{(N+1)}H_{N+2}(\delta _N)`$ making up a projective $`C^{\mathrm{}}(X)`$-module $`𝒞_{(N+1)}`$ of finite rank which is isomorphic to the module of sections of some vector bundle $`E_{N+1}X`$. Let $`\{\mathrm{\Delta }^{r_{N+1}}\}`$ be local bases for this $`C^{\mathrm{}}(X)`$-module. Then any element $`\mathrm{\Phi }H_{N+2}(\delta _N)`$ factorizes
$`\mathrm{\Phi }={\displaystyle \underset{0|\mathrm{\Xi }|}{}}\mathrm{\Phi }_{r_{N+1}}^\mathrm{\Xi }d_\mathrm{\Xi }\mathrm{\Delta }^{r_{N+1}},\mathrm{\Phi }_{r_{N+1}}^\mathrm{\Xi }𝒪_{\mathrm{}}Y,`$ (38)
$`\mathrm{\Delta }^{r_{N+1}}=G^{r_{N+1}}+h^{r_{N+1}}={\displaystyle \underset{0|\mathrm{\Lambda }|}{}}\mathrm{\Delta }_{r_N}^{\mathrm{\Lambda }r_{N+1}}c_\mathrm{\Lambda }^{r_N}+h^{r_{N+1}},`$ (39)
through elements of $`𝒞_{(N+1)}`$. Clearly, this factorization is independent of specification of local bases $`\{\mathrm{\Delta }^{r_{N+1}}\}`$. Let us extend the ring $`𝒫_{\mathrm{}}\{N\}`$ to the ring $`𝒫_{\mathrm{}}\{N+1\}`$ possessing local bases
$`\{y_\mathrm{\Lambda }^i,c_\mathrm{\Lambda }^a,c_\mathrm{\Lambda }^r,c_\mathrm{\Lambda }^{r_1},\mathrm{},c_\mathrm{\Lambda }^{r_N},c_\mathrm{\Lambda }^{r_{N+1}}\},[c_\mathrm{\Lambda }^{r_{N+1}}]=N\mathrm{mod}\mathrm{\hspace{0.17em}2},\mathrm{Ant}[c_\mathrm{\Lambda }^{r_{N+1}}]=N+3.`$
It is provided with the nilpotent graded derivation
$$\delta _{N+1}=\delta _N+\stackrel{}{}{}_{r_{N+1}}{}^{}\mathrm{\Delta }_{}^{r_{N+1}}$$
(40)
of antifield number -1. With this graded derivation, the module $`𝒫_{\mathrm{}}\{N+1\}_{N+4}`$ of graded functions of antifield number Ant$`[f]N+4`$ is split into the chain complex
$`0\mathrm{Im}\delta \stackrel{\delta }{}𝒮_{\mathrm{}}[E;Y]_1\stackrel{\delta _0}{}𝒫_{\mathrm{}}\{0\}_2\stackrel{\delta _1}{}𝒫_{\mathrm{}}\{1\}_3\mathrm{}\stackrel{\delta _{N1}}{}𝒫_{\mathrm{}}\{N1\}_{N+1}`$ (41)
$`\stackrel{\delta _N}{}𝒫_{\mathrm{}}\{N\}_{N+2}\stackrel{\delta _{N+1}}{}𝒫_{\mathrm{}}\{N+1\}_{N+3}\stackrel{\delta _{N+1}}{}𝒫_{\mathrm{}}\{N+1\}_{N+4}.`$
It is readily observed that this complex is $`(N+2)`$-exact. In this case, the $`(N+1)`$-stage Noether identities (36) come from the complete $`(N+1)`$-stage Noether identities
$`{\displaystyle \underset{0|\mathrm{\Lambda }|}{}}\mathrm{\Delta }_{r_N}^{\mathrm{\Lambda }r_{N+1}}d_\mathrm{\Lambda }\mathrm{\Delta }^{r_N}+\delta h^{r_{N+1}}=0,`$
which are reproduced as the nilpotency conditions of the graded derivation (40). $`\mathrm{}`$
The iteration procedure based on Theorem 4 can be prolonged up to an integer $`N_{\mathrm{max}}`$ when the $`N_{\mathrm{max}}`$-stage Noether identities are irreducible, i.e., the homology $`H_{N_{\mathrm{max}}+2}(\delta _{N_{\mathrm{max}}})`$ is trivial. This iteration procedure may also be infinite. It results in the manifested exact Koszul–Tate complex with the Koszul–Tate boundary operator whose nilpotency conditions reproduce all Noether and higher Noether identities of an original differential operator $``$.
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# Generation and Propagation of Nonlinear Waves in Travelling Wave Tubes
## I Introduction
The generation and evolution of nonlinear waves and harmonic distortions in microwave amplifiers such as travelling wave tubes, free electron lasers (FELs) and klystrons have recently attracted much research interest Booske ; Freund ; Bonifacio . In connection with the construction and commissioning of the next generation of FELs and powerful klystrons for accelerating RF cavities in circular machines and linear colliders, this issue has become even more challenging. Of particular importance are the effects of intense self-fields due to space charge and current, as well as wake fields due to interaction impedances. Both of the above influence the propagation of the electron beam in microwave devices, its stability and transport properties. In general, a complete description of collective processes in intense charged particle beams is provided by the Vlasov-Maxwell equations for the self-consistent evolution of the beam distribution function and the electromagnetic fields. Usually, the electron beam in a travelling wave tube can be assumed to be weakly collisional. Hence, the dynamics of electrons is well described by the hydrodynamic equations coupled with the equations for the electromagnetic self-fields, which constitutes a substantial simplification of the model. Although the analytical basis for modelling the dynamics and behaviour of space-charge-dominated beams is well established, a thorough and satisfactory understanding of collective processes, detailed equilibrium and formation of patterns and coherent structures is far from being complete.
While the linear theory of wave generation in microwave amplifiers is relatively well understood Dobson , the nonlinear regime is far from being exhaustively studied. The present paper is aimed at filling this gap. We will be mainly interested in describing the slow evolution of some coarse-grained quantities that are easily measurable, such as the wave amplitudes. Owing to the nonlinear wave interaction contingent on the nonlinear coupling between the hydrodynamic and Maxwell equations, one can expect a formation of nontrivial coherent structure that might be fairly stable in space and time Tzenov ; Tzenov1 . Here, we show that solitary wave patterns in the electron beam density distribution are an irrevocable feature, characteristic of powerful microwave amplifiers.
The paper is organized as follows. In the next section, we state the basic equations which will be the subject of the renormalization group (RG) reduction in section III. Starting from a single equation \[see equation (20)\] for the density distribution of the electron beam, we obtain a formal perturbation expansion of its solution to second order. As expected, it contains secular terms proportional to powers of the time variable which is the only renormalization parameter adopted in our approach. In section IV, the arbitrary constant amplitudes of the perturbation expansion are renormalized such as to eliminate the secular terms. As a result, a set of equations for the renormalized slowly varying amplitudes is obtained, known as the renormalization group equations (RGEs). These equations comprise an infinite system of coupled nonlinear Schrödinger equations. Finally, section V is dedicated to discussion and conclusions.
## II Formulation of the Problem and Basic Equations
The electron beam in a travelling wave tube is assumed to be weakly collisional. Therefore, the dynamics of electrons is well described by the hydrodynamic equations coupled with the equations for the electromagnetic self-fields. We start with the 1D description of a beam of electrons propagating in an external focusing electric field with focusing coefficient $`G`$, which models the bunching of the electron beam in the longitudinal direction. As we will see in the sequel its additional role is to attain a stationary equilibrium by partially compensating the space-charge defocusing. The continuity and the momentum balance equations can be written as
$$\frac{\varrho }{t}+\frac{}{z}\left(\varrho v\right)=0,$$
(1)
$$\frac{v}{t}+v\frac{v}{z}=Gz\frac{k_BT}{m\varrho }\frac{\varrho }{z}\frac{e}{m}\left(\frac{\phi _{sc}}{z}+\frac{\phi _w}{z}\right),$$
(2)
where $`\varrho `$ and $`v`$ are the electron density and the current velocity. Furthermore, $`m`$, $`e`$ and $`T`$ are the mass, the charge and the temperature, respectively, while $`k_B`$ is the Boltzmann constant. The space-charge $`\phi _{sc}`$ and wave $`\phi _w`$ potentials satisfy the Poisson and the wave equation
$$\frac{^2\phi _{sc}}{z^2}=\frac{en\varrho }{\epsilon _0},$$
(3)
$$\mathrm{}\phi _w=\frac{enAZ_0}{c}\frac{^2\varrho }{t^2},\mathrm{}=\frac{^2}{z^2}\frac{1}{c^2}\frac{^2}{t^2},$$
(4)
where $`\mathrm{}`$ denotes the well-known d’Alembert operator. In addition, $`\epsilon _0`$ is the permittivity of free space and $`n=N_e/V_t`$ is the electron number density, where $`N_e`$ is the total number of electrons and $`V_t`$ is the volume occupied by the electron beam in the longitudinal direction. Moreover, the electron beam cross-sectional area is represented by $`A`$, the quantity $`Z_0`$ denotes the interaction impedance, while $`c`$ represents the phase velocity of a cold circuit wave.
Let us introduce the scaling
$$\phi _{sc}=\frac{en}{\epsilon _0}U_{sc},\phi _w=\frac{en}{\epsilon _0}W,$$
(5)
and rewrite equations (2)–(4) as follows
$$\frac{v}{t}+v\frac{v}{z}=Gz\frac{v_T^2}{\varrho }\frac{\varrho }{z}\omega _p^2\left(\frac{U_{sc}}{z}+\frac{W}{z}\right),$$
(6)
$$\frac{^2U_{sc}}{z^2}=\varrho ,\mathrm{}W=𝒵\frac{^2\varrho }{t^2}.$$
(7)
The electron plasma frequency $`\omega _p`$ and the thermal velocity $`v_T`$ of the electron beam are expressed according to
$$\omega _p^2=\frac{e^2n}{\epsilon _0m},v_T^2=\frac{k_BT}{m},$$
(8)
while
$$𝒵=\frac{\epsilon _0Z_0A}{c},$$
(9)
is a shorthand parameter introduced for later convenience. Note that the thermal velocity $`v_T`$ as defined by equation (8) can be alternatively expressed according to the relation
$$v_T=\omega _pr_D,r_D^2=\frac{ϵ_0k_BT}{e^2n_0},$$
(10)
where $`r_D`$ is the electron Debye radius. Equations (1), (6) and (7) possess a stationary equilibrium solution
$$\varrho _0=\frac{G}{\omega _p^2},v_0=0,U_0=\frac{Gz^2}{2\omega _p^2},W_0=0.$$
(11)
Therefore, we can further scale the hydrodynamic and field variables as
$$\varrho =\varrho _0+ϵR,U_{sc}=U_0+ϵU,vϵv,WϵW,$$
(12)
where $`ϵ`$ is a formal small parameter introduced for convenience, which will be set equal to one at the end of the calculations. Thus, the basic equations to be used for the subsequent analysis can be written in the form
$$\frac{R}{t}+\varrho _0\frac{v}{z}+ϵ\frac{}{z}\left(Rv\right)=0,$$
(13)
$$\frac{v}{t}+ϵv\frac{v}{z}=\frac{v_T^2}{\varrho _0+ϵR}\frac{R}{z}\omega _p^2\left(\frac{U}{z}+\frac{W}{z}\right),$$
(14)
$$\frac{^2U}{z^2}=R,\mathrm{}W=𝒵\frac{^2R}{t^2}.$$
(15)
Before we continue with the renormalization group reduction of the system of equations (13)–(15) in the next section, let us assume that the actual dependence of the quantities $`R`$, $`v`$, $`U`$ and $`W`$ on the spatial variables is represented by the expression
$$\widehat{\mathrm{\Psi }}=\widehat{\mathrm{\Psi }}(z,\xi ;t),\widehat{\mathrm{\Psi }}=(R,v,U,W),$$
(16)
where $`\xi =ϵz`$ is a slow spatial variable. Thus, the only renormalization parameter left at our disposal is the time $`t`$ which will prove extremely convenient and simplify tedious algebra in the sequel.
## III Renormalization Group Reduction of the Hydrodynamic Equations
Following the standard procedure of the renormalization group method, we represent $`\widehat{\mathrm{\Psi }}`$ as a perturbation expansion
$$\widehat{\mathrm{\Psi }}=\underset{n=0}{\overset{\mathrm{}}{}}ϵ^n\widehat{\mathrm{\Psi }}_n,$$
(17)
in the formal small parameter $`ϵ`$. The next step consists in expanding the system of hydrodynamic and field equations (13)-(15) in the small parameter $`ϵ`$, and obtaining their naive perturbation solution order by order. It is possible to simplify this system, which will turn out extremely useful in what follows. Differentiating equation (13) with respect to the time $`t`$, differentiating equation (14) with respect to $`z`$ and using equations (15), we can eliminate the electric potentials. As a result of obvious manipulations, we obtain
$$\widehat{}R=ϵ\mathrm{}\left[\frac{\varrho _0}{2}\frac{^2}{z^2}\left(v^2\right)\frac{^2}{tz}(Rv)\right]$$
$$+v_T^2\mathrm{}\frac{^2}{z^2}\left[\frac{\varrho _0}{ϵ}\mathrm{ln}\left(1+\frac{ϵR}{\varrho _0}\right)R\right],$$
(18)
where the linear differential operator $`\widehat{}(_z,_t)`$ is given by the expression
$$\widehat{}(_z,_t)=\mathrm{}\left(\frac{^2}{t^2}v_T^2\frac{^2}{z^2}+G\right)+G𝒵\frac{^4}{t^2z^2}.$$
(19)
Taking into account the expansion of the logarithm, we notice that the right-hand-side of equation (18) is at least of first order in the formal parameter, so that
$$\widehat{}R=ϵ\mathrm{}\left[\frac{\varrho _0}{2}\frac{^2}{z^2}\left(v^2\right)\frac{^2}{tz}(Rv)\right]$$
$$ϵv_T^2\mathrm{}\frac{^2}{z^2}\left(\frac{R^2}{2\varrho _0}\frac{ϵR^3}{3\varrho _0^2}+\frac{ϵ^2R^4}{4\varrho _0^3}\mathrm{}\right),$$
(20)
Equation (20) represents the starting point for the renormalization group reduction, the final goal of which is to obtain a description of the relatively slow dynamics leading to formation of patterns and coherent structures.
Let us proceed order by order. The solution to the zero-order perturbation equation (20) can be written as
$$R_0(z,\xi ;t)=\underset{k}{}A_k\left(\xi \right)\mathrm{e}^{i\psi _k(z;t)},$$
(21)
where
$$\psi _k(z;t)=kz\omega _kt,$$
(22)
and $`A_k`$ is an infinite set of constant complex amplitudes, which will be the subject of the renormalization procedure in the sequel. Here ”constant” means that the amplitudes $`A_k`$ do not depend on the fast spatial variable $`z`$ and on the time $`t`$, however, they can depend on the slow spatial variables $`\xi `$. The summation sign in equation (21) and throughout the paper implies summation over the wave number $`k`$ in the case where it takes discrete values, or integration in the continuous case. From the dispersion equation
$$𝒟(k;\omega _k)=G𝒵k^2\omega _k^2\mathrm{}_k\left(\omega _k^2k^2v_T^2G\right)=0,$$
(23)
it follows that the wave frequency $`\omega _k`$ can be expressed in terms of the wave number $`k`$, where the Fourier-image $`\mathrm{}_k`$ of the d’Alembert operator can be written according to
$$\mathrm{}_k=\frac{\omega _k^2}{c^2}k^2.$$
(24)
It is important to emphasize that
$$\omega _k=\omega _k,A_k=A_k^{},$$
(25)
where the asterisk denotes complex conjugation. The latter assures that the perturbed density distribution as defined by equation (21) is a real quantity. The zero-order current velocity $`v_0(z,\xi ;t)`$ obtained directly from equation (13) can be written as
$$v_0(z,\xi ;t)=\underset{k}{}v_k^{(0)}A_k\left(\xi \right)\mathrm{e}^{i\psi _k(z;t)},v_k^{(0)}=\frac{\omega _k}{\varrho _0k}.$$
(26)
In first order equation (20) acquires the form
$$\widehat{}R_1+\widehat{}_z\frac{R_0}{\xi }=\mathrm{}\left[\frac{\varrho _0}{2}\frac{^2v_0^2}{z^2}\frac{^2}{tz}(R_0v_0)\frac{v_T^2}{2\varrho _0}\frac{^2R_0^2}{z^2}\right],$$
(27)
where by $`\widehat{}_z`$ we have denoted the derivative of the operator $`\widehat{}`$ with respect to $`_z`$. It has now two types of solutions. The first is a secular solution linearly dependent on the time variable in the first-order approximation. As a rule, the highest power in the renormalization parameter of the secular terms contained in the standard perturbation expansion is equal to the corresponding order in the small perturbation parameter. The second solution of equation (20) arising from the nonlinear interaction between waves in the first order, is regular. Taking into account the fact that the Fourier image of the operator $`\widehat{}_z`$ is equal to $`i(𝒟/k)`$, we can write the equation for determining of the secular part of the solution as
$$\widehat{}R_1^{(s)}=i\underset{k}{}\frac{𝒟}{k}\frac{A_k}{\xi }\mathrm{e}^{i\psi _k}.$$
(28)
We note further that
$$\widehat{}t\mathrm{e}^{i\psi _k}=t𝒟\mathrm{e}^{i\psi _k}+i\frac{𝒟}{\omega _k}\mathrm{e}^{i\psi _k}=i\frac{𝒟}{\omega _k}\mathrm{e}^{i\psi _k},$$
(29)
which is a direct consequence of the general relation
$$\widehat{}tG(t)=t\widehat{}G(t)+\widehat{}_tG(t),$$
(30)
holding for a generic function $`G(t)`$. Here $`\widehat{}_t`$ implies differentiation with respect to $`_t`$. To verify equation (30), it suffices to prove the identity by induction in the case, where $`\widehat{}`$ is the monomial operator $`_t^n`$ and then take into account the Taylor expansion of $`\widehat{}`$. With these remarks in hand, it is straightforward to solve equation (28). Combining its solution with the solution of the regular part, we obtain
$$R_1=t\underset{k}{}u_{gk}\frac{A_k}{\xi }\mathrm{e}^{i\psi _k}\frac{1}{2\varrho _0}\underset{k,l}{}\alpha _{kl}A_kA_l\mathrm{e}^{i\left(\psi _k+\psi _l\right)},$$
(31)
where $`u_{gk}`$ is the group velocity defined as
$$u_{gk}=\frac{\mathrm{d}\omega _k}{\mathrm{d}k}=\frac{𝒟}{k}\left(\frac{𝒟}{\omega _k}\right)^1.$$
(32)
In explicit form, the components of the infinite matrix $`\alpha _{kl}`$ are given by the expression
$$\alpha _{kl}=\frac{\gamma _{kl}}{𝒟_{kl}},$$
(33)
where
$$\gamma _{kl}=\mathrm{}_{kl}(k+l)$$
$$\times \left[\left(\omega _k+\omega _l\right)\left(\frac{\omega _k}{k}+\frac{\omega _l}{l}\right)(k+l)\left(v_T^2\frac{\omega _k\omega _l}{kl}\right)\right],$$
(34)
$$\mathrm{}_{kl}=\frac{\left(\omega _k+\omega _l\right)^2}{c^2}(k+l)^2,$$
(35)
$$𝒟_{kl}=𝒟(k+l,\omega _k+\omega _l).$$
(36)
Furthermore, the first-order current velocity can be expressed as
$$v_1=\underset{k}{}v_k^{(1)}\frac{A_k}{\xi }\mathrm{e}^{i\psi _k}\frac{1}{2\varrho _0^2}\underset{k,l}{}\beta _{kl}A_kA_l\mathrm{e}^{i\left(\psi _k+\psi _l\right)},$$
(37)
where
$$v_k^{(1)}=u_{gk}v_k^{(0)}t\frac{i}{\varrho _0k}\left(u_{gk}\varrho _0v_k^{(0)}\right),$$
(38)
$$\beta _{kl}=\frac{\omega _k}{k}+\frac{\omega _l}{l}+\alpha _{kl}\frac{\omega _k+\omega _l}{k+l},\beta _{k,k}=0.$$
(39)
A couple of interesting features of the zero and first-order perturbation solution are noteworthy to be commented at this point. First of all, the zero-order density, current velocity (and electric potentials) are proportional to the arbitrary complex amplitudes $`A_k`$. The second terms in the expressions for the first-order density $`R_1`$ and current velocity $`v_1`$ \[see equations (31) and (37)\] imply contribution from nonlinear interaction between waves. It will be shown in the remainder that these terms give rise to nonlinear terms in the renormalization group equation and describe solitary wave behaviour of a generic mode.
## IV The Renormalization Group Equation
Passing over to the final stage of our renormalization group procedure, we note that particular terms in the second-order perturbation equation (20) will contribute to the secular solution. Since we are interested in this solution to be renormalized later, the contributing terms will be retained only. Hence, we can write
$$\widehat{}R_2+\widehat{}_z\frac{R_1}{\xi }+\frac{\widehat{}_{zz}}{2}\frac{^2R_0}{\xi ^2}=\mathrm{}[\varrho _0\frac{^2\left(v_0v_1\right)}{z^2}$$
$$\frac{^2}{tz}(R_0v_1+R_1v_0)\frac{v_T^2}{\varrho _0}\frac{^2\left(R_0R_1\right)}{z^2}+\frac{v_T^2}{3\varrho _0^2}\frac{^2R_0^3}{z^2}],$$
(40)
or in explicit form
$$\widehat{}R_2=\underset{k}{}\left[u_{gk}\left(it\frac{𝒟}{k}+\frac{^2𝒟}{k\omega _k}\right)+\frac{1}{2}\frac{^2𝒟}{k^2}\right]\frac{^2A_k}{\xi ^2}\mathrm{e}^{i\psi _k}$$
$$+\frac{1}{\varrho _0^2}\underset{k,m}{}\mathrm{\Gamma }_{mk}\left|A_m\right|^2A_k\mathrm{e}^{i\psi _k}.$$
(41)
The components of the infinite matrix $`\mathrm{\Gamma }_{mk}`$ can be expressed according to the relation
$$\mathrm{\Gamma }_{mk}=\mathrm{}_kk^2[\beta _{mk}(\frac{\omega _m}{m}+\frac{\omega _k}{k})$$
$$\alpha _{mk}(v_T^2\frac{\omega _m\omega _k}{mk})v_T^2].$$
(42)
A standard algebra similar to the one outlined in the previous section leads to the second order secular solution
$$R_2=\underset{k}{}\left(\frac{u_{gk}^2t^2}{2}+\frac{G_kt}{2i}\right)\frac{^2A_k}{\xi ^2}\mathrm{e}^{i\psi _k}$$
$$+\frac{t}{i\varrho _0^2}\underset{k,m}{}\mathrm{\Gamma }_{mk}\left(\frac{𝒟}{\omega _k}\right)^1\left|A_m\right|^2A_k\mathrm{e}^{i\psi _k},$$
(43)
where
$$G_k=\left(u_{gk}^2\frac{^2𝒟}{\omega _k^2}+2u_{gk}\frac{^2𝒟}{k\omega _k}+\frac{^2𝒟}{k^2}\right)\left(\frac{𝒟}{\omega _k}\right)^1.$$
(44)
Taking into account the definition (32) of the group velocity, we conclude that
$$\frac{\mathrm{d}u_{gk}}{\mathrm{d}k}=\frac{u_{gk}}{k}+u_{gk}\frac{u_{gk}}{\omega _k}=G_k.$$
(45)
Following the standard procedure Tzenov ; Tzenov1 ; Oono of the RG method, we finally obtain the desired RG equation
$$i\frac{\stackrel{~}{A}_k}{t}+iu_{gk}\frac{\stackrel{~}{A}_k}{z}$$
$$=\frac{G_k}{2}\frac{^2\stackrel{~}{A}_k}{z^2}+\frac{1}{\varrho _0^2}\left(\frac{𝒟}{\omega _k}\right)^1\underset{m}{}\mathrm{\Gamma }_{mk}\left|\stackrel{~}{A}_m\right|^2\stackrel{~}{A}_k,$$
(46)
where now $`\stackrel{~}{A}_k`$ is the renormalized complex amplitude Tzenov . Thus, the renormalized solution for the density perturbation of the electron beam acquires the form
$$R(z;t)=\underset{k}{}\stackrel{~}{A}_k(z;t)\mathrm{e}^{i\psi _k(z;t)},$$
(47)
For the renormalized electric field, we obtain
$$E(z;t)=\frac{ien}{ϵ_0G}\underset{k}{}\frac{k^2v_T^2\omega _k^2}{k}\stackrel{~}{A}_k(z;t)\mathrm{e}^{i\psi _k(z;t)},$$
(48)
Equations (46) represent a system of coupled nonlinear Schrödinger equations for the amplitudes of eigenmodes. Consider a particular mode with a wave number $`k`$. As a first approximation, the contribution of the other modes with $`mk`$ can be neglected, which results in a single nonlinear Schrödinger equation for mode $`k`$. The nonlinearity in the corresponding nonlinear Schrödinger equation describes the nonlinear interaction of the mode $`k`$ with itself.
It should be emphasized that the approach outlined in the present paper is rather general even in the the case of more than one renormalization parameter, where the extension should be straightforward. Note that the unique assumption concerns the knowledge of the dispersion properties of the linear differential operator governing the evolution of the zero-order quantities. The approach can be successfully applied to a wide class of problems of particular physical interest.
## V Discussion and conclusions
We have studied the generation and evolution of nonlinear waves in microwave amplifiers such as travelling wave tubes, free electron lasers and klystrons. The analysis performed in the present paper is based on the hydrodynamic and field equations for the self-consistent evolution of the beam density distribution, the current velocity and the electromagnetic fields. Using further the RG method, a system of coupled nonlinear Schrödinger equations for the slowly varying amplitudes of interacting beam-density waves has been derived. Under the approximation of an isolated mode neglecting the effect of the rest of the modes, this system reduces to a single nonlinear Schrödinger equation for that particular mode.
Since the approach pursued here is rather general, it is presumed that it may find applications to other problems, where of particular interest is the dynamics of slowly varying amplitudes of patterns and coherent structures.
###### Acknowledgements.
It is a pleasure to thank R.C. Davidson for many interesting and useful discussions concerning the subject of the present paper.
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# New distances of unresolved dwarf elliptical galaxies in the vicinity of the Local Group Based on observations collected at the Nordic Optical Telescope
## 1 Introduction
Recent imaging surveys of the Local Group neighbourhood have found and identified a large number of low surface brightness galaxy candidates, which potentially could be nearby dwarf galaxy systems (e.g. Côté et al. 1997; Karachentseva & Karachentsev 1998; Jerjen et al. 2000a). Many are dwarf elliptical galaxies (dEs) located in galaxy groups as satellites of giant galaxies. Whether dEs are local or in the background is an important question to be answered. For example Moore et al. (1999) have numerically studied galactic and cluster halo substructures in a hierarchical universe, and found that simulations predict that the Milky Way galaxy would have many more dwarf satellites than are actually observed. This conflict is a strong motivation for catalogueing and obtaining accurate distances to dwarf galaxy population in the local universe. As a new and numerous target group in extragalactic studies dEs serve also to test the peculiar linearity and smoothness of the Hubble flow in the local volume. As such they may help to define, or at least refine, the dark energy solution to the problem (Baryshev et al. 2001).
Determining distances to galaxies of this type has been a challenge. They have only very little or no neutral hydrogen gas, preventing their detection in the radio at 21 cm and their low surface brightness makes optical spectroscopy feasible only for the few brightest objects (Jerjen et al. 2000b). Instead, distances must be estimated from their stellar content. Taking advantage of the absence of the atmosphere, the Hubble Space Telescope (HST) is used to resolve dEs into stars and to measure the tip of the red giant branch (TRGB) magnitude (Karachentsev et al. 2000). However, this method is expensive in terms of integration time and becomes progressively difficult beyond a few Mpc due to crowding effects.
The Surface Brightness Fluctuation method can be applied to unresolved galaxies and thus offers an alternative to efficiently measure distances of dEs out to 10 Mpc and beyond (e.g. Jerjen 2003). The method was originally introduced by Tonry & Schneider (1988) to measure distances to high surface brightness elliptical galaxies. For low surface brightness dEs the method was developed (Jerjen et al. 1998; Jerjen et al. 2000b) and calibrated (Jerjen et al. 2001) only recently. It is based on the discrete sampling of the unresolved stellar population of a galaxy with the CCD detector and the resulting pixel-to-pixel variance due to the statistical noise in numbers of red giant branch (RGB) stars.
We can report here on new SBF distances for 10 nearby galaxies as part of our continuing project to map the galaxy groups and clouds beyond the Local Group out to $``$10 Mpc. We have studied nine dEs and one S0 type of galaxy in the northern hemisphere. In Table 1 we give a complete list of our galaxy sample including galaxy name, associated environment, morphological type within the extended Hubble classification system (Sandage & Binggeli 1984), and coordinates. UGC 1703 is a field galaxy in the direction of NGC 784, KDG 61 and UGC 5442 are members of the M81 group, UGCA 200 has been assumed to be a companion of NGC 3115 but we cannot confirm this assumption, UGC 5944 is a member of the Leo I group, NGC 4150 (the S0 galaxy) lies in the direction of the Canes Venatici I cloud (CVn I) – probably behind it, BTS 128 is a member of the Coma I group, UGC 7639 is a member of the Canes Venatici II cloud (CVn II), UGC 8882 is a member of the M101 group, and UGC 8799 may lie at the outskirts of the Virgo I cluster.
In Sect. 2, we describe the observations and data reduction. The SBF analysis is presented in Sect. 3. Individual galaxies are discussed in Sect. 4. Finally, we present the summary and draw the conclusions of this work in Sect. 5.
## 2 Observations and Reductions
CCD images were obtained at the 2.56 metre Nordic Optical Telescope on the nights of the 20 January 2002, 3-4 February 2002 and 25-27 February 2003. We used the Andalucia Faint Object Spectrograph and Camera (ALFOSC), which is equipped with a 2048 $`\times `$ 2048 Loral/Lesser CCD detector with a pixel size of 15 $`\mu `$m and a plate scale of 0.188 arcsec, providing a field of view 6.4 arcmin on a side. The gain was set at 1 $`e^{}/`$ADU. A series of six to ten images were taken in $`B`$ and $`R`$ passbands for each of the ten dwarf galaxies, along with bias frames, twilight flats and photometric standard star fields through the nights. The observing log is given in Table 2. The exposure time for individual science frames was 600 seconds, with the exception of NGC 4150 for which we used 180 seconds in $`R`$ to avoid saturation of the central region of the galaxy. The seeing ranged from 0.8 to 1.6 arcsec and all six nights provided photometric conditions with the exception of the night of 4 February 2002, which was partly spectroscopic.
Image reduction was accomplished using routines within the IRAF<sup>1</sup><sup>1</sup>1IRAF, Image Reduction and Analysis Facility, is distributed by the National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc., under contract with the National Science Foundation programme. We removed the bias level from the images by using the bias frames and the overscan region of each image. Images were subsequently trimmed to 2000 $`\times `$ 2000 pixels to remove non-essential data from the borders. Finally, each object image was divided by the corresponding median combined masterflat. Photometric calibration was achieved using the Landolt (1992) standard star fields, which were regularly observed during each night. Thus we determined the photometric zero point (ZP), atmospheric extinction coefficient ($`k`$) and colour term ($`c`$) for each passband and night. Analysis revealed slight variation in extinction coefficients throughout the observation period. The mean $`k`$ value was calculated for each passband and the corresponding values of ZP and $`c`$ were re-evaluated under this constraint. The results are summarised in Table 3. Images taken during the nights of 20 Jan and 4 Feb 2002 (UGC1703, UGC5944) were calibrated with shallow images obtained on 26 Feb 2003.
$`B`$ and $`R`$ images of each galaxy were registered by matching the positions of typically 50 reference stars spread evenly over the image. The alignment was done on a pixel scale in order to avoid dividing galaxy flux in subpixel shifts. The resulting slight degradation of image quality is insignificant in relation to the seeing effects in the images. The sky background level was estimated by fitting a plane to selected star-free areas distributed uniformly over the CCD area but well away from the galaxy. The sky-subtracted images taken in the same passband were cleaned of cosmic rays and median-combined to increase the signal-to-noise. Finally, the resulting master images were flux calibrated.
## 3 R-band SBF Analysis of Selected Early-Type Dwarf Galaxies in the 10 Mpc Range
We applied the Surface Brightness Fluctuation (SBF) method, developed to measure distances to dwarf elliptical galaxies by Jerjen et al. (1998, 2000b) and calibrated in Jerjen et al. (2001). The SBF analysis is done by carefully cleaning the galaxy images from any foreground stars, globular clusters, and background galaxies using procedures that follow the recipes of Jerjen et al. (2000b; 2001). The cleaned galaxy image was then modelled using an isophote fitting routine written in IRAF that allows the centre, ellipticity, and position angle to vary. The best 2D-model was subtracted from the original master image and the residual image divided by the square root of the model for noise normalization.
Any non-radial irregularities in the light distribution of a galaxy such as the detected spiral arms and dust features in NGC 4150 remain visible in the residual image (see Fig. 2). These parts of a galaxy were avoided in the SBF field selection. Only regions where the model follows closely the galaxy light distribution were used.
The largest possible number of slightly overlapping square subimages (hereafter SBF fields) were then defined for each galaxy. The size of the SBF fields was chosen between $`70\times 70`$ and $`100\times 100`$ pixels depending on the apparent diameter of the galaxy. Parts of a SBF field that were contaminated by foreground stars or background galaxies were replaced with randomly selected patches from the fluctuation image, lying outside the field and in the same surface brightness range of the galaxy. The number of patched pixels was 3% or less of the total SBF field area in all cases. In total we defined 61 SBF fields in our 10 sample galaxies. Their positions across the galaxy images are shown in Fig. 1.
All SBF fields were Fourier transformed and the azimuthally averaged power spectra calculated. From isolated bright stars on the galaxy master image we determined the point spread function (PSF) profile. We then fitted a linear combination of the flux normalized and exposure time weighted PSF power spectrum and a constant at the observed galaxy power spectrum $`\text{PS}(k)=P_0\text{PS}_{\text{star}}(k)+P_1`$, demanding a least squares minimization. Data points at low spatial frequencies ($`k5`$) were omitted as they are likely affected by imperfect galaxy model subtraction. Figs. 3 to 9 and 11 to 13 show the power spectrum of each SBF field with the best fitting analytic function indicated as solid lines. Table 4 summarizes the quantities measured in the SBF analysis: Col. 1 – SBF field number and galaxy name, Col. 2 – pixel size of the SBF field, Col. 3 – magnitude $`m_1`$ of a star yielding 1 ADU per second on the CCD, Col. 4 – mean galaxy surface brightness within the SBF field in ADU, Col. 5 – sky brightness in ADU, Col. 6 – exposure time normalized amplitude $`P_0`$ of the best least squares fit at wavenumber $`k=0`$ with fitting error in brackets, Col. 7 – the scale-free white noise component $`P_1`$ in the power spectrum, indicating the ratio of sky to mean galaxy surface brightness within the SBF field.
To estimate the fraction in $`P_0`$ (Col. 6) from unresolved distant background galaxies fainter than the cutoff magnitude $`m_c=24.7`$$`R`$ mag, we made use of a formula that was given in Jensen et al. (1998) and adjusted for the $`R`$-band by Jerjen et al. (2001):
$$P_{\mathrm{BG}}=\frac{p^2}{(0.8\gamma )\mathrm{ln}10}10^{0.8(m_1m_c)\gamma (29.38+RKm_c)},$$
where $`p`$ is the CCD pixel size in arcsec and $`\gamma =0.3`$ the slope of the power-law number distribution for background galaxies in the $`K`$-band (Cowie et al. 1994). Assuming a typical galaxy colour of $`(RK)=2.25`$ (de Jong 1996) we computed $`P_{\mathrm{BG}}`$ and determined the signal-to-noise S/N$`=(P_0P_{\mathrm{BG}})/(P_1+P_{\mathrm{BG}})`$ as well as the relative contribution to the signal $`P_{\mathrm{BG}}/P_0`$ for each individual SBF field. Both numbers are listed in Col. 8 and Col. 9 of Table 4. The contribution from unresolved background galaxies was minimal at the 0–6 per cent level in our SBF fields.
Another potentially significant source of unwanted fluctuations is a rich globular cluster (GC) system in a target galaxy. While this is an important issue for luminous giant ellipticals the expected number of GCs in our dwarf ellipticals is quite small. For instance, the net number of globular cluster candidates for UGC 5944 in the Leo Group is $`3.5\pm 3.9`$ (Miller et al. 1998). The GC frequency ($`S_N`$)–luminosity relation for dE,Ns studied in the Fornax and Virgo clusters (Miller et al. 1998) predicts $`18`$ GCs (assuming $`S_N=2`$) for our brightest dwarf NGC 4150 ($`M_V17.4`$) and $`0`$ GCs (assuming $`S_N=10`$) for the faintest dwarf UGC 5442 ($`M_V10`$). All GCs would be brighter than our cutoff luminosity and thus be excised during the image cleaning process. Therefore, no further corrections were applied to the measured SBF power.
Finally, we calculated the stellar fluctuation magnitude $`\overline{m}_R`$ with the formula $`\overline{m}_R=m_12.5\mathrm{log}(P_0P_{\mathrm{BG}})`$ and measured the $`(BR)`$ colour for each SBF field from the cleaned $`B`$ and $`R`$ galaxy master images. Both quantities were corrected for foreground extinction using the IRAS/DIRBE maps of dust IR emission (Schlegel et al. 1998). The results are listed in Cols. 3 and 4 of Table 5.
The power spectrum fitting error is between 3 and 15%. Other sources of minor errors are the PSF normalization ($``$2%), the shape variation of the stellar PSF over the CCD area (1–2%) and the uncertainty in the photometric calibration ($`0.04`$ mag in $`B`$, 0.03 mag in $`R`$). If we further adopt a 16% error for the foreground extinction (Schlegel et al. 1998), the formal combined error for a single $`\overline{m}_R^0`$ measurement is between 0.05 and 0.20 mag (Col. 3). The error associated with the local colour (Col. 4) has been obtained through the usual error propagation formula from the uncertainties in the sky level determination, the photometry zero points, and Galactic extinction.
## 4 Discussion
### UGC 1703
Neither a velocity nor a distance were known for UGC 1703 to date. This dwarf galaxy is closest in projection (3 degrees) to the spiral galaxy NGC 925 for which a Cepheid distance of $`9.3\pm 0.7`$ Mpc (Silbermann et al. 1996) and a velocity of $`v_{}=553`$ km s<sup>-1</sup> were reported. However, our SBF distance for UGC 1703 of $`4.2\pm 0.3`$ Mpc indicates a much shorter distance from the Milky Way and thus UGC 1703 seems to be only close to NGC 925 in projection. The next nearest major galaxy to UGC 1703 is the SBdm spiral NGC 784 with an angular separation of 4.7 degrees. NGC 784 has a velocity of $`v_{}=198`$ km s<sup>-1</sup> and a reported distance of 5 Mpc (Drozdovsky & Karachentsev 2000). These two results confirm that UGC 1703 is spatially close to NGC 784 with a projected linear distance of only $``$0.4 Mpc at an adopted distance of 4.5 Mpc. Another more qualitative evidence for the short distance of UGC 1703 comes from our deep $`R`$-band image (Fig.1) which shows semi-resolved stars in appearance similar to images of dwarf ellipticals in the Cen A Group observed with a 2.3m ground-based telescope (Jerjen et al. 2000a).
### KDG 61 (\[KK98\] 81)
KDG 61 (\[KK98\] 81) is a member of the M81 group. Karachentsev et al. (2000) reported a distance modulus of $`(mM)=27.78\pm 0.15`$ mag from the measurement of the magnitude of the red giant branch tip and the galaxy has a heliocentric velocity of $`v_{}=135\pm 30`$ km s<sup>-1</sup> (Johnson et al. 1997). We analysed six independent SBF fields across the galaxy’s surface (Fig. 4, left panels). The derived mean distance modulus of $`(mM)_{\mathrm{SBF}}=27.80\pm 0.20`$ mag (see Fig. 4, right panel) is in good agreement with the TRGB result. Due to the small colour range covered by the SBF fields (see Table 5), another SBF distance is technically possible by moving the data points onto the linear branch of the calibration curve. However, that alternative distance modulus of $`(mM)_{\mathrm{SBF}}=30.0\pm 0.3`$ mag is highly inconsistent with the independent TRGB result and thus can be ruled out.
### UGCA 200
Neither a velocity nor a distance were known for UGCA 200 to date. This dE,N was previously photometrically studied by Parodi et al. (2002). The authors reported an outward colour gradient getting redder and an integrated colour of $`(BR)_0=1.38`$. Our five SBF fields were selected from the inner region of the dwarf and thus are slightly bluer, in the range $`1.22<(BR)_0<1.31`$. Nevertheless, the SBF fields of UGCA 200 are the reddest in our sample. The observed correlation between the derived parameters $`\overline{m}_R`$ and $`(BR)`$ colour allowed an unambiguous measurement of the distance. A shift by 29.01 mag yields the best fit to the linear branch of the calibration diagram. UGCA 200 was assumed to be a faint companion of the S0 galaxy NGC 3115 which has a velocity of $`v_{}=720`$ km s<sup>-1</sup> and an accurate SBF distance of $`9.7\pm 0.1`$ Mpc (Tonry et al. 2001). However, our SBF distance for UGCA 200 is significantly shorter with $`6.3\pm 0.8`$ Mpc and thus suggests that this dwarf is actually in the foreground of NGC 3115.
### UGC 5442 (KDG 64)
UGC 5442 (KDG 64) is a member of the M81 group. Karachentsev et al. (2000) reported a distance modulus of $`(mM)=27.84\pm 0.15`$ mag from the measurement of the magnitude of the red giant branch tip. The galaxy has a heliocentric velocity of $`v_{}=18\pm 14`$ km s<sup>-1</sup> (Simien & Prugniel 2002). We analysed six independent SBF fields across the galaxy’s surface (Fig. 6, left panels). The derived mean distance modulus of $`(mM)_{\mathrm{SBF}}=27.74\pm 0.20`$ mag (see Fig. 6, right panel) is in good agreement with the TRGB result. Due to the small colour range covered by the SBF fields (see Table 5), another SBF distance is technically possible by moving the data points onto the linear branch of the calibration curve. However, the alternative distance modulus of $`(mM)_{\mathrm{SBF}}=28.21\pm 0.3`$ mag is inconsistent with the TRGB result.
### UGC 5944 (\[FS90\] 047)
Neither a velocity nor a distance were known for UGC 5944 to date. This dwarf elliptical galaxy was catalogued as member of the Leo Group in Ferguson & Sandage (1990) based on its morphology. With our SBF distance of $`11.1\pm 0.9`$ Mpc we can now confirm this impression. The derived SBF distance is also in good agreement with the mean distance of 11 Mpc for the group (Trentham & Tully 2002). We note that there was not much of colour range found in the five fields we analysed to unambiguously apply the distance calibration. Consequently, there is an alternative distance of $``$22 Mpc for this galaxy. However, this latter distance appears less likely.
### NGC 4150
NGC 4150 is morphologically classified as S0/Sa galaxy (Sandage & Bedke 1994). Once the galaxy light model is subtracted, the residual image shows a prominent, well developed 4-armed spiral structure with non-circular dust pattern in the central region (Fig. 2). The galaxy has a heliocentric velocity of $`v_{}=226\pm 22`$ km s<sup>-1</sup> (Fisher et al. 1995). Tonry et al. (2001) reported a first SBF distance modulus of $`(mM)=30.69\pm 0.25`$ mag (or $`13.7\pm 1.7Mpc`$) and Jensen et al. (2003) another, corrected SBF distance modulus of $`(mM)=30.53\pm 0.24`$ mag (or $`12.8\pm 1.5`$ Mpc). We analysed six independent SBF fields across the galaxy’s surface but well away from the spiral structures (Fig. 2). The derived distance modulus of $`(mM)_{\mathrm{SBF}}=30.79\pm 0.20`$ mag (see Fig. 8) agrees with the earlier published SBF results.
### BTS 128
This dwarf elliptical is located right in the densest region of the Coma I group (Binggeli et al. 1990; Trentham & Tully 2002) about 10 degrees away from the northern boundary of the Virgo cluster. The heliocentric velocity of BTS 128 of $`v_{}=1139\pm 86`$ km s<sup>-1</sup> (Wegner et al. 2001) is in good agreement with the Coma I cluster. BTS 128 is 27.5 arcmin away from the E0 galaxy NGC 4283 for which a SBF distance of $`(mM)=30.98\pm 0.19`$ mag was reported by Tonry et al. (2001). The SBF distance module we derive for BTS 128 is $`(mM)_{\mathrm{SBF}}=31.02\pm 0.25`$ mag ($`16.0\pm 2.0`$ Mpc).
### UGC 7639
This galaxy, classified as dE/Im, is located in the Canes Venatici region. The Im nature of the otherwise featureless dwarf galaxy becomes visible once the galaxy model is subtracted, with a central region of bright young stars and evidence of dust (see Fig. 10). It has a heliocentric velocity of $`v_{}=382`$ km s<sup>-1</sup> (de Vaucouleurs et al. 1991). Makarova et al. (1998) published a rough estimate of the distance of 8.0 Mpc based on brightest blue ($`BV<0.4`$) and red ($`BV>1.6`$) stars. We analysed eight independent SBF fields across the galaxy’s surface (Fig. 10 and 11, left panels). Because of the absence of a steep correlation between colour and SBF magnitude for the different SBF fields, we can use unambiguously the parabolic branch of the calibration diagram. The resulting SBF distance of $`7.1\pm 0.6`$ Mpc or $`(mM)_{\mathrm{SBF}}=29.27\pm 0.16`$ mag (see Fig. 11, right panel) identifies this galaxy as slightly closer than previously thought.
### UGC 8799
No distance was known for this dwarf elliptical to date but Geller et al. (1997) and Wegner et al. (2001) report a velocity of 1132 km s<sup>-1</sup>. Because this galaxy is only 1.5 hours in R.A. from the western boundary of the Virgo cluster and the velocity agrees well with the mean cluster velocity ($`980\pm 60`$ km s<sup>-1</sup>, Tanaka 1985) this may suggest that UGC 8799 has a distance that falls into the observed distance range of the Virgo cluster 14 Mpc $`<`$ D $`<`$ 23 Mpc (Jerjen et al. 2004). The six SBF fields (see Fig. 12) show only a narrow $`BR`$ colour spread and the signal-to-noise in the power spectra is low. This leads to two possible distance moduli: $`(mM)_{\mathrm{SBF}}=30.19\pm 0.23`$ mag or $`(mM)_{\mathrm{SBF}}=30.61\pm 0.26`$ mag (see Fig. 12), suggesting it is in the outskirts of the Virgo cluster.
### UGC 8882
Neither a velocity nor a distance was available for this nucleated dwarf elliptical galaxy. Bremnes et al. (1999) claim the galaxy is part of the M101 group and reported a mean $`BR`$ colour of 1.29. The Cepheid based distance modulus for M101 is $`(mM)_0=29.28\pm 0.14`$ mag (Stetson et al. 1998 and references therein). We analysed eight independent SBF fields across the galaxy’s surface (Fig. 13, left panels). Because of the observed linear relation between colour and SBF magnitude for the different fields we used the linear branch of the calibration diagram. The derived distance modulus is $`(mM)_{\mathrm{SBF}}=29.60\pm 0.20`$ mag (see Fig. 13, right panel). This result is consistent with UGC 8882 being a member of the M101 group.
## 5 Summary and Conclusions
We have analysed $`BR`$-band CCD images of nine nearby dwarf ellipticals and one S0 galaxy. We have measured their stellar $`R`$-band surface brightness fluctuation magnitudes $`\overline{m}_R`$ and $`(BR)_0`$ colours in 61 galaxy fields. The resulting distances were compared with existing distances measured with the TRGB method and SBF method. Agreement between our distances and those of TRGB and previous SBF measurements was very high; all distances agreed within 0.1 magnitudes.
Our SBF distances are given in Table 6 together with heliocentric radial velocities and radial velocities relative to the centre of mass of the Local Group of galaxies. The latter were calculated from heliocentric radial velocities following de Vaucouleurs et al. (1991) and using our distances to the galaxies. The Local Group centre of mass was assumed to lie exactly half way between the Milky Way and M31. The SBF distance to velocity relation positions BTS 128, UGC 7639 and UGC 8799 quite well on the Hubble flow. This is not the case with NGC 4150, KDG 61 and UGC 5442.
Our SBF distance to NGC 4150 of $`14.4\pm 0.7`$ Mpc is in good agreement with earlier SBF distances (Tonry et al. 2001, Jensen et al. 2003) of $`13.7\pm 1.7`$ Mpc and $`12.8\pm 1.5`$ Mpc, and places this galaxy almost 10 Mpc behind the centre of Canes Venatici I cloud. All these estimates are consistent with Karachentsev et al. (2003) lower distance limit of 6.3 Mpc using the TRGB method in a study of CVn I cloud, and their very rough estimate of 20 Mpc, using the globular cluster luminosity function method. The galaxy might then be associated with the outskirts of the Virgo cluster. The velocity of NGC 4150 is rather low ($`v_{}=226`$km s<sup>-1</sup>), but it is difficult to say whether it is discordant with either our distance estimate of $`14.4\pm 0.7`$ Mpc or Karachentsev et al.’s (2003) $`20`$Mpc. Solanes et al. (2002) have found that most Virgo galaxies in the region closest to us, corresponding to our distance to NGC 4150, have high radial velocities outward from the cluster centre (with similarly unusually high velocities away from us for Virgo cluster galaxies in the region behind the centre of the cluster). The angular separation of NGC 4150 from the cluster centre is 19 degrees. If it has a high outward velocity from the cluster centre, its line-of-sight velocity could be quite low. The observed velocity of NGC 4150 can be consistent with membership in the Virgo cluster, at its outer edge.
Our SBF distance to KDG 61 confirms the existing TRGB distance (Karachentsev et al. 2000) and thus the membership in M81 group. Another M81 group member is confirmed as our SBF distance to UGC 5442 agrees well with a previous TRGB distance (Karachentsev et al. 2000). Radial velocities of KDG 61 and UGC 5442 are reasonable relative to M81 velocity of $`v_{}=34\pm 4`$ (de Vaucouleurs et al. 1991) or $`v_{LG}96`$ km s<sup>-1</sup>.
UGC 7639 had only a tentative distance measured with brightest blue and red stars (Makarova et al. 1998) before our SBF distance, which is in relatively good agreement with the earlier distance and confirms the location of this galaxy in the Canes Venatici II cloud.
Distances to galaxies UGC 1703, UGCA 200, UGC 5944 and UGC 8882 were not previously known. We have provided accurate SBF distances to these galaxies. Our distances suggest UGC 1703 may be a distant companion of NGC 784, UGCA 200 may not be a companion of NGC 3115 as was assumed before, UGC 5944 is most certainly a member of Leo I group and UGC 8882 seems to be a member of M101 group. Likewise only radial velocities were known for galaxies BTS 128 and UGC 8799 before our SBF distances, which confirm the membership of BTS 128 in Coma I group and suggest UGC 8799 lies at the outskirts of Virgo I cluster. Its angular separation from the cluster centre, at M87 location of 12<sup>h</sup> 31<sup>m</sup> in RA and +12 23<sup>m</sup> in dec. (J2000.0), is 21 degrees.
The SBF distances we have presented continue to support the understanding of the distribution of dwarf galaxies in galaxy groups and intermediate space in the Local Group neighbourhood. They also demonstrate well the feasibility of the Surface Brightness Fluctuation method in determining accurate distances with 2m class ground-based telescopes out to the near side of the Virgo cluster.
###### Acknowledgements.
The Nordic Optical Telescope is operated on the island of La Palma jointly by Denmark, Finland, Iceland, Norway, and Sweden, in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofísica de Canarias. We thank Kari Nilsson for his help with the observations and the Academy of Finland for support through its funding of the ANTARES programme. Financial support for RR has been provided by Finnish Graduate School in Space Physics and Astronomy and by the Academy of Finland through funding of the project “Calculation of Orbits”. This research has made use of the NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. We are grateful to the referee Enzo Brocato for his useful comments.
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# High Excitation Molecular Gas in the Magellanic Clouds
## 1 Introduction
With their proximity, their unextinguished lines of sight, and their profuse star formation, the Magellanic Clouds are some of the best extragalactic objects in which to study the relationship between molecules and star formation. Because the interstellar medium (ISM) in the Magellanic Clouds is deficient in heavy elements and dust, molecular observations of these objects probe an interesting regime, perhaps more similar to the conditions in early protogalaxies rather than those prevalent today in the Milky Way. Indeed, studies of active, metal–poor, nearby dwarf galaxies such as the Magellanic Clouds should offer insight into the processes at work in primeval sources.
What are the effects of low metallicities on the star–forming molecular ISM? We know that molecules are more difficult to form both in the gas phase (less O, C, and N) and on grain surfaces (fewer grains), and that they are easier to destroy (diminished dust shielding of the UV radiation) in low metallicity environments. Thus, molecules other than $`\mathrm{H}_2`$ are more rare in these sources which in particular translates into a dearth of CO emission in the Magellanic Clouds and other dwarf irregulars (e.g., Israel et al., 1986; Israel, Tacconi & Baas 1995, Taylor, Kobulnicky, & Skillman 1998; Leroy et al. 2005). Furthermore, because the FIR and submillimeter lines of the different forms of carbon and oxygen (C<sup>+</sup>, C, O, and CO) dominate the cooling of the star–forming ISM (e.g., Le Bourlot et al. 1993; Wolfire et al. 1995), the lower abundances of these elements will make the cooling of the molecular gas less efficient. However, at the same time smaller dust–to–gas ratios will yield lower heating of the molecular gas, as photoelectric ejection of electrons from small dust grains is the chief mode by which starlight heats the gas phase of the ISM. If molecular gas temperatures were considerably affected by the metallicity of the ISM, we expect important consequences for star formation in such environments. In particular, if the Jeans criterion is relevant to star–formation, the mass of collapsable clouds grows for decreasing metallicity as the Jeans mass increases $`M_JT^{3/2}`$. Such change could have important effects on the Initial Mass Function of stars in these systems. To a first approximation, models suggest that lowering the metallicity causes a similar decrement in both heating (by diminishing the dust–to–gas ratio) and cooling (by diminishing the C and O abundances; Wolfire et al. 1995).
It is important to realize, however, that there are many possibilities likely to complicate this simple picture. For example, if the dust–to–gas ratio were to decrease faster than the metallicity (as suggested by Lisenfeld & Ferrara 1998), or if there were a lack of very small dust grains in the low metallicity ISM (as suggested by the faintness of the polycyclic aromatic hydrocarbon emission observed toward some of these sources; e.g., Madden 2000), the heating processes may become less efficient and the balance may be shifted toward lower temperatures. Furthermore, because it is also necessary to consider the metallicity threshold below which hitherto secondary heat sources (e.g., chemical heating) become important, the effects of metallicity on the heating and cooling balance of molecular clouds are very difficult to address from a purely theoretical approach.
Answering some of these questions observationally requires studying the physical conditions of the molecular gas in nearby, low metallicity sources. At distances of 55 and 63 kpc, the proximity of the Magellanic Clouds affords single–dish millimeter–wave observations excellent spatial resolution attainable in other galaxies only through the use of interferometers, which permits detailed studies of individual clouds instead of ensemble properties. In particular, the ability to spatially separate the emission from different regions makes the Magellanic Clouds ideal targets to study the excitation of the molecular gas and its relationship with star formation. Multitransition studies of CO and other molecules are very useful tools to determine the physical conditions of the H<sub>2</sub> (e.g., Johansson et al. 1998; Heikkilä 1998; Chin et al. 1998; Heikkilä et al. 1999), but their application is limited if there is no information on the intensities of the higher CO transitions, which are extremely sensitive to density and temperature. Because of their southern declination, however, there is a dearth of submillimeter observations of the Clouds.
We present here a survey of <sup>12</sup>CO ($`J=43`$) in the molecular peaks of the Magellanic Clouds, and the results of an excitation analysis using these and lower $`J`$ observations. These pointings were selected among the brightest <sup>12</sup>CO ($`J=10`$) peaks found by SEST observations (Israel et al. 1993), many of which are associated with star–forming complexes and Henize (1956) H$`\alpha `$ nebulosities, and are thus denoted using the corresponding “N” number. In section §2 we present the observations, in §3 we discuss the LVG analysis and its results, and in §4 we summarize our conclusions.
## 2 Observations and Results
We observed the ($`J=43`$) transition of carbon monoxide (<sup>12</sup>CO) at $`\nu 461.0408`$ GHz (650.7 $`\mu `$m) using AST/RO, the Antarctic Submillimeter Telescope and Remote Observatory located at the Amundsen–Scott South Pole base (Stark et al. 2001). The observations were obtained on the austral winter of 2002, using the lower frequency side of the dual AST/RO SIS waveguide receiver (Walker et al. 1992; Honingh et al. 1997), with system temperatures of $`T_{sys}2000`$ K. The backend was the 2048 channel low resolution (1.07 MHz resolution, 0.68 MHz channels) acousto–optical spectrometer (Schieder, Tolls, and Winnewisser 1989). The spectra were observed in position switching mode, chopping 25′ in Azimuth (which is the same as R.A. at the pole). At 461 GHz, the telescope beam was measured to have a HPBW$`109\mathrm{}`$. The forward efficiency determined from skydips was $`70\%`$, and is assumed to be identical to $`\eta _{mb}`$. Maps of size $`6\mathrm{}\times 6\mathrm{}`$ centered on the SEST coordinates (c.f. Table 1) were obtained for each region, using a 30″ grid (in a few cases the maps were done on a 60″ grid). The data were calibrated using the standard procedure for AST/RO, which includes sky, ambient and cold load measurements every 20–30 minutes, and processed using the COMB astronomical package. The individual maps are shown in Fig. 1, and representative spectra are shown in Figs. 2 and 3. Because the pointing accuracy is estimated to be $`1\mathrm{}`$, we have selected the emission peaks closest to the center of the map within those margins to measure the integrated intensities compiled in Table 1 (indicated by the crosses in Fig. 1). Note that this method may introduce a bias in the direction of obtaining larger CO ($`J=43`$)/($`J=10`$) ratios. Nonetheless we feel that this methodology is justified, as in the Milky Way the positional coincidence between the peaks of both transitions is very often observed. In some cases it is apparent that the structure of the sources is complex (e.g., N 159W; Bolatto et al. 2000), and in particular for N 167 and N 83 the brightest emission peak is well away from the central position. Given the large and variable effects of the atmosphere at submillimeter wavelengths and the pointing accuracy of the telescope, we estimate the overall absolute calibration accuracy of the data to be $`30\%`$.
Table 1 summarizes our observations. Column 4 lists the <sup>12</sup>CO ($`J=10`$) integrated intensities observed by the SEST Magellanic Cloud Key Programme toward our sources, while column 5 shows the intensities convolved to the angular resolution of AST/RO. Finally, column 6 lists the <sup>12</sup>CO ($`J=43`$) integrated intensities used in our LVG analysis, with their statistical errors. These intensities are those measured in the AST/RO pointings marked with crosses in Fig. 1.
## 3 Analysis and Discussion
### 3.1 Modeling of CO
The available observed <sup>12</sup>CO and <sup>13</sup>CO line ratios have been modelled using the large–velocity gradient (LVG) radiative transfer models described by Jansen (1995) and Jansen et al. (1994). These models provide line intensities as a function of three input parameters: gas kinetic temperature $`T_\mathrm{k}`$, molecular hydrogen density $`n(H_2)`$ and CO column density per unit velocity ($`N(\mathrm{CO})`$/d$`V`$). By comparing model line ratios to the observed ratios we determine the physical parameters best describing the conditions in the observed source. In principle, with two isotopes we need to measure five independent line intensities in order to fully determine the conditions of a single molecular gas component (i.e. $`T_\mathrm{k}`$, $`n(H_2)`$, $`N(\text{12}\text{CO})`$/d$`V`$, $`N(\text{13}\text{CO})`$/d$`V`$ and a beam filling–factor). By assuming a fixed isotopical abundance \[<sup>12</sup>CO\]/\[<sup>13</sup>CO\] = 40 (Johansson et al. 1994) we may decrease this requirement to four independent line intensities. As Table 2 shows, this minimum requirement is met by two out of nine LMC objects and five out of six SMC objects. The physical conditions of the remaining eight objects are, in principle, underdetermined.
More realistic, and consequently more complex, models of gas excitation that include more than one component require many more observations to be properly constrained. Full modelling of a two–component molecular cloud using two isotopes requires ten independent measurements, which again are reduced to eight by the introduction of a fixed isotopical abundance. As this is more than we have actually observed in any of the LMC or SMC clouds, it is clear that the solutions may not be unique. As long as the range of possible solutions is not excessive, however, they are still useful to constrain the physical parameters governing the observed emission.
We identified acceptable fits by searching a grid of model parameter combinations (10 K $`T_\mathrm{k}`$ 150 K, $`10^2\mathrm{cm}^3n(\mathrm{H}_2)10^5\mathrm{cm}^3`$, $`6\times 10^{15}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1N(`$CO$`)/`$d$`V\mathrm{\hspace{0.17em}3}\times 10^{18}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1`$) for model line ratios matching the observed values. Although errors in the line ratios increase the range of possible solutions, these ratios tend to define reasonably well–constrained regions of parameter space. The solutions are somewhat degenerate, as variations in the parameters may compensate one another. For instance, a simultaneous increase in kinetic temperature and decrease in $`\mathrm{H}_2`$ densities (or vice versa) yields similar line ratios (see § 3.2.3).
### 3.2 Single–component fits
#### 3.2.1 Objects with two measured intensities
For three objects (LMC-N 48, LMC-N 214C, and SMCB2#6) we only have intensities in the ($`J=10`$) and ($`J=43`$) transitions of <sup>12</sup>CO. The parameters of these clouds are thus poorly constrained, and not summarized in a table. Assuming a single molecular gas component, we find for SMCB2#6 no effective constraints: $`T_\mathrm{k}`$ = 20–150 K, $`n(\mathrm{H}_2)=10^210^4\mathrm{cm}^3`$, and $`N(`$CO$`)/`$d$`V=10^{16}10^{18}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1`$. All <sup>12</sup>CO transitions should be very optically thick. LMC-N 214C is slightly better determined: temperatures below 30 K are not allowed and densities appear high $`n(\mathrm{H}_2)=10^410^5\mathrm{cm}^3`$. The ($`J=10`$) transition should not be very optically thick (although the higher transitions are) and its isotopic intensity ratio should be 10 – 20, as is in fact commonly observed in the LMC (Israel et al. 2003). Finally, the high <sup>12</sup>CO ($`J=43`$)/($`J=10`$) ratio exhibited by LMC-N 48 are, in the single–component approximation, only consistent with low optical depths (isotopic intensity ratios 20 – 40), rather high densities $`n(\mathrm{H}_2)=10^410^5\mathrm{cm}^3`$ and temperatures $`T_\mathrm{k}`$ = 60–150 K, and gradients $`N(`$CO$`)/`$d$`V=10^{16}10^{17}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1`$.
#### 3.2.2 Objects with three measured intensities
For five objects, all in the LMC, (N 55, N 79, N 83A, N 113 and LIRL 648) we have measured <sup>13</sup>CO ($`J=10`$) intensities in addition to the <sup>12</sup>CO ($`J=10`$) and ($`J=43`$) intensities. As is clear from the previous discussion, the molecular gas parameters are still underdetermined, but not fully unconstrained (Table 3). Both N 79 and N 113 fit very well to a hot and fairly dense model cloud with $`T_\mathrm{k}`$ = 100–150 K, $`n(\mathrm{H}_2)=30005000\mathrm{cm}^3`$ and a gradient of about $`6\times 10^{17}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1`$. LIRL 648 is fit, but not very well, by a gas at the somewhat lower temperature of 60 K and the somewhat higher density of $`10^4\mathrm{cm}^3`$. The physical parameters of the molecular gas cloud associated with N 83A are inconsistent with the assumption of a single component. Only a very poor fit is obtained in the high-temperature, high-density, and high-velocity gradient limit.
#### 3.2.3 Objects with four or more measured intensities
Seven objects have a sufficiently large number of measured line intensities to allow a full determination of physical parameters, assuming they can be properly described by a single molecular gas component. This appears indeed to be the case for SMCB1#1, where we find excellent agreement between the observed ratios and those of a model gas characterized by $`T_\mathrm{k}`$ = 50 K, $`n(\mathrm{H}_2)=700\mathrm{cm}^3`$, $`N(`$CO$`)/`$d$`V=5\times 10^{16}\mathrm{cm}^2(\mathrm{km}\mathrm{s}^1)^1`$. To fine-tune the solutions, as well as to gain insight into the various trade-offs, we have run a finer grid of models for this particular source. We summarize the best solutions for the temperature range $`T_\mathrm{k}`$ = 30–60 K in Table 4. The overall best fit is achieved using a kinetic temperature $`T_\mathrm{k}`$ = 50 K, and a density $`n(\mathrm{H}_2)=700\mathrm{cm}^3`$. At constant temperature, the uncertainty in density is about 20$`\%`$. From the Table, it is also clear that we may allow for a similar temperature uncertainty of 20$`\%`$ if we simultaneously increase or decrease the density by 40$`\%`$.
The good fit between model and observed line ratios is of particular significance because the solution is overdetermined with six independent intensity measurements. SMCB1#1 is a small and isolated molecular cloud in the SMC Bar, not associated with a star-forming region (Rubio et al. 1993; Reach, et al. 2000; Rubio et al. 2004), thus a low density of $`700`$ $`\mathrm{cm}^3`$ is not necessarily surprising, although the kinetic temperature is higher than we would have expected for such a cloud in the Milky Way. Table 4, however, illustrates the difficulty in pinning down the physical conditions even in a simple cloud with a simple model.
The measured intensities of N 159-W can also be fit, albeit somewhat poorly, by a single hot ($`T_\mathrm{k}`$ = 150 K) and moderately dense component ($`n(\mathrm{H}_2)=1000\mathrm{cm}^3`$). However, none of the other five objects have intensities consistent with a single component. Typically, a component fitting the observed <sup>12</sup>CO line ratios would fail completely to explain the <sup>12</sup>CO/<sup>13</sup>CO isotopical ratios, indicating that in addition to a warm gas component the presence of a second cooler and dense component should be assumed. The lack of associated luminous objects probably explains the exceptionally homogeneous nature of SMCB1#1, required for such a well-determined single-component fit to be valid.
### 3.3 Dual-component fits
Since (with the exception of SMCB1#1 and perhaps N 159-W) none of the sources for which a lot of constraints are available allows a good fit with a single gas component, we suspect that the succesful single-component fits of the sources in Table 3 result primarily from insufficient information rather than from a simple physical structure. Attempts at more sophisticated modelling achieve little of value for most of these sources. The single exception is LMC-N 83A where only two observed ratios nevertheless conflict with every single component model tried. A dual component model yields acceptable but, not surprisingly, poorly constrained solutions (see Table 5). For all the other well-observed sources where single-component model fits failed, we have also constructed dual-component fits which are likewise listed in Table 5. This table shows that usually one of the two components component is reasonably well-determined, whereas the parameters of the other are generally less tightly constrained.
## 4 Summary and Conclusions
A general result of our survey is the detection of significant <sup>12</sup>CO ($`J=43`$) emission in most molecular peaks in the Magellanic Clouds. By itself this demonstrates the widespread occurrence of significant amounts of warm molecular gas, as tentatively suggested by Israel et al. (2003). Application of LVG model calculations show that molecular gas kinetic temperatures as high as $`T_{\mathrm{kin}}`$ = 100 – 300 K frequently occur. Detailed analysis of the objects for which multiple line ratios are available strongly suggests that the higher temperatures occur in cloud regions that are not very dense ($`n_{\mathrm{H}_2}=10^210^3\mathrm{cm}^3`$), and that this gas is generally associated with colder (typically $`T_{\mathrm{kin}}`$ = 10 – 60 K) but much denser ($`n_{\mathrm{H}_2}=10^410^5\mathrm{cm}^3`$) molecular gas. The simplified analysis possible in cases where fewer line ratios are available strongly suggests that the situation in those objects is similar. We note that recently reported observations of the N 44 complex in the LMC by Kim et al. (2004) show also the presence of strong <sup>12</sup>CO ($`J=43`$) emission, leading the authors to conclude the likely presence of very high densities ($`n_{\mathrm{H}_2}10^5\mathrm{cm}^3`$).
A few clouds are notable in our survey. Already mentioned is the case of SMCB1#1, a quiescent, compact molecular cloud devoid of any star forming activity. It appears to be essentially homogeneous and of modest density ($`n_{\mathrm{H}_2}=700\mathrm{cm}^3`$), but is surprisingly warm ($`T_{\mathrm{kin}}`$ = 50 K). As there are no known embedded heating sources, this temperature must be maintained by the environment of the cloud, located in the south west region of the SMC Bar. The other two notable objects also occur in the SMC. N 12 is likewise located in the southern part of the Bar. Here, both the dense and relatively tenuous molecular phases appear to be surprisingly hot ($`T_{\mathrm{kin}}`$ = 150 K). In N 66, a very hot ($`T_{\mathrm{kin}}`$ = 300 K) and a cooler ($`T_{\mathrm{kin}}=40\pm 20`$ K) phase coexist, both at a rather high density of the order of $`n_{\mathrm{H}_2}=10^4\mathrm{cm}^3`$. This particular molecular cloud is a relatively small remnant in a large and luminous star-forming complex (c.f., Rubio et al. 1996). It is almost certainly in an advanced stage of destructive processing; high densities and temperatures are consistent with such a situation (Rubio et al. 2000; Contursi et al. 2000).
Have we found signs of metallicity effects on the temperature equilibrium of the ISM in metal–poor environments? Given the measurement uncertainties and the dearth of comparable datasets on “normal” clouds in our own galaxy such finding cannot be asserted on the present measurements. It is suggestive, however, that high temperatures seem pervasive even in largely quiescent clouds such as SMCB1#1. This result should be taken with caution, as we suffer from a strong sample selection bias: the brightness of the optically thick <sup>12</sup>CO ($`J=10`$) transition is proportional to the source temperature and beam filling fraction, and we have pointed toward the brightest <sup>12</sup>CO ($`J=10`$) clouds. Thus it may not be surprising that our clouds are warm. Furthermore, studies of Galactic photodissociation regions (e.g., S 140, Draine & Bertoldi 1999; NGC 2023, Draine & Bertoldi 2000) find that H<sub>2</sub> in star–forming clouds is frequently warmer than the theoretical expectation. Similarly, the intensities of the mid–J CO lines tend to be underpredicted by homogeneous PDR calculations (Hollenbach & Tielens 1999 and references therein). It is unclear whether these discrepancies are due to the presence of other heating mechanisms or to shortcomings in the models (e.g., Draine & Bertoldi 1999).
From the observational standpoint, a conclusive study of the effects of metallicity on gas temperature in the Magellanic Clouds must await the availability of more powerful submillimeter instruments observing the southern sky such as the Atacama Large Millimeter Array (ALMA), and the recently deployed single–dish telescopes: the Atacama Pathfinder Experiment (APEX) and the Atacama Submillimeter Telescope Experiment (ASTE).
We wish to thank the AST/RO group, and the anonymous referee. ADB wishes to thank D. Hollenbach for discussions on the subject of warm molecular gas, and for suggestions that helped improve this manuscript. ADB and CLM acknowledge support from NSF grants AST-0228963 and OPP-0126090 respectively.
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# Microscopic calculations of spin polarized neutron matter at finite temperature
## Abstract
The properties of spin polarized neutron matter are studied both at zero and finite temperature within the framework of the Brueckner–Hartree–Fock formalism, using the Argonne v18 nucleon-nucleon interaction. The free energy, energy and entropy per particle are calculated for several values of the spin polarization, densities and temperatures together with the magnetic susceptibility of the system. The results show no indication of a ferromagnetic transition at any density and temperature.
PACS:26.60.+c,21.60.Jz,26.50.+x
Keywords: Neutron matter, Ferromagnetic transition, Finite temperature
Since the suggestion of Pacini and Gold pulsars are generally believed to be rapidly rotating neutron stars with strong surface magnetic fields in the range of $`10^{12}10^{13}`$ Gauss. Despite the great theoretical effort of the last forty years, there is still no general consensus regarding the mechanism to generate such strong magnetic fields in a neutron star. The fields could be a fossil remnant from that of the progenitor star or, alternatively, they could be generated after the formation of the neutron star by some long-lived electric currents flowing in the highly conductive neutron star material. From the nuclear physics point of view, however, one of the most interesting and stimulating mechanisms which have been suggested is the possible existence of a phase transition to a ferromagnetic state at densities corresponding to the theoretically stable neutron stars and, therefore, of a ferromagnetic core in the liquid interior of such compact objects. Such a possibility has been considered since long ago by several authors within different theoretical approaches , but the results are still contradictory. Whereas some calculations, like for instance the ones based on Skyrme-like interactions predict the transition to occur at densities in the range $`(14)\rho _0`$ ($`\rho _0=0.16`$ fm<sup>-3</sup>), others, like recent Monte Carlo and Brueckner–Hartree–Fock calculations using modern two- and three-body realistic interactions exclude such a transition, at least up to densities around five times $`\rho _0`$. This transition could have important consequences for the evolution of a protoneutron star, in particular for the spin correlations in the medium which do strongly affect the neutrino cross sections and the neutrino mean free path inside the star . Therefore, drastically different scenarios for the evolution of protoneutron stars emerge depending on the existence of such a ferromagnetic transition.
Most of the studies of the ferromagnetic transition in neutron and nuclear matter have been done at zero temperature. However, the description of protoneutron stars motivates a study of spin polarized neutron matter at temperature $`T`$ of the order of a few tens of MeV. Recently, the properties of polarized neutron matter both at finite and zero temperature, have been investigated using a large sample of Skyrme-like interactions. The results of Ref. indicate the occurrence of a ferromagnetic phase of neutron matter. However, contrary to what one would intuitively expect, the authors of Ref. have found that the critical density at which ferromagnetism takes place decreases with temperature. This unexpected result was associated to an anomalous behaviour of the entropy of the system which becomes larger for the spin-polarized phase with respect the one for the non-polarized phase, above a certain density. This was shown to be related to the dependence of the effective masses of neutrons with spin up and down on the amount of spin-polarization, and a new constraint on the parameters of the Skyrme force was derived if this anomalous behaviour is to be avoided .
In the present work, we study the bulk and single particle properties of spin-polarized neutron matter at finite temperature. To this aim we make use of a microscopic approach based on the Brueckner–Hartree–Fock (BHF) approximation of the Brueckner–Bethe–Goldstone (BBG) expansion. Here we make use of an extension of the BBG theory (i) to the case in which neutron matter is arbitrarily asymmetric in the spin degree of freedom (i.e., $`\rho _{}\rho _{}`$, where $`\rho _{}`$ ($`\rho _{}`$) is the density of neutron with spin up (down)), and (ii) to the case of finite temperature. In particular, we study the behaviour of the entropy of the system and the effective mass of neutrons as a function of the spin polarization parameter, $`\mathrm{\Delta }=(\rho _{}\rho _{})/(\rho _{}+\rho _{}`$). We show that, contrary to what it is found in Ref. , the entropy of the polarized phase is lower than that of the non-polarized one, according to the idea that the polarized phase is more “ordered” than the non-polarized one.
Our calculation starts with the construction of the neutron-neutron $`G`$-matrix, which describes in an effective way the interaction between two neutrons for each one of the spin combinations $`,,`$ and $``$. This is formally obtained by solving the well known Bethe–Goldstone equation, written schematically as
$$G(\omega )_{\sigma _1\sigma _2,\sigma _3\sigma _4}=V_{\sigma _1\sigma _2,\sigma _3\sigma _4}+\underset{\sigma _i\sigma _j}{}V_{\sigma _1\sigma _2,\sigma _i\sigma _j}\frac{Q_{\sigma _i\sigma _j}}{\omega \epsilon _{\sigma _i}\epsilon _{\sigma _j}+i\eta }G(\omega )_{\sigma _i\sigma _j,\sigma _3\sigma _4},$$
(1)
where the first (last) two sub-indices indicate the spin projection $`\sigma =(`$) of the two neutrons in the initial (final) state, $`V`$ is the bare nucleon-nucleon interaction, $`Q_{\sigma _i\sigma _j}`$ is the Pauli operator which allows only intermediate states compatible with the Pauli principle, and $`\omega `$ is the starting energy defined as the sum of the non-relativistic single-particle energies, $`\epsilon _{()}`$, of the interacting neutrons.
The single-particle energy of a neutron with momentum $`k`$ and spin projection $`\sigma =()`$ is given by
$$\epsilon _\sigma (k)=\frac{\mathrm{}^2k^2}{2m}+Re[U_\sigma (k)],$$
(2)
where the real part of the single-particle potential $`U_\sigma (k)`$ represents the averaged field “felt” by the neutron due to its interaction with the other neutrons of the system. In the BHF approximation it is given by
$$U_\sigma (k)=\underset{\sigma ^{}k^{}}{}n_\sigma ^{}(k^{})\stackrel{}{k}\sigma \stackrel{}{k}^{}\sigma ^{}|G(\omega =\epsilon _\sigma (k)+\epsilon _\sigma ^{}(k^{})|\stackrel{}{k}\sigma \stackrel{}{k}^{}\sigma ^{}_A,$$
(3)
where
$$n_\sigma (k)=\{\begin{array}{cc}1,\text{if }kk_F^\sigma \hfill & \\ 0,\text{otherwise}\hfill & \end{array}$$
(4)
is the corresponding occupation number of a neutron with spin projection $`\sigma `$ and the matrix elements are properly anti-symmetrized. We note here that the so-called continuous prescription has been adopted for the single-particle potential when solving the Bethe–Goldstone equation. As shown by the authors of Refs. , the contribution to the energy per particle from three-body clusters is diminished in this prescription with respect to the one calculated with the gap choice for the single particle potential. We also note that the present calculation has been carried out using the Argonne v18 nucleon-nucleon potential . The momentum dependence of the single-particle spectrum can be characterized by the effective mass $`m_\sigma ^{}(k)`$ defined as:
$$\frac{m_\sigma ^{}(k)}{m}=\frac{k}{m}\left(\frac{d\epsilon _\sigma (k)}{dk}\right)^1,$$
(5)
where $`m`$ is the bare neutron mass.
The total energy per particle is easily obtained once a self-consistent solution of Eqs. (1)–(3) is achieved
$$\frac{E}{A}=\frac{1}{A}\underset{\sigma k}{}n_\sigma (k)\left(\frac{\mathrm{}^2k^2}{2m}+\frac{1}{2}Re[U_\sigma (k)]\right).$$
(6)
The many-body problem at finite temperature has been considered by several authors within different approaches, such as the finite temperature Green’s function method , the thermo field method , or the Bloch–De Domicis (BD) diagrammatic expansion . The latter, developed soon after the Brueckner theory, represents the “natural” extension to finite temperature of the BBG expansion, to which it leads in the zero temperature limit. Baldo and Ferreira showed that the dominant terms in the BD expansion were those that correspond to the zero temperature of the BBG diagrams, where the temperature is introduced only through the Fermi-Dirac distribution
$$f_\sigma (k,T)=\frac{1}{1+exp([\epsilon _\sigma (k,T)\mu _\sigma (T)]/T)},$$
(7)
$`\mu _\sigma (T)`$ being the chemical potential of a neutron with spin projection $`\sigma `$. Therefore, at the BHF level, finite temperature effects can be introduced in a very good approximation just replacing in the Bethe–Goldstone equation: (i) the zero temperature Pauli operator $`Q_{\sigma _i\sigma _j}=(1n_{\sigma _i})(1n_{\sigma _j})`$ by the corresponding finite temperature one $`Q_{\sigma _i\sigma _j}(T)=(1f_{\sigma _i})(1f_{\sigma _j})`$, and (ii) the single-particle energies $`\epsilon _\sigma (k)`$ by the temperature dependent ones $`\epsilon _\sigma (k,T)`$ obtained from Eqs. (2) and (3) when $`n_\sigma (k)`$ is replaced by $`f_\sigma (k,T)`$. These approximations, which are supposed to be valid in the range of densities and temperatures considered here, correspond to the “naive” finite temperature Brueckner–Bethe–Goldstone (NTBBG) expansion discussed in Ref. .
In this case, however, the self-consistent process implies that, together with the Bethe–Goldstone equation and the single-particle potential, the chemical potentials of neutrons with spin up and down must be extracted at each step of the iterative process from the normalization condition
$$\rho _\sigma =\underset{k}{}f_\sigma (k,T).$$
(8)
This is an implicit equation which can be solved numerically. Note that the $`G`$-matrix obtained from the Bethe–Goldstone equation (1) and also the single-particle potentials depend implicitly on the chemical potentials.
Once a self-consistent solution is achieved the total free energy per particle is determined by
$$\frac{F}{A}=\frac{E}{A}T\frac{S}{A},$$
(9)
where $`E/A`$ is evaluated from Eq. (6) replacing $`n_\sigma (k)`$ by $`f_\sigma (k,T)`$ and the total entropy per particle, $`S/A`$, is calculated through the expression
$`{\displaystyle \frac{S}{A}}={\displaystyle \frac{1}{A}}{\displaystyle \underset{\sigma k}{}}[f_\sigma (k,T)\text{ln}(f_\sigma (k,T))+(1f_\sigma (k,T))\text{ln}(1f_\sigma (k,T))].`$ (10)
From the free energy per particle, we can get the remaining macroscopic properties of the system. In our case, we are particularly interested in the magnetic susceptibility $`\chi `$, which characterizes the response of a system to a magnetic field and gives a measure of the energy required to produce a net spin alignment in the direction of the field. It is given by
$$\chi =\frac{\mu ^2\rho }{\left(\frac{^2(F/A)}{\mathrm{\Delta }^2}\right)_{\mathrm{\Delta }=0}}$$
(11)
where $`\mu `$ is the magnetic moment of the neutron.
The single-particle potentials of neutrons with spin up and down have been simultaneously and self-consistently calculated together with their effective interactions. The results at $`\rho =0.16`$ fm<sup>-3</sup> and spin polarization $`\mathrm{\Delta }=0.5`$ are reported for T=0 (left panel) and T=40 MeV (right panel) on the top panels of Fig. 1. The neutron single-particle potential splits up in two different components when a partial spin polarization is assumed. In the case of Fig. 1, the single-particle potential $`Re[U_{}(k)]`$ for neutrons with spin up (the most abundant component) is less attractive than the one for neutrons with spin down, $`Re[U_{}(k)]`$. As demonstrated by the authors of Ref. (see in particular their Eqs. (23) and (24)), this splitting (i) is the result of a phase space effect, i.e. to the change in the number of pairs which the neutron under consideration $`|k,\sigma `$ can form with the remaining neutrons $`|kk_F^\sigma ^{},\sigma ^{}=,`$ of the system as neutron matter is polarized, and (ii) is due to the spin dependence of the neutron-neutron G-matrix in the spin polarized medium (see Eq. (1)). Indeed, as polarization increases, the single particle potential of a spin up neutron is built from a larger number of up-up pairs that form a spin triplet state ($`S=1`$) and, due to the Pauli principle, can only interact through odd angular momentum partial waves. Conversely, the potential of the less abundant species is built from a relatively larger number of up-down pairs which can interact both in the $`S=0`$ and $`S=1`$ two body states. Thus, the potential of the less abundant species receives also contributions from some important attractive channels as e.g. the $`{}_{}{}^{1}S_{0}^{}`$.
The increase of the temperature changes moderately the single-particle potentials. The real part becomes slightly less attractive, whereas the imaginary part increases in size as a consequence of the increase of phase space in the low momentum region.
The momentum dependence of the corresponding effective masses of the two components is also shown in the bottom panels of the figure for the same values of density, spin polarization and temperatures. The general effect of temperature is to smooth out the enhancement of the effective mass near the Fermi surface, as observed in the work of Ref. in symmetric nuclear matter.
In Fig. 2 we show the effective mass $`m_{}^{}`$($`m_{}^{}`$) for neutron with spin up (down) as a function of the spin polarization $`\mathrm{\Delta }`$, for fixed density ($`\rho =0.16`$ fm<sup>-3</sup>) and temperature (T=0 and T=40 MeV). The effective mass is calculated using Eq. (5) taken for each component at the corresponding Fermi momentum. Obviously, for $`\mathrm{\Delta }=0`$ the effective mass of the two components coincides. Once some amount of polarization is considered, the values of the effective masses split in two, the effective mass of the most abundant component being larger than the one of the less abundant. As can be seen the effective masses show an almost linear and symmetric variation with respect to their common value at spin polarization $`\mathrm{\Delta }=0`$, both at T=0 and T=40 MeV. Deviations from this behaviour are only found at the higher polarization values. This behaviour of $`m_\sigma ^{}`$ is a direct consequence of the scissors-like dependence of the single particle potential $`Re[U_\sigma ]`$ as a function of the spin polarization parameter $`\mathrm{\Delta }`$ (see Fig. 2 of Ref. ). A similar qualitative behaviour for the nucleon effective mass, as a function of the isospin asymmetry parameter, $`\beta =(\rho _n\rho _p)/\rho `$, has been found in isospin asymmetric nuclear matter (see in particular Eq. (94) in Ref. ).
The differences of the free energy (F/A), energy (E/A) and entropy (S/A) per particle between the totally polarized and the non-polarized phases are reported in the left, central and right panels of Fig. 3 as a function of the density for several temperatures. The differences in the three quantities increase with density and increase (decrease) with temperature in the case of the free energy (energy and entropy). Contrary to the results of Ref. with the Skyrme interaction, these differences are always positive for the F/A and E/A. This is an indication that the non-polarized phase is energetically preferred in the range of densities explored. Therefore, we can conclude that a phase transition to a ferromagnetic state is not to be expected from our microscopic calculation. If such a transition would exist, the difference in the free energy would become zero at some density, indicating that the ground state of the system would be ferromagnetic from that density on. In addition, the difference in the entropy is always negative indicating, as one intuitively expects, that the totally polarized phase is more “ordered” than the non-polarized one.
In Fig. 4 we show the behaviour of the free energy F/A per particle as a function of the spin polarization for several densities (left panel) and temperatures (right panel). Circles, squares, diamonds and triangles correspond to our BHF results, whereas the solid lines correspond to the parabolic approximation discussed below. As we expected from our previous calculations at zero temperature and , F/A is symmetric in $`\mathrm{\Delta }`$ and it shows a minimum at $`\mathrm{\Delta }=0`$ for all the densities and temperatures considered. This is again an indication that the ground state of neutron matter is paramagnetic, in opposition to what it is found in Ref. for Skyrme-like interactions where, as a consequence of the anomalous behaviour of the entropy, the minimum of F/A is situated at $`0<\mathrm{\Delta }<1`$ and moves to higher polarizations when the temperature increases. It is also interesting to note that the dependence of F/A on the spin polarization is “up to a very good approximation” parabolic. One can try to characterize that dependence in the following simple analytic form:
$$\frac{F}{A}(\rho ,\mathrm{\Delta },T)=\frac{F}{A}(\rho ,0,T)+a(\rho ,T)\mathrm{\Delta }^2$$
(12)
where, assuming the quadratic dependence to be valid up to $`|\mathrm{\Delta }|=1`$ as our results indicate, the value of $`a(\rho ,T)`$ can be easily obtained for each density and temperature as the difference between the total free energies per particle of totally polarized and non-polarized neutron matter
$$a(\rho ,T)=\frac{F}{A}(\rho ,\pm 1,T)\frac{F}{A}(\rho ,0,T).$$
(13)
The magnetic susceptibility can be evaluated then in a very simple way if the parabolic dependence of Eq. (12) is assumed, giving
$$\chi (\rho ,T)=\frac{\mu ^2\rho }{2a(\rho ,T)}.$$
(14)
The ratio $`\chi _F/\chi `$, where $`\chi _F`$ is the magnetic susceptibility of the free Fermi gas, is shown in Fig. 5 as a function of density for several temperatures. Starting from 1, the ratio increases as the density increases at any temperature and no signal of a change of such a trend is expected at higher densities, contrary to the results of Ref. in the case of the Skyrme-like interactions. This is again an indication that a ferromagnetic transition, whose onset would be signaled by the density at which this ratio becomes zero, is not seen and not expected at larger densities either.
Finally, the behaviour of the entropy per particle S/A as a function of the spin polarization at a fixed density $`\rho =0.32`$ fm<sup>-3</sup> for several temperatures is shown in Fig. 6. The entropy, as the free energy, is also symmetric and almost parabolic in $`\mathrm{\Delta }`$. Its maximum is placed at $`\mathrm{\Delta }=0`$ for all the densities and temperatures considered, as one naively expects, contrary to the findings of Ref. . In this reference, it was shown that for a pure parabolic single particle spectrum, as it is the case for the Skyrme interaction, imposing the entropy of the polarized phase to be smaller than the unpolarized one for a given density and temperature, is equivalent to requiring the ratio of the neutron effective masses in the fully polarized and unpolarized phases to be smaller than $`2^{2/3}`$. In the BHF approach, the momentum and temperature dependence of the effective mass prevents from deriving a similar rigorous condition. However, thinking in terms of a value of the effective mass that would characterize the single particle spectrum in average, or considering just the effective mass at the Fermi surface, which is the most relevant for the calculation of the entropy at small temperatures, we can then explore if the BHF calculations respect the condition derived in . In fact, in the case of $`\rho =0.16`$ fm<sup>-3</sup> and $`T=40`$ MeV we find (see Fig. 2) $`m_{}^{}(\mathrm{\Delta }=1)/m_{()}^{}(\mathrm{\Delta }=0)=1.09`$, which is smaller than the limit established in Ref. . This is true for all the densities and temperatures explored in this work and therefore the entropy of the polarized phase is always smaller than that for the unpolarized one.
In summary, we have studied the properties of spin polarized neutron matter both at zero and finite temperature within the framework of the Brueckner–Hartree–Fock formalism. We have determined the single-particle potentials and the effective mass of neutrons with spin up and down for arbitrary values of the density, temperature and spin polarization. We have found that the spin up and spin down effective masses show an almost linear and symmetric variation with respect to their values at spin polarization $`\mathrm{\Delta }=0`$.
We have determined the differences of the free energy (F/A), energy (E/A) and entropy (S/A) per particle between the totally polarized and non-polarized phases. We have found that, in contrast to the results of a similar study with the Skyrme interaction , these differences are always positive for the F/A and E/A which is an indication that the non-polarized phase is energetically favorable, from which we can conclude that a phase transition to a ferromagnetic state is not to be expected. In addition, contrary to the results with the Skyrme interaction, we have found that the difference in the entropy is always negative according to the idea that the totally polarized phase is more “ordered” than the non-polarized one.
Finally, we have seen that both the free energy and the entropy per particle are not only symmetric on the spin polarization but also parabolic in a very good approximation up to $`|\mathrm{\Delta }|=1`$. This finding supports the calculation of the magnetic susceptibility by using only the free energies of the fully polarized and non-polarized phases.
This work is partially supported by DGICYT (Spain) project BFM2002-01868 and by the Generalitat de Catalunya project 2001SGR00064. This research is part of the EU Integrated Infrastructure Initiative Hadron Physics project under contract number RII3-CT-2004-506078. One of the authors (A. Rios) acknowledges the support from DURSI and the European Social Funds.
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# References
PIIPTI, or the
Principle of Increasing Irrelevance of
Preference Type Information
Elemér E Rosinger
Department of Mathematics
University of Pretoria
Pretoria, 0002 South Africa
e-mail : eerosinger@hotmail.com
Abstract
It is shown that in the case of a single decision maker who optimizes several possibly conflicting objectives, the amount of information available in preference relations among pairs of possible decisions, when compared with all other possible information, is tending to zero exponentially with the number of those different objectives. Consequently, in the case of a larger number of conflicting objectives, the only way to obtain a satisfactory amount of information is by the use of non-preference type relations among possible decisions.
1. Introduction
There are three wider areas of decision making, each known to be subject to certain deep structural limitations.
In games, two or more players make decisions, each pursuing his or her own best interest. As is known, Binmore \[1-3\], the complexities involved are not seldom such as to lead to algorithmically unsolvable situations. In this regard it is worth mentioning that, during the late 1940s and early 1950s, when game theory had known a massive interest and development, there was not much awareness about the possibility of the presence of the type of deep difficulties which would more than three decades later be pointed out by Binmore.
In social choice the deeper structural difficulties came to attention relatively early with Arrow’s celebrated impossibility. Here, the issue is the appropriate aggregation of a number of individual preferences. And as it turns out, this in general is not possible, unless there is a ”dictator”.
Social choice can be seen as a partial version of the situation of a single decision maker. Indeed, the aggregation of the set of individual choices amounts to a single decision. On the other hand, even if such a decision is made by a single decision maker, he or she is not supposed to be partial in any way with respect to any of the individual preferences which are aggregated. And yet, Arrow’s impossibility rules the realms of social choice.
The case of a single decision maker facing several conflicting objectives cannot - according to one line of argument - be but more difficult, since he or she is not barred from having specific preferences and expressing them in his or her decision. Consequently, the possible limitations to be faced by a single decision maker may quite likely be more severe than those which lead to Arrow’s impossibility in social choice.
According to another line of argument, however, in the case of one single decision maker, the fact that he or she faces all alone his or her own conflicting objectives gives an easy and natural opportunity for certain cooperative type approaches. After all, cooperation can involve bargaining, and in the case of one single decision maker with several conflicting objectives, he himself, or she herself may end up as if bargaining with himself or herself.
In fact, certain forms of cooperation find a most appropriate context precisely within the thinking of a single decision maker who exhibits a rational behaviour.
As it happens nevertheless, there seems to be little awareness in the literature about the above mentioned issues relating to single decision makers who simultaneously face several conflicting objectives.
Here we present, as one of the two main difficulties facing a single decision maker, what has been named the Principle of Increasing Irrelevance of Preference Type Information, or in short, PIIPTI, see Rosinger .
The other difficulty, mentioned in short in Conclusions, is presented in some detail in Rosinger \[1-5\].
2. A Single Decision Maker with Multiple Conflicting
Objectives
In order to illustrate in detail what is involved, let us consider the following large and practically important class of decision making situations, when the single decision maker SDM has to deal with $`n2`$ typically conflicting objectives given by the utility functions, see von Neumann & Morgenstern, Luce & Raiffa,
(2.1) $`f_1,...,f_n:A`$
and his or her aim is to maximize all of them, taking into account that most often such a thing is not possible simultaneously, due to the conflicts involved.
The main difficulty of this situation is that the SDM is not supposed to have available under any form whatsoever an overall utility function
(2.1) $`f:A`$
which would hopefully synthesize his or her position with respect to those $`n2`$ conflicting objectives in (2.1) taken simultaneously in their totality.
Here, as before, the set $`A`$ describes the available choices, namely, those which the SDM has, and this set $`A`$ may as well be an infinite set, for instance, some open or closed bounded domain in a finite dimensional Euclidean space.
Clearly, the functions $`f_i`$ in (2.1) can be seen as utility functions, and as such, they generate preference relations on the set of choices $`A`$. Namely, the preference relation $`_i`$ corresponding to the utility function $`f_i`$ is defined for $`a,bA`$, by
(2.2) $`a_ibf_i(a)f_i(b)`$
In this way, the problem in (2.1) can be reduced to a choice, according to the natural partial order (2.4) below, of a point in the set of all possible decision outcomes
(2.3) $`B=\{(f_1(a),...,f_n(a))|aA\}^n`$
that is, the set of n-tuples of outcomes $`(f_1(a),...,f_n(a))^n`$ which correspond to various choices $`aA`$ which the SDM can make.
Needless to say, the situation described by (2.1) is not the most general one, since it is possible to encounter cases when the objectives are not given by utility functions, or simply, are not even quantifiable. However, the model in (2.1) can nevertheless offer an edifying enough situation, in order to be able to obtain relevant insights into the nature and extent of the complexities and difficulties which a SDM can face. Furthermore, it can also lead to general enough solution methods, including ways to choose solution concepts, see Rosinger \[1-5\].
Next we give three different arguments supporting PIIPTI. The first and the third ones are of a geometrical nature related to finite dimensional Euclidean spaces. The second argument is of a simple probabilistic-combinatorial kind. Here we should mention that, while the first geometric argument is rather simple and obvious, the other geometric argument, although quite elementary, appears however to be less well known, although it has important connections with Physics.
A First Argument. We start with a very simple geometric fact about finite dimensional Euclidean spaces which can give a good insight into the more involved result in (2.8). On the n-dimensional Euclidean space $`^n`$, with $`n1`$, we consider the natural partial order relation $``$ defined for elements $`x=(x_1,...,x_n),y=(y_1,...,y_n)^n`$, according to
(2.4) $`xyx_iy_i,\text{with}1in`$
Let us denote by
(2.5) $`P_n=\{x^n|x0\}`$
the set of nonnegative elements in $`^n`$, corresponding to the partial order $``$.
Then we can note that, for $`n=1`$, the set $`P_1`$ is half of the space $`^1=`$. Further, for $`n=2`$, the set $`P_2`$ is a quarter of the space $`^2`$. And in general, for $`n1`$, the set $`P_n`$ is $`1/2^n`$ of the space $`^n`$.
It follows that in an n-dimensional Euclidean space $`^n`$, if one is given an arbitrary element $`x^n`$, then the probability for this element $`x`$ to be nonnegative is $`1/2^n`$, thus it tends exponentially to zero with $`n`$.
Consequently, the same happens with the probability that two arbitrary elements $`x,y^n`$ are in the relationship $`xy`$. Indeed, the relationship $`xy`$ is obviously equivalent with $`0yx`$.
This means that one can expect a similar phenomenon to happen when trying to compare points in the set of outcomes in (2.3), which correspond to the $`n`$ conflicting objectives in (2.1).
That very simple geometric fact is, actually, at the root of PIIPTI.
A Second Argument. Let us now assume for the sake of technical simplicity that in (2.1) we have a finite set of choices, namely
(2.6) $`A=\{a_1,...,a_m\},m2`$
A natural single preference relation on $`A`$ corresponding to (2.1), and which may try to synthesize the respective $`n`$ conflicting objectives, should of course be given by a subset
(2.7) $`SA\times A`$
Here, for any $`a,a^{}A`$, the SDM will prefer $`a^{}`$ to $`a`$, in which case we write $`aa^{}`$, or equivalently, $`(a,a^{})S`$, if and only if one has for each objective function $`f_i`$, with $`1in`$, either that $`f_i(a^{})f_i(a)>0`$ and it is not negligible, or $`|f_i(a^{})f_i(a)|`$ is negligible.
Let us therefore see more precisely how much information one single preference relation $`S`$ can carry, when the number $`n`$ of conflicting objectives in (2.1) becomes large, and even if only moderately so. This can be done quite easily by noting that in typical situations, we can have the relation
(2.8) $`\text{car}S/\text{car}(A\times A)=O(1/2^n)`$
where for a finite set $`E`$ we denoted by ”$`\text{car}E`$” the number of its elements.
The proof of (2.8) goes as follows, by using a combinatorial-probabilistic type argument. Let us take any injective function $`g:A`$, and denote by
(2.9) $`S_g=\{(a,a^{})A\times A|g(a)g(a^{})\}`$
which is its corresponding preference relation on $`A`$. Then obviously
(2.10) $`\text{car}S_g=m(m+1)/2`$
Now given any subset $`SA\times A`$, let us denote by $`P(S)`$ the probability that for an arbitrary pair $`(a,a^{})A\times A`$, we have $`(a,a^{})S`$. Then clearly
(2.11) $`P(S_g)=(1+1/m)/2`$
Let us assume about the objective functions in (2.1) the following
(2.12) $`f_1,...,f_n\text{are injective}`$
and furthermore, that their corresponding sets of preferences
(2.13) $`S_{f_1},...,S_{f_n}\text{are probabilisitically independent}`$
Then we obtain, see (2.11)
(2.14) $`P(S_{f_1}...S_{f_n})=P(S_{f_1})...P(S_{f_n})=(1+1/m)^n/2^n`$
And now (2.8) follows, provided that $`n`$ in (2.1) and $`m`$ in (2.6) are such that
(2.15) $`(1+1/m)^n=O(1)`$
which happens in many practical situations.
As for the independence condition (2.13), let us note the following. Let us assume that the objectives $`f_1`$ and $`f_2`$ are such that for $`a,a^{}A`$ we have
(2.16) $`f_1(a)<f_1(a^{})f_2(a)<f_2(a^{})`$
then obviously $`S_{f_1}=S_{f_2}`$, hence (2.14) may fail. But clearly, (2.16) means that $`S_{f_1}`$ and $`S_{f_2}`$ are not independent. In the opposite case, when
(2.17) $`f_1(a)<f_1(a^{})f_2(a^{})<f_2(a)`$
then obviously
(2.18) $`S_{f_1}S_{f_2}=\{(a,a)|aA\}`$
and (2.14) may again fail. However (2.18) once more means that $`S_{f_1}`$ and $`S_{f_2}`$ are not independent, since they are in total conflict with one another.
A Third Argument. For the sake of simplicity, let us assume that the set $`B`$ of outcomes in (2.3) is of the form
(2.19) $`B=\{b=(b_1,...,b_n)^n|\begin{array}{c}b_1,...,b_n0\hfill \\ b_1+...+b_nL\hfill \end{array}\}`$
for a certain $`L>0`$. Then clearly the Pareto maximal, or in other words, the non-dominated subset of $`B`$ is
(2.20) $`B^P=\{b=(b_1,...,b_n)^n|b_1+...+b_n=L\}`$
when considered with the natural partial order (2.4) on $`^n`$.
Now for $`0<ϵ<L`$, the $`ϵ`$-thin shell in $`B`$ corresponding to $`B^P`$ is given by
(2.21) $`B^P(ϵ)=\{b=(b_1,...,b_n)^n|\begin{array}{c}Lϵb_1+...+b_n\hfill \\ L\hfill \end{array}\}`$
And a standard multivariate Calculus argument gives for the volume of $`B`$ in (2.19) the relation
(2.22) $`\text{vol}B=K_nL^n`$
where the constant $`K_n>0`$, involving the Gamma function, does only depend on $`n`$, but not on $`L`$ as well. In this way it is easy to see that
(2.23) $`\text{vol}B^P(ϵ)/\text{vol}B=1(1ϵ/L)^n`$
This leads to a rather counter-intuitive and somewhat paradoxical property of higher dimensional Euclidean spaces.
For instance, in the 20-dimensional case, a shell with a thickness of only 5% of the radius $`L`$ of a sphere will nevertheless contain at least 63% of the total volume of that sphere.
In more simple and direct geometric terms the relation (2.23) means that :
( VOL ) ”The volume of a multidimensional solid is mostly concentrated next to its surface.”
The relevance of this property ( VOL ) to PIIPTI is as follows. The set (2.3), or equivalently (2.19), of outcomes $`B`$ in the multiple objective decision problem (2.1) is the one which determines the choice of the appropriate decision taken in $`A`$, and it does so through the relations (2.2). And obviously, in this respect, only the Pareto maximal, or the non-dominated subset $`B^P`$ of $`B`$, see (2.20), is of relevance. However, within this subset $`B^P`$ no two different points $`uv`$ can be in a relation $`uv`$, see (2.4), this being the very definition of a Pareto maximal, or non-dominated set. In this way all pairs of different elements $`uv`$ in $`B^P`$ are incomparable, thus are outside of being included in a preference relationship.
And as seen in (2.23), the volume of no matter how thin a shell next to $`B^P`$, when compared to that of $`B`$, tends to 1 exponentially with $`n`$ becoming large. Thus, the amount of pairs of different elements $`uv`$ in no matter how small a neighbourhood of $`B^P`$ tends to 1 with $`n`$ becoming large, when compared with all the possible pairs of outcomes in $`B`$. And obviously, any two different elements $`uv`$ in a neighbourhood of $`B^P`$ are comparable, that is, satisfy the relation $`uv`$, see (2.4), only if the are very near to one another, thus they cannot express any kind of a more relevant preference. Otherwise, if two different elements $`uv`$ in a neighbourhood of $`B^P`$ are not near to one another, then they must be incomparable, given the fact that $`B^P`$ is a Pareto maximal, or non-dominated set.
It may be instructive to note that relation (2.23) also has a physical interpretation, as it explains the phenomenon of temperature, see Manin. Indeed, let us assume that a certain simple gas has $`n`$ atoms of unit mass. Then their kinetic energy is given by
(2.24) $`E=\mathrm{\Sigma }_{1in}v_i^2/2`$
where $`v_i`$, with $`1in`$, are the velocities of the respective atoms. Therefore, for a given value of the kinetic energy $`E`$, the state of the gas is described by the vector of $`n`$ velocities, namely
(2.25) $`v=(v_1,...,v_n)S_n(\sqrt{(}2E))`$
where for $`L>0`$, we denoted by
(2.26) $`S_n(L)=\{x=(x_1,...,x_n)^n|x_1^2+...+x_n^2=L^2\}`$
the ( n - 1 )-dimensional surface of the n-dimensional ball with radius $`L`$ in $`^n`$.
Now we can recall that in view of the Avogadro number, under normal conditions for a usual macroscopic volume of gas, one can have
(2.27) $`n>10^{20}`$
Therefore, the above property ( VOL ) which follows from (2.23) is very much manifest. Let us then assume that a small thermometer with a thermal energy $`e`$ negligible compared to $`E`$ is placed in the gas. Then the state (2.25) of the gas will change to a new state
(2.28) $`v=(v_1,...,v_n)S_n(\sqrt{(}2E^{}))`$
However, in view of property ( VOL ), it will follow with a high probability that
(2.29) $`E^{}E`$
And it is precisely this stability or rigidity property (2.29) which leads to the phenomenon of temperature as a macroscopically observable quantity.
3. Conclusions
Situations involving the actions of conscious rational agents are approached in three mathematical theories, namely, the theory of games, the theory of social choice, and decision theory. In games, there are two or more such conscious and rational agents, called players, who are interacting according to the given rules. And except for that, they are free and independent, and there is no overall authority who could influence in any way the players. In social choice, again, there are two or more conscious and rational agents with their given individual preferences. Here however, the issue is to find a mutually acceptable aggregation of those preferences. And such an aggregation is seen as being done by an outsider. Finally, decision theory can be seen as a two person game, in which one of the players is a conscious rational agent, while the other is Nature.
In the case of one single decision maker, who in decision theory is seen as a conscious rational player, playing alone against Nature, one may seem at first two have a situation which enjoys all the advantages that are missing both in games, and in social choice. Indeed, it may at first appear that such a single decision maker does not have to put up with one or more other autonomous players. And also, as the single player, he or she can automatically be seen as a dictator as well, since there is no other conscious agent out there to protest, least of all what is called Nature in such a context.
It would, therefore, appear that in decision theory one has it rather easy.
And yet, in the typical practical situations when the single decision maker is facing multiple and conflicting objectives, all the mentioned seeming advantages are instantly cancelled. Instead, the single decision maker can easily end up feeling as if two or more autonomous agents have moved inside of him or her, and now he or she has to turn into a dictator who, in fact, ends up fighting himself or herself.
In this regard, two facts come to the fore from the beginning in typical situations with a single decision maker facing multiple conflicting objectives, see Rosinger \[1-5\] :
Fact 1. There is no, and there cannot be a unique natural canonical candidate for the very concept of solution. And in fact, the very issue of choosing a solution concept leads to a meta-decision problem which itself has multiple conflicting objectives.
Fact 2. The information contained in the preference structures involved - relative to all other possible, such as for instance, non-preference type information present in the situation - tends exponentially to zero, as the number of conflicting objectives increases. This phenomenon, which in fact is of a very simple higher dimensional geometric nature, can be called the Principle of Increasing Irrelevance of Preference Type Information, or in short PIIPTI.
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# Rapid evaluation of the periodic Green’s function in 𝑑 dimensions
## I Introduction
The Poisson equation is probably one of the most useful equations in physics. In a two-dimensional (2D) space, the periodic solution of this equation corresponds to the solution of particles interacting with the logarithmic interaction, and it has applications in simulations of 2D pancake vortices in high-temperature superconductors tyagilog . In 3D, periodic solutions to the Poisson equation are used in electromagnetism. Here, the solution of the Poisson equation corresponds to a number of charges interacting with the Coulomb potential. This 3D periodic solution is routinely used in most simulations involving charged particles. Recently, the periodic solution of the Poisson equation in higher dimensions has found use in the string theory.
In 1D and 2D, the Green’s function for the Poisson equation for a charge neutral box may be obtained in a closed form. In 3D, one can obtain rapidly converging series representations using well known method by Ewald ewald . The other two approaches for the 3D case were given Lekner lekner and Sperb sperb . However, in higher dimensions, one can either use the Ewald method which has its drawbacks, or use the Jacobi theta function identitiesglasser . In general, there is no efficient way to calculate the Green’s function in a general $`d`$-dimensional space with $`d>3.`$
In this paper, we give an exponentially fast converging series representation for the Green’s function of the Poisson equation in any positive integer dimension. This work will generalize the methods employed for 2D and 3D case tyagipre , and will tie together the different approaches taken by Lekner lekner and Sperb sperb for the especial case of $`d=3`$. The outline of the paper is as follows. In Sec. I we derive expressions giving the Coulomb sum in the $`d`$-dimensional space. In Sec. II we derive recursive relations using the result of the previous section. In Sec. III we discuss the results.
## II Green’s function in $`d`$ dimensions
For simplicity, we consider the case of a unit charge situated within a cubic box in $`d`$ dimensions. The sides of the box are all assumed to be of unit length. From here onwards, we will refer to the box as simulation cell. The basic simulation cell repeats itself in all $`d`$ dimensions. We also assume a charge neutral system. The unit charge interacts with other identical unit charges (for the case of different charges $`q_1`$ and $`q_2`$ one just gets an extra factor of $`q_1q_2`$) situated at the vertices of the periodic structure. The periodic Green’s function in $`d`$ dimensions satisfies the Poisson equation,
$$_d^2G(𝒓)=C_d\underset{𝒍}{}\delta (𝒓+𝒍),$$
(1)
where $`_d^2`$ is the Laplacian operator in $`d`$ dimensions, $`𝒍`$ denotes a $`d`$-dimensional vector, whose components are integers ranging over $`\mathrm{}`$ to $`+\mathrm{},`$ and $`C_d`$ is a dimension-dependent factor. The value of $`C_d`$ for various dimensions is
$$C_d=\{\begin{array}{c}2\text{ }d=1\\ 2\pi \text{ }d=2,\\ 4\pi ^{\nu +1}/\mathrm{\Gamma }\left(\nu \right)\text{ }d>2.\end{array}$$
Here, $`\mathrm{\Gamma }(\nu )`$ stands for the Gamma function, and $`\nu =\left(d2\right)/2`$. We note that with this choice of $`C_d`$ in Eq.(1), the $`G`$ stands for the Coulomb type summation in $`d`$ dimensions. Thus, $`G`$ corresponds to a sum of type $`\left|𝒓\right|`$ in 1D, a logarithmic sum, -$`\mathrm{ln}\left|𝒓\right|`$, in 2D and a sum of type $`\left|𝒓\right|^{(d2)}`$ for a $`d`$-dimensional space with $`d>2`$ . The solution of Eq. (1) diverges, which is a simple consequence of the fact that the interaction energy of a charge with another charge and all its periodic images is infinite. To obtain a meaningful value of $`G`$ we will have to modify Eq. (1) as followsmarshall :
$$^2G_d(𝒓)=C_d\underset{𝒍}{}\delta (𝒓+𝒍)+\frac{C_d}{l_1l_2\mathrm{}l_d}.$$
(2)
The second term in eq.(2) amounts to the presence of a uniform background charge. Thus, for every charge, $`q`$, one may imagine a uniform distribution of charge, such that the total charge per basic simulation cell adds up to $`q`$. For a charge neutral periodic system, imposing these kind of background uniform charge distributions does not matter since the total uniform background charge adds up to zero. However, now a unit charge located within the basic simulation cell at position $`\left\{x_i\right\}`$ not only interacts with a second charge located at the origin and its periodic images, but also interacts with the neutralizing background charge of the second particle. This particular way of introducing the artificial neutralizing background charge leads to only the intrinsic part lekner of the potential energy. We note that once the Green’s function is obtained, the solution of the equation
$$^2V_d=C_d𝝆\left(𝒓\right)$$
under periodic boundary conditions could be simply obtained from
$$V_d=_{\text{cell}}G_d\left(𝒓\mathbf{}𝒓^{}\right)𝝆\left(𝒓^{}\right)𝑑𝒓^{},$$
where $`𝝆`$ is periodic and the simulation cell is overall charge neutral. The rapid evaluation of the $`G_d`$ is discussed in the next section.
The solution of eq.(2) can be written easily in the Fourier space as tyagipre :
$`G_d(x_1,x_2,\mathrm{},x_d)`$ $`={\displaystyle \frac{C_d}{(2\pi )^2}}\times `$
$`\underset{\beta 0}{lim}\left({\displaystyle \underset{\{m\}_d}{}}{\displaystyle \frac{e^{i2\pi (m_1x_1+m_2x_2+\mathrm{}+m_dx_d)}}{\left\{m_1^2+m_2^2+\mathrm{}+m_d^2+\beta ^2/4\pi ^2\right\}}}{\displaystyle \frac{4\pi ^2}{\beta ^2}}\right),`$ (3)
where $`\beta `$ is an infinitesimal parameter which tends to zero. Here , the set $`\{m_{1,d}\}`$ denotes a set of $`d`$ integers $`\{m_1,m_2,\mathrm{},m_d\}`$. Each one of these integers $`m_i`$ runs over $`\mathrm{}`$ to $`+\mathrm{}.`$ Also, $`x_1,x_2,\mathrm{},x_d`$ denote the components of vector $`𝒓_d`$ in $`d`$-dimensions. Due to the periodic boundary conditions, it is sufficient to treat the case where each $`x_i`$ satisfies $`0.5<x0.5`$. The complete expression for the potential has a term arising from the surface contribution. For the 2D case this term turns out to be zero, but for 3D one obtains a contribution from a dipole term deleeuw . At this point, we would recast the Eq. (3) in an alternative form. For that, we use the fact that the solution of
$$\left(^2\beta ^2\right)Q_0(𝒓)=\delta (𝒓)$$
(4)
in $`d`$-dimensional space is given by
$$Q_0(\left|𝒓\right|;\beta )=\frac{1}{\left(2\pi \right)^{\nu +1}}\frac{\beta ^\nu K_\nu \left(\beta \left|𝒓\right|\right)}{r^\nu }.$$
(5)
Thus, the solution of
$$\left(^2\xi ^2\right)Q_d(𝒓;\beta )=C_d\underset{𝒍}{}\delta (𝒓+𝒍)$$
(6)
in $`d`$-dimensional space will be given by
$$Q_d(𝒓;\xi )=\frac{C_d}{\left(2\pi \right)^{\nu +1}}\underset{\left\{m_{1,d}\right\}}{}\left[\xi ^\nu \frac{K_\nu \left(\xi r_{1,d}\right)}{r_{1,d}^\nu }\right],$$
(7)
where
$$r_{1,d}=\left[\underset{i=1}{\overset{d}{}}\left(m_ix_i\right)^2\right]^{1/2}$$
(8)
On the other hand, the solution of Eq. (4) can be written down in the Fourier space easily as
$$Q_d(𝒓;\beta )=\frac{C_d}{\left(2\pi \right)^2}\underset{\{m\}_d}{}\frac{e^{i2\pi (m_1x_1+m_2x_2+\mathrm{}+m_dx_d)}}{\left\{m_1^2+m_2^2+\mathrm{}+m_d^2+\beta ^2/4\pi ^2\right\}}.$$
(9)
Using Eqs.(3) and (4) we see that one can write
$$G_d(x_1,x_2,\mathrm{},x_d)=C_d\underset{\beta 0}{lim}\left(\frac{1}{\left(2\pi \right)^{\nu +1}}\underset{\left\{m_{1,d}\right\}}{}\left[\beta ^\nu \frac{K_\nu \left(\beta r_{1,d}\right)}{r_{1,d}^\nu }\right]\frac{1}{\beta ^2}\right).$$
(10)
A yet another alternative form of $`G_d`$ can be obtained as follows. We can perform one of the $`d`$ sums in Eq.(3) analytically using the formula gradshteyn
$$\underset{i=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{exp}\left(2\pi imx\right)}{m^2+\gamma ^2}=\frac{\pi }{\gamma }\frac{\mathrm{cosh}\left[\pi \gamma \left(12\left|x\right|\right)\right]}{\mathrm{sinh}\left(\pi \gamma \right)}.$$
(11)
Thus, we obtain
$`G_d(x_1,x_2,\mathrm{},x_d)`$ $`={\displaystyle \frac{C_d}{\left(2\pi \right)^2}}\underset{\beta 0}{lim}({\displaystyle \underset{\left\{m_{2,d}\right\}}{}}{\displaystyle \frac{\pi }{\gamma _{\left\{m_{2,d}\right\}}}}{\displaystyle \frac{\mathrm{cosh}\left[\pi \gamma _{\left\{m_{2,d}\right\}}\left(12\left|x_1\right|\right)\right]}{\mathrm{sinh}\left(\pi \gamma _{\left\{m_{2,d}\right\}}\right)}}`$
$`\times \mathrm{exp}\left[2\pi i{\displaystyle \underset{i=2}{\overset{d}{}}}m_ix_i\right]{\displaystyle \frac{1}{\beta ^2}}),`$ (12)
where $`\gamma _{\{m_{2,d}\}}`$ is defined as
$$\gamma _{\left\{m_{2,d}\right\}}=\left(\underset{i=2}{\overset{d}{}}m_i^2+\beta ^2\right)^{1/2}.$$
(13)
For the purpose of taking the limit $`\beta 0,`$ the sum in the first part of Eq. (12) is broken as
$$\underset{\left\{m_{2,d}\right\}}{}=\underset{\left\{m_{2,d}\right\}}{\overset{}{}}+\left(\text{Term with }m_2=0,m_{3=0}\mathrm{},m_d=0\right),$$
(14)
where a prime over the summation sign indicates that the term corresponding to all $`m_i`$ being zero is to be excluded from the summation. This leads to the following representation for $`G_d`$:
$`G_d`$ $`={\displaystyle \frac{C_d}{\left(2\pi \right)^2}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{\overset{}{}}}{\displaystyle \frac{\pi }{\gamma _{\left\{m_{2,d}\right\}}}}{\displaystyle \frac{\mathrm{cosh}\left[\pi \gamma _{\left\{m_{2,d}\right\}}\left(12\left|x_1\right|\right)\right]}{\mathrm{sinh}\left(\pi \gamma _{\left\{m_{2,d}\right\}}\right)}}`$
$`\times \mathrm{exp}\left(2\pi i{\displaystyle \underset{i=2}{\overset{d}{}}}m_ix_i\right)+H_d,`$ (15)
where we have taken the limit $`\beta 0,`$ i.e. we have substituted $`\beta =0`$ in the first part, and $`H_d`$ is given by
$`H_d`$ $`={\displaystyle \frac{C_d}{\left(2\pi \right)^2}}\underset{\beta 0}{lim}\left({\displaystyle \frac{2\pi ^2}{\beta }}{\displaystyle \frac{\mathrm{cosh}\left[\left(1/2\left|x_1\right|\right)\beta \right]}{\mathrm{sinh}\left(\beta /2\right)}}{\displaystyle \frac{4\pi ^2}{\beta ^2}}\right)`$
$`=C_d{\displaystyle \frac{1}{12}}\left(16\left|x_1\right|+6x_1^2\right).`$ (16)
To avoid the bad convergence towards $`x_10,`$ we further modify the summation in the first part of Eq. (15) by using the following trigonometric identity
$$\frac{\mathrm{cosh}(ab)}{\mathrm{sinh}\left(b\right)}=\mathrm{exp}(b)\frac{\mathrm{cosh}(a)}{\mathrm{sinh}(b)}+\mathrm{exp}(a).$$
(17)
Thus, $`G_d`$ can be written as
$$G_d=H_d+J_d+M_d,$$
(18)
where $`H_d`$ is defined in Eq. (16), $`J_d`$ is given by
$`J_d`$ $`={\displaystyle \frac{C_d}{\left(2\pi \right)^2}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{\overset{}{}}}{\displaystyle \frac{\pi }{\gamma _{\left\{m_{2,d}\right\}}}}\mathrm{exp}\left(\pi \gamma _{\left\{m_{2,d}\right\}}\right)`$
$`\times {\displaystyle \frac{\mathrm{cosh}\left[\pi \gamma _{\left\{m_{2,d}\right\}}\left(12\left|x_1\right|\right)\right]}{\mathrm{sinh}\left(\pi \gamma _{\left\{m_{2,d}\right\}}\right)}}\mathrm{exp}\left(2\pi i{\displaystyle \underset{i=2}{\overset{d}{}}}m_ix_i\right),`$ (19)
and
$$M_d=\frac{C_d}{\left(2\pi \right)^2}\underset{\left\{m_{2,d}\right\}}{\overset{}{}}\frac{\pi }{\gamma _{\left\{m_{2,d}\right\}}}\mathrm{exp}\left[2\left|x_1\right|\pi \gamma _{\left\{m_{2,d}\right\}}\right]\mathrm{exp}\left(2\pi i\underset{i=2}{\overset{d}{}}m_ix_i\right).$$
(20)
It is easy to see that Eq. (19) does not have any convergence problem as $`x_1`$tends to zero. Thus, the whole problem has reduced to evaluating the $`M_d`$ term efficiently. This will be done in the next section.
## III Recursive Formulas
In this section we obtain recursive formulas for $`G_d`$ in two different ways, starting with the expressions in Eq. (10) and (18) respectively. The first method, with Eq. (10) as the starting point, will contain Lekner’s results for $`d=3`$ as a special case, while the second method will contain Sperb’s result in 3D as a special case. With the help of Eqs. (10) and (3) we can write
$`G_d(x_1,x_2,\mathrm{},x_d)`$ $`={\displaystyle \frac{C_d}{\left(2\pi \right)^2}}\underset{\beta 0}{lim}(Q_d(x_1,x_2,..,x_d;\beta ){\displaystyle \frac{1}{\beta ^2}})`$
$`={\displaystyle \frac{C_d}{\left(2\pi \right)^2}}\underset{\beta 0}{lim}\left({\displaystyle \underset{\left\{m_{1,d}\right\}}{}}\mathrm{exp}\left(2\pi im_1x\right)\times {\displaystyle \frac{\mathrm{exp}\left(2\pi i_{i=2}^dm_ix_i\right)}{_{i=2}^dm_i^2+\left[\beta ^2+m_1^2\right]}}\right){\displaystyle \frac{1}{\beta ^2}}.`$ (21)
Using the definition of $`G_d`$, Eq. (21) be written as
$`G_d(x_1,x_2,\mathrm{},x_d)`$ $`={\displaystyle \frac{C_d}{C_{d1}}}\underset{\beta 0}{lim}[{\displaystyle \underset{m_1}{}}\mathrm{exp}\left(2\pi im_1x_1\right)`$
$`\times Q_{d1}(x_2,..x_d;\sqrt{\beta ^2+\left(2\pi m_1\right)^2}){\displaystyle \frac{C_{d1}}{\beta ^2}}].`$ (22)
We separate out the term corresponding to $`m_1=0`$ in Eq. (22) so that the limit corresponding to $`\beta `$ can be taken. Thus, we write Eq. (21) as
$`G_d(x_1,x_2,\mathrm{},x_d)`$ $`=2{\displaystyle \frac{C_d}{C_{d1}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)Q_{d1}(x_2,x_3,\mathrm{},x_d;2\pi m_1)`$
$`+{\displaystyle \frac{C_d}{C_{d1}}}\underset{\beta 0}{lim}\left[Q_{d1}(x_2,x_3,\mathrm{},x_d;\beta ){\displaystyle \frac{C_{d1}}{\beta ^2}}\right]`$
$`=2{\displaystyle \frac{C_d}{C_{d1}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)Q_{d1}(x_2,x_3,\mathrm{},x_d;2\pi m_1)`$
$`+{\displaystyle \frac{C_d}{C_{d1}}}G_{d1}(x_2,..,x_d),`$ (23)
where we have taken the limit $`\beta 0`$ in the first term. The Eq. (23) is one of the most important result of this paper. This relates a $`d`$-dimensional sum to a $`\left(d1\right)`$-dimensional sum. This is a recursive relation. If one is able to obtain the Green function for the $`\left(d1\right)`$-dimensional space, one can obtain the Green’s function for the $`d`$-dimensional space. The first term in Eq. (23) can be modified in the following way. We can use a form of $`G_{d1}`$ similar to the one used in Eq. (7) to obtain
$`G_d(x_1,x_2,\mathrm{},x_d)`$ $`=2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{}}\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`\times \left(2\pi m_1\right)^{\nu 1/2}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_1r_{2,d}\right)}{r_{2,d}^{\nu 1/2}}}+{\displaystyle \frac{C_d}{C_{d1}}}G_{d1}(x_2,..,x_d),`$ (24)
where $`\{m_{2,d}\}`$ denotes a sum over sets $`\{m_{2,}m_3,\mathrm{},m_d\}`$ and $`r_{2,d}`$ is defined like Eq. (8)
$$r_{2,d}=\left[\underset{i=2}{\overset{d}{}}\left(m_ix_i\right)^2\right]^{1/2}.$$
(25)
Let us now consider three different cases corresponding to $`d=1,`$ $`d=2`$ and $`d>2.`$ For $`d=1`$ we can evaluate $`G_{d=1}`$ in a closed form:
$`G_1`$ $`={\displaystyle \frac{C_1}{\left(2\pi \right)^2}}\underset{\beta 0}{lim}\left({\displaystyle \underset{m_1}{}}{\displaystyle \frac{\mathrm{exp}\left(2\pi im_1x_1\right)}{\beta ^2+m_1^2}}{\displaystyle \frac{1}{\beta ^2}}\right)`$
$`={\displaystyle \frac{C_1}{\left(2\pi \right)^2}}\underset{\beta 0}{lim}\left({\displaystyle \frac{\pi }{\beta }}{\displaystyle \frac{\mathrm{cosh}\left[\pi \beta \left(12\left|x_1\right|\right)\right]}{\mathrm{sinh}\left[\pi \beta \right]}}{\displaystyle \frac{1}{\beta ^2}}\right)`$
$`=C_1{\displaystyle \frac{1}{12}}\left(16\left|x_1\right|+6x_1^2\right).`$ (26)
Also, the self-energy for this case may be obtained as
$$G_1^{\text{self}}=\underset{x_10}{lim}G_1+\left|x_1\right|=\frac{C_1}{12}.$$
For $`d=2`$ case, we obtain using Eq. (10):
$`G_2(x_1,x_2)`$ $`=2{\displaystyle \frac{C_2}{\left(2\pi \right)^{1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=\mathrm{}}{\overset{+\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)\left(2\pi m_1\right)^{1/2}`$
$`\times {\displaystyle \frac{K_{1/2}\left(2\pi m_1\left|x_2+m_2\right|\right)}{\left|x_2+m_2\right|^{1/2}}}+{\displaystyle \frac{C_2}{C_1}}G_1\left(x_2\right).`$ (27)
Now, using the relation gradshteyn ,
$$K_{1/2}\left(r\right)=\sqrt{\frac{\pi }{2r}}\mathrm{exp}\left(r\right),$$
(28)
we can write
$`G_2(x_1,x_2)`$ $`={\displaystyle \frac{C_2}{2\pi }}{\displaystyle \underset{m_2=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{cos}\left(2\pi m_1x_1\right)}{\left|m_1\right|}}`$
$`\times \mathrm{exp}\left(2\pi m_1\left|x_2+m_2\right|\right)+{\displaystyle \frac{C_2}{C_1}}G_1\left(x_2\right).`$ (29)
The sum over $`m_1`$ can be easily carried out using the identity tyagipre
$`L(x_1,x_2)`$ $`={\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{cos}\left(2\pi m_1x_1\right)}{m_1}}\mathrm{exp}\left(2\pi m_1\left|x_2\right|\right)`$
$`={\displaystyle \frac{1}{2}}\mathrm{ln}\left(12\mathrm{exp}\left[2\pi x_2\right]\mathrm{cos}\left[2\pi x_1\right]+\mathrm{exp}\left[4\pi x_2\right]\right).`$ (30)
Thus, $`G_2`$ can be written as
$`G_2(x_1,x_2)`$ $`={\displaystyle \frac{C_2}{2\pi }}{\displaystyle \underset{m_2=1}{\overset{+\mathrm{}}{}}}L(x_1,\left|x_2+m_2\right|)+L(x_1,\left|x_2m_2\right|)`$
$`+L(x_1,x_2)+{\displaystyle \frac{C_2}{C_1}}G_1\left(x_2\right).`$ (31)
It is also trivial to derive
$$G_2^{\text{self}}=2\frac{C_2}{2\pi }\underset{m_2=1}{\overset{+\mathrm{}}{}}L(0,\left|m_2\right|)\mathrm{ln}2\pi +\frac{C_2}{12}.$$
(32)
Now we consider the case for $`d>2.`$ We can obtain $`G_d`$ from Eq. (24). It is seen that for large arguments the modified Bessel functions decay as
$$K_\nu \left(r\right)\sqrt{\frac{\pi }{2r}}\mathrm{exp}\left(r\right).$$
(33)
As a result, the first term in Eq. (24) decays exponentially. However, one may run into problem if $`r_{2,d\text{ }}`$is very small. In such a case the terms corresponding to $`\{m_{2,d}\}`$ all being zero form a very slowly converging series over $`m_1.`$ This problem of slow convergence when $`r_{2,d}`$ is small can be handled in the following recursive manner. We separate out the particular terms corresponding to $`\{m_{2,d}\}`$ all being zero, and define
$`E_d(x_1,x_2,\mathrm{},x_d)`$ $`=2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{}}\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`\times \left(2\pi m_1\right)^{\nu 1/2}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_1r_{2,d}\right)}{r_{2,d}^{\nu 1/2}}}`$
$`=2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{\overset{}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`\times \left(2\pi m_1\right)^{\nu 1/2}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_1r_{2,d}\right)}{r_{2,d}^{\nu 1/2}}}`$
$`+2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`\times \left(2\pi m_1\right)^{\nu 1/2}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_1r\right)}{r^{\nu 1/2}}}.`$ (34)
Now, we show how to handle the evaluation of $`E_d`$ corresponding to $`d>3.`$ The case for $`d=3`$ will be almost the same. Using the relation gradshteyn (which by the way can be derived from Eq. (21))
$`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left[\left(x+k\right)^2+r^2\right]^{\frac{1}{2}+\nu }}}`$ $`={\displaystyle \frac{\sqrt{\pi }}{\mathrm{\Gamma }\left(\nu +\frac{1}{2}\right)}}\{{\displaystyle \frac{\mathrm{\Gamma }\left(\nu \right)}{r^{2\nu }}}+4\left({\displaystyle \frac{\pi }{r}}\right)^\nu `$
$`\times {\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}l^\nu K_\nu \left(2\pi lr\right)\mathrm{cos}\left(2\pi lx\right)\}\nu >0,`$ (35)
we can write
$`E_d(x_1,x_2,\mathrm{},x_d)`$ $`=2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{\overset{}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`\times \left[\left(2\pi m_1\right)^{\nu 1/2}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_1r_{2,d}\right)}{r_{2,d}^{\nu 1/2}}}\right]`$
$`+{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left[\left(x+k\right)^2+r^2\right]^\nu }}{\displaystyle \frac{\sqrt{\pi }}{\mathrm{\Gamma }\left(\nu \right)}}{\displaystyle \frac{\mathrm{\Gamma }\left(\nu 1/2\right)}{r^{2\nu 1}}}.`$ (36)
Also, the sum over $`k`$ in Eq. (36) can be written as
$`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\left[\left(x+k\right)^2+r^2\right]^\nu }}`$ $`={\displaystyle \frac{1}{\left(x^2+r^2\right)^\nu }}+{\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{1}{\left[\left(x+k\right)^2+r^2\right]^\nu }}+{\displaystyle \frac{1}{\left[\left(xk\right)^2+r^2\right]^\nu }}\right)`$
$`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\nu }{l}}\right)r^{2l}\left[\zeta (2l+2\nu ,N+x)+\zeta (2l+2\nu ,Nx)\right],`$ (37)
where $`N`$ is an arbitrary integertyagipre such that $`N>r+\left|x\right|`$. Using Eqs. (24) , (36) and (37) we can now write
$`G_d(x_1,x_2,..x_d){\displaystyle \frac{1}{\left(x_1^2+r_{}^2\right)^{\nu +1/2}}}`$ $`=2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\left\{m_{2,d}\right\}}{\overset{}{}}}\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`\times \left(2\pi m_1\right)^{\nu 1/2}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_1r_{2,d}\right)}{r_{2,d}^{\nu 1/2}}}`$
$`+{\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{1}{\left[\left(x_1+k\right)^2+r_{}^2\right]^\nu }}+{\displaystyle \frac{1}{\left[\left(x_1k\right)^2+r_{}^2\right]^\nu }}\right)`$
$`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\nu }{l}}\right)r_{}^{2l}\left[\zeta (2l+2\nu ,N+x_1)+\zeta (2l+2\nu ,Nx_1)\right]`$
$`+{\displaystyle \frac{C_d}{C_{d1}}}(G_d(x_2,..x_d){\displaystyle \frac{1}{r_{}^{2\nu }}}).`$ (38)
Note that if $`d=3`$ then instead of Eq.(37) we should use
$`4{\displaystyle \underset{m_1=1}{\overset{\mathrm{}}{}}}K_0\left(2\pi m_1\left(x_2^2+x_3^2\right)^{1/2}\right)\mathrm{cos}\left(2\pi m_1x_1\right)`$
$`=2\left\{\gamma +\mathrm{ln}\left({\displaystyle \frac{\left(x_2^2+x_3^2\right)^{1/2}}{2}}\right)\right\}+{\displaystyle \frac{1}{\sqrt{x_1^2+x_2^2+x_3^2}}}+S(x_1,x_2,x_3),`$ (39)
where
$`S(x_1,x_2,x_3)`$ $`={\displaystyle \underset{n=1}{\overset{N1}{}}}\left({\displaystyle \frac{1}{\sqrt{x_2^2+x_3^2+\left(n+x_1\right)^2}}}+{\displaystyle \frac{1}{\sqrt{x_2^2+x_3^2+\left(nx_1\right)^2}}}\right)`$
$`+{\displaystyle \frac{1}{\sqrt{x_1^2+x_2^2+x_3^2}}}2\gamma \left[\psi (N+x_1)+\psi (Nx_1)\right]`$
$`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{1/2}{l}}\right)\left(x_2^2+x_3^2\right)^l\left[\zeta (2l+1,N+x)+\zeta (2l+1,Nx)\right].`$ (40)
Thus, for the 3D case one would make the following two changes in the expression given in Eq. ( 38). First, there would be an extra term containing $`2\gamma \left[\psi (N+x_1)+\psi (Nx_1)\right]`$ on the right hand side, and second the last term in Eq. (38) would be changed to
$$\frac{C_d}{C_{d1}}\left[G_{d1}(x_2,x_3)+\mathrm{ln}r_{}\right].$$
(41)
Eq. (38) provides us with a general algorithm to calculate $`G_d`$ efficiently in any dimensions. For an example, if we had started out with $`d=10`$, we can obtain $`G_{10}r_8^8`$ by calculating $`G_9r_7^7.`$ Continuing in this fashion we will come down to calculating $`G_2+\mathrm{ln}r_2.`$ Now, this last part $`G_2+\mathrm{ln}r_2`$ has been obtained by several authors . In fact, it can be obtained in a closed form glasser . Thus, we have been able to calculate $`G_{10}r_8^8`$ from which can obtain $`G_{10}`$ by taking the radial part $`r_8^8`$ on the other side. Other forms of $`G_2`$ are given by Gr$`\mathit{}`$nbech-Jensen niels and Tyagi tyagipre . For the sake of completion we write down the result for $`G_2`$:
$`G_2(x_1,x_2)`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{m}{\overset{}{}}}{\displaystyle \frac{\pi }{\left|m\right|}}{\displaystyle \frac{\mathrm{exp}\left(\pi \left|m\right|\right)\mathrm{cosh}\left[2\pi mx_1\right]}{\mathrm{sinh}\left(\pi \left|m\right|\right)}}\mathrm{cos}\left(2\pi mx_2\right)`$
$`{\displaystyle \frac{1}{2}}\mathrm{ln}\left[\mathrm{cosh}\left(2\pi x_1\right)\mathrm{cos}\left(2\pi x_2\right)\right]`$
$`+{\displaystyle \frac{\pi }{6}}\left(1+6x_2^2\right){\displaystyle \frac{\mathrm{ln}\left(2\right)}{2}}.`$ (42)
In the closed form $`G_2`$ is written as logclosed
$$G_2(x_1,x_2)=2\pi \left(\frac{x_2^2}{2}\frac{\mathrm{ln}2}{6\pi }+\frac{1}{2\pi }\mathrm{ln}\left|\frac{\vartheta _1[\pi \left(x_1+ix_2\right),\mathrm{exp}\left(\pi \right)]}{\vartheta _1^{^{}}[0,\mathrm{exp}\left(\pi \right)]^{1/3}}\right|\right),$$
(43)
where $`\vartheta _1`$ represents the Jacobi theta function of the first kind. Also, the self-energy for the 2D case can be obtained from Eq. (42)
$$G_2^{\text{self}}=\frac{1}{\pi }\underset{m=1}{\overset{\mathrm{}}{}}\frac{\pi }{\left|m\right|}\frac{\mathrm{exp}\left(\pi \left|m\right|\right)}{\mathrm{sinh}\left(\pi \left|m\right|\right)}\mathrm{ln}\left(2\pi \right)+\frac{\pi }{6},$$
(44)
or it can be obtained from Eq. (43):
$$G_2^{\text{self}}=\frac{\mathrm{ln}2}{3}\mathrm{ln}\pi \frac{2}{3}\mathrm{ln}\left|\left[\vartheta _1^{^{}}(0,q)\right]\right|.$$
(45)
All three forms Eq. (32) , Eq. (44) and Eq. (45) are equivalent and give numerically the same value for the self-energy. Similarly Eqs. (31), (42) and (43) show perfect agreement.
Now, we give another alternative approach. This time we start with Eq. (18), where $`H_d,`$ $`J_d`$ and $`M_d`$ are defined in Eqs. (16), (19) and (20). $`H_d`$ and $`J_d`$ do not have any convergence problem in the region of interest. We show how to handle $`M_d`$. A recursion formula similar to Eq. (24) can be established for $`M_d`$. It is easy to see just by inspection that $`M_d`$ obeys the following recursion formula:
$`M_d`$ $`={\displaystyle \frac{C_d}{C_{d1}}}M_{d1}+2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_2=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_2x_2\right)\left(2\pi m_2\right)^{\nu 1/2}`$
$`\times {\displaystyle \underset{\{m_{3,d}\}}{}}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_2\sqrt{x_1^2+(m_3x_3)^2+..+(m_dx_d)^2}\right)}{\left[\sqrt{x_1^2+(m_3x_3)^2+..+(m_dx_d)^2}\right]^{\nu 1/2}}},`$ (46)
where $`M_{d1}`$, analogues to Eq. (20), stands for
$$M_{d1}=\frac{C_{d1}}{\left(2\pi \right)^2}\underset{\left\{m_{3,d}\right\}}{\overset{}{}}\frac{\pi }{\gamma _{\left\{m_{3,d}\right\}}}\mathrm{exp}\left[2\left|x_1\right|\pi \gamma _{\left\{m_{3,d}\right\}}\right]\mathrm{exp}\left(2\pi i\underset{i=3}{\overset{d}{}}m_ix_i\right).$$
(47)
In the final step, we break the sum in the second part of Eq. (46) as follows
$$\underset{\{m_{3,d}\}}{}=\underset{\{m_{3,d}\}}{\overset{}{}}+\underset{m_3=0,m_4=0,..}{}.$$
(48)
The term corresponding to $`m_3=0,`$ $`m_4=0\mathrm{}`$ gives rise to a term $`F_d`$ in Eq. (46):
$`F_d`$ $`=2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_2=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_2x_2\right)\left(2\pi m_2\right)^{\nu 1/2}`$
$`\times {\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_2\sqrt{x_1^2+x_3^2+..+x_d^2}\right)}{\left[\sqrt{x_1^2+x_3^2+..+x_d^2}\right]^{\nu 1/2}}}`$
$`={\displaystyle \frac{1}{\sqrt{x_1^2+x_2^2+x_3^2+..+x_d^2}}}{\displaystyle \frac{C_d}{C_{d1}}}{\displaystyle \frac{1}{\sqrt{x_1^2+x_3^2+..+x_d^2}}}`$
$`+{\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{1}{\left[\left(x_2+k\right)^2+r^2\right]^\nu }}+{\displaystyle \frac{1}{\left[\left(x_2k\right)^2+r^2\right]^\nu }}\right)`$
$`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\nu }{l}}\right)r^{2l}\left[\zeta (2l+2\nu ,N+x_2)+\zeta (2l+2\nu ,Nx_2)\right].`$ (49)
Thus, we finally obtain the following recursion relationship for $`M_d`$:
$`\left(M_d{\displaystyle \frac{1}{\left(x_1^2+r_{}^2\right)^{\nu +1/2}}}\right)`$ $`={\displaystyle \frac{C_d}{C_{d1}}}\left(M_{d1}{\displaystyle \frac{1}{r_{}^{2\nu }}}\right)`$
$`+2{\displaystyle \frac{C_d}{\left(2\pi \right)^{\nu +1/2}}}{\displaystyle \underset{m_2=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_2x_2\right)\left(2\pi m_2\right)^{\nu 1/2}`$
$`\times {\displaystyle \underset{\{m_{3,d}\}}{\overset{}{}}}{\displaystyle \frac{K_{\nu 1/2}\left(2\pi m_2\sqrt{x_1^2+\left(m_3x_3\right)^2+\mathrm{}+\left(m_dx_d\right)^2}\right)}{\left[\sqrt{x_1^2+\left(m_3x_3\right)^2+\mathrm{}+\left(m_dx_d\right)^2}\right]^{\nu 1/2}}}`$
$`+{\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{1}{\left[\left(x_2+k\right)^2+r_{}^2\right]^\nu }}+{\displaystyle \frac{1}{\left[\left(x_2k\right)^2+r_{}^2\right]^\nu }}\right)`$
$`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\nu }{l}}\right)r_{}^{2l}\left[\zeta (2l+2\nu ,N+x_2)+\zeta (2l+2\nu ,Nx_2)\right],`$ (50)
where
$$r_{}^2=x_1^2+x_3^2+\mathrm{}+x_d^2.$$
(51)
For $`d=3`$ case, once again, we will have to make two modifications in Eq. (50). With this approach we have obtained Eq. (50), which is analogues to Eq. (38). However, the analysis has become a little bit tedious. The advantage of the second method is that it reduces the computation time, as there is one less summation. The second advantage it can be written down in a product decomposition form. For example, how such a product decomposition form may be written, one may consult Sperb, where a special case corresponding to $`d=3`$ is considered. In general, the procedure of dimensional reduction is to be continued until we have $`M_1`$ on the left hand side. It is clear that $`M_1=0.`$ Let us again consider three special cases. For $`d=1`$ one only has $`H_{d=1}`$ and thus $`G_1=H_1.`$ For $`d=2`$ one obtains
$`J_2`$ $`={\displaystyle \frac{C_2}{\left(2\pi \right)^2}}{\displaystyle \underset{m_2}{\overset{}{}}}{\displaystyle \frac{\pi }{\gamma _{m_2}}}\mathrm{exp}\left(\pi \gamma _{m_2}\right)`$
$`\times {\displaystyle \frac{\mathrm{cosh}\left[\pi \gamma _{m_2}\left(12\left|x_1\right|\right)\right]}{\mathrm{sinh}\left(\pi \gamma _{m_2}\right)}}\mathrm{exp}\left(2\pi im_2x_2\right),`$ (52)
and $`M_2`$ from Eq. (20) and (30) turns out to be just $`L(x_1,x_2)`$
$$M_2=L(x_1,x_2).$$
(53)
Combing $`H_2,`$ $`J_2`$ and $`M_2`$ we obtain the form of $`G_2`$ given in Eq. (42). Considering finally the case for $`d>2`$ case, we can obtain $`G_d`$ again from Eq. (18). Now $`K_d`$ and $`H_d`$ are convergent and $`M_d`$ can be obtained using the recursive relation Eq. (50). For example:
$`\left(M_3{\displaystyle \frac{1}{\left(x_1^2+r_{}^2\right)^{1/2}}}\right)`$ $`={\displaystyle \frac{C_2}{C_1}}\left[M_2+\mathrm{ln}\left(r_{}\right)\right]+2{\displaystyle \frac{C_2}{\left(2\pi \right)^{1/2}}}{\displaystyle \underset{m_2=1}{\overset{\mathrm{}}{}}}\mathrm{cos}\left(2\pi m_2x_2\right)`$
$`\times {\displaystyle \underset{m_3}{\overset{}{}}}K_0\left(2\pi m_2\sqrt{x_1^2+\left(m_3x_3\right)^2}\right)`$
$`+{\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{1}{\left[\left(x_2+k\right)^2+r_{}^2\right]^{1/2}}}+{\displaystyle \frac{1}{\left[\left(x_2k\right)^2+r_{}^2\right]^{1/2}}}\right)`$
$`2\gamma \left[\psi (N+x_2)+\psi (Nx_2)\right]`$
$`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{\nu }{l}}\right)r_{}^{2l}\left[\zeta (2l+2\nu ,N+x_2)+\zeta (2l+2\nu ,Nx_2)\right],`$ (54)
where $`M_2`$ has already been evaluated above. We see that in all the case, expression could be written in a form that the essential Coulomb singularity as the two charges approach each other has been removed.
## IV Conclusions
Using the limiting behavior of the modified Bessel functions, we showed how conditionally convergent Coulomb sums may be handled in an elegant way. We gave two representations of the Green’s function for the Poisson equation in any integer dimensional space. A recursive method was derived that can be applied for wholly periodic cases, as well as for those cases where one may have open boundary conditions along one of the directions. The method may be extended to cover the case where any number of directions may be open. The formulas obtained show rapid convergence in all part of the simulation cell. This method is general enough that it can be easily generalized for a higher dimensional “triclinic” cell. A particular case of the application of this method for a triclinic cell can be seen in a recent papertyagijcp . We have shown that the present work generalizes the work of several authors on periodic and partial periodic systems lekner ; sperb ; mazar . To our knowledge, this treatment is the first of its kind ever taken in a dimension higher than $`d=3.`$
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# USTC-ICTS-05-06 Diquarks in Nonaquark States
## I Introduction
There is growing interest in exotic hadrons, which may open new windows for understanding the hadronic structures and QCD at low energies. Recently, the KEK-PS reported an tribaryon state $`S^0(3115)`$kek1 in the reaction
$$K^{}+^4HeS^0+p$$
(1)
The mass of the state is $`3117_{4.4}^{+1.5}`$MeV , the decay width $`\mathrm{\Gamma }_{S^0}<`$ 21 MeV, and the main decay mode is $`\mathrm{\Sigma }NN`$ rather than $`\mathrm{\Lambda }NN`$. The peak in the proton spectrum is over the background with a significance level 13$`\sigma `$. A strange tribaryon $`S^+(3140)`$ of charge +1 was also reported in the reaction $`K^{}+^4HeS^++n`$ kek2 . The mass and decay width of this exotic state are $`M_{S^+}=3141\pm 3(stat.)_1^{+4}(sys.)`$ MeV and $`\mathrm{\Gamma }_{S^+}23`$ MeV , which is about 25 MeV higher than $`S^+(3115)`$, and its significance is 3.7$`\sigma `$. It also dominantly decay into $`\mathrm{\Sigma }NN`$ rather than $`\mathrm{\Lambda }NN`$.
The $`S^0(3115)`$ was first predicted by Akaishi and Yamazaki first as a deeply-bound kaonic state. Since the discovery of $`S^0(3115)`$ and $`S^+(3140)`$, there has been some theoretical discussion p1 ; p2 ; p3 , in Ref p1 these exotic states are mainly analysized from MIT bag model, and they are identified as kaonic bound state in Ref p2 ; p3 . Since the quark dynamics could be regarded as a cornerstone for hadron physics, it is interesting to investigate the nonaquark states by means of the quark models. Various quark models have been used and proposed in studying the pentaquark baryon state jaffe ; lipkin ; kim . Here we draw the spirit of Jaffe– Wilczek’s workjaffe , since there are evidences for strong diquark correlation in the baryon spectrum diquark , and, especially, in the light nonet-scalar ($`J^{PC}=0^{++}`$) meson spectrum. Their masses are generally below 1000 MeV ($`f_0(600),f_0(980),a_0(980),\kappa (800)`$), and they do not favor the predictions of the q$`\overline{\mathrm{q}}`$-models, but favor the diquark-antidiquark’s quite well. Diquark is a boson with color $`\overline{3}_c`$, flavor $`\overline{3}_\mathrm{f}`$, and spin zero. Diquark correlation is also the basis of color superconductivity in dense quark matter which has not being observed experimentally. This configuration is favored by one gluon exchange georgi ; jaff2 and by instanton interactions hooft ; instanton . It may play important role in the exotic hadron physics. In this paper we try to investigate nonaquark baryons by means of diquarks model, and to learn what happens in the nonaquark case due to the strong diquark correlation. Meanwhile in order to understand the decay of nonaquark states, we suggest a decay mechanism which can qualitatively explain the experiments and give us new predictions. This decay mechanism is quite intuitive. To understand the structure of nonaquark, its mass spectrum and the decay mechanism are the main aims of this paper.
The paper is organized as follows, in the section II we study the direct products of two diquarks states, four diquarks states , and of four diquarks’s plus one quark’s. The irreducible tensors of the allowed nonaquark states are derived. The flavor wave functions are given by identifying the $`SU(3)`$ tensors with the physical tribaryon states. In Section III, the mass spectrum is derived by using the Gell-Mann–Okubo mass formula. The Section IV devotes to study the decays of $`S^0(3115)`$ and $`S^+(3140)`$ under the assumption that the decays caused by a ”fall-apart” mechanism. We find when $`S^0(3115)`$ only belongs to a certain $`\mathrm{𝟐𝟕}`$-plet, its main decay mode is $`\mathrm{\Sigma }NN`$ rather than $`\mathrm{\Lambda }NN`$, and $`S^+(3140)`$ can belongs to either $`\mathrm{𝟐𝟕}`$-plet or $`\overline{\mathrm{𝟑𝟓}}`$-plet. In Section V, we briefly summary the results and give some discussions.
## II The Flavor Wave Function of Nonaquark States
Since the diquark is in the $`\overline{3}_\mathrm{f}`$, it has three configurations in flavor space, which are shown in Fig1. We denote them as
$`QQ^1={\displaystyle \frac{1}{\sqrt{2}}}[d,s]`$
$`QQ^2={\displaystyle \frac{1}{\sqrt{2}}}[s,u]`$
$`QQ^3={\displaystyle \frac{1}{\sqrt{2}}}[u,d]`$ (2)
where $`u`$, $`d`$, $`s`$ are respectively up quark,down quark and strange quark. It is obvious that there are following correspondences: $`QQ^1\overline{u},QQ^2\overline{d},QQ^3\overline{s}.`$
The tensors $`T_{i_1,i_2\mathrm{}i_p}^{j_1,j_2,\mathrm{}j_q}`$ which are the bases for irreducible representations of $`SU(3)`$ are totally symmetric to both all q upper indices and all p low indices, and also are traceless,
$`T_{i_1,i_2\mathrm{}i_p}^{j_1,j_2,\mathrm{}j_q}`$ $`=`$ $`T_{i_1,i_2\mathrm{}i_p}^{j_2,j_1,\mathrm{}j_q}=T_{i_2,i_1\mathrm{}i_p}^{j_1,j_2,\mathrm{}j_q}`$
$`T_{i_1,i_2\mathrm{}i_p}^{i_1,j_2,\mathrm{}j_q}`$ $`=`$ $`0.`$ (3)
Since $`\delta _j^i,\epsilon ^{ijk}\mathrm{and}\epsilon _{\mathrm{ijk}}`$ are tensors, we can use them to raise, low or contract indices when we construct new tensors that are bases of irreducible representation from the direct product tensor. The direct product of two diquarks is
$$QQ^iQQ^j=\frac{1}{\sqrt{2}}S^{ij}+\frac{1}{2\sqrt{2}}\epsilon ^{ijk}T_k,$$
(4)
with $`S^{ij}=\frac{1}{\sqrt{2}}(QQ^iQQ^j+QQ^jQQ^i)`$, $`A^{ij}=\frac{1}{\sqrt{2}}(QQ^iQQ^jQQ^jQQ^i)`$, and $`T_k=\epsilon _{ijk}A^{ij}`$. So the decomposition of the direct product of two diquarks is $`\overline{3}\overline{3}=\overline{6}3`$, which is shown in Fig2. in the Young tabular.
Since the two diquarks can be decomposed into $`\overline{\mathrm{𝟔}}`$ plus $`\mathrm{𝟑}`$, the direct product of four diquarks raises $`\overline{\mathrm{𝟔}}\overline{\mathrm{𝟔}}`$, $`\overline{\mathrm{𝟔}}\mathrm{𝟑}`$, $`\mathrm{𝟑}\overline{\mathrm{𝟔}}`$ and $`\mathrm{𝟑}\mathrm{𝟑}`$. And the corresponding Young tabular is shown in Fig3. It is straightforward that
$`(QQ^iQQ^j)(QQ^mQQ^n)`$
$`={\displaystyle \frac{1}{2}}S^{ij}S^{mn}+{\displaystyle \frac{1}{4}}(\epsilon ^{kij}T_kS^{mn}+\epsilon ^{kmn}S^{ij}T_k)+{\displaystyle \frac{1}{8}}\epsilon ^{kij}\epsilon ^{lmn}T_kT_l`$
$`={\displaystyle \frac{1}{2\sqrt{6}}}T^{ijmn}+{\displaystyle \frac{1}{4\sqrt{2}}}(\epsilon ^{ajm}\delta _b^n\delta _c^i+\epsilon ^{ain}\delta _b^m\delta _c^j)S_a^{bc}+{\displaystyle \frac{1}{2\sqrt{6}}}(\epsilon ^{aim}\epsilon ^{bjn}+\epsilon ^{ajm}\epsilon ^{bin})T_{ab}`$
$`+{\displaystyle \frac{1}{4}}\epsilon ^{kij}[\stackrel{~}{T}_k^{mn}+{\displaystyle \frac{1}{\sqrt{2}}}(\delta _k^m\delta _a^n+\delta _k^n\delta _a^m)\stackrel{~}{Q}^a]+{\displaystyle \frac{1}{4}}\epsilon ^{kmn}[T_k^{ij}+{\displaystyle \frac{1}{\sqrt{2}}}(\delta _k^i\delta _a^j+\delta _k^j\delta _a^i)Q^a]`$
$`+{\displaystyle \frac{1}{4}}\epsilon ^{kij}\epsilon ^{lmn}(\sqrt{2}S_{kl}+\epsilon _{kla}T^a)`$ (5)
where the tensors in the above formula are defined as followings.
$$T^i=\frac{1}{8}\epsilon ^{ijk}(T_jT_kT_kT_j)$$
(6)
$$Q^i=\frac{1}{\sqrt{8}}S^{ij}T_j$$
(7)
$$\stackrel{~}{Q}^i=\frac{1}{\sqrt{8}}T_jS^{ji}$$
(8)
$$S_{ij}=\frac{1}{4\sqrt{2}}(T_iT_j+T_jT_i)$$
(9)
$$T^{ijmn}=\frac{1}{\sqrt{6}}(S^{ij}S^{mn}+S^{mj}S^{in}+S^{in}S^{jm}+S^{mi}S^{jn}+S^{jn}S^{im}+S^{mn}S^{ij})$$
(10)
$$S_i^{jk}=\frac{1}{\sqrt{2}}\epsilon _{imn}(S^{jm}S^{kn}+S^{km}S^{jn})$$
(11)
$$\stackrel{~}{T}_i^{jk}=T_iS^{jk}\frac{1}{\sqrt{2}}(\delta _i^j\delta _m^k+\delta _i^k\delta _m^j)\stackrel{~}{Q}^m$$
(12)
$$T_i^{jk}=S^{jk}T_i\frac{1}{\sqrt{2}}(\delta _i^j\delta _m^k+\delta _i^k\delta _m^j)Q^m$$
(13)
So the four diquarks product can be decomposed into $`\overline{\mathrm{𝟏𝟓}}_1\overline{\mathrm{𝟏𝟓}}_2(3)\overline{\mathrm{𝟑}}(3)\mathrm{𝟔}`$, where the numbers in the parentheses denote the degeneracy in each multiplet. The tensors corresponding to $`\overline{\mathrm{𝟑}}`$ are $`T^i,Q^i,\stackrel{~}{Q}^i`$, the tensor $`S_{ij}`$ form the bases of the irreducible representation $`\mathrm{𝟔}`$, the tensor corresponding to $`\overline{\mathrm{𝟏𝟓}}_1`$ is $`T^{ijmn}`$, and the tensors $`S_i^{jk},\stackrel{~}{T}_i^{jk},T_i^{jk}`$ are respectively the bases of the irreducible representation $`\overline{\mathrm{𝟏𝟓}}_2`$.
### II.1 Nonaquark states
Since quark is in the fundamental representation $`\mathrm{𝟑}`$, when the four diquarks form the irreducible representation $`\overline{\mathrm{𝟑}}`$, the nonaquark state must be in the representation $`\overline{\mathrm{𝟑}}\mathrm{𝟑}=\mathrm{𝟖}+\mathrm{𝟏}`$. This means the nonaquark state can either in the octet or in the singlet. We use $`𝒯^i`$ to stand for $`T^i,Q^i,\stackrel{~}{Q}^i`$, then the tensor product $`𝒯^iq_n`$ can be decomposed as follows
$$𝒯^iq_n=\sqrt{2}(P_n^i+\frac{1}{\sqrt{3}}\delta _n^iS),$$
(14)
where $`S=\frac{1}{\sqrt{6}}𝒯^mq_m`$, $`P_n^i=\frac{1}{\sqrt{2}}(𝒯^iq_n\sqrt{\frac{2}{3}}\delta _n^iS)`$. $`P_n^i`$ stand for the nonaquark octet, and that $`S`$ stands for the nonaquark singlet.
When the four diquarks are in the representation $`\overline{\mathrm{𝟔}}\overline{\mathrm{𝟔}}(=\overline{\mathrm{𝟏𝟓}}_1\overline{\mathrm{𝟏𝟓}}_2\mathrm{𝟔}`$ (see Fig. 3)), they can form the irreducible representative of $`\mathrm{𝟔}`$. Since $`\mathrm{𝟔}\mathrm{𝟑}=\mathrm{𝟏𝟎}\mathrm{𝟖}`$, the nonaquark states can be in decuplet or octet.
$$S_{ij}q_n=\frac{1}{\sqrt{3}}[T_{ijk}+\epsilon _{mjn}P_i^m+\epsilon _{min}P_j^m]$$
(15)
with $`T_{ijk}=\frac{1}{\sqrt{3}}[S_{ij}q_n+S_{in}q_j+S_{jn}q_i]`$, $`P_i^j=\frac{1}{\sqrt{3}}\epsilon ^{jab}S_{ia}q_b`$ and $`T_{ijk},P_i^j`$ respectively correspond to the nonaquark decuplet and octet.
Again, for the four-diquarks states in $`\overline{\mathrm{𝟔}}\overline{\mathrm{𝟔}}`$, they can also form $`\overline{\mathrm{𝟏𝟓}}_1`$. The direct product $`\overline{\mathrm{𝟏𝟓}}_1\mathrm{𝟑}`$ can be reduced as follows
$$T^{ijkl}q_n=T_n^{ijkl}+\frac{1}{\sqrt{6}}(\delta _n^i\delta _b^j\delta _c^k\delta _d^l+\delta _n^j\delta _b^i\delta _c^k\delta _d^l+\delta _n^k\delta _b^i\delta _c^j\delta _d^l+\delta _n^l\delta _b^i\delta _c^j\delta _d^k)D^{bcd}$$
(16)
where $`T_n^{ijkl}=T^{ijkl}q_n\frac{1}{\sqrt{6}}(\delta _n^iD^{jkl}+\delta _n^jD^{ikl}+\delta _n^lD^{ijk}),D^{ijk}=\frac{1}{\sqrt{6}}T^{ijkn}q_n`$, $`T_n^{ijkl}`$ and $`D^{ijk}`$ respectively mean that the nonaquark state belong to $`\overline{\mathrm{𝟑𝟓}}`$plet and $`\overline{\mathrm{𝟏𝟎}}`$-plet.
Finally, we consider the case of that the four-diquarks states form $`\overline{\mathrm{𝟏𝟓}}_2`$. The $`\overline{\mathrm{𝟏𝟓}}_2`$ four-diquarks states have three irreducible representatives: they can be in the product of $`\overline{\mathrm{𝟔}}\mathrm{𝟑}`$, $`\mathrm{𝟑}\overline{\mathrm{𝟔}}`$, or $`\overline{\mathrm{𝟔}}\overline{\mathrm{𝟔}}`$. The corresponding tensors are $`S_i^{jk}`$,$`\stackrel{~}{T}_i^{jk}`$ or $`T_i^{jk}`$ respectively. Using $`𝒯_i^{jk}`$ to denote each one of them, then we have
$$𝒯_i^{jk}q_n=\sqrt{2}T_{in}^{jk}+\sqrt{\frac{2}{3}}\epsilon _{inm}D^{mjk}+\frac{4}{\sqrt{15}}(\delta _n^k\delta _m^j+\delta _m^k\delta _n^j)P_i^m\frac{1}{\sqrt{15}}(\delta _i^k\delta _m^j+\delta _m^k\delta _i^j)P_n^m,$$
(17)
where
$`P_j^i`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{15}}}𝒯_j^{ik}q_k`$
$`D^{ijk}`$ $`=`$ $`{\displaystyle \frac{1}{24}}(\epsilon ^{jab}𝒯_a^{km}q_b+\epsilon ^{kab}𝒯_a^{jm}q_b+\epsilon ^{mab}𝒯_a^{jk}q_b)`$
$`T_{in}^{jk}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}(𝒯_i^{jk}q_n+𝒯_n^{jk}q_i){\displaystyle \frac{\sqrt{30}}{20}}(\delta _i^kP_n^j+\delta _n^kP_i^j+\delta _i^jP_n^k+\delta _n^jP_i^k).`$ (18)
$`T_{in}^{jk}`$ represent the nonaquark $`\overline{35}`$-plet, and $`D^{ijk},P_j^i`$ correspond to nonaquark $`\overline{10}`$-plet and octet. The Young tabular of the tensor decomposit ion are shown in Fig4.
Now we take the color symmetry $`SU(3)^{color}`$ into the account. When the $`SU(3)^{flavor}\times SU(3)^{color}\times SU(2)^{spin}`$ serves as an exact symmetry, then due to the boson statistics the full combined (flavor $`\times `$ color $`\times `$ spin $`\times `$ space) wave-functions of four diquarks states must be symmetric. Since the spin of the diquark is zero, the spin-wave-function is trivially to be symmetric. Thus, the (space $`\times `$ flavor $`\times `$ color) wave function must be symmetric. Under this constraint, only two choices are available:
1)The space-wave-functions are symmetric: In this case, the relative angular momentum $`\mathrm{}`$ is even and starts from $`\mathrm{}=0`$, and then the (flavor $`\times `$ color) wave-functions are symmetric. They must therefore belong to the $`\overline{495}`$ dimensional irreducible representation of SU(9)$``$ SU(3)<sup>flavor</sup>$`\times `$ SU(3)<sup>color</sup>. The reduction of the SU(9) irreducible representation $`\overline{495}`$ with respect to SU(3)<sup>flavor</sup>$`\times `$ SU(3)<sup>color</sup> group is
$$\overline{495}=(\overline{15}_1,\overline{15}_1)+(\overline{15}_2,\overline{15}_2)+(\overline{3},\overline{3})+(6,6).$$
(19)
We can see only the four diquarks’s states $`(flavor,color)=(\overline{3},\overline{3})`$ can combine with the ninth quark with $`(flavor,color)=(3,3)`$ to form a color singlet. Because $`\overline{3}3=81`$, there may exist nonaquark nonet with positive parity. But the states $`S^0(3115)`$ and $`S^+(3140)`$ can not belong to this nonet, because the hypercharge of these two states is $`Y=2`$ , while the maximum of hypercharge of the nonet is 1 .
2) The space-wave-functions are antisymmetric: This means that the relative angular momentum $`\mathrm{}`$ is odd and starts from $`\mathrm{}=1`$, and then the (flavor $`\times `$ color) wave-functions are antisymmetric. And they must belong to the 126 dimensional irreducible representation of SU(9)$``$ SU(3)<sup>flavor</sup>$`\times `$ SU(3)<sup>color</sup>. The reduction of the SU(9) irreducible representation 126 with respect to SU(3)<sup>flavor</sup>$`\times `$ SU(3)<sup>color</sup> is
$$126=(6,6)+(\overline{15}_2,\overline{3})+(\overline{3},\overline{15}_2)$$
(20)
So only when the four diquarks in the state $`(flavor,color)=(\overline{15}_2,\overline{3})`$, they can combine with the ninth quark to form a color singlet hadron. Because $`\overline{15}_23=27\overline{10}8`$, the nonaquark can only in flavor multiplet 27-plet, $`\overline{10}`$-plet, 8-plet, and the nonaquark state can not be flavor 35-plet. This is a rigorous result when the flavor symmetry is exactly. Considering, however, that $`SU(3)^{flavor}`$ is an approximative symmetry which will lead to $`SU(3)^{flavor}\times SU(3)^{color}\times SU(2)^{spin}`$ to be approximate, so we can not completely rule out 35-plet.
From the quantum number of $`S^0(3115)`$ and $`S^+(3140)`$, they possibly belong to the $`\mathrm{𝟐𝟕}`$-plet or $`\overline{\mathrm{𝟑𝟓}}`$-plet. In the exact $`SU(3)^{flavor}\times SU(3)^{color}\times SU(2)^{spin}`$ limit, both $`S^0(3115)`$ and $`S^+(3140)`$ belong to 27-plet whose lowest angular momentum is $`\mathrm{}=1`$ (they are P-wave states), and the corresponding weight diagram is Fig6. The Fig7 is the weight diagram for 35-plet. In the Figures, we show the names of these exotic states, with the subscripts that are the representation-dimensions and the isospin of the particle. Their superscript is the charge of the state.
In tensor representations, the number of lower indices of $`T_{i_1,\mathrm{},i_p}^{j_1,\mathrm{},k_q}`$ is $`p`$ and that of upper indices is $`q`$. Now we suppose that among its lower indices the numbers of 1, 2, and 3 are $`p_1`$, $`p_2`$, and $`p_3`$, respectively, and that among upper indices it has $`q_1`$ 1, $`q_2`$ 2, and $`q_3`$ 3. Then we have $`p_1+p_2+p_3=p`$ and $`q_1+q_2+q_3=q`$. The irreducible tensor is an eigenstate of hypercharge $`Y`$ and the third component of isospin $`I_3`$ with the eigenvalues lowb ; closeb ; inpp
$`Y`$ $`=`$ $`p_1q_1+p_2q_2{\displaystyle \frac{2}{3}}(pq)`$
$`I_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p_1q_1){\displaystyle \frac{1}{2}}(p_2q_2).`$ (21)
The charge of the particle is obtained from the Gell-Mann–Nishijima formula, $`Q=I_3+Y/2`$. By this way, we can match the SU(3) tensors to the physical baryon states.
### II.2 The wavefunction of Nonaquark $`\overline{\mathrm{𝟐𝟕}}`$-plet
It is straightforward to write out the wave functions of nonaquark $`\overline{\mathrm{𝟐𝟕}}`$-plets in the flavor space by means of the irreducible representation tensors. The $`S_{27,1}^+`$ and $`S_{27,1}^0`$ read
$$S_{27,1}^+=\frac{1}{\sqrt{2}}T_{12}^{33}=\frac{1}{4}(𝒯_1^{33}q_2+𝒯_2^{33}q_1)$$
$$S_{27,1}^0=\frac{1}{2}T_{22}^{33}=\frac{1}{2\sqrt{2}}𝒯_2^{33}q_2,$$
(22)
where $`𝒯_k^{ij}`$ stands for the tensors of four-diquarks states. We now provide the explicit expressions of $`S_{27,1}^+`$ and $`S_{27,1}^0`$ for the each irreducible representatives of $`\overline{\mathrm{𝟏𝟓}}_2`$ in order:
1. The case of $`\overline{\mathrm{𝟏𝟓}}_2\overline{\mathrm{𝟔}}\overline{\mathrm{𝟔}}`$: In this case $`𝒯_i^{jk}`$ is $`S_i^{jk}`$ which has been defined in Eq.(11). So,
$`S_1^{33}={\displaystyle \frac{1}{2\sqrt{2}}}([u,d][s,u][u,d][u,d]+[s,u][u,d][u,d][u,d][u,d][u,d][u,d][s,u]`$
$`[u,d][u,d][s,u][u,d]),`$ (23)
$`S_2^{33}={\displaystyle \frac{1}{2\sqrt{2}}}([u,d][u,d][u,d][d,s]+[u,d][u,d][d,s][u,d][u,d][d,s][u,d][u,d]`$
$`[d,s][u,d][u,d][u,d]),`$ (24)
the wave function of $`S_{27_1,1}^+`$ and $`S_{27_1,1}^0`$ are
$`S_{27_1,1}^+`$ $`=`$ $`{\displaystyle \frac{1}{4}}(S_1^{33}q_2+S_2^{33}q_1)`$ (25)
$`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}}}([u,d][s,u][u,d][u,d]d+[s,u][u,d][u,d][u,d]d[u,d][u,d][u,d][s,u]d`$
$`[u,d][u,d][s,u][u,d]d+[u,d][u,d][u,d][d,s]u+[u,d][u,d][d,s][u,d]u`$
$`[u,d][d,s][u,d][u,d]u[d,s][u,d][u,d][u,d]u);`$
$`S_{27_1,1}^0`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}S_2^{33}q_2`$ (26)
$`=`$ $`{\displaystyle \frac{1}{8}}([u,d][u,d][u,d][d,s]d+[u,d][u,d][d,s][u,d]d[u,d][d,s][u,d][u,d]d`$
$`[d,s][u,d][u,d][u,d]d).`$
2. The case of $`\overline{\mathrm{𝟏𝟓}}_2\mathrm{𝟑}\overline{\mathrm{𝟔}}`$: The $`𝒯_i^{jk}`$ is $`\stackrel{~}{T}_i^{jk}`$ which is defined in Eq.(12), and
$$\stackrel{~}{T}_1^{33}=T_1S^{33}=\frac{1}{2}([s,u][u,d][u,d][u,d][u,d][s,u][u,d][u,d])$$
(27)
$$\stackrel{~}{T}_2^{33}=T_2S^{33}=\frac{1}{2}([u,d][d,s][u,d][u,d][d,s][u,d][u,d][u,d])$$
(28)
then, the wave functions of $`S_{27_2,1}^+`$ and $`S_{27_2,1}^0`$ are
$`S_{27_2,1}^+`$ $`=`$ $`{\displaystyle \frac{1}{4}}(\stackrel{~}{T}_1^{33}q_2+\stackrel{~}{T}_2^{33}q_1)`$ (29)
$`=`$ $`{\displaystyle \frac{1}{8}}([s,u][u,d][u,d][u,d]d[u,d][s,u][u,d][u,d]d+[u,d][d,s][u,d][u,d]u`$
$`[d,s][u,d][u,d][u,d]u);`$
$$S_{27_2,1}^0=\frac{1}{2\sqrt{2}}\stackrel{~}{T}_2^{33}q_2=\frac{1}{4\sqrt{2}}([u,d][d,s][u,d][u,d]d[d,s][u,d][u,d][u,d]d).$$
(30)
3. The case of $`\overline{\mathrm{𝟏𝟓}}_2\overline{\mathrm{𝟔}}\mathrm{𝟑}`$: In this case $`𝒯_i^{jk}`$ is the tensor $`T_i^{jk}`$ defined in Eq.(13), obviously
$$T_1^{33}=S^{33}T_1=\frac{1}{2}([u,d][u,d][s,u][u,d][u,d][u,d][u,d][s,u])$$
(31)
$$T_2^{33}=S^{33}T_2=\frac{1}{2}([u,d][u,d][u,d][d,s][u,d][u,d][d,s][u,d])$$
(32)
The wave functions of the two states are
$`S_{27_3,1}^+`$ $`=`$ $`{\displaystyle \frac{1}{4}}(T_1^{33}q_2+T_2^{33}q_1)`$ (33)
$`=`$ $`{\displaystyle \frac{1}{8}}([u,d][u,d][s,u][u,d]d[u,d][u,d][u,d][s,u]d+[u,d][u,d][u,d][d,s]u`$
$`[u,d][u,d][d,s][u,d]u)`$
$`S_{27_3,1}^0`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}T_2^{33}q_2`$ (34)
$`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}}}([u,d][u,d][u,d][d,s]d[u,d][u,d][d,s][u,d]d).`$
### II.3 The wave function of Nonaquark $`\overline{\mathrm{𝟑𝟓}}`$-plet
Since $`\overline{\mathrm{𝟑𝟓}}`$ can not be completely excluded (see the subsection II.A) , we should also discuss it’s wavefunction for completeness. It is easy to identify
$`T_2^{1333}=\sqrt{6}S_{\overline{35},1}^0,T_1^{1333}=\sqrt{3}S_{\overline{35},1}^+\sqrt{2}S_{\overline{35}}^+`$
$`T_2^{2333}=\sqrt{3}S_{\overline{35},1}^+\sqrt{2}S_{\overline{35}}^+,T_3^{3333}=2\sqrt{2}S_{\overline{35}}^+,T_1^{2333}=\sqrt{6}S_{\overline{35},1}^{++}`$ (35)
then
$$\{\begin{array}{c}S_{\overline{35},1}^0=\frac{1}{\sqrt{6}}T_2^{1333}\hfill \\ S_{\overline{35},1}^+=\frac{1}{2\sqrt{3}}(T_2^{2333}T_1^{1333})\hfill \\ S_{\overline{35}}^+=\frac{1}{2\sqrt{2}}T_3^{3333}\hfill \end{array}$$
(36)
and $`T_n^{ijkl}`$ is defined in Eq.(16). It is easy to see
$`T^{1333}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}(S^{13}S^{23}+S^{33}S^{13}+S^{13}S^{33}+S^{31}S^{33}+S^{33}S^{13}+S^{33}S^{13})`$ (37)
$`=`$ $`{\displaystyle \frac{3}{\sqrt{6}}}(S^{13}S^{33}+S^{33}S^{13})`$
$`=`$ $`{\displaystyle \frac{3}{4\sqrt{6}}}([d,s][u,d][u,d][u,d]+[u,d][d,s][u,d][u,d]`$
$`+[u,d][u,d][d,s][u,d]+[u,d][u,d][u,d][d,s])`$
$`T^{2333}`$ $`=`$ $`{\displaystyle \frac{3}{\sqrt{6}}}(S^{23}S^{33}+S^{33}S^{23})`$ (38)
$`=`$ $`{\displaystyle \frac{3}{4\sqrt{6}}}([s,u][u,d][u,d][u,d]+[u,d][s,u][u,d][u,d]`$
$`+[u,d][u,d][s,u][u,d]+[u,d][u,d][u,d][s,u])`$
$$T^{3333}=\sqrt{6}S^{33}S^{33}=\frac{\sqrt{6}}{2}[u,d][u,d][u,d][u,d]$$
(39)
and the wave function of $`S_{\overline{35},1}^0,S_{\overline{35},1}^+,S_{\overline{35}}^+`$ are as followings
$`S_{\overline{35},1}^0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}T_2^{1333}={\displaystyle \frac{1}{\sqrt{6}}}T^{1333}q_2`$ (40)
$`=`$ $`{\displaystyle \frac{1}{8}}([d,s][u,d][u,d][u,d]+[u,d][d,s][u,d][u,d]`$
$`+[u,d][u,d][d,s][u,d]+[u,d][u,d][u,d][d,s])d`$
$`S_{\overline{35},1}^+`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{3}}}(T_2^{2333}T_1^{1333})={\displaystyle \frac{1}{2\sqrt{3}}}(T^{1333}q_1+T^{2333}q_2)`$ (41)
$`=`$ $`{\displaystyle \frac{1}{8\sqrt{2}}}([d,s][u,d][u,d][u,d]u[u,d][d,s][u,d][u,d]u[u,d][u,d][d,s][u,d]u`$
$`[u,d][u,d][u,d][d,s]u+[s,u][u,d][u,d][u,d]d+[u,d][s,u][u,d][u,d]d`$
$`+[u,d][u,d][s,u][u,d]d+[u,d][u,d][u,d][s,u]d)`$
$`S_{\overline{35}}^+`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}T_3^{3333}={\displaystyle \frac{1}{6\sqrt{2}}}(2T^{1333}q_12T^{2333}q_2+T^{3333}q_3)`$ (42)
$`=`$ $`{\displaystyle \frac{1}{8\sqrt{3}}}([d,s][u,d][u,d][u,d]u+[u,d][d,s][u,d][u,d]u+[u,d][u,d][d,s][u,d]u`$
$`+[u,d][u,d][u,d][d,s]u+[s,u][u,d][u,d][u,d]+[u,d][s,u][u,d][u,d]`$
$`+[u,d][u,d][s,u][u,d]+[u,d][u,d][u,d][s,u]2[u,d][u,d][u,d][u,d]s)`$
## III The Mass Spectrum of Nonaquark States
Since all the particles belonging to an irreducible representation of $`SU(3)`$ are degenerate in $`SU(3)`$ symmetry limit, it is necessary to introduce the $`SU(3)`$ symmetry breaking terms into the Hamiltonian in order to obtain the mass splitting. The Hamiltonian that breaks $`SU(3)`$ symmetry but still preserves the isospin symmetry and hypercharge is proportional to the Gell-Mann matrix $`\lambda _8`$, and the baryon mass can be obtained by constructing $`SU(3)`$ singlet term including the hypercharge tensor, in this way we obtain the Gell-Mann-Okubo mass formula:
$$M=M_0+\alpha Y+\beta D_3^3$$
(43)
where $`M_0`$ is a common mass of a given multiplet and $`D_3^3=I(I+1)\frac{Y^2}{4}\frac{C}{6}`$ with $`C=2(p+q)+\frac{2}{3}(p^2+pq+q^2)`$ for the (p,q) representation. $`\alpha `$ and $`\beta `$ are mass constant that are in principle different for different multiplets. Using these constants, we can obtain the masses of all the baryons within the multiplet. Note that in this picture the isospin is conserved.
### III.1 The mass spectrum of Nonaquark $`\mathrm{𝟐𝟕}`$-plet
In the case of $`\mathrm{𝟐𝟕}`$-plet, $`p=q=2`$ and the corresponding weight diagram is Fig.(6). By using the Gell-Mann-Okubo mass formula Eq.(43) we can get all the masses of these states
$`M_{S_{27,1}}=M_{27}+2\alpha _{27}{\displaystyle \frac{5}{3}}\beta _{27},M_{N_{27,\frac{3}{2}}}=M_{27}+\alpha _{27}+{\displaystyle \frac{5}{6}}\beta _{27},M_{N_{27,\frac{1}{2}}}=M_{27}+\alpha _{27}{\displaystyle \frac{13}{6}}\beta _{27}`$
$`M_{\mathrm{\Sigma }_{27,2}}=M_{27}+{\displaystyle \frac{10}{3}}\beta _{27},M_{\mathrm{\Sigma }_{27,1}}=M_{27}{\displaystyle \frac{2}{3}}\beta _{27},M_{\mathrm{\Lambda }_{27}}=M_{27}{\displaystyle \frac{8}{3}}\beta _{27}`$
$`M_{\mathrm{\Xi }_{27,\frac{3}{2}}}=M_{27}\alpha _{27}+{\displaystyle \frac{5}{6}}\beta _{27},M_{\mathrm{\Xi }_{27,\frac{1}{2}}}=M_{27}\alpha _{27}{\displaystyle \frac{13}{6}}\beta _{27}`$
$`M_{\mathrm{\Omega }_{27,1}}=M_{27}2\alpha _{27}{\displaystyle \frac{5}{3}}\beta _{27}`$ (44)
The Gell-Mann-Okubo mass relation for 27–plet nonaquark baryons is
$$3M_{\mathrm{\Lambda }_{27}}+M_{\mathrm{\Sigma }_{27,1}}=2(M_{N_{27,\frac{1}{2}}}+M_{\mathrm{\Xi }_{27,\frac{1}{2}}})=4M_{27}\frac{26}{3}\beta _{27}$$
(45)
The equal mass space relations also exist in the two separate sectors: ($`\mathrm{\Omega }_{27,1},\mathrm{\Xi }_{27,\frac{3}{2}},\mathrm{\Sigma }_{27,2}`$) and ($`\mathrm{\Sigma }_{27,2},N_{27,\frac{3}{2}},S_{27,1}`$). Their mass relations are as follows
$$M_{\mathrm{\Omega }_{27,1}}M_{\mathrm{\Xi }27,\frac{3}{2}}=M_{\mathrm{\Xi }27,\frac{3}{2}}M_{\mathrm{\Sigma }_{27,2}}=\alpha _{27}\frac{5}{2}\beta _{27}$$
(46)
$$M_{\mathrm{\Sigma }_{27,2}}M_{N_{27,\frac{3}{2}}}=M_{N_{27,\frac{3}{2}}}M_{S_{27,1}}=\alpha _{27}+\frac{5}{2}\beta _{27}$$
(47)
From Eq.(44) we can further obtain some relation between mass of these states
$`M_{\mathrm{\Lambda }_{27}}+M_{S_{27,1}}=2M_{N_{27,\frac{1}{2}}},M_{\mathrm{\Sigma }_{27,2}}+M_{S_{27,1}}=2M_{N_{27,\frac{3}{2}}},3M_{\mathrm{\Sigma }_{27,1}}+3M_{S_{27,1}}=4M_{N_{27,\frac{1}{2}}}+2M_{N_{27,\frac{3}{2}}}`$
$`3M_{\mathrm{\Xi }_{27,\frac{3}{2}}}+6M_{S_{27,1}}=5M_{N_{27,\frac{1}{2}}}+4M_{N_{27,\frac{3}{2}}},3M_{\mathrm{\Xi }_{27,\frac{1}{2}}}+6M_{S_{27,1}}=8M_{N_{27,\frac{1}{2}}}+M_{N_{27,\frac{3}{2}}}`$
$`3M_{\mathrm{\Omega }_{27,1}}+9M_{S_{27,1}}=10M_{N_{27,\frac{1}{2}}}+2M_{N_{27,\frac{3}{2}}}`$ (48)
Both the masses of $`\overline{\mathrm{𝟑𝟓}}`$-plets and the masses of $`\mathrm{𝟐𝟕}`$-plets contain three parameters: $`M_{\overline{35}},\alpha _{\overline{35}},\beta _{\overline{35}}`$ or $`M_{27},\alpha _{27},\beta _{27}`$, but we only known experimentally the mass of $`S^0(3115)`$ and $`S^+(3140)`$. So we can not fix the masses of other nonaquark states which are predicted by us.
It mostly seems that $`S^0(3115)`$ and $`S^+(3140)`$ would belong to the same isospin multiplet and the mass difference between them are mainly due to electromagnetic interaction and the mass differences between u quark’s and d quark’s .
### III.2 The mass spectrum of the Nonaquark $`\overline{\mathrm{𝟑𝟓}}`$-plet
By means of the Gell-Mann–Okubo mass formula eq.(43), and noting the $`\overline{\mathrm{𝟑𝟓}}`$-plet with $`p=1,q=4`$ whose weight diagram is shown in Fig7., the masses of all the exotic nonaquark states are as follows.
$`M_{X_{\overline{35},\frac{1}{2}}}=M_{\overline{35}}+3\alpha _{\overline{35}}{\displaystyle \frac{11}{2}}\beta _{\overline{35}},M_{S_{\overline{35},1}}=M_{\overline{35}}+2\alpha _{\overline{35}}3\beta _{\overline{35}},M_{S_{\overline{35}}}=M_{\overline{35}}+2\alpha _{\overline{35}}5\beta _{\overline{35}}`$
$`M_{N_{\overline{35},\frac{3}{2}}}=M_{\overline{35}}+\alpha _{\overline{35}}{\displaystyle \frac{\beta _{\overline{35}}}{2}},M_{N_{\overline{35},\frac{1}{2}}}=M_{\overline{35}}+\alpha _{\overline{35}}{\displaystyle \frac{7}{2}}\beta _{\overline{35}},M_{\mathrm{\Sigma }_{\overline{35},2}}=M_0+2\beta _{\overline{35}}`$
$`M_{\mathrm{\Sigma }_{\overline{35},1}}=M_{\overline{35}}2\beta _{\overline{35}},M_{\mathrm{\Xi }_{\overline{35},\frac{5}{2}}}=M_{\overline{35}}\alpha _{\overline{35}}+{\displaystyle \frac{9}{2}}\beta _{\overline{35}},M_{\mathrm{\Xi }_{\overline{35},\frac{3}{2}}}=M_{\overline{35}}\alpha _{\overline{35}}{\displaystyle \frac{1}{2}}\beta _{\overline{35}}`$
$`M_{\mathrm{\Omega }_{\overline{35},2}}=M_{\overline{35}}2\alpha _{\overline{35}}+\beta _{\overline{35}}`$ (49)
We can find that the equal space rule holds for two sectors of nonaquark baryons: ($`X_{\overline{35},\frac{1}{2}},S_{\overline{35},1},N_{\overline{35},\frac{3}{2}},\mathrm{\Sigma }_{\overline{35},2},\mathrm{\Xi }_{\overline{35},\frac{5}{2}}`$) and ($`S_{\overline{35}},N_{\overline{35},\frac{1}{2}},\mathrm{\Sigma }_{\overline{35},1},\mathrm{\Xi }_{\overline{35},\frac{3}{2}},\mathrm{\Omega }_{\overline{35},2}`$). They satisfy the mass relations as follows
$$M_{X_{\overline{35},\frac{1}{2}}}M_{S_{\overline{35},1}}=M_{S_{\overline{35},1}}M_{N_{\overline{35},\frac{3}{2}}}=M_{N_{\overline{35},\frac{3}{2}}}M_{\mathrm{\Sigma }_{\overline{35},2}}=M_{\mathrm{\Sigma }_{\overline{35},2}}M_{\mathrm{\Xi }_{\overline{35},\frac{5}{2}}}=\alpha _{\overline{35}}\frac{5}{2}\beta _{\overline{35}}$$
(50)
$$M_{S_{\overline{35}}}M_{N_{\overline{35},\frac{1}{2}}}=M_{N_{\overline{35},\frac{1}{2}}}M_{\mathrm{\Sigma }_{\overline{35},1}}=M_{\mathrm{\Sigma }_{\overline{35},1}}M_{\mathrm{\Xi }_{\overline{35},\frac{3}{2}}}=M_{\mathrm{\Xi }_{\overline{35},\frac{3}{2}}}M_{\mathrm{\Omega }_{\overline{35},2}}=\alpha _{\overline{35}}\frac{3}{2}\beta _{\overline{35}}$$
(51)
And we can also derive the mass relation between these nonaquark states
$`M_{X_{\overline{35},\frac{1}{2}}}+M_{N_{\overline{35},\frac{3}{2}}}=2M_{S_{\overline{35},1}},3(M_{S_{\overline{35}}}M_{S_{\overline{35},1}})=2(M_{N_{\overline{35},\frac{1}{2}}}M_{N_{\overline{35},\frac{3}{2}}})`$
$`M_{\mathrm{\Sigma }_{\overline{35},2}}+M_{S_{\overline{35},1}}=2M_{N_{\overline{35},\frac{3}{2}}},M_{\mathrm{\Sigma }_{\overline{35},1}}+3M_{S_{\overline{35},1}}=2(M_{S_{\overline{35}}}+M_{N_{\overline{35},\frac{3}{2}}})`$
$`M_{\mathrm{\Xi }_{\overline{35},\frac{5}{2}}}+2M_{S_{\overline{35},1}}=3M_{N_{\overline{35},\frac{3}{2}}},2M_{\mathrm{\Xi }_{\overline{35},\frac{3}{2}}}+9M_{S_{\overline{35},1}}=5M_{S_{\overline{35}}}+6M_{N_{\overline{35},\frac{3}{2}}}`$
$`M_{\mathrm{\Omega }_{\overline{35},2}}+6M_{S_{\overline{35},1}}=3M_{S_{\overline{35}}}+4M_{N_{\overline{35},\frac{3}{2}}}`$ (52)
## IV Decay of The Nonaquark States
We assume the nonaquark state decays arise by a ”fall-apart” mechanism closecon ; Jm ; carlson ; bs without need for gluon exchange to trigger the decay. There are some discussions on this mechanism in the pentaquark spectrum studies. By this mechanism, a diquark in the pentaquark must be so clever that its two quarks are detached to two isolated quarks, and one of them enters into the adjacent diquarks to form a baryon separatively, and another combines with the residual anti-quark to form a meson, and then a pentaquark baryon decays into a usual baryon plus a meson. In this mechanism, the dynamics from the color coupling contributes a common factor to the decay amplitude for a certain flavor multiplet which is irrelevant to the discussions on its decay branch fractions. Extending this ”fall-apart” mechanism to nonaquark state decays is natural and straightforward: 1, one diquark is detached into two quarks; 2, these two quarks enter the adjacent diquaks, and form two baryons separatively; 3, the ninth quark also enter a diquark to form a baryon, and then the nonaquark state consequently decays into three baryons. We show this mechanism in the Fig5.
To the baryon octet, its quark content and the corresponding tensor are well-known, which is listed in the Table I. The tensor basis of the baryon octet is
$$B_r^l=\frac{1}{\sqrt{3}}\epsilon ^{lmn}S_{rn}q_m$$
(53)
with $`S_{rn}=\frac{1}{\sqrt{2}}(q_rq_n+q_nq_r)`$.
And the baryon singlet $`\mathrm{\Lambda }_1^0`$ is given by
$`\mathrm{\Lambda }_1^0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\sqrt{6}}}\epsilon ^{klm}(q_lq_mq_mq_l)q_k`$ (54)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}([d,s]u+[s,u]d+[u,d]s)`$
From these wave function we can see:
$$\{\begin{array}{c}\text{[d,s]}u=\mathrm{\Sigma }^0+\frac{\sqrt{6}}{3}\mathrm{\Lambda }_1^0\frac{\sqrt{3}}{3}\mathrm{\Lambda }_8^0\hfill \\ \text{[s,u]}d=\mathrm{\Sigma }^0+\frac{\sqrt{6}}{3}\mathrm{\Lambda }_1^0\frac{\sqrt{3}}{3}\mathrm{\Lambda }_8^0\hfill \\ \text{[u,d]}s=\frac{\sqrt{6}}{3}\mathrm{\Lambda }_1^0+\frac{2\sqrt{3}}{3}\mathrm{\Lambda }_8^0\hfill \end{array}$$
(55)
### IV.1 The decay of $`\mathrm{𝟐𝟕}`$-plet
Since the $`\mathrm{𝟐𝟕}`$-plet is degenerate, i.e., $`\mathrm{𝟐𝟕}_\mathrm{𝟏},\mathrm{𝟐𝟕}_\mathrm{𝟐}`$ and $`\mathrm{𝟐𝟕}_\mathrm{𝟑}`$ are the three irreducible representatives of $`\mathrm{𝟐𝟕}`$dimension, we discuss the decays for each case following the method of dealing with the ”fall-apart” decay close . Starting with the wave function of $`S_{27_1,1}^+`$ ( or $`S_{27_1,1}^0`$ etc.) which is given in Eq.(25)(Eq.(26) etc.), we rewrite the wave function in the form of (qqq)(qqq)(qqq) which is similar to Ref.close , then using Eq.(55) and Table I and considering fermion statistics, we can map $`S_{27_1,1}^+`$ ( or $`S_{27_1,1}^0`$ etc.) into the ground state of three baryons, from which we can learn what particles are the decay products of the $`S_{27_1,1}^+`$ (or $`S_{27_1,1}^0`$ etc.), and derive further the ratios of the branch fractions of these channels. We discuss the decays for each case respectively by means of the above spirit of dealing with the decays in the follows:
1. The first case: The wave function of $`S_{27_1,1}^+`$ and $`S_{27_1,1}^0`$ is given by Eq.(25) and Eq.(26) separately , and they are mapped into
$$S_{27_1,1}^+\frac{1}{\sqrt{2}}\mathrm{\Sigma }^0pn$$
(56)
$$S_{27_1,1}^0\frac{1}{2}(\sqrt{2}\mathrm{\Sigma }^{}pn+\frac{\sqrt{6}}{3}\mathrm{\Lambda }_1^0nn+\frac{2\sqrt{3}}{3}\mathrm{\Lambda }_8^0nn).$$
(57)
Then, we obtain that the main decay channel of $`S_{27_1,1}^+`$ is $`\mathrm{\Sigma }NN`$, and the ratio of coupling constants is
$$g(S_{27_1,1}^+\mathrm{\Sigma }^{}pn):g(S_{27_1,1}^+\mathrm{\Lambda }_1^0nn):g(S_{27_1,1}^+\mathrm{\Lambda }_8^0nn)=\sqrt{2}:\frac{\sqrt{6}}{3}:\frac{2\sqrt{3}}{3}.$$
(58)
Considering phase space effect, we get
$$\frac{BR[S_{27_1,1}^+\mathrm{\Sigma }NN]}{BR[S_{27_1,1}^+\mathrm{\Lambda }NN]}=\frac{2}{((\frac{\sqrt{6}}{3})^2+\frac{2\sqrt{3}}{3})^2)\times 4.164}0.24$$
(59)
Therefore, $`S_{27_1,1}^0`$ dominantly decays into $`\mathrm{\Lambda }NN`$, and $`S_{27_1,1}^0`$ can not be $`S^0(3115)`$.
2. The second case: Eq.(29) and Eq.(30) give the wave function of $`S_{27_2,1}^+`$ and $`S_{27_2,1}^0`$ in this case, and the map is as following
$$S_{27_2,1}^+\frac{1}{2\sqrt{2}}(\mathrm{\Sigma }^+nn+\mathrm{\Sigma }^{}pp+\sqrt{2}\mathrm{\Sigma }^0pn),$$
(60)
$$S_{27_2,1}^0\frac{1}{2\sqrt{2}}(\mathrm{\Sigma }^0nn+\sqrt{2}\mathrm{\Sigma }^{}pn\sqrt{6}\mathrm{\Lambda }_1^0nn\sqrt{3}\mathrm{\Lambda }_8^0nn).$$
(61)
We note that in this case $`S_{27_2,1}^+`$ also can only decay into $`\mathrm{\Sigma }NN`$, and the ratio of the branch fractions
$$BR[S_{27_2,1}^+\mathrm{\Sigma }^+nn]:BR[S_{27,1}^+\mathrm{\Sigma }^{}pp]:BR[S_{27,1}^+\mathrm{\Sigma }^0pn]=1:1:2$$
(62)
and the main decay modes of $`S_{27_2,1}^0`$ is $`\mathrm{\Lambda }NN`$
$$\frac{BR[S_{27_2,1}^0\mathrm{\Sigma }NN]}{BR[S_{27_2,1}^0\mathrm{\Lambda }NN]}0.08$$
(63)
3. The third case: The wave function of $`S_{27_3,1}^+`$ and $`S_{27_3,1}^0`$ is given by Eq.(33) and Eq.(34), similarly
$$S_{27_3,1}^+\frac{3}{2}\mathrm{\Sigma }^0pn+\frac{1}{2\sqrt{2}}\mathrm{\Sigma }^+nn+\frac{1}{2\sqrt{2}}\mathrm{\Sigma }^{}pp$$
(64)
$$S_{27_3,1}^0\frac{1}{4\sqrt{2}}(6\sqrt{2}\mathrm{\Sigma }^{}pn2\mathrm{\Sigma }^0nn\frac{2\sqrt{6}}{3}\mathrm{\Lambda }_1^0nn+\frac{2\sqrt{3}}{3}\mathrm{\Lambda }_8^0nn).$$
(65)
From the above two formula, we know that in this case $`S_{27_3,1}^+`$ decay mainly into $`\mathrm{\Sigma }NN`$ not to $`\mathrm{\Lambda }NN`$,and
$$BR[S_{27_3,1}^+\mathrm{\Sigma }^0pn]:BR[S_{27_3,1}^+\mathrm{\Sigma }^+nn]:BR[S_{27_3,1}^+\mathrm{\Sigma }^{}pp]=18:1:1$$
(66)
and the ratio of branch fractions
$$\frac{BR[S_{27_3,1}^0\mathrm{\Sigma }NN]}{BR[S_{27_3,1}^0\mathrm{\Lambda }nn]}2.8$$
(67)
We note that in this case the main decay mode of $`S_{27_3,1}^0`$ is $`\mathrm{\Sigma }NN`$ which is consistent with the decay of $`S^0(3115)`$.
### IV.2 The decay of $`\overline{\mathrm{𝟑𝟓}}`$-plet
Starting with the wave function of $`S_{\overline{35},1}^+`$ which is given in Eq.(41), we discuss the decays by means of the method used in the previous sub-section to deal with the decays. We map $`S_{\overline{35},1}^+`$ into the ground state of three baryons
$$S_{\overline{35},1}^+\frac{1}{2}(\mathrm{\Sigma }^{}pp+\mathrm{\Sigma }^+nn)$$
(68)
So, $`S_{\overline{35},1}^+`$ can only decay into $`\mathrm{\Sigma }^{}NN`$, but can not decay into $`\mathrm{\Lambda }^{}NN`$ (N stands for nucleon,i.e., proton or neutron). And the ratio of branch fractions is
$$\frac{BR[S_{\overline{35},1}^+\mathrm{\Sigma }^{}pp]}{BR[S_{\overline{35},1}^+\mathrm{\Sigma }^+nn]}=\frac{1}{1}=1$$
(69)
Similarly $`S_{\overline{35}}^+`$ whose wave function is defined in Eq.(42) is mapped onto
$$S_{\overline{35}}^+\frac{1}{\sqrt{6}}(\mathrm{\Sigma }^{}pp\mathrm{\Sigma }^+nn)$$
(70)
also $`S_{\overline{35}}^+`$ decay only into $`\mathrm{\Sigma }^{}NN`$, can not decay into $`\mathrm{\Lambda }^{}NN`$, and the ratio of branch fractions is
$$\frac{BR[S_{\overline{35}}^+\mathrm{\Sigma }^{}pp]}{BR[S_{\overline{35}}^+\mathrm{\Sigma }^+nn]}=\frac{1}{1}=1$$
(71)
We also obtain that the $`S_{\overline{35}}^+\mathrm{\Sigma }^{}pp`$ and $`S_{\overline{35}}^+\mathrm{\Sigma }^+nn`$ interactions have different phase, while the $`S_{\overline{35},1}^+\mathrm{\Sigma }^{}pp`$ and $`S_{\overline{35},1}^+\mathrm{\Sigma }^+nn`$ interactions have the same phase (see Eqs.(70) and (68)).
In the same way, starting from Eq.(40) which is the wave function of $`S_{\overline{35},1}^0`$, we have:
$$S_{\overline{35},1}^0\frac{1}{2}(\mathrm{\Sigma }^0+\sqrt{3}\mathrm{\Lambda }_8^0)nn$$
(72)
So the main decay channel of $`S_{\overline{35},1}^0`$ is $`\mathrm{\Lambda }NN`$, but not $`\mathrm{\Sigma }NN`$, i.e., ratio of the effective couplings reads
$$\frac{g(S_{\overline{35},1}^0\mathrm{\Sigma }^0nn)}{g(S_{\overline{35},1}^0\mathrm{\Lambda }_8^0nn)}=\frac{1}{\sqrt{3}}.$$
(73)
Considering the different three body phase space, we can obtain the the ratio of branch fractions
$$\frac{BR[S_{\overline{35},1}^0\mathrm{\Sigma }^0nn]}{BR[S_{\overline{35},1}^0\mathrm{\Lambda }_8^0nn]}0.05.$$
(74)
So the observed $`S^0(3115)`$ kek1 can not be $`S_{\overline{35},1}^0`$.
From the decay of $`\overline{\mathrm{𝟑𝟓}}`$-plet and $`\mathrm{𝟐𝟕}`$-plet, we see that the nonaquark state $`S^0(3115)`$ can only possibly belong to $`\mathrm{𝟐𝟕}`$-plet, and its flavor configuration is Eq.(34), or its main component is $`S_{27_3,1}^0`$ with small mixing of $`S_{27_1,1}^0,S_{27_2,1}^0,S_{35,1}^0`$ . While $`S^+(3140)`$ maybe belong to $`\overline{\mathrm{𝟑𝟓}}`$-plet or $`\mathrm{𝟐𝟕}`$-plet, it possibly is $`S_{27_1,1}^+,S_{27_2,1}^+,S_{27_3,1}^+,S_{35,1}^+`$ or the mixing of all of them, since these states have same quantum numbers. From the discussion of subsection A in section II , we know that in the exact $`SU(3)^{flavor}\times SU(3)^{color}\times SU(2)^{spin}`$ limit, $`S^0(3115)`$ and $`S^+(3140)`$ possibly belong to 27-plet, this further give support to our suggestion that $`S^0(3115)`$ can only possibly belong to $`\mathrm{𝟐𝟕}`$-plet. Furthermore in this case the nonaquark states have negative parity, this is a unusual results, since ”standard” nonaquark state which involve 9 quarks in relative S-wave have positive parity, there maybe really exists exotica nonaquark 27-plet. We can obtain useful information about $`S^+(3140)`$ by experimentally measuring the branch fractions of its decay channels.
## V Discussion and Conclusion
In summary, we have obtain the wave functions of the nonaquark states in $`SU(3)`$ quark model with diquark correlation using standard direct tensor decomposition, we predict the existence of other nonaquark states. It would be helpful for constructing the effective interaction Lagrangian to describe the nonaquark decays with rational $`SU(3)`$flavor structure. We obtain some interesting mass sum rules for the nonaquark $`\overline{\mathrm{𝟑𝟓}}`$-plet and $`\mathrm{𝟐𝟕}`$-plet, but it is still open to fix the spectrum of $`\overline{\mathrm{𝟑𝟓}}`$-plet and $`\mathrm{𝟐𝟕}`$-plet due to the scarcity of experiments. More data are expected. Under the assumption of the ”fall-apart” decay mechanism which has been subtly used in studying pentaquark decays, we find out that the $`S^0(3115)`$ belongs to a $`\mathrm{𝟐𝟕}`$-plet. Its main component is $`S_{27_3,1}^0`$, and the mixing with $`S_{27_1,1}^0,S_{27_2,1}^0,S_{35,1}^0`$ must be few. It is possible that $`S^0(3115)`$ and $`S^+(3140)`$ belong to the same isospin multiplet, since their mass difference is about 25 MeV which can be interpreted by electromagnetic interaction and u,d quark mass difference. Further in the exact $`SU(3)^{flavor}\times SU(3)^{color}\times SU(2)^{spin}`$ limit, $`S^0(3115)`$ and $`S^+(3140)`$ are belong to the 27-plet, and its parity is negative. We suggest to study the decays of $`S^+(3140)`$ more in experiment, especially the ratios of the various decay mode’s branch fractions, through which we can learn the structure of $`S^+(3140)`$ and maybe discover new nonaquark state .
There maybe other types of quark correlation in nonaquark state, such as (qqq)-(qqq)-(qqq), since either (qq$`\overline{\mathrm{q}}`$)-(qq) quark correlation lipkin or diquark correlation jaffe is possible in pentaquark, they have different flavor and color structures from the diquark correlation, so the masses, the decay property, and the products etc of the states may be different from the prediction of quark model with diquark correlation.
It would be interesting to study the mixing between the nonaquark states with same quantum numbers. It is necessary to considering the mixing in order to compare theoretical prediction with experimental data more precisely. This would be highly nontrivial.
Finally, we might image that this multi-baryon problem could be studied in other quark models, e.g., the model in swart , and in the chiral soliton modelskyrme .
ACKNOWLEDGMENTS
We acknowledge Dr. J. Deng for his help in plotting the diagrams. We would like also to thank the referee for his suggestion to discuss the color symmetry in the nonaquark wavefunction. This work is partially supported by National Natural Science Foundation of China under Grant Numbers 90403021, and by the PhD Program Funds of the Education Ministry of China and KJCX2-SW-N10 of the Chinese Academy.
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# Acknowledgements
## Acknowledgements
J. M. M., N. S. and V. T. are supported by CNRS. N. K., J. M. M., V. T. are supported by the European network EUCLID-HPRNC-CT-2002-00325. J. M. M. and N.S. are supported by INTAS-03-51-3350. N.S. is supported by the French-Russian Exchange Program, the Program of RAS Mathematical Methods of the Nonlinear Dynamics, RFBR-05-01-00498, Scientific Schools 2052.2003.1. N. K, N. S. and V. T. would like to thank the Theoretical Physics group of the Laboratory of Physics at ENS Lyon for hospitality, which makes this collaboration possible.
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# References
On Integrating Out Heavy Fields in SUSY Theories
S. P. de Alwis
Physics Department, University of Colorado,
Boulder, CO 80309 USA
Abstract
We examine the procedure for integrating out heavy fields in supersymmetric (both global and local) theories. We find that the usual conditions need to be modified in general and we discuss the restrictions under which they are valid. These issues are relevant for recent work in string compactification with fluxes.
PACS numbers: 11.25. -w, 98.80.-k; COLO-HEP 508
e-mail: dealwis@pizero.colorado.edu
In a theory containing both heavy (mass $`M`$) and light (mass $`mM`$) fields, one may derive an effective field theory Weinberg (1980) valid for energy scales $`EM`$, by doing the functional integral over the heavy fields (as well as light field modes with frequencies greater than $`M`$). If one is just considering the classical approximation then this just means solving the classical equation for the heavy field in terms of the light field and substituting back into the action. So if the potential for the heavy ($`\mathrm{\Phi }`$) and light ($`\varphi `$) fields is $`V(\mathrm{\Phi },\varphi )=\frac{1}{2}M^2\mathrm{\Phi }^2+\stackrel{~}{V}(\mathrm{\Phi },\varphi )`$ the equation of motion for the heavy field gives,
$$\mathrm{\Phi }=\frac{1}{\mathrm{}M^2}\frac{\stackrel{~}{V}}{\mathrm{\Phi }}=\frac{1}{M^2}\frac{\stackrel{~}{V}}{\mathrm{\Phi }}+\frac{1}{M^2}O(\frac{\mathrm{}}{M^2}\frac{\stackrel{~}{V}}{\mathrm{\Phi }}).$$
In other words up to terms in the derivative expansion that maybe ignored at energy $`EM`$, $`\mathrm{\Phi }`$ is a solution of the equation
$$\frac{V}{\mathrm{\Phi }}=0$$
(1)
Consider now a globally supersymmetric theory of chiral scalar fields with superspace action
$$S=d^4xd^4\theta \overline{\mathrm{\Phi }^i}\mathrm{\Phi }^i+(d^4xd^2\theta W(\mathrm{\Phi }^i)+c.c.].$$
(2)
The superfield equations of motion are (with $`D^2=D^\alpha D_\alpha `$where $`D_\alpha `$ is the spinor covariant derivative)<sup>1</sup><sup>1</sup>1We use the conventions of Wess and Bagger (1992). For the derivations below see also Gates et al. (1983).,
$$\frac{1}{4}\overline{D}^2\overline{\mathrm{\Phi }^i}+\frac{W}{\mathrm{\Phi }^i}=0.$$
(3)
The component form of this is obtained by successive application of the spinorial covariant derivative and setting $`\theta =\overline{\theta }=0`$ (indicated below by the symbol $`|`$). In particular the auxiliary field equation (from $`D^0`$) and the scalar field (from $`\frac{1}{4}D^2`$) become (after setting the fermions to zero and using $`D^2\overline{D^2}\overline{\mathrm{\Phi }}=16\mathrm{}\overline{\mathrm{\Phi }}`$)
$`\overline{F^i}+{\displaystyle \frac{W}{\mathrm{\Phi }^i}}|`$ $`=0`$ (4)
$`\mathrm{}\overline{\mathrm{\Phi }^i|}+{\displaystyle \underset{j}{}}{\displaystyle \frac{W}{\mathrm{\Phi }^i\mathrm{\Phi }^j}}|F^j`$ $`=0`$ (5)
Now suppose that we wish integrate out a heavy field with say $`i=H`$ to get an effective theory for the light fields with $`i=l`$. For example one might have $`W=\frac{1}{2}M\mathrm{\Phi }^{H2}+\frac{1}{2}\lambda _{ll^{}}\mathrm{\Phi }^H\mathrm{\Phi }_l\mathrm{\Phi }_l^{}+W_L(\mathrm{\Phi }_l)`$. Then the analog of (1) (see for example Intriligator and Seiberg (1996)) is to require that
$$\frac{W}{\mathrm{\Phi }^H}=0.$$
(6)
Taking components of this equation, (and ignoring the fermionic terms)
$`{\displaystyle \frac{W}{\mathrm{\Phi }^H}}|`$ $`=0`$ (7)
$`{\displaystyle \underset{j}{}}{\displaystyle \frac{^2W}{\mathrm{\Phi }^H\mathrm{\Phi }^l}}|D^2\mathrm{\Phi }^l+{\displaystyle \frac{^2W}{\mathrm{\Phi }^H\mathrm{\Phi }^H}}|{\displaystyle \frac{\overline{W}}{\overline{\mathrm{\Phi }}^H}}|`$ $`=0`$ (8)
where in the second (obtained by acting with $`D^2`$ and setting fermions to zero) we have used (4) for $`i=H`$. The second condition is of course the requirement that the potential $`V=_iF^i\overline{F^i}`$ be extremized with respect to $`\mathrm{\Phi }_H,`$ which is what one would impose in a non-supersymmetric theory (see (1)). However here we have the additional condition (7) which in conjunction with (8) and the equatinon of motion for the light field leads to
$$\underset{j}{}\frac{^2W}{\mathrm{\Phi }^H\mathrm{\Phi }^l}|\frac{W}{\mathrm{\Phi }^l}|=0$$
But (after solving for the heavy field using (6) this is a constraint on the light fields in theories where there is a non-zero coupling between the heavy and light fields as in the above example. In particular if there is only one light field it imposes the condition $`\frac{W}{\mathrm{\Phi }^l}|=0`$, meaning that the light scalar is also at the minimum of the potential. In other words one does not get an effective potential for the light field - the only consistent result of integrating out the heavy field is that all fields are sitting at the SUSY minimum (if it exists). When there is more than one light field this is not necessarily the case but nevertheless the light field space is constrained. If we had kept the fermionic terms then this constraint would be a relation between the bosonic components and squares of fermionic components of the light fields.
To see where this comes from let us write the superpotential for the heavy light theory as
$$W(\mathrm{\Phi }^H,\mathrm{\Phi }^l)=\frac{1}{2}M\mathrm{\Phi }^{H2}+\stackrel{~}{W}(\mathrm{\Phi }^H,\mathrm{\Phi }^l).$$
(9)
Operating on the (conjugate of the) equation of motion (3) with ($`\frac{1}{4}D^2`$), using (3) again and rearranging we have,
$$\mathrm{\Phi }^H=\frac{1}{\mathrm{}M^2}(M\frac{\stackrel{~}{W}}{\mathrm{\Phi }^H}+\frac{\overline{D}^2}{4}\frac{\overline{\stackrel{~}{W}}}{\overline{\mathrm{\Phi }}^H})$$
Expanding the inverse Klein-Gordon operator as before we can rewrite this as,
$$M\mathrm{\Phi }^H+\frac{\stackrel{~}{W}}{\mathrm{\Phi }^H}=\frac{\overline{D}^2}{4M}\frac{\overline{\stackrel{~}{W}}}{\overline{\mathrm{\Phi }}^H}+O(\frac{\mathrm{}}{M^2}(\mathrm{})).$$
So to the lowest order in the space-time derivative (momentum) expansion what we get for the equation determining the heavy field in terms of the light is,
$$\frac{W}{\mathrm{\Phi }^H}=\frac{\overline{D}^2}{4M}\frac{\overline{\stackrel{~}{W}}}{\overline{\mathrm{\Phi }}^H},$$
(10)
rather than (6). To get the latter one needs the additional assumption that the possible values of $`\overline{D^2}\frac{\overline{\stackrel{~}{W}}}{\overline{\mathrm{\Phi }}^H}`$ are small compared to $`M`$. For instance in the above example (in the paragraph after (5) this means that we need $`\mathrm{\Phi }_l<<M`$. This of course is what one would expect. However as we pointed out earlier in this same approximation the light field space appears to be constrained.
Let us be even more specific and consider the model (9) with $`\stackrel{~}{W}=\frac{1}{2}HL^2`$. (We’ve relabelled the heavy field as $`H`$ and the light field as $`L`$.) Then (10) becomes,
$`{\displaystyle \frac{W}{H}}`$ $`=MH+{\displaystyle \frac{1}{2}}L^2`$ $`={\displaystyle \frac{\overline{D}^2}{4M}}({\displaystyle \frac{1}{2}}\overline{L}^2)`$ (11)
$`={\displaystyle \frac{1}{4M}}(\overline{D}^{\dot{\alpha }}\overline{\overline{L}D_{\dot{\alpha }}}\overline{L}`$ $`+\overline{L}\overline{D}^2\overline{L)}`$
Again we see that the strict imposition of $`W/H=0`$ leads to the constraint on the light field space that we found above. To see in what approximation this equation is valid let us solve it for $`H`$ (giving $`H=L^2/M`$) and then plug it back into the RHS of the last equality of (11) after using the light field equation. The bosonic term is then $`L^2\frac{|L|^2}{M^2}`$and is small when $`|L|<<M`$.
It is perhaps worthwhile looking at this example in component form. The superspace Lagrangian in the above example is
$$d^4\theta (\overline{H}H+\overline{L}L)+[d^2\theta \frac{1}{2}(MH^2+HL^2)+c.c.]$$
(12)
In components one has, writing the scalar and F components of $`H=(A,`$F) and of $`L=(a,f)`$ and ignoring the fermion terms,
$`L`$ $`=`$ $`\overline{A}\mathrm{}A+\overline{a}\mathrm{}a+\overline{F}F+\overline{f}f+{\displaystyle \frac{1}{2}}(Fa^2+\overline{F}\overline{a}^2)+(Aaf+\overline{A}\overline{a}\overline{f})`$ (13)
$`+`$ $`M(AF+\overline{A}\overline{F})`$
The heavy field equations are
$`\mathrm{}\overline{A}`$ $`=`$ $`MF+af`$ (14)
$`\overline{F}`$ $`=`$ $`{\displaystyle \frac{a^2}{2}}MA`$ (15)
Solving them we get after expanding in powers of $`\mathrm{}/M^2`$,
$`\overline{A}`$ $`=`$ $`{\displaystyle \frac{1}{2M^2}}(M\overline{a}^2+2af)+{\displaystyle \frac{1}{M^2}}O({\displaystyle \frac{\mathrm{}}{M^2}})`$
$`\overline{F}`$ $`=`$ $`{\displaystyle \frac{\overline{a}\overline{f}}{M}}+{\displaystyle \frac{1}{M^2}}O({\displaystyle \frac{\mathrm{}}{M^2}})`$
It is easily checked that these are precisely the equations that would be obtained by looking at the components of the superfield equation (10). Plugging these equations into (13) we get the light field potential,
$$V=\overline{f}f(1+\frac{\overline{a}a}{M^2})\frac{1}{2M}(fa^3+\overline{f}\overline{a}^3)$$
Eliminating the light auxiliary field we have
$$V=\frac{|a|^6}{4M^2}(1+\frac{|a|^2}{M^2})^1$$
If we had just imposed the usual condition $`W/H=0`$ we would not have got the $`|a|^2/M^2`$term in the parenthesis. This means that this condition gives the correct result for the potential only for small values of the field $`|a|<<M`$.
To derive the analog of (10) in supergravity we use the formalism in chapter 8 of Gates et al. (1983) (though we remain with the conventions of Wess and Bagger (1992)). In this formalism the matter equations can be derived by replacing the supervielbein determinant by the so-called chiral compensator field $`\varphi `$ and treating the coupling of the matter fields as in flat superspace. Thus the supercovariant derivative $`D_\alpha `$is the flat space one and satisfies $`D^3=0`$. Also acting on chiral fields we have $`(D^2\overline{D}^2/16)\mathrm{\Phi }=\mathrm{}\mathrm{\Phi }`$. The action is (with $`M_p^2=1`$)
$$S=3d^4xd^4\theta \overline{\varphi }\varphi e^{K/3}+(d^4xd^2\theta \varphi ^3W+h.c.)$$
(16)
where the superpotential is a holomorphic function of the chiral scalar fields $`W=W(\chi ^i)`$ and the Kaehler potential is a real function $`K=K(\chi ,\overline{\chi })`$. From this action one obtains the following equations of motion.
$`{\displaystyle \frac{1}{4}}\overline{D^2}\overline{\chi }^{\overline{k}}+{\displaystyle \frac{1}{4}}K^{\overline{k}i}(K_{i\overline{j}l}{\displaystyle \frac{2}{3}}K_{i\overline{j}}K_{\overline{l}})\overline{D}^{\dot{\alpha }}\overline{\chi }^{\overline{j}}\overline{D}_{\dot{\alpha }}\overline{\chi }^{\overline{l}}`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\overline{D}^{\dot{\alpha }}\overline{\varphi }}{\overline{\varphi }}}\overline{D}_{\dot{\alpha }}\overline{\chi }^{\overline{k}}`$
$`=e^{K/3}{\displaystyle \frac{\varphi ^2}{\overline{\varphi }}}K^{\overline{k}i}D_iW`$
$`{\displaystyle \frac{1}{4}}\overline{D^2}(\overline{\varphi }e^{K/3})=\varphi ^2W(\chi )`$
Using the identity above (16) we get,
$$\mathrm{}\overline{\chi }^{\overline{k}}=\frac{D^2}{16}[K^{\overline{k}i}(K_{i\overline{j}l}\frac{2}{3}K_{i\overline{j}}K_{\overline{l}})]\overline{D}^{\dot{\alpha }}\overline{\chi }^{\overline{j}}\overline{D}_{\dot{\alpha }}\overline{\chi }^{\overline{l}}\overline{\varphi }^1\frac{D^2}{4}(e^{K/3}\varphi ^2D^{\overline{k}}W)$$
(17)
Let us now consider the case with one heavy superfield $`H`$ which we will take to be canonically normalized. So the Kahler potential becomes,
$$K=\overline{H}H+K^l(L,\overline{L})$$
(18)
where $`L`$ stands for the light fields. Also the superpotential is taken to be
$$W=\frac{1}{2}MH^2+\stackrel{~}{W}(H,L)$$
where $`M`$is a large mass parameter.
Then from (17) we have for the heavy field
$`4\overline{\varphi }\mathrm{}\overline{H}`$ $`=D^2(e^{K/3M_p^2}\varphi ^2)D_HW+4Me^{2K/M_p^2}\varphi ^2\overline{\varphi }^2(1+{\displaystyle \frac{\overline{H}H}{M_p^2}})D_{\overline{H}}\overline{W}`$
$`+`$ $`(1+{\displaystyle \frac{\overline{HH}}{M_p^2}})[e^{K/3M_p^2}\varphi ^2(({\displaystyle \frac{2}{3}}{\displaystyle \frac{K_l}{M_p^2}}D^\alpha \chi ^l{\displaystyle \frac{2D^\alpha \varphi }{\overline{\varphi }}})+{\displaystyle \frac{\overline{H}D^\alpha H}{M_p^2}})`$ $`+2D^\alpha (e^{K/3M_p^2}\varphi ^2)]D_\alpha H`$
$``$ $`{\displaystyle \frac{\overline{\varphi }}{6M_p^2}}D^2K_{\overline{l}}\overline{D}^{\dot{\alpha }}\overline{H}\overline{D}_{\dot{\alpha }}\overline{\chi }^{\overline{l}}+2D^\alpha (e^{K/3M_p^2}\varphi ^2)D_\alpha D_H\stackrel{~}{W}`$ $`+e^{K/3M_p^2}\varphi ^2D^2D_H\stackrel{~}{W}`$ (19)
In the above $`D_HW=_HW+K_HW/M_p^2`$ is the Kaehler derivative of the superpotential with respect to the heavy field and we have restored the dependence on the Planck mass. To integrate out a heavy field $`H`$ we have to set $`\mathrm{}\overline{H}p^2\overline{H}0`$ and the condition for that is that the right hand side of the above equation is set to zero. Note that this condition reduces to (10) in the global limit $`M_p\mathrm{}`$ and $`\varphi 1`$. Here up to terms involving a factor of $`D_\alpha H`$ we have
$$D^2(e^{K/3M_p^2}\varphi ^2)D_HW+4Me^{2K/M_p^2}\varphi ^2\overline{\varphi }^2(1+\frac{\overline{H}H}{M_p^2})D_{\overline{H}}\overline{W}=e^{K/3M_p^2}\varphi ^2D^2D_H\stackrel{~}{W}$$
As in the case of the global SUSY discussion one may expect that with some restriction on the light field space (so that the right hand side of the equation is small compared to $`M`$) the relevant condition would be the natural generalization of (6)
$$D_HW=_HW+\frac{1}{M_p^2}W_HK=0$$
(20)
Note that (by taking its spinor derivative) this condition implies
$$W\overline{D}_{\dot{\alpha }}\overline{H}=0,$$
so that the other $`O(M)`$ terms which all have a factor of $`D_\alpha H`$ also vanish (since the superpotential should not vanish at a generic point). So (20) is certainly a sufficient condition in the sense that it implies $`\mathrm{}H=0`$ .
However it is easy to see that the strict implementation of the condition (20) leads to the conclusion that the light fermion fields would have to be set to zero. Let us look at this in somewhat more general terms than above.
We assume that the Kaehler potential is a sum of terms as in the string theory examples discussed in de Alwis (2005). In particular if we call the heavy superfields $`H^I`$ and the light superfields $`L^i`$ assume that
$$K=K^h(H,\overline{H)}+K^l(L,\overline{L})$$
(21)
Thus in the example in section 4 below, $`H^I=z^i`$and $`L^i=S,T`$. Then the generalization of (20) becomes (note that capital letters $`I,J`$ go over the heavy fields)
$$K^{h\overline{I}J}D_JW=0$$
(22)
Using the non-degeneracy of the metric on the heavy fields this becomes,
$$_IW+K_I^hW=0$$
(23)
Taking the anti-chiral derivative of this equation and using the chirality of the superpotential we get,
$$WK_{I\overline{J}}^h\overline{𝒟}_{\dot{\alpha }}\overline{H}^{\overline{J}}=0$$
(24)
implying $`\overline{𝒟}_{\dot{\alpha }}\overline{H}^{\overline{J}}=0`$ since the metric is $`K_{I\overline{J}}^h`$ is non-degenerate and $`W0`$ at generic points.
Now let us assume that there is a solution $`H=H(L,\overline{L})`$ of equation (23). Using the chirality of $`H`$ and $`L`$ we get by differentiating this solution,
$$\overline{𝒟}_{\dot{\alpha }}H=\frac{H^I}{L^j}\overline{𝒟}_{\dot{\alpha }}L^j+\frac{H^I}{\overline{L}^{\overline{j}}}\overline{𝒟}_{\dot{\alpha }}\overline{L}^{\overline{j}}=\frac{H^I}{\overline{L}^{\overline{j}}}\overline{𝒟}_{\dot{\alpha }}\overline{L}^{\overline{j}}=0.$$
(25)
Similarly from the result (see (24)) that the chiral derivative of $`H`$ also vanishes,
$$𝒟_\alpha H^I=\frac{H^I}{L^j}𝒟_\alpha L^j+\frac{H}{\overline{L}^{\overline{j}}}𝒟_\alpha \overline{L}^{\overline{j}}=\frac{H^I}{L^j}𝒟_\alpha L^j=0$$
(26)
These two equations tell us that the chiral derivative of $`L`$ should be zero - in other words the light fermions should also be set to zero. This of course means that the light field theory is not supersymmetric! However the problem is that the condition (20) or its generalization is really too strong. It was obtained by ignoring the $`D_\alpha H`$ terms as well as the $`O(1/M)`$ terms. Of course as was pointed out earlier, these terms are zero if one imposes (20) so that one gets $`\mathrm{}H=0`$, but the condition itself is not necessary. The actual necessary condition which follows from (19) is a relation between the Kahler derivative terms and the fermionic superfield terms. This condition would then express the bosonic component of the heavy superfield in terms of the light bosonic field as well as squares of light fermionic fields. However if we are interested only in the scalar potential we do not need to keep these terms. Thus the condition (20) can clearly be still used if one is just interested in computing the potential for the light chiral scalars.
Unlike the case of rigid supersymmetry (20) is not a holomorphic equation since the Kaehler derivative involves the real function $`K(\mathrm{\Phi },\overline{\mathrm{\Phi })}`$. This means that the solution for the heavy field will not in general be a holomorphic function of the light fields and hence the light field theory in general will have a superpotential that is just one (or a constant) and the whole effect of integrating out the heavy field will be accounted for by changing the Kaehler potential. In fact the original potential should be expressed in terms of the Kaehler invariant function $`G=\mathrm{ln}K+\mathrm{ln}|W|^2`$before integrating out the heavy fields. In a companion paper de Alwis (2005) we show this explicitly in some examples that come from type IIB string theory compactifications with fluxes.
Now one might worry that the restriction on the range of the light field essentially forces us back to the global case. This would indeed be the case in an example such as (12). Evaluating (20) in this case we have,
$$D_HW=MH+\frac{1}{2}L^2+\overline{H}\frac{1}{2}(MH^2+HL^2)=0$$
Solving this for $`H`$ we may compute the scalar potential using the standard supergravity formula (expressed in terms of $`G`$) but now the question is whether it is consistent to keep the supergravity corrections given that the integrating out formula above, is valid only for $`|a|^2/M^2<<1`$ as in the global case discussed earlier. The point is that necessarily $`MM_p`$ so that supergravity corrections which in this model are $`O(|a|^2/M_p^2)`$ should also be ignored for consistency of the approximation.
However very often in string theoretic examples (such as those with flux compactifications) there is a constant in the superpotential $`W=W_0+\mathrm{}`$ where the ellipses denote field dependent terms. In these cases the supergravity corrections are indeed significant since $`D_LW=_LW+K_LW/M_p^2_LW+K_LW_0/M_P^2`$ and the second term may even be of $`(O(1)`$ even thought light field space is restricted to $`|a|^2<<M^2`$. So this does not necessarily force us to the global limit since in many examples of interest in string theory there would be a constant in the superpotential which is generically of the order of the Planck/String scale. Thus the extra piece in the Kaehler derivative (as compared to the ordinary derivative) of the superpotential has to be kept. So we may use the condition (22) with the understanding that it is to be used only for the scalar components of the superfields for the purpose of calculating the light scalar field potential, still remaining within the context of supergravity.
Acknowledgments:
I’m very grateful to Martin Rocek for several useful suggestions.
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# An Efficient Approximation Algorithm for Point Pattern Matching Under Noise**footnote *A preliminary version was presented at the 7th International Symposium, Latin American Theoretical Informatics (LATIN 2006) [15].
## 1 Introduction
The general problem of finding large similar common substructures in two point sets arises in many areas ranging from computer vision to structural bioinformatics. In this paper, we study one of the more general problems, known as the *largest common point set problem* (LCP), which has several variants to be discussed below.
#### Problem Statement.
Given two point sets in $`^3`$, $`P=\{p_1,\mathrm{},p_m\}`$ and $`Q=\{q_1,\mathrm{},q_n\}`$, and an error parameter $`ϵ0`$, we want to find a rigid motion $`\mu `$ that maximizes the cardinality of subset $`IQ`$, such that $`\mathrm{𝖽𝗂𝗌𝗍}(P,\mu (I))ϵ`$. For an optimal set $`I`$, denote $`|I|`$ by $`\mathrm{𝖫𝖢𝖯}(P,Q)`$. There are two commonly used distance measures between point sets: *Hausdorff distance* and *bottleneck distance*. The Hausdorff distance $`\mathrm{𝖽𝗂𝗌𝗍}(P,Q)`$ between two point sets $`P`$ and $`Q`$ is given by $`\mathrm{max}_{qQ}\mathrm{min}_{pP}pq`$. The bottleneck distance $`\mathrm{𝖽𝗂𝗌𝗍}(P,Q)`$ between two point sets $`P`$ and $`Q`$ is given by $`\mathrm{min}_f\mathrm{max}_{qQ}f(q)q`$, where $`f:QP`$ is an injection. Thus we get two versions of the LCP depending on which distance is used.
Another distinction that is made is between the *exact*-LCP and the *threshold*-LCP. In the former we have $`ϵ=0`$ and in the latter we have $`ϵ>0`$. The exact-LCP is computationally easier than the threshold-LCP; however, it is not useful when the data suffers from round-off and sampling errors, and when we wish to measure the resemblance between two point sets and do not expect exact matches. These problems are better modeled by the threshold-LCP, which turns out to be harder, and various kinds of approximation algorithms have been considered for it in the literature (see below). A special kind of threshold-LCP in which one assumes that the minimum interpoint distance is greater than the error parameter $`2ϵ`$ is called *tolerant*-LCP. tolerant-LCP more accurately captures many problems arising in practice, and it appears that it is algorithmically easier than threshold-LCP. Notice that for the tolerant-LCP, the Hausdorff and bottleneck distances are essentially the same in the sense that the problem has a solution of Hausdorff distance $`ϵ`$ if and only if the solution is of bottleneck distance $`ϵ`$. Thus, for the tolerant-LCP, there is no need to specify which distance is in use.
In practice, it is often the case that the size of the solution set $`I`$ to the LCP is required to be at least a certain fraction of the minimum of the sizes of the two point sets: $`|I|\frac{1}{\alpha }\mathrm{min}(|P|,|Q|),`$ where $`\alpha `$ is a positive constant. This version of the LCP is known as the $`\alpha `$-LCP. A special case of the LCP which requires matching the entire set $`Q`$ is called Pattern Matching (PM) problem. Again, we have exact-PM, threshold-PM, and tolerant-PM versions.
In this paper, we focus on approximation algorithms for tolerant-LCP and tolerant-$`\alpha `$-LCP. There are two natural notions of approximation. (1) *Distance approximation:* The algorithm finds a transformation that brings a set $`IQ`$ of size at least $`\mathrm{𝖫𝖢𝖯}(P,Q)`$ within distance $`ϵ^{}`$ for some constant $`ϵ^{}>ϵ`$. (2) *Size-approximation:* The algorithm guarantees that $`|I|(1\delta )\mathrm{𝖫𝖢𝖯}(P,Q)`$, for constant $`\delta [0,1)`$.
#### Previous work.
The LCP has been extensively investigated in computer vision (e.g. ), computational geometry (e.g. ), and also finds applications in computational structural biology (e.g. ). For the exact-LCP problem, there are four simple and popular algorithms: alignment (e.g. ), pose clustering (e.g. ), geometric hashing (e.g. ) and generalized Hough transform (GHT) (e.g. ). These algorithms are often confused with one another in the literature. For convenience of the reader, we include brief descriptions of these algorithms in the appendix. Among these four algorithms, the most efficient algorithm is GHT.
#### Exact algorithms for tolerant-LCP.
As we mentioned above, the tolerant-LCP (or more generally, threshold-LCP) is a better model of many situations that arise in practice. However, it turns out that it is considerably more difficult to solve the tolerant-LCP than the exact-LCP. Intuitively, a fundamental difference between the two problems lies in the fact that for the exact-LCP the set of rigid motions, that may potentially correspond to the solution, is discrete and can be easily enumerated. Indeed, the algorithms for the exact-LCP are all based on the (explicit or implicit) enumeration of rigid motions that can be obtained by matching triplets to triplets. On the other hand, for the tolerant-LCP this set is continuous, and hence the direct enumeration strategies do not work. Nevertheless, the optimal rigid motions can be characterized by a set of high degree polynomial equations as in . A similar characterization was made by Alt and Guibas in for the 2D tolerant-PM problem and by the authors in for the 3D tolerant-PM. All known algorithms for the threshold-LCP use these characterizations and involve solving systems of high degree equations which leads to “numerical instability problem” . Note that exact-LCP and the exact solution for tolerant-LCP are two distinct problems. (Readers are cautioned not to confuse these two problems as in Gavrilov et al. .) Ambühl et al. gave an algorithm for tolerant-LCP with running time $`O(m^{16}n^{16}\sqrt{m+n})`$. The algorithm in for threshold-PM can be adapted to solve the tolerant-LCP in $`O(m^6n^6(m+n)^{2.5})`$ time. Both algorithms are for bottleneck distances. These algorithms can be modified to solve threshold-LCP under Hausdorff distance with a better running time by replacing the maximum bipartite graph matching algorithm which runs in $`O(n^{2.5})`$ with the $`O(n\mathrm{log}n)`$ time algorithm for nearest neighbor search. Both of these algorithms are for the general threshold-LCP, but to the best of our knowledge, these algorithms are the only known exact algorithms for the tolerant-LCP also.
#### Approximation algorithms for tolerant-LCP.
Like threshold-LCP, the exact algorithm for threshold-PM is difficult, even in 2D (see ). Two types of approximation algorithms were studied. First, Goodrich et al showed that there is a small discrete set of rigid motions which contains a rigid motion approximating (in distance) the optimal rigid motion for the threshold-PM problem, and thus the threshold-PM problem can be solved approximately by an enumeration strategy. Based on this idea and the alignment approach of enumerating all possible such discrete rigid motions, Akutsu , and Biswas and Chakraborty gave distance-approximation algorithms with running time $`O(m^4n^4\sqrt{m+n})`$ for the threshold-LCP under bottleneck distance, which can be modified to give $`O(m^3n^4\mathrm{log}m)`$ time algorithm for the tolerant-LCP. Second, Heffernan and Schirra introduced approximate decision algorithms to approximate the minimum Hausdorff distance between two point sets. Given $`ϵ>0`$, their algorithm answers correctly (YES/NO) if $`ϵ`$ is not too close to the optimal value $`ϵ^{}`$ (which is the minimum Hausdorff distance between the two point sets) and DON’T KNOW if the answer is too close to the optimal value. Notice that this approximation framework can not be “similarly” adopted to the LCP problem because in the LCP case there are two parameters – size and distance – to be optimized. This appears to be mistaken by Indyk et al. in where their approximation algorithm for tolerant-LCP is not well defined. Cardoze and Schulman gave an approximation algorithm (with possible false positives) but the transformations are restricted to translations for the LCP problem. Given $`\alpha `$, let $`ϵ_{min}(\alpha )`$ denote the smallest $`ϵ`$ for which $`\alpha `$-LCP exists; given $`ϵ`$, let $`\alpha _{min}(ϵ)`$ denote the smallest $`\alpha `$ for which $`\alpha `$-LCP exists. Biswas and Chakraborty combined the idea from Heffernan and Schirra and the algorithm of Akutsu to give a size-approximation algorithm which returns $`\alpha _u>\alpha _l`$ such that $`\mathrm{min}\{\alpha :ϵ>8ϵ(\alpha )\}\alpha _u\alpha _{min}(ϵ)`$ and $`\alpha _{min}(ϵ)>\alpha _l\mathrm{max}\{\alpha :ϵ<\frac{1}{8}ϵ_{min}(\alpha )\}`$. However, all these approximation algorithms still take high running time of $`\stackrel{~}{O}(m^3n^4)`$ (the notation $`\stackrel{~}{O}`$ hides poly log factors in $`m`$ and $`n`$).
#### Heuristics for tolerant-LCP.
In practice, the tolerant-LCP is solved heuristically by using the geometric hashing and GHT algorithms for which rigorous analyses are only known for the exact-LCP. For example, the algorithms in are for tolerant-LCP but the analyses are for exact-LCP only. Because of its practical performance, the exact version of GHT was carefully analyzed by Akutsu et al. , and a randomized version of the exact version of geometric hashing in 2D was given by Irani and Raghavan . The tolerant version of GHT (and geometric hashing) is based on the corresponding exact version by replacing the exact matching with the approximate matching which requires a distance measure to compare the keys. We can no longer identify the optimal rigid motion by the maximum votes as in the exact case. Instead, the tolerant version of GHT *clusters* the rigid motions (which are points in a six-dimensional space) and heuristically approximates the optimal rigid motion by a rigid motion in the largest cluster. Thus besides not giving any guarantees about the solution, this heuristic requires clustering in six dimensions, which is computationally expensive.
#### Other Related Work.
There is some closely related work that aims at computing the minimum Hausdorff distance for PM (see, e.g., and references therein). Also, the problems we are considering can be thought of as the point pattern matching problem under uniform distortion. Recently, there has been some work on point pattern matching under non-uniform distortion .
#### Our results.
There are three results in this paper. First, we introduce a new distance-approximation algorithm for tolerant-LCP algorithm, called Diheda (because our algorithm is based on Dihedral Angle comparisons).
###### Theorem 1.1
Let $`P,Q^3`$ of size $`m`$ and $`n`$, with $`mn`$, and $`ϵ>0`$. Suppose that interpoint distances in $`P`$ and in $`Q`$ be $`>2ϵ`$ (this is the condition for tolerant-LCP). Diheda (see Algorithm 1) finds a rigid motion $`\mu `$ and a subset $`I`$ of $`Q`$ such that
* $`|I|\mathrm{𝖫𝖢𝖯}(P,Q)`$ and
* $`\mathrm{𝖽𝗂𝗌𝗍}(P,\mu (I))4ϵ`$
in $`O(m^3n^3\mathrm{log}m)`$ time.
Diheda is simple and more efficient than the known distance-approximation algorithms (which are alignment-based) for tolerant-LCP. The running time of Diheda is $`O(m^3n^3\mathrm{log}m)`$ in the worst case. For general input, we expect the algorithm to be much faster because it is simpler and more efficient than the previous heuristics that are known to be fast in practice. This is because our clustering step is simple (sorting linearly ordered data) while the clustering step in those heuristics requires clustering high-dimensional data.
Second, based on a combinatorial observation, we improve the algorithms for exact-$`\alpha `$-LCP by a linear factor for pose clustering or GHT and a quadratic factor for alignment or geometric hashing. This also corrects a mistake by Irani and Raghavan .
Finally, we achieve a similar speed-up for Diheda using a sampling approach based on expander graphs at the expense of approximation in the matched set size. We remark that this result is mainly of theoretical interest because of the large constant factor involved. Expander graphs have been used before in geometric optimization for fast deterministic algorithms ; however, the way we use these graphs appears to be new. Our results also hold when we extend the set of transformations to scaling; for simplicity we restrict ourselves to rigid motions in this paper.
#### Outline.
The paper is organized as follows. The rest of this section contains some preliminaries. In Section 2 we introduce our new distance-approximation algorithm for tolerant-LCP. In Section 3 we show how a simple deterministic sampling strategy based on the pigeonhole principle yields speed-ups for the exact-$`\alpha `$-LCP algorithms. In Section 4 we show how to use expander graphs to further speed up the Diheda algorithm for tolerant-$`\alpha `$-LCP at the expense of approximation in the matched set size. Section 5 is the conclusion. In the appendix, we recall and compare the existing four basic algorithms for exact-LCP: pose clustering, alignment, GHT and geometric hashing.
#### Terminology and Notation.
For a transformation $`\mu `$, denote by $`I_\mu `$ the set of points in $`\mu (Q)`$ that are within distance $`ϵ`$ of some point in $`P`$. We call $`I_\mu `$ the matched set of $`\mu `$ and say that $`\mu `$ is an $`|I_\mu |`$-matching. We call the transformation $`\mu `$ that maximizes $`|I_\mu |`$ the maximum matching transformation. A basis is a minimal (for containment relation) ordered tuple of points which is required to uniquely define a rigid motion. For example, in 2D every ordered pair is a basis; while in 3D, every non-collinear triplet is a basis. In Figure 1, a rigid motion in 3D is specified by mapping a basis $`(q_1,q_2,q_3)`$ to another basis $`(p_1,p_2,p_3)`$.
We call a key used to represent an ordered tuple $`S`$ a rigid motion invariant key if it satisfies the following: (1) the key remains the same for all $`\mu (S)`$ where $`\mu `$ is any rigid motion, and (2) for any two ordered tuples $`S`$ and $`S^{}`$ with the same rigid motion invariant key there is a unique rigid motion $`\mu `$ such that $`\mu (S)=S^{}`$. For example, as rigid motion preserves orientation and distances among points, given a non-degenerate triangle $`\mathrm{\Delta }`$, the 3 side lengths of $`\mathrm{\Delta }`$ together with the orientation (the sign of the determinant of the ordered triplet) form a rigid motion invariant key for $`\mathrm{\Delta }`$ in $`^3`$. Henceforth, for simplicity of exposition, in the description of our algorithms we will omit the orientation part of the key.
## 2 Diheda
In this section, we introduce a new distance-approximation algorithm, called Diheda , for tolerant-LCP. The algorithm is based on a simple geometric observation. It can be seen as an improvement of a known GHT-based heuristic such that the output has theoretical guarantees.
### 2.1 Review of GHT
First, we review the idea of the pair-based version of GHT for exact-LCP. See the appendix or for more details. For each congruent pair, say $`(p_1,p_2)`$ in $`P`$ and $`(q_1,q_2)`$ in $`Q`$, and for each of the remaining points $`pP`$ and $`qQ`$, if $`(q_1,q_2,q)`$ is congruent to $`(p_1,p_2,p)`$, compute the rigid motion $`\mu `$ that matches $`(q_1,q_2,q)`$ to $`(p_1,p_2,p)`$. We then cast one vote for $`\mu `$. The rigid motion that receives the maximum number of votes corresponds to the maximum matching transformation sought. See Figure 1 for an example.
### 2.2 Comparable rigid motions by dihedral angles
For the exact-LCP, one only needs to compare rigid motions by equality (for voting). For the tolerant-LCP, one needs to measure how close two rigid motions are. In $`^3`$, each rigid motion can be described by 6 parameters (3 for translations and 3 for rotations). How to define a distance measure between rigid motions? We will show below that the rigid motions considered in our algorithm are related to each other in a simple way that enables a natural notion of distance between the rigid motions.
#### Observation.
In the pair-based version of GHT as described above, the rigid motions to be compared have a special property: the rigid motions transform a common pair — they all match $`(q_1,q_2)`$ to $`(p_1,p_2)`$ in Figure 1. Two such transformations no longer differ in all 6 parameters but differ in only one parameter. To see this, we first recall that a dihedral angle is the angle between two intersecting planes; see Figure 2 for an example.
In general, we can decompose the rigid motion for matching $`(q_1,q_2,q_3)`$ to $`(p_1,p_2,p_3)`$ into two parts: first, we transform $`(q_1,q_2)`$ to $`(p_1,p_2)`$ by a transformation $`\varphi _1`$; then we rotate the point $`\varphi _1(q_3)`$ about $`\stackrel{}{p_1p_2}`$ by an angle $`\theta `$, where $`\theta `$ is the dihedral angle between the planes $`(p_1,p_2,p_3)`$ and $`(\varphi _1(q_1),\varphi _1(q_2),\varphi _1(q_3))`$. This will bring $`q_3`$ to coincide with $`p_3`$. Thus, we have the following lemma:
###### Lemma 2.1
Let $`(p_1,p_2,p_3)`$ and $`(q_1,q_2,q_3)`$ be two congruent non-collinear triplets, and let $`\varphi _1`$ be a rigid motion that takes $`q_i`$ to $`p_i`$ for $`i=1,2`$. Let $`\varphi _2`$ be the rotation about $`\stackrel{}{p_1p_2}`$ by an angle $`\theta `$, where $`\theta `$ is the dihedral angle between the planes $`(p_1,p_2,p_3)`$ and $`(\varphi _1(q_1),\varphi _1(q_2),\varphi _1(q_3))`$. Then the unique rigid motion that takes $`(p_1,p_2,p_3)`$ to $`(q_1,q_2,q_3)`$ is equal to $`\varphi _2\varphi _1`$.
We now state another lemma that will be useful in the description and proof of correctness of Diheda . Let $`(p_1,p_2,p)`$ and $`q`$ be four points as shown in Figure 2. Consider the rotations about $`\stackrel{}{p_1p_2}`$ that take $`q`$ to within $`ϵ`$ of $`p`$. The rotation angles of these transformations form a subinterval of $`[0,2\pi )`$. This is because a circle $`C`$ (corresponding to the trajectory of $`p`$) intersects with the sphere $`B`$ (around $`p`$ with radius $`ϵ`$) at at most two points (corresponding to a subinterval of $`[0,2\pi )`$), as shown in Figure 2. That is, we have the following lemma:
###### Lemma 2.2
Let $`p_1,p_2,p,q^3`$ be four points (not necessarily non-collinear), then the rotation angles of transformations that rotate $`q`$ about $`\stackrel{}{p_1p_2}`$ to within $`ϵ`$ of $`p`$ form a subinterval of $`[0,2\pi )`$.
### 2.3 Approximating the optimal rigid motion by the “diametric” rigid motion
For a point set $`S^3`$, we call a pair of points $`\{p,q\}S^2`$ *diameter-pair* if $`pq=\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S)`$. A rigid motion of $`Q`$ that takes $`q_1`$ to $`p_1`$ and $`q_2`$ on the line $`p_1p_2`$ and closest possible to $`p_2`$ is called a $`(p_1,p_2,q_1,q_2)`$-rigid motion. Based on an idea similar to the one behind Lemma 2.4 in Goodrich et al. , we have the following lemma:
###### Lemma 2.3
Let $`\mu `$ be a rigid motion such that each point of $`\mu (S)`$, where $`SQ`$, is within distance $`ϵ`$ of a point in $`P`$. Let $`\{q_1,q_2\}`$ be a diameter-pair of $`S`$. Let $`p_iP`$ be the closest point to $`\mu (q_i)`$ for $`i=1,2.`$ Then we have a $`(p_1,p_2,q_1,q_2)`$-rigid motion $`\mu ^{}`$ of $`Q`$ such that each point of $`\mu ^{}(S)`$ is within $`4ϵ`$ of a point in $`P`$.
Proof Sketch. Translate $`\mu (q_1)`$ to $`p_1`$; this translation shifts each point by at most $`ϵ`$. Next, rotate about $`p_1`$ such that $`\mu (q_2)`$ is closest to $`p_2`$ (which implies $`\mu ^{}(q_1),\mu ^{}(q_2)`$ and $`p_2`$ are collinear). Since $`\{q_1,q_2\}`$ is a diameter-pair, this rotation moves each point by at most $`2ϵ`$. Thus, each point is at most $`ϵ+ϵ+2ϵ=4ϵ`$ from its matched point.
### 2.4 Approximation algorithm for tolerant-LCP
We first describe the idea of our algorithm Diheda . Input is two point sets in $`^3`$, $`P=\{p_1,\mathrm{},p_m\}`$ and $`Q=\{q_1,\mathrm{},q_n\}`$ with $`mn`$, and $`ϵ0`$. Suppose that the optimal rigid motion $`\mu _0`$ was achieved by matching a set $`I_{\mu _0}=\{q_1,q_2,\mathrm{},q_k\}Q`$ to $`J_{\mu _0}=\{p_1,p_2,\mathrm{},p_k\}P`$. WLOG, assume that $`\{q_1,q_2\}`$ is the diameter pair of $`I_{\mu _0}`$. Then by Lemma 2.3, there exists a $`(p_1,p_2,q_1,q_2)`$-rigid motion $`\mu `$ of $`Q`$ such that $`\mu (I_{\mu _0})`$ is within $`4ϵ`$ of a point in $`P`$. Since we do not know the matched set, we do not know a diameter-pair for the matched set either. Therefore, we exhaustively go through each possible pair. Namely, for each pair $`(q_1,q_2)Q`$ and each pair $`(p_1,p_2)P`$, if they are approximately congruent then we find a $`(p_1,p_2,q_1,q_2)`$-rigid motion $`\mu `$ of $`Q`$ that matches as many remaining points as possible. Note that $`(p_1,p_2,q_1,q_2)`$-rigid motions are determined up to a rotation about the line $`p_1p_2`$. By Lemma 2.2, the rotation angles that bring $`\mu (q_i)`$ to within $`4ϵ`$ of $`p_i`$ form a subinterval of $`[0,2\pi )`$. And the number of non-empty intersection subintervals corresponds to the size of the matched set. Thus, to find $`\mu `$, for each pair $`(p,q)P\{p_1,p_2\}\times Q\{q_1,q_2\}`$, we compute the dihedral angle interval according to Lemma 2.2. The rigid motion $`\mu `$ sought corresponds to an angle $`\varphi `$ that lies in the maximum number of dihedral intervals. The details of the algorithm are described in Algorithm 1.
Time Complexity. For each triplet in $`Q`$, using kd-tree for range query, it takes $`O(m^{3(1\frac{1}{3})}+m^3+m^3\mathrm{log}m^2)=O(m^3\mathrm{log}m)`$ for lines 11–20. For each pair $`(q_1,q_2)`$ and $`(p_1,p_2)`$, we spend time $`O(mn)`$ to find the subintervals for the dihedral angles, and time $`O(mn\mathrm{log}m)`$ to sort these subintervals and do the scan to find an angle that lies in the maximum number of subintervals. Thus the total time is $`O(m^3n^3\mathrm{log}m)`$.
## 3 Improvement by pigeonhole principle
In this section we show how a simple deterministic sampling strategy based on the pigeonhole principle yields speed-ups for the four basic algorithms for exact-$`\alpha `$-LCP. Specifically, we get a linear speed-up for pose clustering and GHT, and quadratic speed-up for alignment and geometric hashing. It appears to have been erroneously concluded previously that no such improvements were possible deterministically .
In pose clustering or GHT, suppose we know a pair $`(q_1,q_2)`$ in $`Q`$ that is in the sought matched set, then the transformation sought will be the one receiving the maximum number of votes among the transformations computed for $`(q_1,q_2)`$. Thus if we have chosen a pair $`(q_1,q_2)`$ that lies in the matched set, then the maximum matching transformation will be found. We are interested in the question “can we find a pair in the matched set without exhaustive enumeration”? The answer is yes: we only need to try a linear number of pairs $`(q_1,q_2)`$ to find the maximum matching transformation or conclude that there is none that matches at least $`\frac{n}{\alpha }`$ points.
We are given a set $`Q=\{q_1,\mathrm{},q_n\}`$, and let $`IQ`$ be an unknown set of size $`\frac{n}{\alpha }`$ for some constant $`\alpha >1`$. We need to discover a pair $`(p,q)`$ with $`p,qI`$ by using queries of the following type. A query consist of a pair $`(a,b)`$ with $`a,bQ`$. If we have $`a,bI`$, the answer to the query is YES, otherwise the answer is NO. Thus our goal is to devise a deterministic query scheme such that as few queries are needed as possible in the worst case (over the choice of $`I`$) before a query is answered YES. Similarly, one can ask the question about querying triplets to discover a triplet entirely in $`I`$.
###### Theorem 3.1
For an unknown set $`IQ`$ with $`|I|\frac{n}{\alpha }`$ and $`|Q|=n`$ using queries as described above,
(1) it suffices to query $`O(\alpha n)`$ pairs to discover a pair in $`I`$;
(2) it suffices to query $`O(\alpha ^2n)`$ triplets to discover a triplet in $`I`$.
Proof. The proof is based on the pigeonhole principle. To prove (1), we assume for simplicity that $`\alpha `$ and $`\frac{n}{\alpha }`$ are both integers. Partition the set $`Q`$ into $`\frac{n}{\alpha }`$ subsets of size $`\alpha `$ each. Since the size of $`I`$ is more than $`\frac{n}{\alpha }`$, by the pigeonhole principle, there is a pair of points in $`I`$ that lies in one of the above chosen subsets. Thus querying all pairs in these subsets will discover $`I`$. This gives that $`\frac{n}{\alpha }\left(\genfrac{}{}{0pt}{}{\alpha }{2}\right)\alpha n`$ queries are sufficient to discover $`I`$.
Similarly, to prove (2), partition $`Q`$ into $`\frac{n}{2\alpha }`$ subsets $`P_1,\mathrm{},P_{\frac{n}{2\alpha }}`$ of size $`2\alpha `$ each (we assume, as before, that $`2\alpha `$ and $`\frac{n}{2\alpha }`$ are both integers). Now we test all triplets that lie in the $`P_i`$’s. Any set $`IQ`$ that intersects with each of the $`P_i`$’s in at most $`2`$ points has size $`\frac{n}{\alpha }`$. Hence if $`|I|>\frac{n}{\alpha }`$ then it must intersect with one of the sets above in at least $`3`$ points. Thus testing the triplets from the $`P_i`$’s is sufficient to discover $`I`$. The number of triplets tested is $`\frac{n}{2\alpha }\left(\genfrac{}{}{0pt}{}{2\alpha }{3}\right)\alpha ^2n`$.
Remark: It can be shown that the schemes in the proof above are the best possible in requiring the smallest number of queries (up to constant factors).
In alignment and geometric hashing algorithms if we have chosen a triplet $`(q_1,q_2,q_3)`$ from the maximum matching set $`IQ`$ then we will discover $`I`$. The question, as before, is how many triplets in $`Q`$ need to be queried to discover a set $`I`$ of size $`>\frac{n}{\alpha }`$. By Theorem 3.1 (2), we only need to query $`O(\alpha ^2n)`$ triplets. Thus the running times of both alignment and geometric hashing are improved by a factor of $`\mathrm{\Theta }(n^2)`$.
See Table 1 for the time complexity comparison of deterministic algorithms for exact-$`\alpha `$-LCP in $`^3`$.
Finally, our approximation algorithm for tolerant-LCP adapts naturally for exact-$`\alpha `$-LCP with pigeonhole sampling. We analyze the running time of our algorithm for exact-$`\alpha `$-LCP with the pigeonhole sampling of pairs. In the exact case, each exact matched pair of points $`(q,p)`$ corresponds to a single dihedral angle. We thus find the dihedral angle that occurs the maximum number of times by sorting all the dihedral angles. For a fixed pair $`(q_1,q_2)`$ and a point $`q`$ in $`Q`$ the number of triplets $`((p_1,p_2),p_3)`$ in $`P`$ that match $`((q_1,q_2),q_3)`$ is bounded above by $`3H_2(m)`$, where $`H_2(m)`$ is the maximum possible number of the congruent triangles in a point set of size $`m`$ in $`R^3`$. Total time spent for pair $`(q_1,q_2)`$ then is $`O(nH_2)`$. Since we use $`O(\alpha n)`$ pairs, the overall running time is $`O(\alpha n^2H_2)`$. Agarwal and Sharir show that $`H_2(m)m^{\frac{5}{3}}g(m)`$, where $`g(m)`$ is a very slowly growing function of $`m`$ of inverse-Ackermann type.
As is often the case for algorithms for LCP, analysis involves determining quantities such as $`H_2(m)`$, which is a difficult problem. In the above table we have tried to give references for the first four algorithms including the tightest analyses rather than the original sources. Note that our algorithm is simpler than the others in the first column which involve checking for congruent simplices in a dictionary.
## 4 Expander-based sampling
While for the exact-$`\alpha `$-LCP the simple pigeonhole sampling served us well, for the tolerant-$`\alpha `$-LCP we do not know any such simple scheme for choosing pairs. The reason is that now we not only need to guarantee that each large set contain some sampled pairs, but also that each large set contain a sampled pair with large length (diameter-pair) as needed for the application of Lemma 2.3 in the Diheda algorithm. Our approach is based on expander graphs (see, e.g., ). Informally, expander graphs have linear number of edges but the edges are “well-spread” in the sense that there is an edge between any two sufficiently large disjoint subsets of vertices. Let $`G`$ be an expander graph with $`Q`$ as its vertex set. We show that for each $`SQ`$, if $`|S|`$ is not too small, then there is an edge $`(u,v)`$ in $`G`$ such that $`(u,v)S^2`$ and $`uv`$ approximates the diameter of $`S`$.
By choosing the pairs for the Diheda algorithm from the edge set of $`G`$ (the rest of the algorithm is same as before), we obtain a bicriteria – distance and size – approximation algorithm as stated in Theorem 4.4 below. We first give a few definitions and recall a result about expander graphs that we will need to prove the correctness of our algorithm.
###### Definition 4.1
Let $`S`$ be a finite set of points of $`^r`$ for $`r1`$, and let $`0kn`$. Define $`\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,k)=\mathrm{min}_{T:|T|=k}\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(ST).`$
That is, $`\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,k)`$ is the minimum of the diameter of the sets obtained by deleting $`k`$ points from $`S`$. Clearly, $`\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,0)=\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S)`$.
Let $`U`$ and $`V`$ be two disjoint subsets of vertices of a graph $`G`$. Denote by $`e(U,V)`$ the set of edges in $`G`$ with one end in $`U`$ and the other in $`V`$. We will make use of the following well-known theorem about the eigenvalues of graphs (see, e.g. , for the proof and related background).
###### Theorem 4.2
Let $`G`$ be a $`d`$-regular graph on $`n`$ vertices. Let $`d=\lambda _1\lambda _2\mathrm{}\lambda _n`$ be the eigenvalues of the adjacency matrix of $`G`$. Denote $`\lambda =\mathrm{max}_{2in}|\lambda _i|.`$ Then for every two disjoint subsets $`U,WV`$,
$$\left||e(U,W)|\frac{d|U||W|}{n}\right|\lambda \sqrt{|U||W|}.$$
(1)
###### Corollary 4.3
Let $`U,WV`$ be two disjoint sets with $`|U|=|W|>\frac{\lambda n}{d}`$. Then $`G`$ has an edge in $`U\times W`$.
Proof. It follows from (1) that if $`\frac{d|U||W|}{n}>\lambda \sqrt{|U||W|}`$ then $`|e(U,W)|>0`$, and since $`|e(U,W)|`$ is integral, $`|e(U,W)|1`$. But the above condition is clearly true if we take $`U`$ and $`W`$ as in the statement of Corollary 4.3.
There are efficient constructions of graph families known with $`\lambda <2\sqrt{d}`$ (see, e.g., ). Let us call such graphs *good expander graphs*. We can now state our main result for this section.
###### Theorem 4.4
For an $`\alpha `$-LCP instance $`(P,Q)`$ with $`\mathrm{𝖫𝖢𝖯}(P,Q)>\frac{n}{\alpha }`$, the Diheda algorithm with expander-based sampling using a good expander graph of degree $`d>2500\alpha ^2`$ finds a rigid motion $`\mu `$ in time $`O(m^3n^2\mathrm{log}m)`$ such that there is a subset $`I`$ satisfying the following criteria:
(1) size-approximation criterion: $`|I|\mathrm{𝖫𝖢𝖯}(P,Q)\frac{50}{\sqrt{d}}n`$;
(2) distance-approximation criterion: each point of $`\mu (I)`$ is within distance $`6ϵ`$ from a point in $`P`$.
Thus by choosing $`d`$ large enough we can get as good size-approximation as desired. The constants in the above theorem have been chosen for simplicity of the proof and can be improved slightly.
For the proof we first need a lemma showing that choosing the query pairs from a graph with small $`\lambda (G)`$ (the second largest eigenvalue of $`G`$) gives a long (in a well-defined sense) edge in every not too small subset of vertices.
###### Lemma 4.5
Let $`G`$ be a $`d`$-regular graph with vertex set $`Q^3`$, and $`|Q|=n`$. Let $`SQ`$ be such that $`|S|>\frac{25\lambda (G)n}{d}`$. Then there is an edge $`\{s_1,s_2\}E(G)S^2`$ such that $`s_1s_2\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,\frac{25\lambda (G)}{d}n)}{2}`$.
Proof. For a positive constant $`c`$ to be chosen later, remove $`cn`$ pairs from $`S`$ as follows. First remove a diameter pair, then from the remaining points remove a diameter pair, and so on. Let $`T`$ be the set of points in the removed pairs and $`T^p`$ the set of removed pairs. The remaining set $`ST`$ has diameter $`\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)`$ by the definition of $`\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)`$, and hence each of the removed pairs has length $`\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)`$. For $`B,CS`$ let $`||B,C||=\mathrm{min}_{bB,cC}||bc||`$.
###### Claim 1
The set $`T`$ defined above can be partitioned into three sets $`B`$, $`C`$, $`E`$, such that $`|B|,|C|\frac{cn}{6}`$, and $`||B,C||\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)}{2}`$.
Proof. Fix a Cartesian coordinate system and consider the projections of the pairs in $`T^p`$ on the $`x`$-, $`y`$\- ,and $`z`$-axes. It is easy to see that for at least one of these axes, at least $`\frac{cn}{3}`$ pairs have projections of length $`\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)}{\sqrt{3}}`$. Suppose without loss of generality that this is the case for the $`x`$-axis, and denote the set of projections of pairs on the $`x`$-axis with length $`\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)}{\sqrt{3}}`$ by $`T_x^p`$, and the set of points in the pairs in $`T_x^p`$ by $`T_x`$. We have $`|T_x|2cn/3`$. Now consider a sliding window $`W`$ on the $`x`$-axis of length $`\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)}{2}`$, initially at $`\mathrm{}`$, and slide it to $`+\mathrm{}`$. At any position of $`W`$, each pair in $`T_x^p`$ has at most $`1`$ point in $`W`$, as the length of any pair is more than the length of $`W`$. Thus at any position, $`W`$ contains $`|T_x^p|=|T_x|/2`$ points. It is now easy to see by a standard continuity argument that there is a position of $`W`$, call it $`\overline{W}`$, where there are $`\frac{|T_x|}{4}\frac{cn}{6}`$ points of $`T_x`$ both to the left and to the right of $`\overline{W}`$.
Now, $`B`$ is defined to be the set of points in $`T`$ whose projection is in $`T_x`$ and is to the left of $`\overline{W}`$; similarly $`C`$ is the set of points in $`T`$ whose projection is in $`T_x`$ and is to the right of $`\overline{W}`$. Clearly any two points, one from $`B`$ and the other from $`C`$, are $`\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,2cn)}{2}`$-apart.
Coming back to the proof of Lemma 4.5, the property that we need from the query-graph is that for any two disjoint sets $`B,CS`$ of size $`\delta |S|`$, where $`\delta `$ is a small positive constant, the query-graph should have an edge in $`B\times C`$.
By Corollary 4.3 if $`|B|\frac{cn}{6}>\frac{\lambda n}{d}`$, and $`|C|\frac{cn}{6}>\frac{\lambda n}{d}`$, that is, if $`c>\frac{6\lambda }{d}`$, then $`G`$ has an edge in $`B\times C`$. Taking $`c=\frac{12.5\lambda }{d}`$ completes the proof of Lemma 4.5.
Proof of Theorem 4.4. If we take $`G`$ to be a good expander graph then Lemma 4.5 gives that $`G`$ has an edge of length $`\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,\frac{50}{\sqrt{d}}n)}{2}`$. Let $`S`$ also be a solution to tolerant-LCP for input $`(P,Q)`$ with error parameter $`ϵ>0`$. We have that one of the sampled pairs has length at least $`\frac{\mathrm{𝖽𝗂𝖺𝗆𝖾𝗍𝖾𝗋}(S,\frac{50}{\sqrt{d}}n)}{2}`$. Thus applying an appropriate variant (replacing the diameter pair by the sampled pair with large length as guaranteed by Lemma 4.5) of Lemma 2.3, we get a rigid motion $`\mu `$ such that there is a subset $`I`$ satisfying the following:
(1) $`|I||S|\frac{50}{\sqrt{d}}n`$ for any $`d>2500\alpha ^2`$;
(2) Each point of $`I`$ is within $`6ϵ(=ϵ+ϵ+4ϵ)`$ of a point in $`M`$.
## 5 Discussion
We have presented a new practical algorithm for point pattern matching. Our Diheda algorithm is the fastest known distance-approximation algorithm for tolerant-LCP, and is simple compared to other known distance-approximation algorithms and heuristics which involve 6-dimensional clustering. Our analysis of Diheda is not tight, and perhaps better bounds can be obtained if the interpoint distance is greater than $`ϵ`$ by a sufficiently large constant factor.
Our technique of pigeonhole sampling yields speed-ups for all four popular algorithms and also the fastest known deterministic algorithm for the exact-LCP. Again, our algorithms are simpler than the previous best algorithms. Akutsu et al. give a tighter analysis for GHT in terms of the function $`\lambda ^{3,2}(m,n)`$. Our analysis of Diheda (and GHT) with pigeonhole sampling was based on $`H_2(m)`$. Presumably, a better analysis similar to the idea in is possible.
Point pattern matching is of fundamental importance for computer vision and structural bioinformatics. Indeed, this investigation stemmed from research in structural bioinformatics. Current software, which uses either geometric hashing or generalized Hough transform, can immediately benefit from this work. We have implemented a randomized version of Diheda for molecular common substructure detection and the results were reported in .
Acknowledgment. We thank S. Muthukrishnan and Ali Shokoufandeh for the helpful comments and advice.
Appendix
## Appendix A Voting Algorithms for Exact-LCP
In this appendix, we review and compare four popular algorithms for exact-LCP: pose clustering, alignment, generalized Hough transform(GHT), and geometric hashing. These algorithms are all based on a voting idea and are sometimes confused in the literature. Please see Algorithms 2, 3, 4, 5) for a full description of the algorithms in their generic form independent of the search data structure used. In particular, geometric hashing algorithms need not use a hash-table as a search data structure. We describe all the algorithms in terms of a dictionary of objects (which are either transformations or a set of points and can be ordered lexicographically). Denote the query time for this dictionary by $`S(x)+O(k)`$ where $`x`$ is the size of the dictionary, and $`k`$ is the size of the output depending on the query. For example, if the dictionary is implemented by a search tree we have $`S(x)=O(\mathrm{log}x)`$.
Pose clustering and alignment are the basic methods. GHT and geometric hashing can be regarded as their respective efficient implementations. Efficiency is achieved by preprocessing of the point sets using their rigid motion invariant keys which speeds-up the searches.
In pose clustering, for each pair of triplets $`(q_1,q_2,q_3)Q`$ and $`(p_1,p_2,p_3)P`$, we check if they are congruent. If they are then we compute the rigid motion $`\mu `$ such that $`\mu (q_1,q_2,q_3)=(p_1,p_2,p_3)`$. We then cast one vote for $`\mu `$. The rigid motion which receives the maximum number of votes corresponds to the maximum matching transformation sought. The running time of pose clustering is $`O(m^3n^3S(m^3n^3))`$ as the size of the dictionary of transformations can be as large as $`O(m^3n^3)`$.
In alignment, for each pair of triplets $`(q_1,q_2,q_3)Q`$ and $`(p_1,p_2,p_3)P`$ we check if they are congruent. If they are then we compute the rigid motion $`\mu `$ such that $`\mu (q_1,q_2,q_3)=(p_1,p_2,p_3)`$. Then we count the number of points in $`\mu (Q)`$ that coincide with points in $`P`$. This number gives the number of votes the rigid motion $`\mu `$ gets. The rigid motion which receives the maximum number of votes corresponds to the maximum matching transformation sought. The running time is $`O(m^3n^4S(m))`$.
The difference between pose clustering and alignment is the voting space: in pose clustering voting is done for transformations while in alignment it is for bases (triplets of points). In both pose clustering and alignment algorithms, each possible triplet in $`Q`$ is compared with each possible triplet in $`P`$. However, by representing each triplet with its rigid motion invariant key, only triplets with the same key (rigid motion invariant) are needed to be compared. This provides an efficient implementation. For example, the GHT algorithm is an efficient implementation of pose clustering. Here we preprocess $`P`$ by storing the triplets of points with the rigid motion invariant keys in a dictionary. Now for each triplet $`(q_1,q_2,q_3)`$ in $`Q`$ we find congruent triplets in $`P`$ by searching for the rigid motion invariant key for $`(q_1,q_2,q_3)`$. The rest of the algorithm is the same as pose clustering. Similarly the geometric hashing algorithm is an efficient implementation of the alignment method.
GHT is faster than geometric hashing, however geometric hashing has the advantage that algorithm can stop as soon as it has found a good match. Depending on the application this gives geometric hashing advantage over GHT.
As observed by Olson and Akutsu et al. , pose clustering and GHT can be further improved. This is because a $`k`$-matching transformation can be identified by matching $`(k2)`$ bases which match a common pair. We call this version of the generalized Hough transform the pair-based version; it is described below in Algorithm 6. Although the worst case time complexity of the pair-based version and the original version are the same, this will serve as a basis for our new scheme, called Diheda . The pair-based version also allows efficient random sampling of pairs .
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# Photon Echoes Produced by Switching Electric Fields
## Abstract
We demonstrate photon echoes in Eu<sup>3+</sup>:Y<sub>2</sub>SiO<sub>5</sub> by controlling the inhomogeneous broadening of the Eu<sup>3+</sup> <sup>7</sup>F$`{}_{0}{}^{}_{}^{5}`$D<sub>0</sub> optical transition. This transition has a linear Stark shift and we induce inhomogeneous broadening by applying an external electric field gradient. After optical excitation, reversing the polarity of the field rephases the ensemble, resulting in a photon echo. This is the first demonstration of such a photon echo and its application as a quantum memory is discussed.
In the emerging area of quantum information science the ability to store and recall quantum states of light it is highly desirable. To date all proposals to achieve this rely on the mapping of light states onto the states of atom-like systems. In spite of the significant effort that has been directed towards this goal, high fidelity reversible mapping of the state of a light field onto atom-like systems has not yet been demonstrated.
The use of atomic ensembles for quantum memory applications is attractive because it allows strong coupling without the need for the high finesse cavities of earlier proposals Cirac et al. (1997). Much attention has be been given to using the phenomena of electromagnetically induced transparency (EIT) to build quantum memories using atomic ensembles Fleischhauer and Lukin (2000, 2002). EIT has led to some dramatic experimental results including ultra-slow group velocities and light storage in trapped atomic systems Hau et al. (1999); Liu et al. (2001), atomic vapors Kash et al. (1999); Budker et al. (1999); Philips et al. (2001) and the solid state Turukin et al. (2002); Longdell et al. (2005). While EIT can lead to very slow group velocities it has proven difficult to achieve delays with large time-bandwidth products and without high losses Matsko et al. (2005). Low time-bandwidth products for delays make high fidelity light storage using EIT difficult. This is because the time-bandwidth product is a measure of how-many distinct pulses fit inside the delay medium and for high fidelity storage the entire pulse must remain inside the medium as the group velocity is slowly reduced to zero. In this work we explore an alternative ensemble based quantum memory on photon echoes.
Coherent manipulation and storage of classical light states using photon echo techniques dates back to the 1980’s Mossberg (1982). In contrast to EIT based techniques where the required slow group velocities are directly linked to narrow bandwidths, recent photon echo based experiments have demonstrated the ability to store thousands of pulses Lin et al. (1995) and do signal processing at gigahertz bandwidths Merkel et al. (2004). However there is no way currently known to use a standard photon echo, where the rephasing of the atomic coherence is achieved with an optical $`\pi `$-pulse, in a quantum memory.
In 2001 Moiseev and Kröll Moiseev and Kroll (2001) published a proposal for a quantum memory based on modified photon echos. Instead of rephasing the atomic coherence with intense optical pulses, it used reversible inhomogeneous broadening. Here we report the first demonstration of photon echoes produced by reversible inhomogeneous broadening. In this case this is achieved using optical centers in a solid which have a linear stark shift and macroscopic electric field gradients.
A quantum memory based on controlled inhomogeneous broadening can be understood as follows: consider a coherent light pulse entering a medium of two-level atoms. The Hamiltonian for the system is of the form:
$$H=H_{\text{field}}+\underset{n}{}\omega _n\sigma _n^z+d(\sigma _n^+E(z_n,t)+\sigma _n^{}E^{}(z_n,t))$$
(1)
Here the $`\sigma `$ are Pauli operators and $`d`$ the transition dipole moment. We assume that the area of the incoming pulse is small and that each individual atom is never driven far from its ground state, enabling the approximation $`\sigma _z=1/2`$. Working in this small-pulse regime linearizes the equations of motion allowing simple analytic expressions to be derived. It is also the regime of interest for a quantum memory.
The treatment here is semi-classical but because the equations of motion are linear the results should carry straight over to a fully quantum analysis. Treating the atoms as a continuous field $`\sigma _n^{}(t)\phi (\delta ,z,t)dzd\delta `$ one obtains the following equations of motion:
$`({\displaystyle \frac{^2}{z^2}}+{\displaystyle \frac{1}{c^2}}{\displaystyle \frac{^2}{t^2}})E(z,t)=\eta {\displaystyle 𝑑\delta g(\delta ,z)\frac{^2}{t^2}\phi (\delta ,z,t)}`$ (2)
$`{\displaystyle \frac{}{t}}\phi (\delta ,z,t)=i(\omega _0+\delta )\phi (\delta ,z,t)+dE(z,t)`$ (3)
Here $`g(\delta ,z)`$ describes the atom density as a function of detuning and position. Changing variables by setting
$`E(z,t)`$ $`=`$ $`E_f(z,t)\mathrm{exp}(i(kz\omega _0t))`$ (4)
$`\phi (\delta ,z,t)`$ $`=`$ $`\alpha (\delta ,z,t)\mathrm{exp}(i(kz\omega _0t))`$ (5)
and assuming that the amplitudes $`E_f`$ and $`\alpha `$ are slowly varying functions of both $`z`$ and $`t`$. We obtain
$`2ik\left(_z{\displaystyle \frac{1}{c}}_t\right)E_f`$ $`=`$ $`\eta \omega _0^2{\displaystyle 𝑑\delta g(\delta ,z)\alpha (\delta ,z,t)}`$
$`_t\alpha (\delta ,z,t)`$ $`=`$ $`i\delta \alpha (\delta ,z,t)+E_f(z,t)`$ (6)
If we narrow our focus briefly to the situation where we have a spatially homogeneous system with large inhomogeneous broadening ($`g(\delta ,z)=1`$) these have the analytical solution Crisp (1970)
$`E_f(z,t)z`$ $`=`$ $`e^{\eta z}E_f(0,tz/c)`$ (7)
$`\alpha (\delta ,z,t)`$ $`=`$ $`e^{(i\delta /c\eta )z}{\displaystyle _{\mathrm{}}^t}𝑑\tau E_f(0,\tau )e^{i\delta \tau }`$ (8)
As it propagates, the optical pulse exponentially decays as the atoms absorb the energy. Because of the inhomogeneous broadening the macroscopic coherence of the atomic ensemble also decays. However the process should not be seen as dissipative, each individual atom has not lost any coherence and the absorption process is reversible as described below. For the experiments we performed the the detuning is a function of position and the decay of the field is not exponential. The different spectral components of the input field are absorbed at different positions. However as long light pulse is could be totally absorbed the argument for how the memory works is unchanged.
Defining backwards propagating atomic ($`\beta `$) and optical ($`E_b`$) fields via
$`E(z,t)`$ $`=`$ $`E_b(z,t)\mathrm{exp}(i(kz\omega _0t))`$ (9)
$`\phi (\delta ,z,t)`$ $`=`$ $`\beta (\delta ,z,t)\mathrm{exp}(i(kz\omega _0t))`$ (10)
one arrives at the equations of motion
$`2ik\left(_z+{\displaystyle \frac{1}{c}}_t\right)E_b(z,t)=\eta \omega _0^2{\displaystyle 𝑑\delta g(\delta ,z)\beta (\delta ,z,t)}`$
$`_t\beta (\delta ,z,t)=i\delta \beta (\delta ,z,t)E_b(z,t)`$ (11)
Comparing Eqns. (11) with the time reversed versions of (6) it can be seen that the two coincide if the sign of $`\delta `$ is reversed. Thus all that is required to make the pulse come out again as a time reversed copy of itself is to flip the detunings of the atoms and at the same time to apply an phase matching operation such that the value of $`\beta `$ after the operation is equal to the value that $`\alpha `$ was before the operation. As $`\beta (\delta ,z,t)=\alpha (\delta ,z,t)\mathrm{exp}(i2kz)`$ this phase matching operation is a position dependent phase shift for the atomic states. This can be achieved by driving from the excited state to an auxiliary ground state with a $`\pi `$ pulse and driving back up again with another $`\pi `$ pulse such that the wavevector difference is $`2k`$. These two $`\pi `$ pulses can be separated in time and between them the coherence is stored in the hyperfine transitions. Coherence times of many seconds have been demonstrated for hyperfine transitions Fraval et al. (2005) in rare earth ion doped systems.
In Moiseev and Kroll (2001) it was proposed that the reversible inhomogeneous broadening be achieved using Doppler broadening in an atomic gas. Since then using impurity ions to achieve this controlled inhomogeneous broadening has been proposed Moiseev et al. (2003) and independent of this work the use of electric field gradients has also been proposed Nilsson and Kroll (2005). A theoretical treatment of the general case has also appeared Kraus et al. (2005).
In order for a quantum memory based on controlled inhomogeneous broadening to be practical it is necessary for the induced broadening to be larger than the unbroadened linewidth of the transition. A further requirement is that the field polarity be switched in a time short compared to the inverse of this linewidth. Here we show that these conditions can be achieved in rare earth ion doped systems by successfully demonstrating a Stark echo.
The optical transition used in this experiment was the <sup>7</sup>F$`{}_{0}{}^{}`$ <sup>5</sup>D<sub>0</sub> in <sup>151</sup>Eu, at 579.879 nm in 0.1 at% Eu<sup>3+</sup>:Y<sub>2</sub>SiO<sub>5</sub>. FIG. 1 shows the hyperfine structure of the two electronic singlet states. The transition was excited with linearly polarized light propagating along the C<sub>2</sub> axis of the crystal, with the polarization chosen to maximize the absorption. The length of the crystal in the direction of propagation was 4 mm. The crystal was cooled to below 4 K in a liquid helium bath cryostat. A quadrupole electric field was applied to the sample using four 10 mm long, 2 mm diameter rods in a quadrupolar arrangement as shown in FIG. 1. Two amplifiers with 1 MHz bandwidth supplied the voltage across the electrodes. These amplifiers had two opposite polarity outputs and voltage rails of $`\pm `$35 V. This configuration provided an electric field that varied linearly across the sample in the direction of light propagation with a maximum field gradient of approximately 300 Vcm<sup>-2</sup>.
The optical setup was essentially the same as in previous work Pryde et al. (2000); Longdell and Sellars (2004); Longdell et al. (2004). A highly stabilized dye laser was used with an established stability of better than 200 Hz over timescales of 0.2 s. The light incident on the sample was gated with two acousto-optic modulators (AOMs) in series. These allowed pulses with an arbitrary amplitude and phase envelope to be applied to the sample. A Mach-Zehnder interferometer arrangement with the AOMs and sample in one arm was employed to enable heterodyne detection of the coherent emission from the sample. The overall frequency shift introduced by the AOMs was 51 MHz. The intensity of the beat signal was detected with a photo-diode.
The linear Stark shift for the <sup>7</sup>F$`{}_{0}{}^{}`$ <sup>5</sup>D<sub>0</sub> transition in Eu<sup>3+</sup>:Y<sub>2</sub>SiO<sub>5</sub> has not been reported but from a study in YAlO<sub>3</sub> it is expected to be of the order of 35 kHzVcm<sup>-1</sup> Meixner et al. (1992). With the current experimental setup the anticipated Stark-induced spectral broadening was therefore 2 MHz. Although this broadening is large compared to the 122 Hz homogeneous linewidth of the optical transition it is significantly smaller than the inhomogeneous linewidth of our sample which was 3 GHz. To create an optical feature which was narrow compared to the induced broadening, the same optical pumping procedure as used in Longdell and Sellars (2004); Longdell et al. (2004) was employed. This consisted of burning a relatively wide ($``$3 MHz) spectral hole in the absorption line by scanning the laser frequency. A narrow anti-hole was placed in the middle of this region by applying RF excitation at 80.7 MHz as well as light at a different frequency. This frequency was given by the combination of ground and excited state hyperfine splittings as shown in FIG. 1. The spectral width of the anti-hole was then reduced by optically pumping, out of resonance, ions more than 12.5 kHz from the center frequency of the anti-hole. The peak absorption of the feature was approximately 40%.
Figure 2 (dotted) shows the free induction decay (FID) resulting from excitation of the spectral feature created as described above, with a 3 $`\mu `$s long optical pulse at $`\omega _p`$. The coherent emission was found to be consistent with that from a 25 kHz wide feature with a top hat profile.
In order to determine the degree of induced broadening, the FID measurement was repeated with an electric field gradient applied after the creation of the spectral feature. By measuring the length of FIDs as the voltage was increased, we estimate the rate of induced broadening to be 42 kHzV<sup>-1</sup>.
To observe an echo from a feature broadened by applying 25 V to the electrodes, we excited the feature using a 1 $`\mu `$s optical pulse. The polarity of the field was reversed after a time $`\tau `$, and after a further delay of $`\tau `$ the echo was observed.
FIG. 3 shows the echoes created using a varying delay between the input pulse and the electric field reversal. The intensity of the echo as a function of the delay is plotted in FIG. 4. From this it can be seen that the envelope of the echo amplitude has the same profile as the FID of the unbroadened spectral feature. The time-bandwidth product, or the number of distinct pulses that can be stored is four.
The intensity of the echo for a 1.8 $`\mu `$s long input pulse as a function of the input pulse intensity is shown in FIG. 5. At the highest input power of 7 mW the pulse area was $`\pi `$/2. The size of the pulse corresponding to a $`\pi `$/2 pulse was determined by nutation measurements. The output amplitude is seen to be linear for low input intensities but saturates for pulse areas approaching $`\pi `$/2.
In these initial experiments the efficiency of the echo was limited by the low level of absorption of the broadened spectral feature. The 40% absorption of the unbroadened line was reduced to 1% on the application of the electric field gradient. For quantum memory applications it will be necessary for the broadened feature to be optically thick.
One way to increase the optical thickness is to increase the interaction length, using longer samples or multi-pass cells. Alternately the absorption could be enhanced by placing the sample an optical cavity.
Another way to increase the optical depth is by increasing the spectral density of the ions. With doped samples it may not be possible to increase the spectral density of the ions by simply increasing the dopant level Sellars et al. (2004). This is due to the random strain in the crystal introduced by the dopant ions, producing inhomogeneous broadening that increases linearly with the dopant concentration Cone . An alternative strategy for achieving a high spectral density is to use stoichiometric materials such as EuCl<sub>3</sub>.6H<sub>2</sub>O. The spectral density in EuCl<sub>3</sub>.6H<sub>2</sub>O is high due to the high concentration of Eu<sup>3+</sup> ions and the low level of strain induced broadening.
In conclusion we have shown that it is possible to rephase optical coherence through reversing an external electric field gradient. We demonstrate a time-bandwidth product of four and linear operation at low input intensities. This rephasing using controlled inhomogeneous broadening is the first demonstration of this key effect which has potential applications as a quantum memory.
An elegant aspect of this memory scheme, when compared to other coherence optical memories, is that the only optical excitation is the signal to be stored. The limitation of the current work was the low optical thickness resulting in very low efficiency. Methods for overcoming this limitation were discussed.
The authors would like to thank D. Freeman for his careful reading of the manuscript. This work was supported by the Australian Defence Science and Technology Organisation and the Australian Research Council.
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# Additivity for the parametrized topological Euler characteristic and Reidemeister torsion
## 1. Introduction
In recent years various attempts have been made to generalize the classical Reidemeister torsion (which is an invariant of non-simply connected finite CW complexes) to a parametrized version, i.e. to an invariant of fiber bundles. One such generalization is the notion of an analytic torsion introduced by Bismut and Lott . It is defined for bundles of smooth manifolds satisfying some additional conditions. Another definition was proposed by Igusa and Klein , who construct parametrized Reidemeister torsion - also for bundles of smooth manifolds - using generalized Morse functions.
Our main interest in this paper lies in yet another definition of torsion developed by Dwyer, Weiss and Williams . One of the main features of their construction is that it is described in the language of the homotopy theory which makes it relatively simple. Briefly, it proceeds as follows. Given a manifold $`M`$ and a locally constant sheaf of $`R`$-modules $`\rho :VM`$ such that the homology groups $`H_{}(M,\rho )`$ vanish, the classical Reidemeister torsion of $`M`$ can be defined as an element of a group $`\mathrm{Wh}(\rho )`$ which is a certain quotient of $`K_1(R)`$. Let $`Q(M_+)`$ denote the infinite loops space associated with the suspension spectrum of $`M`$. The starting point for the construction of Dwyer, Weiss and Williams is the observation that the group $`\mathrm{Wh}(\rho )`$ can be identified with $`\pi _0\mathrm{\Phi }_\rho ^s(M)`$ where where $`\mathrm{\Phi }_\rho ^s(M)`$ is the homotopy fiber of a certain map $`\lambda :Q(M_+)K(R)`$ defined in terms of the sheaf $`\rho `$. It follows that the classical torsion of $`M`$ determines a connectedness component of $`\mathrm{\Phi }_\rho ^s(M)`$. In fact, more is true: one can construct the torsion invariant as a unique point $`\tau _\rho ^s(M)\mathrm{\Phi }_\rho ^s(M)`$. Assume now that we are given a smooth bundle $`p:EB`$ and a sheaf of $`R`$-modules $`\rho :VE`$ such that the homology groups of the fibers $`M_b`$ of $`p`$ with coefficients in $`\rho |_{M_b}`$ vanish. For every $`bB`$ we obtain a point $`\tau _\rho ^s(M_b)\mathrm{\Phi }_\rho ^s(M_b)`$. The spaces $`\mathrm{\Phi }_\rho ^s(M_b)`$ can be glued together using the topology of the bundle $`p`$ in such way that we obtain a fibration $`\mathrm{\Phi }_\rho ^s(p)B`$. The assignment $`b\tau _\rho ^s(M_b)`$ defines then a section $`\tau _\rho ^s(p)`$ of this fibration which can be interpreted as a parametrized smooth torsion of the bundle $`p`$.
An advantage of the construction of torsion sketched above is its flexibility: while $`\tau _\rho ^s(p)`$ is defined for smooth bundles of manifolds, its variants can be used to extend the notion of parametrized torsion to more general settings. In fact, Dwyer, Weiss and Williams gave two additional versions of their definition. The homotopy Reidemeister torsion $`\tau _\rho ^h(p)`$ is defined for any fibration $`p:EB`$ with homotopy finitely dominated fibers and is obtained just as the smooth torsion, but with $`Q(M_+)`$ replaced by $`A(M)`$ – the Waldhausen $`A`$-theory of $`M`$. The topological Reidemeister torsion $`\tau _\rho ^t(p)`$ exists whenever $`p:EB`$ happens to be a bundle of compact topological manifolds and is constructed using $`A^\%(M)`$ \- the excisive version of the $`A`$-theory.
Each of these invariants takes values in a different infinite loop space; thus topological torsion of a smooth bundle is not the same as its smooth torsion, but it can be seen as an approximation of the smooth torsion. Similarly, homotopy torsion can be considered as an approximation of the topological torsion when both are defined.
The work of Goette , and Igusa extended our understanding of the relationship between the analytical torsion of Bismut–Lott and the torsion of Igusa–Klein. It is far less clear, however, how these notions relate to the smooth torsion of Dwyer–Weiss–Williams. Our goal here is to bring these constructions closer together. The starting point is an axiomatization of parametrized Reidemeister torsion proposed by Igusa . He showed that if a cohomological version of torsion of smooth bundles satisfies two conditions, then it is unique up to a scalar multiple. The first condition is the product formula relating the torsion of a composition of two fibrations $`p_2p_1`$ to the torsion of $`p_1`$ and $`p_2`$. The second condition is additivity, which describes the torsion of a pushout of bundles over a space $`B`$. In the second author showed that a homotopy theoretical analog of additivity is satisfied by the homotopy torsion of Dwyer, Weiss and Williams. The main result of our present paper shows that additivity holds also for the topological torsion $`\tau _\rho ^t`$.
###### 1.1 Definition.
Let $`p:EB`$ be a bundle of closed topological manifolds. We say that $`p`$ admits a fiberwise codimension one splitting if there are subbundles of manifolds $`p_i:E_iB`$ ($`i=0,1,2`$) such that $`E=E_1_{E_0}E_2`$, fibers of $`p_1`$, $`p_2`$ are compact submanifolds (with boundary) of the fibers of $`p`$, and the fibers of $`p_0`$ are the common boundary of the fibers of $`p_1`$ and $`p_2`$.
###### 1.2 Theorem.
Suppose that $`p:EB`$ is a bundle of closed manifolds which admits a fiberwise codimension one splitting into subbundles $`p_i:E_iB`$ for $`i=0,1,2`$. Let $`R`$ be a ring and $`\rho :VE`$ be a locally constant sheaf of finitely generated projective left $`R`$-modules. Finally, assume that $`H_{}(p^1(b);\rho )=0`$ and $`H_{}(p_i^1(b);\rho |_{E_i})=0`$ for $`i=0,1,2`$ and all $`bB`$. Then there exists a preferred homotopy class of paths in the topological Whitehead space $`\mathrm{\Phi }_\rho ^t(p)`$ joining $`\tau _\rho ^t(p)`$ with $`j_1\tau _\rho ^t(p_1)+j_2\tau _\rho ^t(p_2)j_0\tau _\rho ^t(p_0)`$, where $`j_i`$ is the map induced by the inclusion $`j_i:E_iE`$.
The definitions of $`\mathrm{\Phi }_\rho ^t(p)`$ and $`\tau _\rho ^t(p)`$ are recalled in Section 2. We note that this property of $`\tau _\rho ^t`$ parallels the additivity of the classical combinatorial Reidemeister torsion relating torsion of a finite CW–complex $`X_ZY`$ to the torsions of $`X`$, $`Y`$ and $`Z`$.
While the proof of Theorem 1.2 is more subtle than in the case of homotopy torsion, the essential idea is to reduce the problem to additivity of $`\tau _\rho ^h`$ and then use the arguments of . We expect that, similarly, a proof of the additivity for smooth torsion may be obtained by reduction to topological case and application of Theorem 1.2.
The main component of the proof of Theorem 1.2 is the additivity theorem for the topological Euler characteristic:
###### 1.3 Theorem.
Let $`p:EB`$ be a fiber bundle of closed topological manifolds admitting a fiberwise codimension one splitting as in Theorem 1.2. There exists a preferred homotopy class of paths in $`A^\%(p)`$ joining $`\chi ^t(p)`$ with $`j_1\chi ^t(p_1)+j_2\chi ^t(p_2)j_0\chi ^t(p_0)`$.
The space $`A^\%(p)`$ denotes here the parametrized excisive $`A`$-theory of $`p`$, and $`\chi ^t(p)A^\%(p)`$ is the topological Euler characteristic of $`p`$. Again, definitions of these notions are sketched in Section 2.
In the present paper Theorem 1.3 serves as a step in establishing the additivity of topological torsion, but it has also other potential applications. One of the main results of says that a bundle of compact manifolds $`p:EB`$ is fiber homotopy equivalent to a smooth bundle if and only if $`\chi ^t(p)`$ can be lifted to the infinite loop space associated with the parametrized suspension spectrum of $`p`$. From this perspective additivity of the topological Euler characteristic provides a tool for computing an obstruction for smoothing of the bundle $`p`$.
###### 1.4.
Organization of the paper. In Section 2 we recall the constructions of Dwyer, Weiss, and Williams. In their setting the topological (resp. homotopy) Reidemeister torsion can be thought of as a lift of the parametrized topological (resp. homotopy) Euler characteristic. In order to prove additivity for $`\tau _\rho ^t`$ it is then enough to verify that additivity holds for the topological Euler characteristic, and then show that we can lift the resulting path. As the first step we demonstrate (§3) that excisive Euler characteristic is additive in non-parametrized case, that is for bundles over a one-point space. In Section 4 we show how to extend this result to bundles of manifolds with a discrete structure group. Subsequently in §5 we show that additivity of the Euler characteristic for arbitrary bundles follows from additivity for certain universal bundles. Then, in Section 6 we show that additivity for the universal bundles follows from additivity of the topological Euler characteristic for bundles with a discrete structure groups. This completes the proof of Theorem 1.3. Finally, in Section 7 we prove Theorem 1.2.
###### Acknowledgments.
The authors wish to thank Bruce Williams for conversations which contributed to this work. Comments of the referee helped to clarify several passages.
## 2. Dwyer-Weiss-Williams constructions
The purpose of this section is to provide a quick review of constructions leading to the definition of the topological Reidemeister torsion. We refer to for a detailed treatment of this subject. We also set here the notation which we will use throughout the paper. In general we tried to preserve the notation of , although some differences occur.
###### 2.1.
Non-parametrized Euler characteristics. By a Waldhausen category we will mean a category $`𝒞`$ together with a choice of two subcategories: a subcategory of cofibrations and a subcategory of weak equivalences satisfying the axioms of (in the terminology of such a category $`𝒞`$ is called a category with cofibrations and weak equivalences). Applying the $`S_{}`$-construction to $`𝒞`$ one obtains $`\mathrm{\Omega }|wS_{}𝒞|`$ – the $`K`$-theory space of $`𝒞`$ \[22, 1.3\]. This is an infinite loop space which we will denote by $`K(𝒞)`$. Every object $`c𝒞`$ represents a point $`[c]K(𝒞)`$, and a weak equivalence $`\phi :cc^{}`$ determines a path from $`[c]`$ to $`[c^{}]`$.
An exact functor of Waldhausen categories $`F:𝒞𝒟`$ is a functor preserving the distinguished subcategories and all other relevant structures. Any such functor induces a map of the associated infinite loop spaces $`F_{}:K(𝒞)K(𝒟)`$. The functor $`𝒞\times 𝒞𝒞`$ which assigns to a pair of objects $`(c,c^{})`$ their coproduct $`cc^{}`$ in $`𝒞`$ is exact and defines a map $`K(𝒞)\times K(𝒞)K(𝒞)`$. This equips $`K(𝒞)`$ with an $`H`$-space structure such that $`[c]+[c^{}]=[cc^{}]`$. If we have defined a suspension functor $`\mathrm{\Sigma }:𝒞𝒞`$ \[22, p. 349\] then the map $`K(𝒞)K(𝒞)`$ induced by $`\mathrm{\Sigma }`$ represents a homotopy inverse with respect to the $`H`$-space structure on $`K(𝒞)`$. Thus, for $`c𝒞`$ we can write $`[c]:=[\mathrm{\Sigma }c]`$. All Waldhausen categories considered here come equipped with suspension functors.
If $`cc^{}c^{\prime \prime }`$ is a cofibration sequence in $`𝒞`$ then there is a path in $`K(𝒞)`$ joining $`[c^{}]`$ with $`[c]+[c^{\prime \prime }]`$. One way to get such a path is to use Waldhausen’s additivity theorem \[22, Prop. 1.3.2\]. Another way is to observe that (in the notation of \[22, 1.3\]) a cofibration sequence $`cc^{}c^{\prime \prime }`$ defines a point $`[cc^{}c^{\prime \prime }]|wS_2C|`$. The restriction of the map $`|wS_2C|\times \mathrm{\Delta }^2|wS_{}𝒞|`$ to $`[cc^{}c^{\prime \prime }]\times \mathrm{\Delta }^2`$ yields the desired path \[22, 1.3.3\]. This second construction is more explicit and easily adapts to the parametrized setting (see Lemma 4.3). This is the construction we are using throughout the paper.
Following \[5, p. 40\] we will denote by $`^{fd}(X)`$ the category of homotopy finitely dominated retractive spaces over $`X`$. The objects of $`^{fd}(X)`$ are diagrams
such that $`rs=\mathrm{id}_X`$, $`s`$ is a cofibration, and $`Y`$ is a homotopy finitely dominated space over $`X`$. The category $`^{fd}(X)`$ can be equipped with a Waldhausen category structure where a morphism in $`^{fd}(X)`$ is a weak equivalence or a cofibration if its underlying map of spaces is a homotopy equivalence or, respectively, a map with the homotopy extension property. Its $`K`$-theory space is denoted $`A(X)`$ and called the $`A`$-theory of $`X`$.
###### 2.2 Definition.
Let $`X`$ be a finitely dominated space. The characteristic object $`X^h^{fd}(X)`$ is the retractive space
where $`s(X)=X\times \{1\}`$, and $`r`$ is the projection map. The homotopy Euler characteristic of $`X`$ is the point $`\chi ^h(X)A(X)`$ represented by $`X^h`$.
The assignment $`XA(X)`$ defines a functor on the category of finitely dominated spaces. It is not a homology theory since it does not satisfy the excision axiom. By there exists a functor $`XA^\%(X)`$ which in a certain sense is the best possible approximation of $`A()`$ by an excisive functor. In \[5, §7\] the authors show that if $`X`$ is an Euclidean neighborhood retract (ENR) then $`A^\%(X)`$ can be explicitly constructed using Waldhausen categories. We outline this construction next.
For an ENR space $`X`$ we have a category $`^{ld}(𝕁X)`$ the objects of which are diagrams
such that $`rs=\mathrm{id}_{X\times [0,\mathrm{})}`$ and where $`Y`$ is homotopy locally finitely dominated as a space over $`X\times [0,\mathrm{})`$ \[5, p.48\]. Morphisms in $`^{ld}(𝕁X)`$ are retractive maps. Given objects $`r_i:Y_iX\times [0,\mathrm{}):s_i`$ for $`i=1,2`$ and morphisms $`f,g:Y_1Y_2`$ we have the notion of controlled homotopy between $`f`$ and $`g`$. By this we mean a map $`H:Y_1\times [0,1]Y_2`$ which gives a homotopy between $`f`$ and $`g`$ in the usual sense and which commutes with the maps $`s_1,s_2`$, but commutes with the retractions $`r_1,r_2`$ only in a relaxed, controlled way \[5, p. 47\]. The category $`^{ld}(𝕁X)`$ can be equipped with a Waldhausen category structure where weak equivalences are controlled homotopy equivalences and cofibrations are the maps with the controlled homotopy extension property. The $`K`$-theory space of $`^{ld}(𝕁X)`$ will be denoted by $`A^𝕁(X)`$.
We have a functor
$$I:^{fd}(X)^{ld}(𝕁X)$$
which assigns to a retractive space $`r:YX:s`$ over $`X`$ a retractive space over $`X\times [0,\mathrm{})`$
where $`\overline{Y}`$ is a pushout in the diagram
The functor $`I`$ is an embedding of categories, and – considered as a functor of Waldhausen categories – it is exact so it induces a map of the $`K`$-theory spaces
$$I_{}:A(X)A^𝕁(X)$$
Next, let $`𝒱(X)`$ denote the category of proper retractive ENRs over $`X\times [0,\mathrm{})`$. This is a subcategory of $`^{ld}(𝕁X)`$ whose objects $`r:YX\times [0,\mathrm{}):s`$ satisfy the conditions that $`Y`$ is an ENR and that $`r`$ is a proper map \[5, 7.8\]. Let $`J:𝒱(X)^{ld}(𝕁X)`$ be the inclusion functor. The category $`𝒱(X)`$ admits a Waldhausen category structure such that $`J`$ becomes an exact functor \[5, p. 51\]. As a consequence we obtain an infinite loop space $`V(X):=K(𝒱(X))`$ and a map $`J_{}:V(X)A^𝕁(X)`$.
###### 2.3 Definition.
Let $`X`$ be a compact ENR. The space $`A^\%(X)`$ is the homotopy limit
$$\begin{array}{ccccc}A^\%(X):=\underset{}{\mathrm{holim}}(A(X)& \stackrel{I_{}}{}& A^𝕁(X)& \stackrel{J_{}}{}& V(X))\end{array}$$
We call $`A^\%(X)`$ the excisive $`A`$-theory of $`X`$. The natural map $`\alpha :A^\%(X)A(X)`$ is called the assembly map.
###### 2.4 Remark.
1) In \[5, §7\] the space $`A^\%(X)`$ is defined in somewhat different manner, as a homotopy fiber of a map $`V(X)K(𝒢^{ld}(X))`$ where $`𝒢^{ld}(X)`$ is the Waldhausen category with the same objects as $`^{ld}(𝕁X)`$ but with germs of retractive maps as morphisms. The resulting infinite loop space is however homotopy equivalent to the one described above (see \[5, Lemma 8.7\]).
2) The space $`V(X)`$ is in fact contractible \[5, p.52\], so $`A^\%(X)`$ is equivalent to a homotopy fiber of the map $`I_{}`$. The above description of $`A^\%(X)`$ allows us however to perform certain constructions in $`A^\%(X)`$ combinatorially, on the level of Waldhausen categories as follows. For a compact ENR $`X`$ let $`^\%(X)`$ denote the pullback of the diagram of categories
$$\begin{array}{ccccc}^{fd}(X)& \stackrel{I}{}& ^{ld}(𝕁X)& \stackrel{J}{}& 𝒱(X)\end{array}$$
Thus, objects of $`^\%(X)`$ are pairs $`(a,b)`$ where $`a^{fd}(X)`$, $`b𝒱(X)`$, and $`I(a)=J(b)`$, and morphisms are defined similarly. The category $`^\%(X)`$ has the obvious structure of a Waldhausen category such that the functors $`^\%(X)^{fd}(X)`$, $`^\%(X)𝒱(X)`$ are exact. It follows that we have a commutative diagram of infinite loop spaces:
As a consequence we obtain a map $`K(^\%(X))A^\%(X)`$. In particular any object $`(a,b)^\%(X)`$ determines a point $`[a,b]A^\%(X)`$, a weak equivalence $`(a,b)(a^{},b^{})`$ defines a path from $`[a,b]`$ to $`[a^{},b^{}]`$, and a cofibration sequence
$$(a,b)(a^{},b^{})(a^{\prime \prime },b^{\prime \prime })$$
in $`^\%(X)`$ determines a path in $`A^\%(X)`$ joining $`[a^{},b^{}]`$ with $`[a,b]+[a^{\prime \prime },b^{\prime \prime }]`$. It will be sometimes convenient to describe points and paths in $`A^\%(X)`$ in this way.
We are now ready to define the topological Euler characteristic of a space.
###### 2.5 Definition.
Let $`X`$ be an ENR. The characteristic object of $`X`$ in $`𝒱(X)`$ is the object $`X^v`$ given by the retractive space
The topological Euler characteristic of $`X`$ is the point $`\chi ^t(X)A^\%(X)`$ determined by the object $`X^t:=(X^h,X^v)^\%(X)`$.
Notice that we have $`\alpha (\chi ^t(X))=\chi ^h(X)`$.
###### 2.6 Remark.
Let $`f:XY`$ be a map of spaces. By abuse of notation by $`f_{}`$ we will denote each of the functors induced by $`f`$: $`^{fd}(X)^{fd}(Y)`$, $`𝒱(X)𝒱(Y)`$, $`^{ld}(𝕁X)^{ld}(𝕁Y)`$, and $`^\%(X)^\%(Y)`$, as well as the maps of infinite loop spaces: $`A(X)A(Y)`$, $`V(X)V(Y)`$, $`A^𝕁(X)A^𝕁(Y)`$, and $`A^\%(X)A^\%(Y)`$.
Euler characteristics are not preserved in general by the maps $`f_{}`$. However, we have the following
###### 2.7 Lemma.
Let $`f:XY`$ be a homotopy equivalence of finitely dominated spaces. We have a canonical weak equivalence $`f^h:f_{}(X^h)Y^h`$ in $`^{fd}(Y)`$. Moreover, if $`g:YZ`$ is another homotopy equivalence, then $`(gf)^h=g^hg_{}(f^h)`$.
###### Proof.
Define
$$f^h:=f\mathrm{id}:f_{}(X^h)=XYYY$$
The properties of $`f^h`$ are straightforward to check. ∎
###### 2.8 Corollary.
If $`f:XY`$ is a homotopy equivalence of finitely dominated spaces then there is a canonical path $`\omega _f`$ from $`f_{}\chi ^h(X)`$ to $`\chi ^h(Y)`$.
###### 2.9 Lemma.
Let $`f:XY`$ be a cell-like map , , of compact ENRs. We have a canonical weak equivalence $`f^v:f_{}(X^v)Y^v`$ in $`𝒱(X)`$. Moreover the weak equivalences $`f^h`$ and $`f^v`$ satisfy the equations $`I(f^h)=J(f^v)`$, and so they define a weak equivalence $`f^t:=(f^h,f^v):f_{}(X^t)Y^t`$ in $`^\%(Y)`$. The weak equivalences $`f^v`$ and $`f^t`$ satisfy a cochain condition analogous to the one described in Lemma 2.7.
###### Proof.
The map $`f^v`$ is given by
$$f^v:=f\mathrm{id}:f_{}(X^v)=XY\times [0,\mathrm{})YY\times [0,\mathrm{})$$
It is a weak equivalence in $`𝒱(X)`$ by \[5, p. 53\]. Verification of the remaining properties of $`f^v`$ is straightforward. ∎
###### 2.10 Corollary.
If $`f:XY`$ is a cell-like map of compact ENRs then there is a canonical path $`\sigma _f`$ joining $`f_{}(\chi ^t(X))`$ with $`\chi ^t(Y)`$. Moreover, if $`a:A^\%(X)A(X)`$ is the assembly map, then $`a(\sigma _f)=\omega _f`$ where $`\omega _f`$ is the path from Lemma 2.8.
We will refer to the properties of $`\chi ^h()`$ and $`\chi ^t()`$ described in Corollaries 2.8 and 2.10 as the lax naturality of Euler characteristics.
###### 2.11.
Parametrization. The constructions sketched above can be generalized to the setting where the space $`X`$ is replaced by a fibration $`E\stackrel{𝑝}{}B`$ with a fiber $`F`$. Intuitively, one can construct in this case a fibration $`A_B\left(E\right)B`$ the fiber of which is $`A(F)`$. An analog of the homotopy Euler characteristic in this context is a section of this fibration which restricts to $`\chi ^h(F)`$ over every point of $`B`$. A similar idea underlies the notions of the parametrized excisive $`A`$-theory and the parametrized topological Euler characteristic. For technical reasons it is more convenient, however, to define the parametrized Euler characteristics in different terms. We give these formal definitions first, and then explain how they relate to the above idea.
Let $`𝒞`$ be a small category, and let $`F:𝒞𝒮paces`$ be a functor. Recall that $`\mathrm{holim}_𝒞F`$ is the space of natural transformations
$$\underset{𝒞}{\mathrm{holim}}F:=\mathrm{Map}_𝒞(|𝒞/|,F)$$
where $`|𝒞/|`$ is the functor which assigns to $`c𝒞`$ the nerve of the overcategory $`C/c`$. The following fact is implicitly present in :
###### 2.12 Lemma.
Let $`𝒲Cat`$ denote the category of Waldhausen categories with exact functors as morphisms. Assume that for a functor $`F:𝒞𝒲Cat`$ we have a rule which assigns to $`c𝒞`$ an object $`c^!F(c)`$, and to $`f\mathrm{Mor}_𝒞(c,d)`$ a weak equivalence $`f^!:F(f)(c^!)d^!`$ in $`F(d)`$ in such way that $`(gf)^!=g^!F(f)(f^!)`$. Then the assignment $`|C/c|[c^!]`$ defines a point $`[c^!,f^!]\mathrm{holim}_𝒞K(F(c))`$.
Indeed, if $`wF(c)`$ denotes the subcategory of weak equivalences in the Waldhausen category $`F(c)`$, then the assignments $`cc^!`$, $`ff^!`$ as in Lemma 2.12 define a point in $`\mathrm{holim}_𝒞|wF(c)|`$ (where $`|wF(c)|`$ is the nerve of $`wF(c)`$). Using the natural transformation of functors $`|wF(c)|K(F(c))`$ we obtain a point in $`\mathrm{holim}_𝒞K(F(c))`$ as claimed.
Now, let $`𝔅:\mathrm{\Delta }^{op}𝒮ets`$ be a simplicial set. Denote by $`𝒮(𝔅)`$ the category whose objects are all simplices $`x𝔅`$, and where morphisms $`xy`$ in $`𝒮(𝔅)`$ come from morphisms $`\phi `$ in $`\mathrm{\Delta }^{op}`$ satisfying $`𝔅(\phi )(y)=x`$ (thus $`𝒮(𝔅)`$ is the opposite category of the Grothendieck construction on the functor $`𝔅`$ \[4, p.22\]). Let $`B`$ be the geometric realization of $`𝔅`$, and for $`x𝔅`$ let $`\phi _x:\mathrm{\Delta }^{|x|}B`$ denote the characteristic map of $`x`$. Assume that we have a fibration $`E\stackrel{𝑝}{}B`$ with a homotopy finitely dominated fiber $`F`$. We can define a functor
$$𝒮(𝔅)𝒮paces,xE_x$$
where $`E_x:=lim(E\stackrel{𝑝}{}B\stackrel{\phi _x}{}\mathrm{\Delta }^{|x|})`$. As a result we get a functor
$$F:𝒮(𝔅)𝒲Cat,x^{fd}(E_x)$$
For $`x𝒮(𝔅)`$ consider the assignment $`xE_x^h`$ where $`E_x^h^{fd}(X)`$ is the characteristic object of $`E_x`$ (2.2). Notice that for any morphism $`f:xy`$ in $`𝒮(𝔅)`$ the map $`F(f):E_xE_y`$ is a homotopy equivalence, so by Lemma 2.7 it defines a weak equivalence $`F(f)^h:F(f)_{}(E_x^h)E_y^h`$ in $`^{fd}(E_y)`$ . Lemma 2.7 also shows that the assignments $`xE_x^h`$, $`fF(f)^h`$ satisfy the conditions Lemma 2.12, and so they define a point $`[E_x^h,F(f)^h]\underset{𝒮(𝔅)}{\mathrm{holim}}A(E_x)`$.
###### 2.13 Definition.
The homotopy Euler characteristic of a fibration $`E\stackrel{𝑝}{}B`$ is the point $`[E_x^h,F(f)^h]\underset{𝒮(𝔅)}{\mathrm{holim}}A(E_x)`$. We will write $`\chi ^h(p):=[E_x^h,F(f)^h]`$.
Next, notice that the constant maps $`A(E_x)`$ define a map of homotopy colimits
$$p_{}:\underset{𝒮(𝔅)}{\mathrm{hocolim}}A(E_x)\underset{𝒮(𝔅)}{\mathrm{hocolim}}=B$$
which by \[7, p.180\] is a quasi-fibration with the fiber $`A(F)`$. Let $`p_{}:A_B\left(E\right)B`$ be the fibration associated to this quasi-fibration. For all morphisms $`xy`$ in $`𝒮(𝔅)`$ the map $`E_xE_y`$ is a homotopy equivalence, and thus so is the map $`A(E_x)A(E_y)`$. Therefore as an application of \[3, Prop. 3.12\] we obtain
###### 2.14 Proposition.
The map
$$\underset{x𝒮(𝔅)}{\mathrm{holim}}A(E_x)=\mathrm{Map}_{𝒮(𝔅)}(|𝒮(𝔅)/|,A(E_{()}))\mathrm{Map}_B(\underset{x𝒮(𝔅)}{\mathrm{hocolim}}|𝒮(𝔅)/x|,\underset{x𝒮(𝔅)}{\mathrm{hocolim}}A(E_x))$$
$$f\underset{}{\mathrm{hocolim}}f$$
is a weak equivalence.
Here $`\mathrm{Map}_B(,)`$ denotes the mapping space of spaces over $`B`$. Also, in the category of spaces over $`B`$ we have weak equivalences
$$\underset{x𝒮(𝔅)}{\mathrm{hocolim}}|𝒮(𝔅)/x|B\mathrm{and}\underset{x𝒮(𝔅)}{\mathrm{hocolim}}A(E_x)A_B\left(E\right)$$
which combined with Proposition 2.14 give
$$\underset{x𝒮(𝔅)}{\mathrm{holim}}A(E_x)\mathrm{Map}_B(B,A_B\left(E\right))$$
It follows that $`\chi ^h(p)\underset{}{\mathrm{holim}}A(E_x)`$ defines a point (unique up to a contractible space of choices) in $`\mathrm{\Gamma }(p_{})`$ – the space of sections of the fibration $`p_{}:A_B\left(E\right)B`$. This brings us back to the intuitive construction of $`\chi ^h(p)`$ sketched at the beginning of this section.
As we have mentioned above the idea behind the definition of the parametrized topological characteristic $`\chi ^t(p)`$ is similar. The details, however, are more involved. The problem is that one cannot (mimicking Definition 2.13) define $`\chi ^t(p)`$ using Lemmas 2.9 and 2.12 since the maps $`E_xE_y`$ are usually not cell-like. In the authors overcome this difficulty by replacing $`𝔅`$ with a new simplicial set $`p𝔅`$, and the functor $`F:𝒮(𝔅)𝒮paces`$ by a new functor $`tF:𝒮(p𝔅)𝒮paces`$, which sends every morphism in $`𝒮(p𝔅)`$ to a homeomorphism. Then the assumptions of Lemma 2.9 are satisfied and we can apply 2.12 in order to define a point in $`\chi ^t(p)\mathrm{holim}_{𝒮(p𝔅)}A^\%(tF)`$. The details follow.
###### 2.15 Definition (Compare \[5, p.13\] ).
Let $`E\stackrel{𝑝}{}B`$ be a (locally trivial) fiber bundle, where $`B`$ is the geometric realization of a simplicial set $`𝔅`$. By $`p𝔅`$ we will denote the simplicial set whose $`k`$-simplices are pairs $`(x,\theta )`$ where $`x`$ a $`k`$-simplex in $`𝔅`$, and $`\theta `$ is an equivalence relation on $`E_x`$ such that the quotient map $`E_xE_x^\theta `$ and the projection $`E_x\mathrm{\Delta }^k`$ give a homeomorphism $`E_xE_x^\theta \times \mathrm{\Delta }^k`$.
Notice that if $`F`$ is the fiber of the fiber bundle $`E\stackrel{𝑝}{}B`$ then for every $`(x,\theta )p𝔅`$ we have $`E_x^\theta F`$ and the assignment $`(x,\theta )E_x^\theta `$ defines a functor $`tF:𝒮(p𝔅)𝒮paces`$ which sends every morphism in $`𝒮(p𝔅)`$ to a homeomorphism.
###### 2.16 Definition.
Let $`E\stackrel{𝑝}{}B`$ be a fiber bundle whose fiber $`F`$ is a compact ENR. The assignments $`(x,\theta )(E_x^\theta )^t`$ (see 2.5) , and $`fF(f)^t`$ (2.9), where $`f`$ is a morphism in $`p𝔅`$, satisfy the conditions of Lemma 2.12. Therefore they define a point
$$\chi ^t(p):=[(E_x^\theta )^t,F(f)^t]\underset{(x,\theta )p𝔅}{\mathrm{holim}}A^\%(E_x^\theta )$$
We call $`\chi ^t(p)`$ the topological Euler characteristic of the bundle $`p`$.
The following fact lets us compare the homotopy and topological Euler characteristics of bundles.
###### 2.17 Proposition.
Let $`E\stackrel{𝑝}{}B`$ be a locally trivial fiber bundle whose fiber is a compact topological manifold $`M`$ (perhaps with boundary). We have a commutative diagram
The vertical maps are induced by assembly maps. All horizontal maps are weak equivalences. The horizontal maps on the left are induced by the forgetful functor $`𝒮(p𝔅)𝒮(𝔅)`$, while the horizontal maps on the right come from the natural transformation $`E_xE_x^\theta `$ of functors over $`𝒮(p𝔅)`$.
The maps on the right are weak equivalences by homotopy invariance of homotopy limits. The maps on the left are weak equivalences by \[5, Corollary 2.7\].
###### 2.18.
The implications of Proposition 2.17 are twofold. On one hand it lets us think about $`\chi ^t(p)`$ as an element of $`\mathrm{holim}_{x𝒮(𝔅)}A^\%(E_x)`$ defined uniquely up to a contractible space of choices. Similarly as in 2.14 we then get $`\mathrm{holim}_{x𝒮(𝔅)}A^\%(E_x)\mathrm{\Gamma }(p_{}^\%)`$ where $`\mathrm{\Gamma }(p_{}^\%)`$ is the space of sections of the fibration $`p_{}^\%:A_B^\%\left(E\right)B`$ associated to the quasi-fibration $`\mathrm{hocolim}_{x𝒮(𝔅)}A^\%(E_x)B`$ with the fiber $`A^\%(M)`$. In particular we can think of $`\chi ^t(p)`$ as a section of $`p_{}^\%`$. On the other hand consider the image of $`\chi ^h(p)`$ under the composition of the bottom maps in the diagram above. One can see that the point $`\mathrm{holim}_{(x,\theta )}A(E_x^\theta )`$ defined in this way can be explicitly described by means of Lemma 2.12 as coming from the assignment $`(x,\theta )(E_x^\theta )^{ah}:=E_xE_x^\theta `$ for any object $`(x,\theta )𝒮(p𝔅)`$, and $`ff^{ah}`$ for any morphism $`f:(x,\theta )(y,\theta ^{})`$ in $`𝒮(p𝔅)`$ where $`f^{ah}`$ is the map
$$f^{ah}:tF(f)(E_xE_x^\theta )=E_yE_x^\theta \stackrel{\mathrm{id}tF(f)}{}E_yE_y^\theta ^{}$$
We can consider the element $`[(E_x^\theta )^{ah},f^{ah}]\mathrm{holim}_{(x,\theta )}A(E_x^\theta )`$ as a re-definition of $`\chi ^h(p)`$. Unlike the non-parametrized case this element is not equal of the image of $`\chi ^t(p)`$ under the assembly map. However, the homotopy equivalences $`E_xE_x^\theta E_x^\theta E_x^\theta `$ define a canonical path $`\sigma _p`$ in $`\mathrm{holim}_{(x,\theta )}A(E_x^\theta )`$, joining the image of $`\chi ^t(p)`$ with the image of $`\chi ^h(p)`$.
###### 2.19 Notation.
We will denote by $`A(p)`$ the homotopy limit $`\mathrm{holim}_xA(E_x)`$. Similarly, $`A^\%(p)`$ will denote $`\mathrm{holim}_{(x,\theta )}A^\%(E_x^\theta )`$.
###### 2.20.
Reidemeister Torsions. Let $`R`$ be a (discrete) ring with identity, and let $`Ch^{fd}(R)`$ denote the category of chain complexes of left projective $`R`$-modules which are chain homotopy equivalent to finitely generated complexes. The category $`Ch^{fd}(R)`$ can be equipped with a Waldhausen category structure where weak equivalences are chain homotopy equivalences and cofibrations are chain maps which are split injective on every level. The infinite loop space $`K(Ch^{fd}(R))`$ is homotopy equivalent to the $`K`$-theory space of the ring $`R`$ \[5, p. 43\]. Thus, from now on we will denote $`K(Ch^{fd}(R))`$ by $`K(R)`$.
Let $`\rho :VE`$ be a locally constant sheaf of projective modules. The sheaf $`\rho `$ induces a functor $`L_\rho :^{fd}(E)Ch^{fd}(R)`$ which assigns to any object $`(XE)^{fd}(E)`$ the relative singular chain complex $`C(X,E,\rho )`$ with local coefficients given by $`\rho `$. The functor is $`L_\rho `$ is not exact. It is however close enough to being exact that it still defines a map $`L_\rho :A(E)K(R)`$ if we slightly modify the construction of $`A(E)`$ and $`K(R)`$ using a variant of the $`S_{}`$-construction proposed by Thomason (see \[5, p. 43\]). Assume that $`H_{}(E,\rho )=0`$. In this case the complex $`L_\rho (E^h)`$ is acyclic, so the map $`0L_\rho (E^h)`$ (where $`0`$ is the zero chain complex) is a weak equivalence in $`Ch^{fd}(E)`$. This gives a canonical path $`\sigma _\rho (E)`$ in $`K(R)`$ joining $`L_\rho (\chi ^h(E))`$ with $``$–the basepoint of $`K(R)`$ which is represented by the chain complex $`0`$.
###### 2.21 Definition.
The pair $`(\chi ^h(E),\sigma _\rho (E))`$ defines a point in $`\mathrm{\Phi }_\rho ^h(E)`$ \- the homotopy fiber of the map $`L_\rho `$. This point is called the homotopy Reidemeister torsion of the space $`E`$ and is denoted by $`\tau _\rho ^h(E)`$.
If $`E`$ is a compact ENR then let $`\mathrm{\Phi }_\rho ^t(E)`$ denote the homotopy fiber of the map $`L_\rho \alpha `$ where $`\alpha `$ is the assembly map. The point $`\tau _\rho ^t(E)\mathrm{\Phi }_\rho ^t(E)`$ defined by the pair $`(\chi ^t(E),\sigma _\rho (E))`$ is the topological Reidemeister torsion of $`E`$.
We will call $`\mathrm{\Phi }_\rho ^h(E)`$ and $`\mathrm{\Phi }_\rho ^t(E)`$ the homotopy (resp. topological) Whitehead spaces.
Next, let $`E\stackrel{𝑝}{}B`$ be a fibration with a homotopy finitely dominated fiber $`F`$, where as before $`B=|𝔅|`$, and and let $`\rho :VE`$ be a locally constant sheaf of finitely generated projective left $`R`$-modules. In such case following \[5, p. 66\] we can define fibrations $`\mathrm{\Phi }_{\rho ,B}^h(E)B`$ and $`\mathrm{\Phi }_{\rho ,B}^t(E)B`$ with fibers $`\mathrm{\Phi }_\rho ^h(F)`$ and $`\mathrm{\Phi }_\rho ^t(F)`$ respectively. If $`H_{}(E_x,\rho |_{E_x})=0`$ for all $`x𝔅`$ these fibrations admit sections which assign $`\tau _{\rho |_F}^h(F)`$ (resp. $`\tau _{\rho |_F}^t(F)`$) to every point $`bB`$. We can think of these sections as parametrized versions of the Reidemeister torsions. However, similarly as it was the case for the parametrized Euler characteristics (2.11), this describes torsions only up to a contractible space of choices. We will then again need other, more precise definitions.
For $`x𝒮(𝔅)`$ consider the maps $`E_xE`$. The induced functors of Waldhausen categories $`^{fd}(E_x)^{fd}(E)`$ define a natural transformation $`\eta `$ from the functor $`F:𝒮(𝔅)𝒲Cat`$, $`F(x)=^{fd}(E_x)`$ to the constant functor over $`𝒮(𝔅)`$ with the value $`^{fd}(E)`$. Therefore we obtain a map
$$\eta _{}:A(p)=\underset{x𝒮(𝔅)}{\mathrm{holim}}A(E_x)\underset{x𝒮(𝔅)}{\mathrm{holim}}A(E)$$
Recall the the homotopy Euler characteristic $`\chi ^h(p)`$ was defined as the point in $`A(p)=\mathrm{holim}_xA(E_x)`$ represented by the assignments $`xE_x^h`$, $`fF(f)^h`$ (2.13). The image of $`\chi ^h(p)`$ under $`\eta _{}`$ is in turn represented by the assignments $`xE_xE`$, $`f(F(f)\mathrm{id}:E_xEE_yE)`$. Assume that $`H_{}(E_x,\rho |_{E_x})=0`$ for all $`x𝒮(𝔅)`$. In this case the relative chain complexes $`C(E_xE,E,\rho )`$ are acyclic, and the weak equivalences $`0C(E_xE,E,\rho )`$ define a canonical path $`\sigma _\rho (p)`$ joining the basepoint and $`L_\rho _{}\eta _{}(\chi ^h(p))`$ in $`\mathrm{holim}_xK(R)`$.
###### 2.22 Definition.
Let $`\mathrm{\Phi }_\rho ^h(p)`$ be the homotopy fiber of the map $`L_\rho \eta _{}:A(p)\mathrm{holim}_xK(R)`$ over the basepoint of $`\mathrm{holim}_xK(R)`$. The homotopy Reidemeister torsion of $`p`$ is the point $`\tau _\rho ^h(p)\mathrm{\Phi }_\rho ^h(p)`$ given by the pair $`(\chi ^h(p),\sigma _\rho (p))`$.
Assume now that $`E\stackrel{𝑝}{}B`$ is a bundle of compact manifolds. Let $`\stackrel{~}{A}^\%(p)`$ be the pullback of the diagram
$$\underset{𝑥}{\mathrm{holim}}A(E_x)\stackrel{}{}\underset{(x,\theta )}{\mathrm{holim}}A(E_x^\theta )\stackrel{𝛼}{}\underset{(x,\theta )}{\mathrm{holim}}A^\%(E_x^\theta )$$
where the maps are as in Proposition 2.17. Recall (2.18) that we have a canonical path $`\sigma _p`$ in $`\mathrm{holim}_{(x,\theta )}A(E_x^\theta )`$ joining the images of $`\chi ^h(p)`$ and $`\chi ^t(p)`$. Thus, the triple $`(\chi ^t(p),\sigma _p,\chi ^h(p))`$ defines a point $`\stackrel{~}{\chi }^t(p)\stackrel{~}{A}^\%(p)`$.
###### 2.23 Definition.
Let $`\mathrm{\Phi }_\rho ^t(p)`$ denote the homotopy fiber of the map
$$\stackrel{~}{A}^\%(p)\underset{𝑥}{\mathrm{holim}}A(E_x)\stackrel{L_\rho \eta _{}}{}\underset{𝑥}{\mathrm{holim}}K(R)$$
If $`H_{}(E_x,\rho |_{E_x})=0`$ for all $`x𝒮(𝔅)`$ then the pair $`(\stackrel{~}{\chi }^t(p),\sigma _\rho (p))`$ (where $`\sigma _\rho (p)`$ is the path as in Definition 2.22) defines a point $`\tau _\rho ^t(p)\mathrm{\Phi }_\rho ^t(p)`$. We call it the topological parametrized Reidemeister torsion of the bundle $`p`$.
Notice that we have a pullback diagram
and that $`\gamma _1(\tau _\rho ^t(p))=\tau _\rho ^h(p)`$, $`\gamma _2(\tau _\rho ^t(p))=\stackrel{~}{\chi }^t(p)`$.
## 3. Non-parametrized additivity theorem
The goal of this section is to prove the following
###### 3.1 Theorem (Additivity for the topological Euler characteristic).
Let $`M`$ be a closed topological manifold which admits a splitting along a compact codimension one submanifold $`M_0`$:
$$M=M_1_{M_0}M_2$$
There exists a preferred path $`\omega `$ in $`A^\%(M)`$ from $`\chi ^t(M)`$ to $`k_1\chi ^tM_1+k_2\chi ^tM_2k_0\chi ^tM_0`$, where $`k_i:M_iM`$ is the inclusion map ($`i=0,1,2`$).
Applying the assembly map $`\alpha :A^\%(M)A(M)`$ to the path $`\omega `$ we obtain
###### 3.2 Corollary.
If $`M`$ is a manifold as in Theorem 3.1, then there exists a path in $`A(M)`$ joining $`\chi ^h(M)`$ with $`k_1\chi ^hM_1+k_2\chi ^hM_2k_0\chi ^hM_0`$.
Thus we recover the additivity theorem for the homotopy Euler characteristic which was proved in by the second author.
###### Proof of Theorem 3.1.
Consider a manifold
$$\overline{M}=M_1_{M_0\times \{1\}}M_0\times [1,1]_{M_0\times \{1\}}M_2$$
For $`i=1,2`$ let $`\overline{k}_i:M_i\overline{M}`$ denote the inclusion map. We have a map $`f:\overline{M}M`$ which restricts to the inclusions $`k_1`$, $`k_2`$ on $`M_1`$, $`M_2`$ respectively, and which sends $`M_0\times [1,1]\overline{M}`$ to $`M_0M`$ via the projection map onto the first factor.
We will construct the path $`\omega `$ as a concatenation of three paths in $`A^\%(M)`$:
* a path $`\omega _1`$ from $`\chi ^t(M)`$ to $`f_{}\chi ^t(\overline{M})`$,
* a path $`\omega _2`$ from $`f_{}\chi ^t(\overline{M})`$ to $`f_{}(\overline{k}_1\chi ^t(M_1)+\overline{k}_2\chi ^t(M_1)+C)`$ for some $`CA^\%(\overline{M})`$,
* a path $`\omega _3`$ from $`f_{}\overline{k}_1\chi ^t(M_1)+f_{}\overline{k}_2\chi ^t(M_1)+f_{}C`$ to $`k_1\chi ^hM_1+k_2\chi ^hM_2k_0\chi ^hM_0`$.
Construction of $`\omega _1`$. Since $`f`$ is a cell-like map, the path $`\omega _1`$ exists by lax naturality of $`\chi ^t`$ (2.10).
Construction of $`\omega _2`$. We will construct a point $`CA^\%(\overline{M})`$ and a path $`\sigma `$ in $`A^\%(\overline{M})`$ joining $`\chi ^t(\overline{M})`$ with $`(\overline{k}_1\chi ^t(M_1)+\overline{k}_2\chi ^t(M_1)C`$. Then we will have $`\omega _2=f_{}(\sigma )`$.
First we define the point $`C`$ as follows. Let $`C^{}^{fd}(\overline{M})`$ be the retractive space over $`\overline{M}`$ given by
$$C^{}:=\overline{M}_{M_0\times \{1,1\}}M_0\times [1,1]$$
and let $`C^{\prime \prime }𝒱(\overline{M})`$ be given by the retractive space over $`\overline{M}\times [0,\mathrm{})`$
$$C^{\prime \prime }:=\overline{M}\times [0,\mathrm{})_{M_0\times \{1,1\}}M_0\times [1,1]$$
where the embedding $`M_0\times \{1,1\}\overline{M}\times [0,\mathrm{})`$ is given by $`(x,\pm 1)((x,\pm 1),0)`$. One can check that (in the notation of Remark 2.4) we have $`I(C^{})=J(C^{\prime \prime })`$, so $`(C^{},C^{\prime \prime })`$ is an object of $`^\%(\overline{M})`$. Take $`C`$ to be the point in $`A^\%(\overline{M})`$ represented by this object.
Next, recall (2.5) that for any compact ENR $`X`$ the Euler characteristic $`\chi ^t(X)`$ is represented by the object $`X^t=(X^h,X^v)^\%(X)`$ . Also, notice that in our case the object $`\overline{k}_iM_i^h^{fd}(\overline{M})`$ is the retractive space $`M_i\overline{M}`$ over $`\overline{M}`$. We have a cofibration sequence in $`^{fd}(\overline{M})`$
$$\overline{k}_1M_1^h(M_1M_2\overline{M})\overline{k}_2M_2^h$$
Similarly, the object $`\overline{k}_iM_i^v𝒱(\overline{M})`$ is given by the retractive space $`M_i\overline{M}\times [0,\mathrm{})`$ over $`\overline{M}\times [0,\mathrm{})`$, which gives a cofibration sequence in $`𝒱(\overline{M})`$
$$\overline{k}_1M_1^v(M_1M_2\overline{M}\times [0,\mathrm{}))\overline{k}_2M_2^v$$
These two cofibration sequences lift to a cofibration sequence in $`^\%(M)`$:
$$\overline{k}_1M_1^t(M_1M_2\overline{M},M_1M_2\overline{M}\times [0,\mathrm{}))\overline{k}_2M_2^t$$
which gives us a path in $`A^\%(\overline{M})`$ joining $`\overline{k}_1\chi ^t(M_1)+\overline{k}_2\chi ^t(M_2)`$ with the point $`[M_1M_2\overline{M},M_1M_2\overline{M}\times [0,\mathrm{})]`$. As a consequence we only need to construct a path from $`\chi ^t(\overline{M})`$ to $`[M_1M_2\overline{M},M_1M_2\overline{M}\times [0,\mathrm{})]+C`$. In order to accomplish this notice that $`C^{}`$ fits into a cofibration sequence in $`^{fd}(\overline{M})`$
$$\begin{array}{ccc}(M_1M_2)\overline{M}& \stackrel{(\overline{k}_1\overline{k}_2)1}{}& \overline{M}^h=\overline{M}\overline{M}C^{}\end{array}$$
while $`C^{\prime \prime }`$ is the cofiber in the following cofibration sequence in $`𝒱(\overline{M})`$:
$$\begin{array}{ccc}(M_1M_2)\overline{M}\times [0,\mathrm{})& \stackrel{(\overline{k}_1\overline{k}_2)1}{}& \overline{M}^v=\overline{M}\overline{M}\times [0,\mathrm{})C^{\prime \prime }\end{array}$$
As before these sequences yield a cofibration sequence in $`^\%(\overline{M})`$
$$(M_1M_2\overline{M},M_1M_2\overline{M}\times [0,\mathrm{}))\overline{M}^t=(\overline{M}^h,\overline{M}^v)(C^{},C^{\prime \prime })$$
Passing from $`^\%(\overline{M})`$ to $`A^\%(\overline{M})`$ we obtain the desired path.
Construction of $`\omega _3`$. We will show that in $`A^\%(M)`$ we have paths $`\delta _i`$ ($`i=1,2`$) joining $`f_{}\overline{k}_i\chi ^tM_i`$ with $`k_i\chi ^tM_i`$, and a path $`\delta _0`$ from $`f_{}C`$ (where $`CA^\%(\overline{M})`$ is defined as above) to $`k_0\chi ^tM_0`$. Then we can take $`\omega _3=\delta _1+\delta _2+\delta _0`$.
For $`i=1,2`$ we have $`f\overline{k}_i=k_i`$, so $`f_{}\overline{k}_i\chi ^tM_i=k_i\chi ^tM_i`$, and we can choose $`\delta _1,\delta _2`$ to be the constant paths.
The construction of the path $`\delta _0`$ resembles the construction of $`\omega _2`$ above. We have $`f_{}C=([f_{}C^{}],[f_{}C^{\prime \prime }])`$. Notice that $`f_{}C^{}`$ fits into the following pushout diagram in $`^{fd}(M)`$:
Similarly, $`f_{}C^{\prime \prime }`$ can be represented as a pushout in $`𝒱(M)`$:
where the map $`M_0\times \{1,1\}M\times [0,\mathrm{})`$ is given by $`(x,\pm 1)(x,0)`$. As a consequence we have a cofibration sequence in the category $`^{fd}(M)`$:
$$\begin{array}{c}M_0MM_0\times [1,1]_{M_0\times \{1\}}Mf_{}C^{}\end{array}$$
as well as a cofibration sequence in $`𝒱(M)`$:
$$\begin{array}{c}M_0M\times [0,\mathrm{})M_0\times [1,1]_{M_0\times \{1\}}M\times [0,\mathrm{})f_{}C^{\prime \prime }\end{array}$$
(we identify here $`M_0\times \{1\}`$ with a subspace of $`M\times [0,\mathrm{})`$ via the embedding $`(x,1)(x,0)`$). Notice that $`M_0M=k_0M_0^h`$ in $`^{fd}(M)`$, and $`M_0M\times [0,\mathrm{})=k_0M_0^v`$ in $`𝒱(M)`$. Again, we can lift these two sequence to a cofibration in $`^\%(M)`$
$$k_0(M^h,M^v)(M_0\times [1,1]_{M_0\times \{1\}}M,M_0\times [1,1]_{M_0\times \{1\}}M\times [0,\mathrm{}))f_{}(C^{},C^{\prime \prime })$$
As a consequence we obtain a path $`\delta _0^{}`$ in $`A^\%(M)`$ joining $`k_0\chi ^t(M_0)+f_{}C`$ with the point
$$[M_0\times [1,1]_{M_0\times \{1\}}M,M_0\times [1,1]_{M_0\times \{1\}}M\times [0,\mathrm{})]A^\%(M)$$
Finally notice that the retractive spaces $`M_0\times [1,1]_{M_0\times \{1\}}M`$ and $`M_0\times [1,1]_{M_0\times \{1\}}M\times [0,\mathrm{})`$ are weakly equivalent to the trivial retractive spaces $`MM`$, and (respectively) $`M\times [0,\mathrm{})M\times [0,\mathrm{})`$, which implies that $`\delta _0^{}`$ can be further extended to the basepoint of $`A^\%(M)`$. The path $`\delta _0`$ can be now obtained shifting $`\delta _0^{}`$ by the element $`k_0\chi ^t(M_0)A^\%(M)`$.
## 4. Flat bundles
The argument which led us to the proof of Theorem 3.1 can be generalized to give a proof of Theorem 1.3 for a certain class of fiber bundles. Namely, assume that $`M`$ is a closed topological manifold with a codimension one splitting $`M=M_1_{M_0}M_2`$, and let $`G`$ be a discrete group acting on $`M`$ on the right by homeomorphisms which preserve the splitting. In such case the bundle $`p^G:EG\times _GMBG`$ splits into subbundles $`p_i^G:EG\times _GM_iBG`$. For $`i=0,1,2`$ let $`j_i:EG\times _GM_iEG\times _GM`$ denote the inclusion map. We have
###### 4.1 Proposition.
For the bundle $`p^G`$ as above we have a path in $`A^\%(p^G)`$ joining $`\chi ^t(p^G)`$ with $`j_i\chi ^t(p_1^G)+j_i\chi ^t(p_1^G)j_i\chi ^t(p_1^G)`$.
The proof of this fact will rely on two lemmas. The first lemma generalizes lax naturality of the topological Euler characteristic (2.10). Recall (2.9) that if $`𝒞`$ is a small category and if $`F:𝒞𝒮paces`$ is a diagram of compact ENRs then we have assignments $`cF(c)^t`$ for $`c𝒞`$ and $`fF(f)^t`$ for a morphism $`f`$ in $`𝒞`$, which by Lemma 2.12 define a point $`[F(c)^t,F(f)^t]\mathrm{holim}_{c𝒞}A^\%(F(c))`$.
###### 4.2 Lemma.
Let $`𝒞`$ be a small category. Assume that $`F,G:𝒞𝒮paces`$ are diagrams of compact ENRs and cell-like maps, and let $`\eta :FG`$ be a natural transformation such that for every $`c𝒞`$ the map $`\eta _c:F(c)G(c)`$ is cell-like. Then there is a path in $`\mathrm{holim}_{c𝒞}A^\%(G(c))`$ joining $`\eta _{}[F(c)^t,F(f)^t]`$ with $`[G(c)^t,G(f)^t]`$. Here
$$\eta _{}:\mathrm{holim}_{c𝒞}A^\%(F(c))\mathrm{holim}_{c𝒞}A^\%(G(c))$$
is the map induced by $`\eta `$.
Proof of this fact resembles justification for Lemma 2.12. The natural transformation $`\eta `$ defines a path in $`\mathrm{holim}_{c𝒞}|w^\%(G(c))|`$. Using the map $`\mathrm{holim}_{c𝒞}|w^\%(G(c))|\mathrm{holim}_{c𝒞}A^\%(G(c))`$ we obtain the required path.
The second lemma describes how one can construct paths in homotopy limits of diagrams of $`K`$-theory spaces using cofibration sequences.
###### 4.3 Lemma.
Let $`F:𝒞𝒲Cat`$ be a functor, and for $`i=1,2,3`$ let $`cc_i^!`$, $`ff_i^!`$ be assignments as in the Lemma 2.12. Assume also that for every $`c𝒞`$ we have a cofibration sequence in $`F(c)`$:
$$c_1^!\stackrel{\phi _c}{}c_2^!\stackrel{\varphi _c}{}c_3^!$$
such that for any morphism $`cd`$ in $`𝒞`$ the following diagram commutes:
Then there is a path in $`\mathrm{holim}_𝒞K(F(c))`$ joining the point $`[c_2^!,f_2^!]`$ with $`[c_1^!,f_1^!]+[c_3^!,f_3^!]`$.
Indeed, in the notation of the cofibration sequences $`c_1^!\stackrel{\phi _c}{}c_2^!\stackrel{\varphi _c}{}c_3^!`$ define a point $`[c_1^!\stackrel{\phi _c}{}c_2^!\stackrel{\varphi _c}{}c_3^!]\mathrm{holim}_{c𝒞}|wS_2F(c)|`$. Also, we have a map
$$\underset{c𝒞}{\mathrm{holim}}|wS_2F(c)|\times \mathrm{\Delta }^2\underset{c𝒞}{\mathrm{holim}}|wS_{}F(c)|$$
Restricting this map to $`[c_1^!\stackrel{\phi _c}{}c_2^!\stackrel{\varphi _c}{}c_3^!]\times \mathrm{\Delta }^2`$ we obtain the desired path.
###### Proof of Proposition 4.1.
Consider the group $`G`$ as a category with one object $``$. The action of $`G`$ on $`M`$ defines functors $`F:G^{op}Spaces`$ and $`F_i:G^{op}Spaces`$ (where $`G^{op}`$ is the opposite category of $`G`$) such that $`F()=A^\%(M)`$ and $`F_i()=A^\%(M_i)`$ for $`i=0,1,2`$. One can check that we have weak equivalences.
$$A^\%(p^G)\underset{G^{op}}{\mathrm{holim}}A^\%(M)\mathrm{and}A^\%(p_i^G)\underset{G^{op}}{\mathrm{holim}}A^\%(M_i)$$
Moreover, the maps
$$k_i:\underset{G^{op}}{\mathrm{holim}}A^\%(M_i)\underset{G^{op}}{\mathrm{holim}}A^\%(M)$$
induced by the inclusions $`k_i:M_iM`$ correspond under these weak equivalences to the maps $`j_i:A^\%(p_i^G)A^\%(p^G)`$. Also, the point of $`\mathrm{holim}_{G^{op}}A^\%(M)`$ corresponding to $`\chi ^t(p^G)`$ is (in the notation of 2.5 and 2.9) the point $`[M^t,F(g)^t]`$, and similarly the Euler characteristics $`\chi ^t(p_i^G)`$, $`i=0,1,2`$ correspond to $`[M_i^t,F_i(g)^t]\mathrm{holim}_{G^{op}}A^\%(M_i)`$. As a consequence it is enough to construct a path $`\omega `$ in $`\mathrm{holim}_{G^{op}}A^\%(M)`$ joining the point $`[M^t,F(g)^t]`$ with $`k_1[M_1^t,F(g)^t]+k_2[M_2^t,F(g)^t]k_0[M_0^t,F(g)^t]`$.
The construction of this path follows the same steps as the proof of Theorem 3.1. As in that proof we construct a manifold $`\overline{M}=M_1_{M_0\times \{1\}}M_0\times [1,1]_{M_0\times \{1\}}M_2`$. Any homeomorphism of $`M`$ extends to a homeomorphism of $`\overline{M}`$ (which is a product map on $`M_0\times [1,1]`$), thus the group $`G`$ acts on $`\overline{M}`$, and we have a functor $`\overline{F}:G^\%(\overline{M})`$ such that $`\overline{F}()=^\%(\overline{M})`$. The map $`f:\overline{M}M`$ contracting $`M_0\times [1,1]\overline{M}`$ to $`M_0M`$ is $`G`$-equivariant, so it defines a natural transformation $`f_{}:\overline{F}F`$ which in turn induces a map of homotopy limits
$$f_{}:\underset{G^{op}}{\mathrm{holim}}A^\%(\overline{M})\underset{G^{op}}{\mathrm{holim}}A^\%(M)$$
Similarly as in the proof of Theorem 3.1 we can now build the path $`\omega `$ as a concatenation of three paths. The first path in 3.1, (joining $`f_{}\chi ^t(\overline{M})`$ with $`\chi ^t(M)`$) was obtained using lax naturality of $`\chi ^t`$. In our present context we can construct a path in $`\mathrm{holim}_{G^{op}}A^\%(M)`$ joining $`f_{}[\overline{M}^t,\overline{F}(g)^t]`$ with $`[M^t,F(g)^t]`$ by applying Lemma 4.2 to the natural transformation $`f_{}:\overline{F}F`$. The other two paths in the proof of Theorem 3.1 were constructed using certain cofibration sequences in $`^\%(\overline{M})`$ and $`^\%(M)`$. One can start with the same cofibration sequences and then for each retractive space appearing in them choose assignments $`\{gg^!\}_{gG}`$ in such way that the conditions of Lemma 4.3 are satisfied. As a consequence one obtains the desired paths in $`\mathrm{holim}_{G^{op}}A^\%(\overline{M})`$ and $`\mathrm{holim}_{G^{op}}A^\%(M)`$. ∎
## 5. Universal bundles
In the last section we verified that Theorem 1.3 holds for some bundles whose structure group is discrete. Our strategy of proving Theorem 1.3 in its whole generality is as follows. In this section we show that the bundles which admit the required splitting into subbundles are induced by some universal bundle $`p^U`$. Moreover, additivity of the topological Euler characteristic for an arbitrary bundle follows from its additivity for this universal bundle. Thus, we only need to show that $`\chi ^t`$ is additive for the bundle $`p^U`$. We prove this special case, using the results of the last section, in §6.
Let $`M`$ be a closed topological manifold which admits a splitting along a codimension one submanifold $`M_0`$:
$$MM_1_{M_0}M_2$$
Let $`TOP(M)`$ be the simplicial group of homeomorphisms of $`M`$ and let $`T`$ denote the subgroup of $`TOP(M)`$ consisting of homeomorphisms which preserve the splitting. Consider the bundle
$$p^U:ET\times _TMBT$$
The bundle $`p^U`$ admits a fiberwise codimension one splitting into sub-bundles $`p_i^U:ET\times _TM_iBT`$ for $`i=0,1,2`$, such that the fiber of $`p_i^U`$ is $`M_i`$. Moreover, $`p^U`$ is the universal bundle for bundles $`E\stackrel{𝑝}{}B`$ with fiber $`M`$ which admit such a splitting. Thus, if $`p`$ splits into subbundles $`p_i`$ with fibers $`M_i`$ for $`i=0,1,2`$ then we have a map $`c:BBT`$ which fits into pullback diagrams:
(1)
Moreover, these maps commute with all inclusions of sub-bundles. Our goal will be to show that the statement of Theorem 1.3 holds for the bundle $`p^U`$.
###### 5.1 Proposition.
Let $`\iota _i:ET\times _TM_iET\times _TM`$ be the inclusion map ($`i=0,1,2`$). Then there is a path $`\sigma _{p^U}`$ in $`A^\%(p^U)`$ joining $`\chi ^t(p^U)`$ with $`\iota _1\chi ^t(p_1^U)+\iota _2\chi ^t(p_2^U)\iota _0\chi ^t(p_0^U)`$.
We postpone the proof of this fact until §6. Meanwhile we will show that, as indicated at the beginning of this section, Theorem 1.3 can be obtained from this special case. We will need the following lemmas which follow directly from the constructions of $`A^\%(p)`$, $`\chi ^t(p)`$, and from Lemma 4.2.
###### 5.2 Lemma.
Assume that we have a pullback square
where $`p`$, $`p^{}`$ are fiber bundles of compact topological manifolds. Let $`\overline{f}^{}:A^\%(p)A^\%(p^{})`$ denote the map induced by the pullback. Then there exists a path in $`A^\%(p)`$ joining $`\chi ^t(p^{})`$ with $`\overline{f}^{}\chi ^t(p)`$.
###### 5.3 Lemma.
Assume that we have a commutative diagram
where for $`i=1,2`$ the maps $`p_i`$ are bundles of compact topological manifolds, $`p_i^{}`$ is a sub-bundle of $`p_i`$, $`j_i:E_i^{}E_i`$ is the inclusion map, and both squares in the middle are pullbacks. Then $`j_1\overline{f}^{}=\overline{f}^{}j_2`$, where $`j_i:A^\%(p_i^{})A^\%(p_i)`$, $`\overline{f}^{}:A^\%(p_2)A^\%(p_1)`$, $`\overline{f}^{}:A^\%(p_2^{})A^\%(p_1^{})`$ are the maps induced by $`j_i,\overline{f},\overline{f^{}}`$ respectively.
###### Proof of Theorem 1.3.
Let $`p:EB`$ be a fiber bundle as in the statement of Theorem 1.3, and let $`c,c_i,\overline{c},\overline{c}_i`$ denote the maps as in the pullback squares (1) above. Applying the map $`\overline{c}^{}`$ to the path $`\sigma _{p^U}`$ from Proposition 5.1 we obtain a path $`\overline{c}^{}\sigma _{p^U}`$ in $`A^\%(p)`$ joining $`\overline{c}^{}\chi ^t(p^U)`$ with $`\overline{c}^{}\iota _1\chi ^t(p_1^U)+\overline{c}^{}\iota _2\chi ^t(p_2^U)\overline{c}^{}\iota _0\chi ^t(p_0^U)`$. Lemma 5.2 implies existence of a path $`\eta `$ joining $`\chi ^t(p)`$ with $`\overline{c}^{}\chi ^t(p^U)`$. Similarly for $`i=0,1,2`$ we have paths $`\eta _i`$ in $`A^\%(p_i)`$ joining $`\overline{c}_i^{}\chi ^t(p_i^U)`$ with $`\chi ^t(p_i)`$. Take $`\eta ^{}:=j_1\eta _1+j_2\eta _2j_0\eta _0`$. This is a path in $`A^\%(p)`$ with endpoints $`j_1\overline{c}_1^{}\chi ^t(p_1^U)+j_2\overline{c}_2^{}\chi ^t(p_2^U)j_0\overline{c}_0^{}\chi ^t(p_0^U)`$ and $`j_1\chi ^t(p_1)+j_1\chi ^t(p_2)j_1\chi ^t(p_0)`$. By Lemma 5.3 we have $`j_i\overline{c}_i^{}=\overline{c}^{}\iota _{}`$. Therefore we can concatenate $`\eta `$, $`\overline{c}^{}\sigma _{p^U}`$ and $`\eta ^{}`$ and obtain a path in $`A^\%(p)`$ which joins $`\chi ^t(p)`$ with $`j_1\chi ^t(p_1)+j_2\chi ^t(p_2)j_0\chi ^t(p_0)`$. ∎
## 6. From topological to discrete structure group
As the previous section demonstrated, checking additivity for the parametrized topological Euler characteristic reduces to showing that it holds for the universal bundle $`p^U:ET\times _TMBT`$. The structure group of this bundle is the simplicial group $`T`$. Our next goal is to show that it is enough to verify additivity for a certain flat bundle, i.e. a bundle whose structure group is discrete.
Recall that a map $`f:B^{}B`$ is a homology equivalence if it induces isomorphisms on homology groups with arbitrary local coefficients.
###### 6.1 Lemma.
Let $`E\stackrel{𝑝}{}B`$ be a fiber bundle with a compact topological manifold $`M=M_1_{M_0}M_2`$ as a fiber. Assume that $`p`$ admits a decomposition into sub-bundles as in the statement in Theorem 1.3, and that we have a pullback diagram
If the map $`f:B^{}B`$ is a homology equivalence, then the additivity path for $`\chi ^t`$ exists for the bundle $`p`$ if it exists for the bundle $`p^{}`$ (where $`p^{}`$ comes with decomposition into sub-bundles induced from $`p`$).
Lemma 6.1 follows from the following
###### 6.2 Proposition.
Assume that we have a homotopy pullback square
such that $`q`$, $`q^{}`$ are fibrations whose fibers are nilpotent spaces. Assume also that $`f:BB^{}`$ is a homology equivalence. Then the induced map of the spaces of sections
$$\mathrm{\Gamma }(q)\mathrm{\Gamma }(q^{})$$
is a weak equivalence.
###### Proof.
See \[5, Proof of Cor. 2.7\]
###### Proof of Lemma 6.1.
Let $`p_i:E_iB`$ be the sub-bundles in the decomposition of $`p`$, and let $`j_i:E_iE`$ be the inclusion maps. Using the assumption that we have an additivity path for $`p^{}`$ and Lemma 5.2 we obtain a path in $`A(p^{})`$ joining $`\overline{f}^{}\chi ^t(p)`$ with $`\overline{f}^{}(j_1\chi ^t(p_1)+j_2\chi ^t(p_2)j_0\chi ^t(p_0))`$ Thus it suffices to prove that the map $`\overline{f}^{}:A^\%(p)A^\%(p^{})`$ is a weak equivalence. In order to show that consider the homotopy pullback square
where the fibrations $`p_{}^{}`$, $`p_{}`$ are defined as in 2.18. The fiber of $`p_{}`$ and $`p_{}^{}`$ is $`A^\%(M)`$ which is an infinite loop space, thus in particular a nilpotent space. By Lemma 6.2 we have a weak equivalence of the spaces of sections
$$\mathrm{\Gamma }(p_{}^{})\mathrm{\Gamma }(p_{})$$
On the other hand we have $`\mathrm{\Gamma }(p_{})A^\%(p)`$, and $`\mathrm{\Gamma }(p_{}^{})A^\%(p^{})`$ (see 2.17, 2.18). It follows that $`A^\%(p)A^\%(p^{})`$. ∎
Our application of Lemma 6.1 is as follows. Consider the universal bundle $`p^U:ET\times _TMBT`$ as in Section 5. By the collaring theorem the submanifold $`M_0`$ has a bicollar in $`M`$ i.e. we have an embedding $`c:M_0\times (1,1)M`$ such that $`c(m,0)=m`$ for $`mM_0`$. Let $`T_c`$ denote the subgroup of $`T`$ consisting of all these splitting preserving homeomorphisms of $`M`$ which are product maps on the bicollar $`c`$. In other words, $`fT_c`$ if there is a homeomorphism $`f^{}:M_0M_0`$ such that the following diagram commutes:
Let $`T_\epsilon :=\mathrm{colim}_\mathrm{c}T_c`$ where the colimit is taken over all bicollar neighborhoods of $`M_0`$, and let $`T_\epsilon ^\delta `$ denote the group $`T_\epsilon `$, but equipped with the discrete topology. We have
###### 6.3 Lemma.
The homomorphism $`T_\epsilon ^\delta T`$ induces a homology equivalence of classifying spaces $`BT_\epsilon ^\delta BT`$.
###### Proof.
The homomorphism $`T_\epsilon ^\delta T`$ is a composition of the inclusion $`T_\epsilon T`$ and the map $`T_\epsilon ^\delta T_\epsilon `$. The map of classifying spaces induced by the first of these homomorphisms $`BT_\epsilon BT`$ is a homotopy equivalence by Siebenmann’s isotopy extension theorem . The map $`BT_\epsilon ^\delta BT_\epsilon `$ induced by the second homomorphism is a homology equivalence by the results of McDuff , Thurston , Segal , and Mather . ∎
We are now in position to give a proof of Lemma 5.1, and thus complete the proof of Theorem 1.3.
###### Proof of Lemma 5.1.
We have a pullback diagram
By Lemma 6.3 the map $`BT_\epsilon ^\delta BT`$ is a homology equivalence. Moreover, by Proposition 4.1 the additivity of $`\chi ^t`$ holds for the bundle $`p^{T_\epsilon ^\delta }`$. Therefore using Lemma 6.1 we obtain additivity of $`\chi ^t`$ for the bundle $`p^U`$. ∎
## 7. Additivity for topological Reidemeister torsion
Our final task is to give the proof of additivity of the topological Reidemeister torsion, i.e. Theorem 1.2. Let then $`E\stackrel{𝑝}{}B`$ be a bundle of compact topological manifolds and let $`V\stackrel{𝜌}{}E`$ be a locally constant sheaf of finitely generated projective left $`R`$-modules such that the assumptions of Theorem 1.2 are satisfied. From the pullback diagram below Definition 2.23 it follows that in order to prove additivity for $`\tau _\rho ^t(p)`$ we need to construct
* a path $`\stackrel{~}{\omega }_p^t`$ in $`\stackrel{~}{A}^\%(p)`$ joining $`\stackrel{~}{\chi }^t(p)`$ with $`j_1\stackrel{~}{\chi }^t(p_1)+j_2\stackrel{~}{\chi }^t(p_2)j_0\stackrel{~}{\chi }^t(p_0)`$;
* a path $`\omega _\rho ^h`$ in $`\mathrm{\Phi }_\rho ^h(p)`$ joining $`\tau _\rho ^h(p)`$ with $`j_1\tau _\rho ^h(p_1)+j_2\tau _\rho ^h(p_2)j_0\tau _\rho ^h(p_0)`$ and such that $`\delta _1(\stackrel{~}{\omega }^\%)=\delta _2(\omega _\rho ^h)`$.
Construction of the path $`\stackrel{~}{\omega }_p^t`$. Recall that the space $`\stackrel{~}{A}^\%(p)`$ was obtained as a homotopy pullback of the diagram
$$\underset{𝑥}{\mathrm{holim}}A(E_x)\stackrel{𝛽}{}\underset{(x,\theta )}{\mathrm{holim}}A(E_x^\theta )\stackrel{𝛼}{}\underset{(x,\theta )}{\mathrm{holim}}A^\%(E_x^\theta )$$
and that the point $`\stackrel{~}{\chi }^t(p)\stackrel{~}{A}(p)`$ was represented by the triple $`(\chi ^h(p),\sigma _p,\chi ^t(p))`$ where $`\sigma _p`$ is the canonical path in $`\mathrm{holim}_{(x,\theta )}A(E_x^\theta )`$ joining the images of $`\chi ^h(p)`$ and $`\chi ^t(p)`$. As a consequence in order to describe the path $`\stackrel{~}{\omega }_p^t`$ we need to construct
* a path $`\omega _p^t`$ in $`\mathrm{holim}_{(x,\theta )}A^\%(E_x^\theta )`$ joining $`\chi ^t(p)`$ with $`j_1\chi ^t(p_1)+j_2\chi ^t(p_2)j_0\chi ^t(p_0)`$;
* a path $`\omega _p^h`$ in $`\mathrm{holim}_xA(E_x)`$ joining $`\chi ^h(p)`$ with $`j_1\chi ^h(p_1)+j_2\chi ^h(p_2)j_0\chi ^h(p_0)`$;
* a homotopy $`H:[0,1]\times [0,1]\mathrm{holim}_{(x,\theta )}A(E_x^\theta )`$ such that $`H(,0)=\alpha \omega _p^t`$, $`H(,1)=\beta \omega _p^h`$, $`H(0,)=\sigma _\rho `$, and $`H(1,)=j_1\sigma _{p_1}+j_2\sigma _{p_2}j_0\sigma _{p_0}`$.
We take $`\omega _p^t`$ to be the additivity path for the topological Euler characteristic which we constructed in that proof of Theorem 1.3. In order to construct a suitable path $`\omega _p^h`$ we will use $`A^{biv}(p)`$ – the bivariant $`A`$-theory of the bundle $`p`$ . The space $`A^{biv}(p)`$ is the $`K`$-theory space of the Waldhausen category $`^{fd}(p)`$ which is a full subcategory of the category of retractive spaces over $`E`$. The objects of $`^{fd}(p)`$ are these retractive spaces $`r:XE:s`$ which satisfy the condition that $`s`$ is a cofibrations and that for every point $`bB`$ the homotopy fiber $`F_{pr}^b`$ of the map $`pr`$ over $`b`$ is a homotopy finitely dominated space. Notice that $`F_{pr}^b`$ is in a natural way a retractive space over $`F_p^b`$ – the homotopy fiber of $`p`$, thus the above assumption says that for every $`bB`$ the fiber $`F_{pr}^b`$ is an object of the Waldhausen category $`^{fd}(F_p^b)`$ (2.1). The Waldhausen category structure on $`^{fd}(p)`$ is defined by taking a morphism to be a weak equivalence or a cofibrations if its underlying map of spaces is a homotopy equivalence or respectively a cofibration.
The retractive space $`EEE`$ is an object of $`^{fd}(p)`$, and so it defines a point $`\chi ^{biv}(p)A^{biv}(p)`$. Let $`j_i:A^{biv}(p_i)A^{biv}(p)`$, $`i=0,1,2`$, be the maps induced by inclusions of subbundles. Using the same constructions as in the proof of Theorem 3.1 with $`E`$, $`E_i`$ taken in place of $`M`$, $`M_i`$ we can construct a path $`\omega _p^{biv}`$ in $`A^{biv}(p)`$ joining $`\chi ^{biv}(p)`$ with $`j_1\chi ^{biv}(p_1)+j_2\chi ^{biv}(p_2)j_0\chi ^{biv}(p_0)`$.
For every simplex $`x𝒮(𝔅)`$ we have an exact functor $`^{fd}(p)^{fd}(E_x)`$ which yields a map of infinite loop spaces $`A^{biv}(p)A(E_x)`$. These maps can be combined to give a map
$$\underset{𝑥}{\mathrm{holim}}A^{biv}(p)\underset{𝑥}{\mathrm{holim}}A(E_x)=A(p)$$
where the first homotopy limit is taken over the constant functor with the value $`A^{biv}(p)`$. Composing this map with the Bousfield-Kan map $`A^{biv}(p)lim_xA^{biv}(p)\mathrm{holim}_xA^{biv}(p)`$ we obtain
$$\alpha ^{}:A^{biv}(p)A(p)$$
(this is the generalized coassembly map of ). The image of $`\chi ^{biv}(p)`$ under the map $`\alpha ^{}`$ does not coincide with $`\chi ^h(p)`$, but there is a canonical path in $`\mathrm{holim}_xA(E_x)`$ which joins these two points. As a consequence the path $`\alpha ^{}\omega _p^{biv}`$ defines an additivity path $`\omega _p^h`$ for the homotopy Euler characteristic.
In order to see that we have the required homotopy $`H`$ one needs to retrace our construction of the additivity path for the topological Euler characteristic. First, one considers bundles with a discrete structure group where the homotopy $`H`$ can described using constructions on the level of Waldhausen categories. For more general bundles one uses the fact that the additivity path $`\omega _p^t`$ was obtained from the above special case by means of pullbacks. This is enough to verify that the homotopy $`H`$ will still exist.
Construction of the path $`\omega _\rho ^h`$. Notice that the path $`\omega _\rho ^h`$ is an additivity path for the homotopy Reidemeister torsion of the bundle $`p`$. The existence of such path was proved in . In our present setting we need the path $`\omega _\rho ^h`$ which is compatible with $`\stackrel{~}{\omega }_p^t`$, thus it is not enough to quote that result. The path $`\omega _\rho ^h`$ can be obtained however by mimicking the constructions of . Briefly, one starts by constructing $`A^{ac}(p)`$ \- the acyclic bivariant $`A`$-theory of the bundle $`p`$. The space $`A^{ac}(p)`$ is obtained from a Waldhausen category $`^{ac}(p)`$ which is a full subcategory of $`^{fd}(p)`$. The objects of $`^{ac}(p)`$ are these retractive spaces which satisfy the condition that for every $`bB`$ the relative chain complex of the pair $`(F_{pr},F_r)`$ with local coefficients in $`\rho |_{F_r}`$ is acyclic (as before $`F_{pr}`$ and $`F_p`$ denote here the homotopy fibers over $`b`$ of the maps $`pr`$ and $`p`$ respectively). The inclusion of Waldhausen categories $`^{ac}(p)^{fd}(p)`$ induces a map $`A^{ac}(p)A^{biv}(p)`$. Directly from the construction of the path $`\omega _p^{biv}`$ it follows that it admits a lift $`\omega ^{ac}`$ to $`A^{ac}(p)`$. The space $`A^{ac}(p)`$ in turn maps into the homotopy Whitehead space $`\mathrm{\Phi }_\rho ^h(p)`$ in such way that we obtain a commutative diagram
We can take $`\omega _\rho ^h`$ to be the image of $`\omega ^{ac}`$ under this map.
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# Yet Another Normalisation Proof for Martin-Löf’s Logical Framework
## 1 Introduction
The normalisation proofs for dependently typed systems are known to be notoriously difficult. For example, if we have a task to prove strong normalisation for Martin-Löf’s Logical Framework (MLF) (in the Appendix), and if we use typed operational semantics as in \[Gog94\], the proof would be more than one hundred pages long. When a proof is long and complicated, it is likely found to contain mistakes and bugs \[Coq85, CG90, Alt94\]. This paper presents an elegant and comprehensible proof of strong normalisation for MLF.
We often associate well-typedness with strong normalisation in type systems. But this paper suggests that well-typedness may have little to do with strong normalisation in essence, and proves that terms with correct arities are strongly normalising. The condition of “correct arity” is weaker than that of well-typedness (*i.e.* well-typed terms have correct arities). This paper will also demonstrate the difference between types and arities when we extend MLF with some inductive data types and their computation rules. New reduction rules will not increase the set of terms with correct arities, but they will usually increase the set of well-typed terms. One of the reasons is that there are reductions inside types (*i.e.* one type can be reduced to another type) in a dependently typed system but there is no reduction for arities.
Our goal is to prove the strong normalisation w.r.t. $`\beta `$ and $`\eta `$-reduction. But it is very difficult to prove it directly. An important technique in the paper is that, we extend the definition of terms and kinds, and introduce a new reduction rule $`\beta _2`$ for kinds. Then, we prove a stronger and more general property, that is, strong normalisation w.r.t. $`\beta `$, $`\eta `$ and $`\beta _2`$-reduction. In this way, the proof becomes easier although the property is stronger. Without the $`\beta _2`$-reduction, the proof of soundness in Section 4 is impossible to go through.
In Section 2, we give some basic definitions that are used throughout the paper. In Section 3, the inference rules of arities are formally presented. In Section 4, we give more definitions such as saturated sets, and prove the strong normalisation for the arity system. In Section 5, the computation rules for the type of dependent pairs and finite types and simple computation rules for universes are introduced. The strong normalisation for a dependently typed system is proved by the commutation property between these rules and $`\beta `$-reduction. The conclusions and future work are discussed in the last section.
#### Related work
Logical frameworks arise because one wants to create a single framework, which is a kind of meta-logic or universal logic. The Edinburgh Logical Framework \[HHP87, HHP92\] presents logics by a *judgements-as-types* principle, which can be regarded as the meta-theoretical analogue of the well-known *propositions-as-types* principles \[CF58, dB80, How80\]. Martin-Löf’s logical framework \[ML84, NPS90\] has been developed by Martin-Löf to present his intensional type theory. In UTT \[Luo94\], Luo proposed a typed version of Martin-Löf’s logical framework, in which untyped functional operations of the form $`(x)k`$ are replaced by typed $`[x:K]k`$.
There are many normalisation proofs for simply typed systems and dependently typed systems in literature \[Bar92, Luo90, Alt93\] \[MW96, Gog94\] \[Geu93, Wer92\]. The techniques employed in this paper such as the interpretation of arities and saturated sets are inspired by and closely related to the proof for simply typed calculus in \[Bar92\]. The concept of arity is well-known in mathematics and it is often defined as the maximum number of arguments that a function can have. But in this paper, the definition of arity and the concept of “correct arity” are different. The complexity of the normalisation proof for MLF is dramatically decreased because of this concept and other techniques such as a new case of kinds and the corresponding $`\beta _2`$-reduction. The commutation property was also studied in literature such as \[Bar84, Cos96\]. The properties of Church-Rosser and strong normalisation for finite types in simply typed systems are also studied in \[SC04\].
## 2 Basic definitions
In this section, we give some basic definitions that will be used later, and give the redice and the corresponding reduction rules.
###### Definition 1
(Terms and Kinds)
* Terms
1. a variable is a term,
2. $`\lambda x:K.M`$ is a term if $`x`$ is a variable, $`K`$ is a kind and $`M`$ is a term,
3. $`MN`$ is a term if $`M`$ and $`N`$ are terms.
* Kinds
1. $`Type`$ is a kind,
2. $`El(M)`$ is a kind if $`M`$ is a term,
3. $`(x:K_1)K_2`$ is a kind if $`K_1`$ and $`K_2`$ are kinds,
4. $`KN`$ is a kind if $`K`$ is a kind and $`N`$ is a term.
###### Remark 1
Terms and kinds are mutually and recursively defined. This definition allows more terms and kinds than that of MLF since the forth case for the definition of kinds is not included in MLF (see Appendix for details).
Following the tradition, $`\mathrm{\Lambda }`$ denotes the set of all terms and $`\mathrm{\Pi }`$ the set of all kinds. We sometimes write $`f(a)`$ for $`fa`$, $`f(a,b)`$ for $`(fa)b`$ and so on. $`[N/x]M`$ stands for the expression obtained from $`M`$ by substituting $`N`$ for the free occurrences of variable $`x`$ in $`M`$. $`FV(M)`$ is the set of free variables in $`M`$.
### Redice and reduction rules
There are three different forms of redice: $`(\lambda x:K.M)N`$, $`((x:K_1)K_2)N`$ and $`\lambda x:K.Mx`$ when $`xFV(M)`$. The reduction rules for these redice are the following.
$$(\lambda x:K.M)N_\beta [N/x]M$$
$$((x:K_1)K_2)N_{\beta _2}[N/x]K_2$$
$$\lambda x:K.Mx_\eta MxFV(M)$$
###### Remark 2
The second rule $`_{\beta _2}`$ is new and is not included in MLF. This rule will make the soundness proof go through easily although the property is stronger and more general.
$`_R`$ represents one-step $`R`$-reduction, precisely, $`M_RN`$ if a sub-term $`P`$ of $`M`$ is a $`R`$-redex and $`N`$ is obtained by replacing $`P`$ by the result after applying the reduction rule $`R`$. $`M_RN`$ means there is $`0`$ or more but finite steps of $`R`$-reduction from $`M`$ to $`N`$. $`M_R^+N`$ means there is at least one but finite steps of $`R`$-reduction from $`M`$ to $`N`$.
###### Definition 2
(Arities)
* $`Zero`$ is an arity,
* $`(a_1,a_2)`$ is an arity if $`a_1`$ and $`a_2`$ are arities.
$`\mathrm{\Omega }`$ denotes the set of all arities.
## 3 Inference rules
In this section, we formally present the inference rules of arities.
The judgement form will be the following form,
$$AM:a$$
where $`A<x_1:a_1,\mathrm{},x_n:a_n>`$ is a finite sequence of $`x_i:a_i`$, $`x_i`$ is a variable and $`a_i`$ is an arity; $`M`$ is a term or kind; and $`a`$ is an arity. We shall read this judgement like “under the context $`A`$, the term or kind $`M`$ has arity $`a`$”.
For a context $`Ax_1:a_1,\mathrm{},x_n:a_n`$, $`FV(A)`$ represents the set $`\{x_1,\mathrm{},x_n\}`$.
All of the inference rules of arities are in Figure 1.
###### Definition 3
We say that a term or kind $`M`$ has a correct arity if $`AM:a`$ is derivable for some $`A`$ and $`a`$.
###### Remark 3
We have the following remarks:
* A well-typed term has a correct arity (a proof will be given later), but a term which has a correct arity is not necessarily well-typed. For instance, under the context
$$A:Type,B:Type,C:Type,f:(x:A)C,b:B$$
the term $`f(b)`$ is not well-typed, but it has a correct arity $`Zero`$ under the following context
$$A:Zero,B:Zero,C:Zero,f:(Zero,Zero),b:Zero$$
Another example with dependent type is that, under the context
$$\begin{array}{c}A:Type,B:(x:A)Type,f:(x:A)(y:B(x))Type,\\ x_1:A,x_2:A,b:B(x_2)\end{array}$$
the term $`f(x_1,b)`$ is not well-typed, but it has a correct arity $`Zero`$ in the following context
$$\begin{array}{c}A:Zero,B:(Zero,Zero),f:(Zero,(Zero,Zero)),\\ x_1:Zero,x_2:Zero,b:Zero\end{array}$$
* For any judgement $`AM:a`$, $`M`$ must be either a kind or a term. A derivation such as $`\frac{AType:Zero}{AEl(Type):Zero}`$ is not possible, because $`El(Type)`$ is neither a term nor a kind.
###### Lemma 1
If both $`AM:a`$ and $`AM:b`$ are derivable then $`a`$ and $`b`$ are syntactically the same ($`ab`$). And $`AMM:a`$ is not derivable for any $`A`$, $`M`$ and $`a`$.
###### Proof
By induction on the derivations of $`AM:a`$ and $`AM:b`$.
###### Remark 4
One may recall that the non-terminating example $`\omega \omega `$ where $`\omega \lambda x.xx`$. It is impossible that $`\omega `$ is well-typed in a simply typed calculus \[Bar92\]. By Lemma 1, it is also impossible to have a correct arity for $`\omega `$.
## 4 Normalisation proof
In this section, we give more definitions such as saturated sets to prove the strong normalisation for the arity system.
###### Definition 4
(Interpretation of arities)
* $`SN^\mathrm{\Lambda }=_{df}\{M\mathrm{\Lambda }|Misstronglynormalising\}`$.
* $`SN^\mathrm{\Pi }=_{df}\{M\mathrm{\Pi }|Misstronglynormalising\}`$.
* $`Zero^\mathrm{\Lambda }=_{df}SN^\mathrm{\Lambda }`$.
* $`Zero^\mathrm{\Pi }=_{df}SN^\mathrm{\Pi }`$.
* $`(a_1,a_2)^\mathrm{\Lambda }=_{df}\{M\mathrm{\Lambda }|Na_1^\mathrm{\Lambda },MNa_2^\mathrm{\Lambda }\}`$.
* $`(a_1,a_2)^\mathrm{\Pi }=_{df}\{K\mathrm{\Pi }|Na_1^\mathrm{\Lambda },KNa_2^\mathrm{\Pi }\}`$.
###### Remark 5
$`a^\mathrm{\Lambda }`$ is a set of terms, while $`a^\mathrm{\Pi }`$ is a set of kinds for any arity $`a`$.
We shall write $`\overline{R}`$ for $`R_1,R_2,\mathrm{},R_n`$ for some $`n0`$, and $`M\overline{R}`$ for $`(\mathrm{}((MR_1)R_2)\mathrm{}R_n)`$.
###### Definition 5
(Saturated sets)
* A subset $`XSN^\mathrm{\Lambda }`$ is called saturated if
1. $`\overline{R}SN^\mathrm{\Lambda }`$, $`x\overline{R}X`$ where $`x`$ is any term variable,
2. $`\overline{R}SN^\mathrm{\Lambda }`$, $`QSN^\mathrm{\Lambda }`$ and $`KSN^\mathrm{\Pi }`$,
$$([Q/x]P)\overline{R}X(\lambda x:K.P)Q\overline{R}X$$
* A subset $`YSN^\mathrm{\Pi }`$ is called saturated if $`\overline{R}SN^\mathrm{\Lambda }`$, $`NSN^\mathrm{\Lambda }`$ and $`K_1SN^\mathrm{\Pi }`$,
$$([N/x]K_2)\overline{R}Y((x:K_1)K_2)N\overline{R}Y$$
* $`SAT^\mathrm{\Lambda }=_{df}\{XSN^\mathrm{\Lambda }|Xissaturated\}`$
* $`SAT^\mathrm{\Pi }=_{df}\{YSN^\mathrm{\Pi }|Yissaturated\}`$
###### Lemma 2
(Arities and saturated sets)
* $`SN^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$ and $`SN^\mathrm{\Pi }SAT^\mathrm{\Pi }`$.
* $`a\mathrm{\Omega }a^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$ and $`a^\mathrm{\Pi }SAT^\mathrm{\Pi }`$.
###### Proof
By the definition of saturated sets and by induction on arities.
* Let’s prove $`SN^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$ first. We have $`SN^\mathrm{\Lambda }SN^\mathrm{\Lambda }`$ and $`x\overline{R}SN^\mathrm{\Lambda }`$ if $`\overline{R}SN^\mathrm{\Lambda }`$. Now we need to prove for $`Q,\overline{R}SN^\mathrm{\Lambda }`$ and $`KSN^\mathrm{\Pi }`$,
$$([Q/x]P)\overline{R}SN^\mathrm{\Lambda }(\lambda x:K.P)Q\overline{R}SN^\mathrm{\Lambda }$$
Since $`([Q/x]P)\overline{R}SN^\mathrm{\Lambda }`$, we have $`PSN^\mathrm{\Lambda }`$ and after any finitely many steps reducing inside $`P`$, $`Q`$ and $`\overline{R}`$, $`([Q^{}/x]P^{})\overline{R^{}}SN^\mathrm{\Lambda }`$ with $`P_{\beta \eta }P^{}`$ , $`Q_{\beta \eta }Q^{}`$ and $`\overline{R}_{\beta \eta }\overline{R^{}}`$.
From $`(\lambda x:K.P)Q\overline{R}`$, after any finitely many steps reducing inside $`P`$, $`Q`$, $`\overline{R}`$ and $`K`$, and we get $`(\lambda x:K^{}.P^{})Q^{}\overline{R^{}}`$. From here, we may have two choices.
+ $`(\lambda x:K^{}.P^{})Q^{}\overline{R^{}}_\beta ([Q^{}/x]P^{})\overline{R^{}}`$
+ $`P^{}Fx`$ and $`xFV(F)`$ and
$$(\lambda x:K^{}.P^{})Q^{}\overline{R^{}}_\eta FQ^{}\overline{R^{}}([Q^{}/x]P^{})\overline{R^{}}$$
For both cases, because $`([Q^{}/x]P^{})\overline{R^{}}SN^\mathrm{\Lambda }`$, we have $`(\lambda x:K.P)Q\overline{R}SN^\mathrm{\Lambda }`$.
* The proof of $`SN^\mathrm{\Pi }SAT^\mathrm{\Pi }`$ is similar to that of $`SN^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$.
* Now, let’s prove $`a^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$ by induction on $`a`$. The base case (*i.e.* $`Zero^\mathrm{\Lambda }=SN^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$) has been proved. So we only need to prove $`(a_1,a_2)^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$. By induction hypothesis, we have $`a_1^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$ and $`a_2^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$.
Then we have $`xa_1^\mathrm{\Lambda }`$ for all variable $`x`$. Therefore
$`F(a_1,a_2)^\mathrm{\Lambda }`$ $``$ $`Fxa_2^\mathrm{\Lambda }`$
$``$ $`FxSN^\mathrm{\Lambda }`$
$``$ $`FSN^\mathrm{\Lambda }`$
So, we have $`(a_1,a_2)^\mathrm{\Lambda }SN^\mathrm{\Lambda }`$.
Now, we need to prove that for any variable $`x`$ and $`\overline{R}SN^\mathrm{\Lambda }`$, we have $`x\overline{R}(a_1,a_2)^\mathrm{\Lambda }`$. This means
$$Na_1^\mathrm{\Lambda }x\overline{R}Na_2^\mathrm{\Lambda }$$
which is true since $`a_1^\mathrm{\Lambda }SN^\mathrm{\Lambda }`$ and $`a_2^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$.
Finally, we need to prove that for $`\overline{R}SN^\mathrm{\Lambda }`$, $`QSN^\mathrm{\Lambda }`$ and $`KSN^\mathrm{\Pi }`$,
$$([Q/x]P)\overline{R}(a_1,a_2)^\mathrm{\Lambda }(\lambda x:K.P)Q\overline{R}(a_1,a_2)^\mathrm{\Lambda }$$
Since $`([Q/x]P)\overline{R}(a_1,a_2)^\mathrm{\Lambda }`$, we have $`([Q/x]P)\overline{R}Na_2^\mathrm{\Lambda }`$ for $`Na_1^\mathrm{\Lambda }`$. And since $`a_1^\mathrm{\Lambda }SN^\mathrm{\Lambda }`$ and $`a_2^\mathrm{\Lambda }SAT^\mathrm{\Lambda },`$ we have $`(\lambda x:K.P)Q\overline{R}Na_2^\mathrm{\Lambda }`$ and hence
$$(\lambda x:K.P)Q\overline{R}(a_1,a_2)^\mathrm{\Lambda }$$
* The proof of $`a^\mathrm{\Pi }SAT^\mathrm{\Pi }`$ is similar to that of $`a^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$
We often use $`SN`$ for $`SN^\mathrm{\Lambda }SN^\mathrm{\Pi }`$ and $`a`$ for $`a^\mathrm{\Lambda }a^\mathrm{\Pi }`$.
###### Definition 6
(Valuation)
* A *valuation* is a map $`\rho :V\mathrm{\Lambda }`$, where $`V`$ is the set of all term variables.
* Let $`\rho `$ be a valuation. Then
$$M_\rho =_{df}[\rho (x_1)/x_1,\mathrm{},\rho (x_n)/x_n]M$$
where $`x_1,\mathrm{},x_n`$ are all of the free variable in $`M`$.
* Let $`\rho `$ be a valuation. Then
+ $`\rho `$ satisfies $`M:a`$, notation $`\rho M:a,`$ if $`M_\rho a`$;
+ $`\rho `$ satisfies $`A`$, notation $`\rho A`$, if $`\rho x:a`$ for all $`x:aA`$;
+ $`A`$ satisfies $`M:a`$, notation $`AM:a,`$ if
$$\rho (\rho A\rho M:a)$$
###### Remark 6
For any valuation $`\rho `$, if $`M`$ is a term, $`M_\rho `$ is also a term, and similarly, if $`M`$ is a kind, $`M_\rho `$ is also a kind. If a valuation $`\rho `$ satisfies that $`\rho (x)=x`$ then $`M_\rho M`$.
###### Lemma 3
(Soundness) $`\mathrm{A}\mathrm{M}:\mathrm{a}\mathrm{A}\mathrm{M}:\mathrm{a}`$ where $`\mathrm{M}`$ is a term or kind.
###### Proof
By induction on the derivations of $`AM:a`$.
1. The last rule is
$$\frac{Avalid}{AType:Zero}$$
Since $`Type_\rho =Type`$ for any $`\rho `$ and $`TypeSN=Zero`$, we have $`Type_\rho Zero`$.
2. The last rule is
$$\frac{AM:Zero}{AEl(M):Zero}$$
Since $`El(M)_\rho =El(M_\rho )`$ for any $`\rho `$ and $`M_\rho Zero=SN`$, we have $`El(M)_\rho SN=Zero`$.
3. The last rule is
$$\frac{AK_1:a_1A,x:a_1K_2:a_2}{A(x:K_1)K_2:(a_1,a_2)}$$
We must show that
$$\rho (\rho A\rho (x:K_1)K_2:(a_1,a_2))$$
That is, we must show that $`(x:K_1)K_2_\rho (a_1,a_2)^\mathrm{\Pi }`$. By the definition of $`(a_1,a_2)^\mathrm{\Pi }`$, we must show that, for all $`Na_1^\mathrm{\Lambda }`$,
$$(x:K_1)K_2_\rho Na_2^\mathrm{\Pi }$$
Note that
$`(x:K_1)K_2_\rho N`$ $``$ $`((x:K_1^{})K_2^{})N`$
$`_{\beta _2}`$ $`[N/x]K_2^{}`$
$``$ $`K_2_{\rho (N/x)}`$
where $`K_1^{}K_1_\rho [\rho (y_i)/y_i\mathrm{}]K_1`$ and $`K_2^{}K_2_\rho [\rho (y_i)/y_i\mathrm{}]K_2`$
Now, let’s consider the induction hypothesis. Since $`\rho (N/x)A,x:a_1`$, we have $`K_1_\rho a_1^\mathrm{\Pi }`$ and $`K_2_{\rho (N/x)}a_2^\mathrm{\Pi }`$. So, we have $`[N/x]K_2^{}a_2^\mathrm{\Pi }`$, and because $`a_2^\mathrm{\Pi }`$ is saturated, we have $`((x:K_1^{})K_2^{})Na_2^\mathrm{\Pi }`$, i.e. $`(x:K_1)K_2_\rho Na_2^\mathrm{\Pi }`$. Note that, since $`a_1^\mathrm{\Lambda }SN^\mathrm{\Lambda }`$ and $`a_1^\mathrm{\Pi }SN^\mathrm{\Pi }`$, we know that $`NSN^\mathrm{\Lambda }`$ and $`K_1^{}SN^\mathrm{\Pi }`$.
4. The last rule is
$$\frac{AK:(a_1,a_2)AN:a_1}{AKN:a_2}$$
We must show that
$$\rho (\rho A\rho KN:a_2)$$
By induction hypothesis, we have $`K_\rho (a_1,a_2)^\mathrm{\Pi }`$ and $`N_\rho a_1^\mathrm{\Lambda }`$.
By the definition of $`(a_1,a_2)^\mathrm{\Pi }`$, we have $`K_\rho N_\rho a_2^\mathrm{\Pi }`$, *i.e.* $`KN_\rho a_2^\mathrm{\Pi }`$.
5. The last rule is
$$\frac{A,x:a,A^{}valid}{A,x:a,A^{}x:a}$$
Easy.
6. The last rule is
$$\frac{AK:a_1A,x:a_1M:a_2}{A\lambda x:K.M:(a_1,a_2)}$$
Similar to case 3.
7. The last rule is
$$\frac{AM:(a_1,a_2)AN:a_1}{AMN:a_2}$$
Similar to case 4. ∎
###### Theorem 4.1
If $`AM:a`$, then $`M`$ is strongly normalising.
###### Proof
By Lemma 3 and take the evaluation $`\rho _0`$ that satisfies $`\rho _0(x)=x`$.
By Lemma 3, we have $`AM:a`$. So, by definition, we have
$$\rho _0A\rho _0M:a$$
Suppose $`Ax_1:a_1,\mathrm{},x_n:a_n.`$ Since $`a_i^\mathrm{\Lambda }SAT^\mathrm{\Lambda }`$, we have $`x_ia_i^\mathrm{\Lambda }`$. Hence $`\rho _0A`$. So, we have $`\rho _0M:a`$ and hence $`M=M_{\rho _0}aSN`$.∎
### Translation from kinds to arities
Now, we define a map to translate kinds to arities, and prove that well-typed terms have correct arities.
###### Definition 7
A map $`arity:\mathrm{\Pi }\mathrm{\Omega }`$ is inductively defined as follows.
* $`arity(Type)=Zero`$,
* $`arity(El(A))=Zero`$,
* $`arity((x:K_1)K_2)=(arity(K_1),arity(K_2))`$.
Suppose a context $`\mathrm{\Gamma }x_1:K_1,\mathrm{},x_n:K_n`$, then $`arity(\mathrm{\Gamma })x_1:arity(K_1),\mathrm{},x_n:arity(K_n)`$.
###### Theorem 4.2
*(Well-typed terms have correct arities)* If $`\mathrm{\Gamma }M:K`$ is derivable in MLF, then $`arity(\mathrm{\Gamma })M:arity(K)`$ is derivable.
###### Proof
By induction on the derivations of $`\mathrm{\Gamma }M:K`$ (see the inference rules of MLF in Appendix).
###### Theorem 4.3
If $`\mathrm{\Gamma }M:K`$ is derivable in MLF, then $`M`$ is strongly normalising.
###### Proof
By Theorem 4.1 and Theorem 4.2.
## 5 Computation rules
In this section, we shall introduce computation rules for the type of dependent pairs and finite types and simple computation rules for universes. The strong normalisation is proved in a way that no one has ever take before in dependently typed systems, to the author’s best knowledge. Recall that adding new computation (or reduction) rules will not increase the set of terms with correct arities. The basic strategy we adopt is to prove strong normalisation one reduction rule after another. That is, if we have already proved strong normalisation for a set of reduction rules, after adding one new reduction rule, can we still prove strong normalisation? This strategy will not work for dependently typed systems if we want to prove the statement that “well-typed terms are strongly normalising”, because whenever we add a single computation rule, the set of well-typed terms may increase.
### 5.1 The type of dependent pairs
In MLF, the constants and computation rules for the type of dependent pairs can be specified as follows:
$`\mathrm{\Sigma }`$ $`:`$ $`(A:Type)(B:(A)Type)Type`$
$`pair`$ $`:`$ $`(A:Type)(B:(A)Type)(a:A)(b:B(a))\mathrm{\Sigma }(A,B)`$
$`\pi _1`$ $`:`$ $`(A:Type)(B:(A)Type)(z:\mathrm{\Sigma }(A,B))A`$
$`\pi _2`$ $`:`$ $`(A:Type)(B:(A)Type)(z:\mathrm{\Sigma }(A,B))B(\pi _1(A,B,z))`$
$`\pi _1(A,B,pair(A,B,a,b))`$ $`=`$ $`a:A`$
$`\pi _2(A,B,pair(A,B,a,b))`$ $`=`$ $`b:B(a)`$
In the arity system of the paper, we change the kinds to arities and the constants and the reduction rules are introduced as the following:
$`\mathrm{\Sigma }`$ $`:`$ $`(Zero,((Zero,Zero),Zero))`$
$`pair`$ $`:`$ $`(Zero,((Zero,Zero),(Zero,(Zero,Zero))))`$
$`\pi _1`$ $`:`$ $`(Zero,((Zero,Zero),(Zero,Zero)))`$
$`\pi _2`$ $`:`$ $`(Zero,((Zero,Zero),(Zero,Zero)))`$
$`\pi _1(A,B,pair(A,B,a,b))`$ $`_{\pi _1}`$ $`a:Zero`$
$`\pi _2(A,B,pair(A,B,a,b))`$ $`_{\pi _2}`$ $`b:Zero`$
### 5.2 Finite types
In type systems, a finite type $`𝒯`$ can be represented by following constants
$`𝒯`$ $`:`$ $`Type`$
$`c_1`$ $`:`$ $`𝒯`$
. . .
$`c_n`$ $`:`$ $`𝒯`$
$`_𝒯`$ $`:`$ $`(P:(𝒯)Type)`$
$`(P(c_1))\mathrm{}(P(c_n))`$
$`(z:𝒯)(P(z))`$
and the following computation rules
$`_𝒯(P,p_1,\mathrm{},p_n,c_1)`$ $`=`$ $`p_1:P(c_1)`$
$`\mathrm{}\mathrm{}`$
$`_𝒯(P,p_1,\mathrm{},p_n,c_n)`$ $`=`$ $`p_n:P(c_n)`$
In the arity system of the paper, we change the kinds to arities and the constants and the computation rules are introduced as follows.
$`𝒯`$ $`:`$ $`Zero`$
$`c_1`$ $`:`$ $`Zero`$
. . .
$`c_n`$ $`:`$ $`Zero`$
$`_𝒯`$ $`:`$ $`((Zero,Zero),`$
$`(Zero,(Zero,\mathrm{}(Zero,`$
$`(Zero,Zero)\mathrm{})`$
and the following reduction rules
$`_𝒯(P,p_1,\mathrm{},p_n,c_1)`$ $``$ $`p_1:Zero`$
$`\mathrm{}\mathrm{}`$
$`_𝒯(P,p_1,\mathrm{},p_n,c_n)`$ $``$ $`p_n:Zero`$
Now, let’s consider a concrete example, boolean type. Its representation in type systems and in the arity system are the following.
$`Bool`$ $`:`$ $`Type`$
$`true`$ $`:`$ $`Bool`$
$`false`$ $`:`$ $`Bool`$
$`_{Bool}`$ $`:`$ $`(P:(Bool)Type)`$
$`(p_1:P(true))(p_2:P(false))`$
$`(z:Bool)P(z)`$
$`_{Bool}(P,p_1,p_2,true)`$ $`=`$ $`p_1:P(true)`$
$`_{Bool}(P,p_1,p_2,false)`$ $`=`$ $`p_2:P(false)`$
$`Bool`$ $`:`$ $`Zero`$
$`true`$ $`:`$ $`Zero`$
$`false`$ $`:`$ $`Zero`$
$`_{Bool}`$ $`:`$ $`((Zero,Zero),`$
$`(Zero,(Zero,`$
$`(Zero,Zero))))`$
$`_{Bool}(P,p_1,p_2,true)`$ $`_{b_1}`$ $`p_1:Zero`$
$`_{Bool}(P,p_1,p_2,false)`$ $`_{b_2}`$ $`p_2:Zero`$
### 5.3 Universe operator
We consider some simple case, for example,
$`U`$ $`:`$ $`Type`$
$`Bool`$ $`:`$ $`Type`$
$`bool`$ $`:`$ $`U`$
$`uo`$ $`:`$ $`(U)Type`$
$`uo(bool)`$ $`=`$ $`Bool`$
$`U`$ $`:`$ $`Zero`$
$`Bool`$ $`:`$ $`Zero`$
$`bool`$ $`:`$ $`Zero`$
$`uo`$ $`:`$ $`(Zero,Zero)`$
$$uo(bool)_uBool:Zero$$
### 5.4 Strong normalisation w.r.t. $`\beta \eta \pi _1`$-reduction
We have proved strong normalisation w.r.t. $`\beta \eta `$-reduction in Section 4. Now, we add the reduction rule $`\pi _1`$ and prove strong normalisation w.r.t. $`\beta \eta \pi _1`$-reduction. As mentioned before, the strategy is to prove strong normalisation one reduction rule after another. So after proving it w.r.t. $`\beta \eta \pi _1`$-reduction, we can add another rule (*eg, $`\pi _2`$*-reduction), and so on. In this section, we demonstrate the proof techniques through the proof w.r.t. $`\beta \eta \pi _1`$-reduction. For other reduction rules such as $`\pi _2`$, $`b_1`$, $`b_2`$ and $`u`$, the proof methods are the same.
###### Theorem 5.1
If $`M`$ doesn’t have a correct arity under a context $`A`$ without the $`\pi _1`$-reduction then $`M`$ still doesn’t have a correct arity under the context $`A`$ with the $`\pi _1`$-reduction.
###### Proof
The arities of the left hand side and the right hand side of the reduction rule $`\pi _1`$ are the same, and there is no reduction for arities. So, $`\pi _1`$-reduction becomes irrelevant whether $`M`$ has a correct arity.
###### Remark 7
As mentioned before, in dependently typed systems, a term that is not well-typed can become a well-typed term after adding new reduction rules. For instance, under a context $`f:(x:B(a))C`$ and $`y:B(\pi _1(pair(a,b)))`$, the term $`f(y)`$ is not well-typed (some details are omitted here). However, if we add the $`\pi _1`$-reduction rule, then it becomes a well-typed term. This example shows that, after adding new reduction rules, well-typed terms may increase. This is one of the difficulties to prove the statement that “well-typed terms are strongly normalising”.
Now, in order to prove strong normalisation, we prove some lemmas first.
###### Lemma 4
(Substitution for $`\mathrm{\eta }`$) If $`\mathrm{M}_1_\mathrm{\eta }\mathrm{M}_2`$ then $`[\mathrm{N}/\mathrm{x}]\mathrm{M}_1_\mathrm{\eta }[\mathrm{N}/\mathrm{x}]\mathrm{M}_2`$. And if $`\mathrm{N}_1_\mathrm{\eta }\mathrm{N}_2`$ then $`[\mathrm{N}_1/\mathrm{x}]\mathrm{M}_\mathrm{\eta }[\mathrm{N}_2/\mathrm{x}]\mathrm{M}`$.
###### Proof
For the first part, we proceed the proof by induction on $`M_1`$, and for the second part, by induction on $`M`$. In the case that $`M`$ is a variable, we consider two sub-cases: $`Mx`$ and $`Mx`$.
###### Lemma 5
If $`M_1_\beta M_2`$ and $`xFV(M_1)`$ then $`xFV(M_2)`$.
###### Proof
By induction on $`M_1`$.
###### Lemma 6
If $`M_1_\eta \lambda x:K_2.M_2`$ then there are three and only three possibilities as the following:
* $`M_1\lambda y:K_1.(\lambda x:K_2.M_2)y`$ for some $`y`$ and $`K_1`$, and $`yFV(\lambda x:K_2.M_2)`$.
* $`M_1\lambda x:K_2.N`$ for some $`N`$ and $`N_\eta M_2`$.
* $`M_1\lambda x:K_1.M_2`$ for some $`K_1`$ and $`K_1_\eta K_2`$.
###### Proof
By the understanding of one-step reduction.
###### Lemma 7
(Commutation for $`\mathrm{\eta }\mathrm{\beta }`$) If $`\mathrm{M}_1_\mathrm{\eta }\mathrm{M}_2`$ and $`\mathrm{M}_2_\mathrm{\beta }\mathrm{M}_3`$ then there exists a $`\mathrm{M}_2^{}`$ such that $`\mathrm{M}_1_\mathrm{\beta }^+\mathrm{M}_2^{}`$ and $`\mathrm{M}_2^{}_\mathrm{\eta }\mathrm{M}_3`$.
###### Proof
By induction on $`M_1`$ and Lemma 4, 5 and 6.
###### Lemma 8
(Substitution for $`\mathrm{\pi }_1`$) If $`\mathrm{M}_1_{\mathrm{\pi }_1}\mathrm{M}_2`$ then $`[\mathrm{N}/\mathrm{x}]\mathrm{M}_1_{\mathrm{\pi }_1}[\mathrm{N}/\mathrm{x}]\mathrm{M}_2`$. And if $`\mathrm{N}_1_{\mathrm{\pi }_1}\mathrm{N}_2`$ then $`[\mathrm{N}_1/\mathrm{x}]\mathrm{M}_{\mathrm{\pi }_1}[\mathrm{N}_2/\mathrm{x}]\mathrm{M}`$.
###### Proof
Similar to the proof of Lemma 4.
###### Lemma 9
If $`M_1_{\pi _1}\lambda x:K_2.M_2`$ then there are two and only two possibilities as the following:
* $`M_1\lambda x:K_2.N`$ for some $`N`$ and $`N_{\pi _1}M_2`$.
* $`M_1\lambda x:K_1.M_2`$ for some $`K_1`$ and $`K_1_{\pi _1}K_2`$.
###### Proof
By the understanding of one-step reduction and the arity of $`M_1`$ is not $`Zero`$.
###### Lemma 10
(Commutation for $`\mathrm{\pi }_1\mathrm{\beta }`$) If $`\mathrm{M}_1_{\mathrm{\pi }_1}\mathrm{M}_2`$ and $`\mathrm{M}_2_\mathrm{\beta }\mathrm{M}_3`$ then there exists a $`\mathrm{M}_2^{}`$ such that $`\mathrm{M}_1_\mathrm{\beta }\mathrm{M}_2^{}`$ and $`\mathrm{M}_2^{}_{\mathrm{\pi }_1}\mathrm{M}_3`$.
###### Proof
By induction on $`M_1`$ and Lemma 8 and 9.
###### Theorem 5.2
If $`AM:a`$, then $`M`$ is strongly normalising w.r.t. $`\beta \eta \pi _1`$-reduction.
###### Proof
We proceed the proof by contradiction, and by Theorem 4.1 and Lemma 7 and 10.
Suppose there is an infinite reduction sequence for $`M`$ and it is called $`S`$. By Theorem 4.1, $`M`$ is strongly normalising w.r.t. $`\beta \eta `$-reduction. So, $`S`$ must contain infinite times of $`\pi _1`$-reduction. Every time when $`\eta `$-reduction or $`\pi _1`$-reduction rule is applied, terms become smaller. So, $`M`$ is strongly normalising w.r.t. $`\eta \pi _1`$-reduction. And hence $`S`$ must also contain infinite times of $`\beta `$-reduction. In fact, $`S`$ must be like the following,
$$M_{\eta \pi _1}^+M_1_\beta ^+M_2_{\eta \pi _1}^+M_3_\beta ^+M_4_{\eta \pi _1}^+\mathrm{}$$
or
$$M_\beta ^+M_1_{\eta \pi _1}^+M_2_\beta ^+M_3_{\eta \pi _1}^+M_4_\beta ^+\mathrm{}$$
where $`_\beta ^+`$ means one or more but finite reduction steps of $`\beta `$, and similarly, $`_{\eta \pi _1}^+`$ means one or more but finite reduction steps of $`\eta `$ or $`\pi _1`$.
Now, by Lemma 7 and Lemma 10, for the infinite sequence $`S`$, we can always move the $`\beta `$-reduction steps forward and build an infinite sequence of $`\beta `$-reduction. This is a contradiction to that $`M`$ is strongly normalising w.r.t. $`\beta `$-reduction.∎
## 6 Conclusions and future work
Strong normalisation for MLF has been proved in the paper, but we did not follow the traditional understanding, that is, well-typed terms are strongly normalising. Instead, a weaker condition has been proposed, which says terms with correct arities are strongly normalising. The author hopes this new understanding will inspire us to think the question “why is a term strongly normalising?” again, and to simplify the proofs for dependently typed systems.
Another important technique employed in the paper is that, in order to prove what we want, we prove a more general and stronger property. In the paper, the definition of terms and kinds is extended and a new reduction rule $`\beta _2`$ is introduced. And we proved strong normalisation w.r.t. $`\beta \eta \beta _2`$-reduction instead of w.r.t. $`\beta \eta `$-reduction only. This generalisation is quite different from the traditional idea of generalising induction hypothesis.
We only studied the computation rules for some inductive data types and these rules have commutation property. However, some computation rules do not have such property, for instance, the computation rule for the type of function space. How to prove strong normalisation for such rules needs further study. The question of how to develop weaker conditions to simplify the normalisation proofs for other type systems is also worth being taken into our consideration.
#### Acknowledgements
Thanks to Zhaohui Luo, Sergei Soloviev, James McKinna and Healfdene Goguen for discussions on the issue of strong normalisation, and for reading the earlier version of the paper, and for their helpful comments and suggestions.
## Appendix
Terms and Kinds in MLF
* Terms
1. a variable is a term,
2. $`\lambda x:K.M`$ is a term if $`x`$ is a variable, $`K`$ is a kind and $`M`$ is a term,
3. $`MN`$ is a term if $`M`$ and $`N`$ are terms.
* Kinds
1. $`Type`$ is a kind,
2. $`El(M)`$ is a kind if $`M`$ is a term,
3. $`(x:K_1)K_2`$ is a kind if $`K_1`$ and $`K_2`$ are kinds.
Reduction rules in MLF
$$(\lambda x:K.M)N_\beta [N/x]M$$
$$\lambda x:K.Mx_\eta MxFV(M)$$
Inference rules for MLF
Contexts and assumptions
$$\frac{}{<>valid}\frac{\mathrm{\Gamma }KkindxFV(\mathrm{\Gamma })}{\mathrm{\Gamma },x:Kvalid}\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}valid}{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}x:K}$$
Equality rules
$$\frac{\mathrm{\Gamma }Kkind}{\mathrm{\Gamma }K=K}\frac{\mathrm{\Gamma }K=K^{}}{\mathrm{\Gamma }K^{}=K}\frac{\mathrm{\Gamma }K=K^{}\mathrm{\Gamma }K^{}=K^{\prime \prime }}{\mathrm{\Gamma }K=K^{\prime \prime }}$$
$$\frac{\mathrm{\Gamma }k:K}{\mathrm{\Gamma }k=k:K}\frac{\mathrm{\Gamma }k=k^{}:K}{\mathrm{\Gamma }k^{}=k:K}\frac{\mathrm{\Gamma }k=k^{}:K\mathrm{\Gamma }k^{}=k^{\prime \prime }:K}{\mathrm{\Gamma }k=k^{\prime \prime }:K}$$
$$\frac{\mathrm{\Gamma }k:K\mathrm{\Gamma }K=K^{}}{\mathrm{\Gamma }k:K^{}}\frac{\mathrm{\Gamma }k=k^{}:K\mathrm{\Gamma }K=K^{}}{\mathrm{\Gamma }k=k^{}:K^{}}$$
Substitution rules
$$\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}valid\mathrm{\Gamma }k:K}{\mathrm{\Gamma },[k/x]\mathrm{\Gamma }^{}valid}$$
$$\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}K^{}kind\mathrm{\Gamma }k:K}{\mathrm{\Gamma },[k/x]\mathrm{\Gamma }^{}[k/x]K^{}kind}\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }K^{}kind\mathrm{\Gamma }k=k^{}:K}{\mathrm{\Gamma },[k/x]\mathrm{\Gamma }^{}[k/x]K^{}=[k^{}/x]K^{}}$$
$$\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}k^{}:K^{}\mathrm{\Gamma }k:K}{\mathrm{\Gamma },[k/x]\mathrm{\Gamma }^{}[k/x]k^{}:[k/x]K^{}}\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}k^{}:K^{}\mathrm{\Gamma }k_1=k_2:K}{\mathrm{\Gamma },[k_1/x]\mathrm{\Gamma }^{}[k_1/x]k^{}=[k_2/x]:[k_1/x]K^{}}$$
$$\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}K^{}=K^{\prime \prime }\mathrm{\Gamma }k:K}{\mathrm{\Gamma },[k/x]\mathrm{\Gamma }^{}[k/x]K^{}=[k/x]K^{\prime \prime }}\frac{\mathrm{\Gamma },x:K,\mathrm{\Gamma }^{}k^{}=k^{\prime \prime }:K^{}\mathrm{\Gamma }k:K}{\mathrm{\Gamma },[k/x]\mathrm{\Gamma }^{}[k/x]k^{}=[k/x]k^{\prime \prime }:[k/x]K^{}}$$
The kind type
$$\frac{\mathrm{\Gamma }valid}{\mathrm{\Gamma }Typekind}\frac{\mathrm{\Gamma }A:Type}{\mathrm{\Gamma }El(A)kind}\frac{\mathrm{\Gamma }A=B:Type}{\mathrm{\Gamma }El(A)=El(B)}$$
Dependent product kinds
$$\frac{\mathrm{\Gamma }Kkind\mathrm{\Gamma },x:KK^{}kind}{\mathrm{\Gamma }(x:K)K^{}kind}\frac{\mathrm{\Gamma }K_1=K_2\mathrm{\Gamma },x:K_1K_1^{}=K_2^{}}{\mathrm{\Gamma }(x:K_1)K_1^{}=(x:K_2)K_2^{}}$$
$$\frac{\mathrm{\Gamma },x:Kk:K^{}}{\mathrm{\Gamma }\lambda x:K.k:(x:K)K^{}}(\xi )\frac{\mathrm{\Gamma }K_1=K_2\mathrm{\Gamma },x:K_1k_1=k_2:K}{\mathrm{\Gamma }\lambda x:K_1.k_1=\lambda x:K_2.k_2:(x:K_1)K}$$
$$\frac{\mathrm{\Gamma }f:(x:K)K^{}\mathrm{\Gamma }k:K}{\mathrm{\Gamma }f(k):[k/x]K^{}}\frac{\mathrm{\Gamma }f=f^{}:(x:K)K^{}\mathrm{\Gamma }k_1=k_2:K}{\mathrm{\Gamma }f(k_1)=f^{}(k_2):[k_1/x]K^{}}$$
$$(\beta )\frac{\mathrm{\Gamma },x:Kk^{}:K^{}\mathrm{\Gamma }k:K}{\mathrm{\Gamma }(\lambda x:K.k^{})(k)=[k/x]k^{}:[k/x]K^{}}(\eta )\frac{\mathrm{\Gamma }f:(x:K)K^{}xFV(f)}{\mathrm{\Gamma }\lambda x:K.f(x)=f:(x:K)K^{}}$$
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# The cotangent space at a monomial ideal of the Hilbert scheme of points of an affine space
## 1. Introduction
### 1.1. Summary of results
Let $`𝗄`$ be an algebraically closed field of any characteristic, and
$$𝔸_𝗄^r=\mathrm{Spec}(𝗄[x_1,\mathrm{},x_r])=\mathrm{Spec}(𝗄[𝐱])$$
the affine space of dimension $`r`$ over $`𝗄`$. The Hilbert scheme
$$\mathrm{Hilb}_{𝔸_𝗄^r}^n=\text{H}^n$$
parameterizes the 0-dimensional closed subschemes
$$\mathrm{Spec}(𝗄[𝐱]/I)𝔸_𝗄^r$$
having length $`n`$, that is,
$$dim_𝗄(𝗄[𝐱]/I)=\text{colength of }I=n.$$
In this paper we study the cotangent space of a point of $`\text{H}^n`$ that corresponds to a monomial ideal $`I`$ (that is, $`I`$ is generated by monomials). If $`I`$ is a monomial ideal of colength $`n`$, and we let
$$\beta =\{\text{monomials }mmI\},$$
then it is clear that $`\beta `$ is a $`𝗄`$-basis of the quotient $`𝗄[𝐱]/I`$; furthermore, $`\beta `$ has the property that for monomials $`m_1`$, $`m_2`$,
$$m_1\beta \text{ and }m_2|m_1m_2\beta ;$$
we shall call any set of monomials $`\beta `$ with this property a basis set of monomials. Let $`U_\beta `$ denote the (affine) open subscheme of $`\text{H}^n`$ whose $`𝗄`$-points $`t`$ are associated to ideals $`I_t`$ such that $`\beta `$ is a $`𝗄`$-basis of the quotient $`𝗄[𝐱]/I_t`$. In particular, the monomial ideal $`I`$ $`=`$ $`I_\beta `$ that we started with has this property, so
$$t_\beta U_\beta ,\text{ where }t_\beta \text{ is the point corresponding to }I_\beta =I_{t_\beta };$$
we can therefore identify the cotangent space of $`t_\beta `$ with the $`𝗄`$-vector space $`M/M^2`$, where $`M`$ is the maximal ideal of $`t_\beta `$ in the coordinate ring $`R`$ of $`U_\beta `$. Since $`t_\beta `$ lies in the closure of the locus corresponding to closed subschemes supported at $`n`$ distinct points of $`𝔸_𝗄^r`$, which is an $`(rn)`$-dimensional component of $`\text{H}^n`$, we have that (Proposition 2.4.1)
(1)
dimk(/MM2) rn, and
tβ is nonsingular
dimk(/MM2) = rn.
dimk(/MM2) rn, and
tβ is nonsingular
dimk(/MM2) = rn.
missing-subexpression\left.\begin{array}[]{cc}{}\parbox{289.07999pt}{
\begin{itemize}
\par\itemize@item@$\dim_{\mathsf{k}}(M/M^{2})$ $\geq$ $rn$, and
\par\itemize@item@$t_{\beta}$ is nonsingular $\Leftrightarrow$
$\dim_{\mathsf{k}}(M/M^{2})$ $=$ $rn$.
\par\end{itemize}
}&{}\hfil\end{array}\right.
Fortunately, we have a concrete description of $`R`$, based on the observation that for every point $`t`$ $``$ $`U_\beta `$ and every monomial
$$x_1^{d_1}x_2^{d_2}\mathrm{}x_r^{d_r}=𝐱^𝐝I_\beta ,$$
there is a unique polynomial
$$F_𝐝(t)=𝐱^𝐝\underset{𝐱^𝐣\beta }{}c_𝐣^𝐝(t)𝐱^𝐣I_t,c_𝐣^𝐝(t)𝗄,$$
since the quotient $`𝗄[𝐱]/I_t`$ has $`𝗄`$-basis $`\beta `$. As $`t`$ ranges over $`U_\beta `$, the coefficients $`c_𝐣^𝐝(t)`$ define functions $`c_𝐣^𝐝`$ $``$ $`R`$ that generate $`R`$ as a $`𝗄`$-algebra; moreover, the relations among these functions can be completely described. Note in particular that the point $`t_\beta `$ is the “origin” of $`U_\beta `$, the point at which all the functions $`c_𝐣^𝐝`$ vanish.
The maximal ideal $`M`$ $``$ $`R`$ that cuts out the point $`t_\beta `$ is therefore generated by the functions $`c_𝐣^𝐝`$. Each of these functions can be viewed as an arrow pointing from $`𝐱^𝐝`$ to $`𝐱^𝐣`$ in the lattice of monomials obtained by identifying each monomial $`𝐱^𝐛`$ with its $`r`$-tuple $`𝐛`$ of exponents (for example, see Figures 1 and 2). It then turns out that translation of arrows (keeping the head inside $`\beta `$ and the tail outside $`\beta `$) corresponds to equivalence modulo $`M^2`$; more precisely (see Section 3),
* If $`c_{𝐣_1}^{𝐝_1}`$ can be translated to $`c_{𝐣_2}^{𝐝_2}`$, then $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_{𝐣_2}^{𝐝_2}(modM^2)`$.
* If $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`0(modM^2)`$, then $`c_{𝐣_1}^{𝐝_1}`$ can be translated so that its head exits $`\beta `$ across a hyperplane ($`x_i`$-degree $`=`$ $`0`$), and conversely.
* If $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_{𝐣_2}^{𝐝_2}`$ $``$ $`0(modM^2)`$, then $`c_{𝐣_1}^{𝐝_1}`$ can be translated to $`c_{𝐣_2}^{𝐝_2}`$.
Our first main result, Theorem 3.2.1, states that if $`S`$ is a finite set of $`d`$ arrows such that no member of $`S`$ can be translated so that its head exits $`\beta `$ across a hyperplane ($`x_i`$-degree $`=`$ $`0`$), and no two distinct members of $`S`$ can be translated one to the other, then the $`𝗄`$-span of $`S`$ in $`M/M^2`$ has dimension $`d`$, and conversely; in this case, we say that $`S`$ has maximal rank (mod $`M^2`$). Using this criterion, we construct in Theorem 4.1.3 a maximal rank (mod $`M^2`$) set $`𝒮`$ of $`rn`$ arrows $`c_𝐣^𝐝`$ of a special kind that we call minimal standard arrows: *minimal* means that the tail $`𝐱^𝐝`$ is a minimal generator of the monomial ideal $`I_\beta `$, and *standard* means that the vector $`𝐣`$ $``$ $`𝐝`$ has only one negative component. (Since $`𝐱^𝐣`$ $``$ $`\beta `$ and $`𝐱^𝐝`$ $``$ $`\beta `$, the vector of any arrow must always have at least one negative component, to avoid the contradiction $`𝐱^𝐝`$ $`|`$ $`𝐱^𝐣`$.) Therefore, in view of (1),
$$t_\beta \text{ is nonsingular }𝒮\text{ spans }M/M^2𝒮\text{ is a }𝗄\text{-basis of }M/M^2.$$
When any (and hence all) of these equivalent conditions holds, we say that $`\beta `$ is a smooth basis set. (The set $`𝒮`$ is in general not unique.)
###### Remark 1.1.1.
The idea of arrow translation and its connection to equivalence modulo $`M^2`$ is due to Haiman (in the case of two variables); it is a pleasure to acknowledge the inspiration of that beautiful article.
The remainder of the paper explores conditions on $`\beta `$ sufficient to imply that $`\beta `$ either is or is not smooth. The basic result is Theorem 5.1.1, which in essence says that
(2)
$$\beta \text{ smooth }\{\begin{array}{c}c_𝐣^𝐝\text{ non-standard }c_𝐣^𝐝0(modM^2),\text{ and}\hfill \\ c_𝐣^{}^𝐝^{}\text{ standard and }0(modM^2)c_𝐣^{}^𝐝^{}\hfill \\ \text{can be translated to an arrow in }𝒮.\hfill \end{array}$$
In Section 5.2, we present a simple condition on $`\beta `$ that ensures the existence of non-standard arrows that are *not* $``$ 0 modulo $`M^2`$. One first checks that any minimal arrow $`c_𝐣^𝐝`$ with head $`𝐱^𝐣`$ *maximal* in $`\beta `$ (for the partial ordering defined by divisibility of monomials) cannot be translated except to itself (Lemma 5.2.1); we call such an arrow rigid. It is then clear that any non-standard rigid arrow will violate the first bulleted condition in (2), forcing $`\beta `$ to be non-smooth (Corollary 5.2.2). Figure 1 illustrates such an arrow; a specific example is discussed in Section 5.3.
We go on to present three ways to construct smooth basis sets:
In Section 6, we show that a basis set $`\beta `$ in $`r`$ variables that is the “thickening” (18) of a smooth basis set $`\beta _0`$ in $`r1`$ variables will be smooth (Theorem 6.4.1). As a consequence, we obtain that every box (a basis set that has the shape of a rectangular parallelepiped) is smooth (Proposition 6.5.3).
In Section 9, we show that a basis set $`\beta `$ that is obtained by adding a box to a smooth basis set in an appropriate way (see Subsection 9.1) is smooth (Theorem 9.3.1). Therefore, if we start with a box, and add a finite sequence of boxes to obtain a compound box (see Figure 6), we obtain a smooth basis set (Corollary 9.4.1). In particular, *every* basis set in two variables is a compound box (see Figure 7) and therefore smooth, so we have obtained a variant proof of Haiman’s result \[6, Prop. 2.4, p. 209\], using a lemma valid in arbitrary dimension. In addition, we prove in Section 10 that every smooth basis set in three variables is a compound box (Theorem 10.3.1).
In Section 11, we show that a basis set that is the union of two boxes is a smooth basis set (Theorem 11.3.2). We present this result to show that the characterization of smooth basis sets in three variables (as compound boxes) does not carry over to higher dimensions, because for $`r`$ $``$ $`4`$ there are two-box unions that are not compound boxes (Example 11.4).
Sections 7 and 8 are devoted to technical preparations needed for Sections 9 and 10. Section 7 introduces the operation of *truncation* of a basis set (see Figure 4), and Section 8 establishes conditions under which a truncation of a smooth basis set will be smooth.
### 1.2. Table of contents
1. Introduction.
2. 1. Summary of results.
2. Table of contents.
3. The Hilbert scheme of points $`\text{H}^n`$ and its open subschemes $`U_\beta `$.
4. 1. The Hilbert scheme $`\text{H}^n`$.
2. The open subschemes $`U_\beta `$.
3. The coordinate ring of $`U_\beta `$.
4. The locus of reduced subschemes meets $`U_\beta `$.
5. The cotangent space of the point $`t_\beta `$ $``$ $`\text{H}^n`$.
6. 1. Arrows, translation, and congruence modulo $`M^2`$.
2. $`𝗄`$-Linearly independent sets of arrows in $`M/M^2`$.
7. A linearly independent set in $`M/M^2`$ of cardinality $`rn`$.
8. 1. Standard and minimal arrows.
2. Advancement of minimal standard arrows.
3. Shadow promotion.
4. Iterated shadow promotion.
5. Proof of Theorem 4.1.3
6. Summary and terminology.
9. Consequences of Theorem 4.1.3.
10. 1. Necessary and sufficient conditions for $`t_\beta `$ $``$ $`\text{H}^n`$ to be nonsingular.
2. A sufficient condition for $`t_\beta `$ $``$ $`\text{H}^n`$ to be singular.
3. Example: $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_3\}`$.
11. Thickening of basis sets.
12. 1. Definition of thickening.
2. Minimal generators of $`I_\beta `$.
3. A standard bunch $`𝒮`$ for $`\beta `$.
4. Thickenings of smooth basis sets are smooth.
5. Example: “Boxes” are smooth basis sets.
13. Truncation of basis sets.
14. 1. Definition of truncation.
2. Minimal generators of $`I_{\beta _t}`$.
3. Lifting arrows from $`\beta _t`$ to $`\beta `$.
4. Non-standard arrows on $`\beta `$ and $`\beta _t`$.
15. Sufficient conditions for a truncation to be smooth.
16. 1. The additional hypothesis.
2. $`x_j`$-sub-bunches of arrows for $`\beta `$ and $`\beta _t`$.
3. $`x_k`$-standard arrows of $`x_j`$-height $``$ $`h`$.
4. Linear independence of lifts of $`x_k`$-sub-bunches.
5. $`x_k`$-sub-bunches of arrows for $`\beta `$ and $`\beta _t`$, $`x_k`$ $``$ $`x_j`$.
6. Main theorem on truncations.
17. Addition of boxes to basis sets.
18. 1. Definition of box addition.
2. Minimal generators of $`I_\beta ^{}`$.
3. Main theorem on box additions.
4. Compound boxes.
5. Example: Basis sets in two variables.
6. Example: The lexicographic point.
7. Example: $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_1x_2,x_3\}`$.
19. Smooth basis sets in three variables are compound boxes.
20. 1. The main lemma.
2. Proof of Lemma 10.1.1.
3. The main theorem.
21. The union of two boxes.
22. 1. Notation.
2. Minimal generators of $`I_\beta `$.
3. Two-box unions are smooth basis sets.
4. Example: $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_1x_2,x_3,x_4,x_3x_4\}`$.
## 2. The Hilbert scheme of points $`\text{H}^n`$ and its open subschemes $`U_\beta `$
In this section we briefly recall the definition and some properties of the Hilbert scheme of points $`\text{H}^n`$, and the open subschemes $`U_\beta `$ that form an open covering of $`\text{H}^n`$. Much of our subsequent work is based on the explicit representation of the coordinate ring of $`U_\beta `$ that is obtained in .
### 2.1. The Hilbert scheme $`\text{H}^n`$
The scheme $`\text{H}^n`$ can be defined functorially as the parameter scheme of a universal flat and proper family $`Z_n`$ of zero-dimensional closed subschemes of $`𝔸_𝗄^r`$ having length $`n`$. In other words, $`Z_n`$ $``$ $`\text{H}^n\times 𝔸_𝗄^r`$ is a closed subscheme that is finite and flat of degree $`n`$ over $`\text{H}^n`$ and satisfies the following universal property:
(3) Let $`T`$ be a scheme over $`𝗄`$. Then the set of maps $`f:T\text{H}^n`$ is in natural bijective correspondence with the set of closed subschemes $`Z_fT\times 𝔸_𝗄^r`$ that are finite and flat of degree $`n`$ over $`T`$; the bijection $`fZ_f`$ is defined by $`Z_f`$ = $`T\times _{\text{H}^n}Z_n`$.
In particular, the inclusion of the $`𝗄`$-point $`t\text{H}^n`$ corresponds to a unique closed subscheme
$$Z_tt\times 𝔸_𝗄^r𝔸_𝗄^r;$$
the map $`tZ_t`$ defines a bijection from the set of $`𝗄`$-points of $`\text{H}^n`$ to the set of 0-dimensional closed subschemes of length $`n`$, or, equivalently, to the set of ideals
$$I𝗄[x_1,\mathrm{},x_r]=𝗄[𝐱]$$
of colength $`n`$. We write $`I_t`$ to denote the ideal corresponding to the $`𝗄`$-point $`t`$ $``$ $`\text{H}^n`$.
The existence of $`\text{H}^n`$ can be established as follows: It is the open subscheme of $`\mathrm{Hilb}_{_𝗄^r}^n`$ (a projective scheme the existence of which is a consequence of Grothendieck’s general construction given in ) arising from the inclusion of $`𝔸_𝗄^r`$ in $`_𝗄^r`$ as a standard affine. Alternatively, $`\text{H}^n`$ is a special case of the multigraded Hilbert scheme constructed by Haiman and Sturmfels . Elementary constructions of $`\text{H}^n`$ are given in and .
### 2.2. The open subschemes $`U_\beta `$
Let $`\beta `$ denote a nonempty set of $`n`$ monomials in the indeterminates $`x_1,\mathrm{},x_r`$ that satisfies the following property: if $`m_1`$ $``$ $`\beta `$ and $`m_2`$ is a monomial dividing $`m_1`$, then $`m_2`$ $``$ $`\beta `$. We will call such a set $`\beta `$ a basis set, and its members basis monomials. As a set of $`𝗄`$-points, we define $`U_\beta `$ as follows:
(4)
$$U_\beta =\{t\text{H}^n𝗄[𝐱]/I_t\text{ has }𝗄\text{-basis }\beta \}\text{H}^n.$$
Every point $`t`$ $``$ $`\text{H}^n`$ belongs to at least one set of the form $`U_\beta `$, a fact that is, in M. Haiman’s words, “often regarded as a part of Gröbner basis theory but actually goes back to Gordan \[6, p. 207\]. Therefore, the sets $`U_\beta `$ form an open covering of $`\text{H}^n`$. In fact, the $`U_\beta `$ form an open *affine* covering of $`\text{H}^n`$; this is proved in , where the coordinate ring of $`U_\beta `$ is explicitly obtained as a quotient of a polynomial ring. Note that Haiman introduced the open subschemes $`U_\beta `$ in for the case of the affine plane $`(r=2)`$; he showed that they are affine using other means.
###### Remark 2.2.1.
Let $`Z_\beta `$ denote the restriction of the universal closed subscheme $`Z_n`$ to $`U_\beta `$ $``$ $`\text{H}^n`$. Then the direct image of the structure sheaf of $`Z_\beta `$ on $`U_\beta `$ is free with basis $`\beta `$, because $`\beta `$ is everywhere a local basis of the (locally free) direct image. This also follows directly from the construction of $`(U_\beta ,Z_\beta )`$ given in \[9, Sec. 7.3\].
Note that the monomials that do *not* belong to a basis set $`\beta `$ generate a monomial ideal
$$I=I_\beta 𝗄[𝐱];$$
it is immediate that the quotient $`𝗄[𝐱]/I_\beta `$ has $`\beta `$ as $`𝗄`$-basis, that is, $`I_\beta `$ corresponds to a point
(5)
$$t=t_\beta U_\beta .$$
Conversely, if $`I`$ is a monomial ideal such that $`𝗄[𝐱]/I`$ has $`𝗄`$-dimension $`n`$, then the quotient has as $`𝗄`$-basis the set of monomials not in $`I`$, and this is clearly a basis set $`\beta `$ of $`n`$ members.
### 2.3. The coordinate ring of $`U_\beta `$
Let $`R`$ denote the coordinate ring of the affine scheme $`U_\beta `$, and let $`I_{U_\beta }`$ $``$ $`R[𝐱]`$ be the ideal that cuts out the universal closed subscheme $`Z_\beta `$ $``$ $`U_\beta \times 𝔸_𝗄^r`$. By Remark 2.2.1, we have that the module $`R[𝐱]/I_{U_\beta }`$ is $`R`$-free with basis $`\beta `$. Therefore, every monomial
$$x_1^{d_1}x_2^{d_2}\mathrm{}x_r^{d_r}=𝐱^𝐝R[𝐱]$$
is congruent modulo $`I_{U_\beta }`$ to a unique $`R`$-linear combination of the monomials $`𝐱^𝐣`$ $``$ $`\beta `$, or, in other words, for every $`𝐱^𝐝`$ there is a unique polynomial of the form
(6)
$$F_𝐝=𝐱^𝐝\underset{𝐣\beta }{}c_𝐣^𝐝𝐱^𝐣I_{U_\beta },c_𝐣^𝐝R,$$
where we abuse the notation (here and elsewhere) by writing $`𝐣\beta `$ for $`𝐱^𝐣\beta `$. When writing monomials in this way, we will reserve $`𝐣`$ for basis monomials. Note that
(7)
$$𝐝\beta \{\begin{array}{c}c_𝐣^𝐝=1\text{ if }𝐝=𝐣,\text{ and}\hfill \\ c_𝐣^𝐝=0\text{ if }𝐝𝐣.\hfill \end{array}$$
Following Haiman \[6, p. 210\], we multiply the polynomial (6) by the variable $`x_i`$, and then expand each monomial $`x_i𝐱^𝐣`$ using (6) to obtain a polynomial of the form
$$x_i𝐱^𝐝+(R\text{-linear combination of basis monomials})I_{U_\beta },$$
which must therefore be equal to $`F_𝐝^{}`$, where $`𝐱^𝐝^{}`$ = $`x_i𝐱^𝐝`$. Equating coefficients, we obtain the relations
(8)
$$c_{𝐣_0}^𝐝^{}\underset{𝐣\beta }{}c_𝐣^𝐝c_{𝐣_0}^𝐣^{}=0,$$
where $`𝐱^𝐣^{}`$ $`=`$ $`x_i𝐱^𝐣`$.
For each coefficient $`c_𝐣^𝐝`$ such that $`𝐝`$ $``$ $`\beta `$ we introduce an indeterminate $`C_𝐣^𝐝`$, and let
(9)
$$\delta :𝗄[(C_𝐣^𝐝)]R,C_𝐣^𝐝c_𝐣^𝐝,$$
be the natural map. We have the following
###### Proposition 2.3.1.
The coordinate ring $`R`$ of $`U_\beta `$ is generated as a $`𝗄`$-algebra by the coefficients $`c_𝐣^𝐝`$ such that $`𝐝`$ $``$ $`\beta `$; that is, the map $`\delta `$ is surjective. Furthermore, the kernel of $`\delta `$ is generated by the polynomials
$$\tau _{𝐣_0}^{(𝐝,x_i)}=C_{𝐣_0}^𝐝^{}\underset{𝐣\beta }{}C_𝐣^𝐝C_{𝐣_0}^𝐣^{},𝐝\beta ,𝐣_0\beta ,1ir,$$
that are obtained from (8) by replacing each coefficient $`c_𝐣^𝐛`$ with the corresponding indeterminate $`C_𝐣^𝐛`$, if $`𝐛`$ $``$ $`\beta `$, or the appropriate constant — $`0`$ or $`1`$, according to *(7)*, and denoted $`C_𝐣^𝐛`$ — if $`𝐛`$ $``$ $`\beta `$.
*Proof:* The construction of $`R`$ given in \[9, Sec. 7\] shows that $`R`$ is generated by a certain finite subset of the coefficients $`c_𝐣^𝐝`$, and the kernel of the map $`\delta `$ is generated by a finite set of polynomials (there denoted $`\rho _𝐣^{(𝐛_1,𝐛_2)}`$), each of which is in fact either equal to one of the polynomials $`\tau `$ or to a $`𝗄`$-linear combination of two of them. ∎
###### Remark 2.3.2.
It is clear that the point $`t_\beta `$ (5) is the “origin” of $`U_\beta `$; that is, $`t_\beta `$ is the point at which all the coefficient functions $`c_𝐣^𝐝`$ with $`𝐝`$ $``$ $`\beta `$ vanish.
### 2.4. The locus of reduced subschemes meets $`U_\beta `$
Let $`\text{H}_{}`$ $``$ $`\text{H}^n`$ denote the open subscheme parameterizing the closed subschemes of $`𝔸_𝗄^r`$ that are supported at $`n`$ distinct points (and are therefore reduced). Since $`\text{H}_{}`$ is irreducible and of dimension $`rn`$, its closure $`\overline{\text{H}}_{}`$ is a component of $`\text{H}^n`$ of dimension $`rn`$. It is well-known that the point $`t_\beta `$ corresponding to the monomial ideal $`I_\beta `$ $``$ $`𝗄[𝐱]`$ lies in $`\overline{\text{H}}_{}`$; in fact, one can exhibit a one-parameter family of length-$`n`$ closed subschemes of $`𝔸_𝗄^r`$ such that the fiber over 0 is the subscheme defined by $`I_\beta `$ and the general fiber is a reduced subscheme. This family is constructed using Hartshorne’s concept of distraction as adapted by Geramita et al. in . The construction works whenever the ground field is infinite (which is true for us, since $`𝗄`$ is assumed algebraically closed) or sufficiently large.
Briefly, the construction goes like this. First, choose an infinite (or sufficiently long) sequence of distinct elements $`a_0,a_1,a_2,\mathrm{}`$ of $`𝗄`$. Then, for every minimal generator
$$x_1^{d_1}x_2^{d_2}\mathrm{}x_r^{d_r}=𝐱^𝐝$$
of the monomial ideal $`I_\beta `$, form the homogeneous polynomial
$$f_𝐝=\underset{i=1}{\overset{r}{}}\left(\underset{j=0}{\overset{d_i1}{}}(x_ia_jw)\right)𝗄[w][𝐱].$$
One then proves that the homogeneous ideal $``$ generated by the polynomials $`f_𝐝`$ is the ideal of all functions vanishing on the $`n`$ (distinct projective) points
$$[1:a_{e_1}:a_{e_2}:\mathrm{}:a_{e_r}]_𝗄^r,x_1^{e_1}x_2^{e_2}\mathrm{}x_r^{e_r}=𝐱^𝐞\beta $$
(\[2, Theorem 2.2\]). In particular, if we view $`w`$ as a parameter, then for any nonzero value $`w_1`$ of $`w`$, the ideal $`(w_1)`$ $``$ $`𝗄[𝐱]`$ obtained by replacing the variable $`w`$ by the constant $`w_1`$ is supported at the $`n`$ distinct points
$$(w_1a_{e_1},w_1a_{e_2},\mathrm{},w_1a_{e_r}),𝐱^𝐞\beta $$
(when $`w_1`$ $`=`$ $`0`$, clearly $`(0)`$ $`=`$ $`I_\beta `$, the original monomial ideal of colength $`n`$). One then checks that the quotient $`𝗄[w][𝐱]/`$ is $`𝗄[w]`$-free with basis $`\beta `$ as a $`𝗄[w]`$-module. (Indeed, from the form of the polynomials $`f_𝐝`$, it is clear that every minimal generator of the monomial ideal $`I_\beta `$ is congruent (mod $``$) to a $`k[w]`$-linear combination of basis monomials *that divide the minimal generator*. Induction on the degree then shows that every monomial $`m`$ $``$ $`I_\beta `$ is congruent (mod $``$) to a $`k[w]`$-linear combination of basis monomials that divide $`m`$; therefore, the basis monomials span the quotient $`𝗄[w][𝐱]/`$. Since the $`𝗄`$-dimension of each quotient $`𝗄[𝐱]/(w_1)`$ is $``$ $`n`$, we find that $`\beta `$ is locally a basis everywhere, and hence a global basis, as desired.) In this way we obtain a finite and flat family of subschemes which defines a map $`\phi :\mathrm{Spec}(𝗄[t])U_\beta `$ $``$ $`\text{H}^n`$; one sees easily that $`\phi (0)`$ $`=`$ $`t_\beta `$ and $`\phi (w_1)`$ $``$ $`\text{H}_{}`$ $``$ $`U_\beta `$ for all $`w_1`$ $``$ $`0`$. It follows at once that $`t_\beta `$ $``$ $`\overline{\text{H}}_{}`$, which yields
###### Proposition 2.4.1.
The $`𝗄`$-dimension of the (Zariski) tangent and cotangent spaces of the point $`t_\beta `$ is $``$ $`dim(\overline{\text{H}}_{})`$ $`=`$ $`rn`$, and $`t_\beta `$ is a smooth point of $`\text{H}^n`$ if and only if
$$dim_𝗄(\text{cotangent space of }t_\beta )=rn.\text{}$$
###### Remark 2.4.2.
We will give another proof that the $`𝗄`$-dimension of the cotangent space of $`t_\beta `$ is at least $`rn`$ in Section 4.
###### Remark 2.4.3.
The foregoing demonstrates that each open subscheme $`U_\beta `$ meets the locus $`\text{H}_{}`$ of reduced subschemes nontrivially. This is shown in a different way in \[9, Sec. 2.4\].
## 3. The cotangent space of the point $`t_\beta `$ $``$ $`\text{H}^n`$
Recall that the cotangent space of the point $`t_\beta `$ $``$ $`\text{H}^n`$ is the $`𝗄`$-vector space $`𝔪_{t_\beta }/𝔪_{t_\beta }^2`$, where $`𝔪_{t_\beta }`$ is the maximal ideal of the local ring $`𝒪_{t_\beta }`$. In light of Remark 2.3.2, we see that we can identify the cotangent space with the $`𝗄`$-vector space $`M/M^2`$, where $`M`$ is the maximal ideal generated by the functions $`c_𝐣^𝐝`$ ($`𝐝`$ $``$ $`\beta `$) in the coordinate ring $`R`$ of $`U_\beta `$. In , Haiman observed, in the two-variable case, that congruence (mod $`M^2`$) can be visualized as translation-equivalence of arrows defined by the $`c_𝐣^𝐝`$. We now extend Haiman’s idea to any number $`r`$ of variables.
### 3.1. Arrows, translation, and congruence modulo $`M^2`$
As stated by Haiman, with symbols adjusted, we have that \[6, p. 210\]
> \[m\]odulo $`M^2`$, the terms -$`c_𝐣^𝐝c_{𝐣_0}^𝐣^{}`$ \[in\] equation (8) reduce to zero for $`𝐣^{}`$ $``$ $`\beta `$ and for $`𝐣^{}`$ $``$ $`\beta `$, $`𝐣^{}`$ $``$ $`𝐣_0`$ \[recall (7) — here we are assuming that $`𝐝`$ $``$ $`\beta `$, so that $`c_𝐣^𝐝`$ $``$ $`M`$\]. The remaining term is -$`c_{𝐣_1}^𝐝`$, where
>
> $$𝐣_1^{}=𝐣_0,\text{ that is, }𝐱^{𝐣_1}x_i=𝐱^{𝐣_0},$$
>
> or zero if $`𝐱^{𝐣_0}`$ is not divisible by $`x_i`$. Thus in $`M/M^2`$ we have
>
> (10)
> $$c_{𝐣_0}^𝐝^{}=\{\begin{array}{c}c_{𝐣_1}^𝐝,\text{ if }x_i\text{ divides }𝐱^{𝐣_0},\text{ and}\hfill \\ 0,\text{ otherwise}.\hfill \end{array}$$
>
> It is convenient to depict each $`c_𝐣^𝐝`$ by an arrow from $`𝐝`$ to $`𝐣`$ (see Figure 2). Equation (10) says that we may move these arrows \[in the $`x_i`$-direction, $`1`$ $``$ $`i`$ $``$ $`r`$,\] without changing their values modulo $`M^2`$, provided we keep the head inside $`\beta `$ and the tail outside. More generally, as long as we keep the tail in the first \[orthant\] and outside $`\beta `$, we may even move the head across the \[hyperplane ($`x_i`$-degree $`=`$ $`0`$)\]. When this is possible, the value of the arrow (mod $`M^2`$) is zero.
Henceforth when we speak of an “arrow” $`c_𝐣^𝐝`$, we will assume that $`𝐝`$ $``$ $`\beta `$, so that there is a corresponding indeterminate $`C_𝐣^𝐝`$ (as defined in Section 2.3). We will say that two arrows $`c_{𝐣_1}^{𝐛_1}`$ and $`c_{𝐣_2}^{𝐛_2}`$ are translation-equivalent if the first can be moved to the second by a sequence of discrete steps in the various variable directions such that the head of the arrow remains inside $`\beta `$ and the tail remains outside $`\beta `$. This clearly defines an equivalence relation on the set of arrows, which we denote $`c_{𝐣_1}^{𝐛_1}`$ $``$ $`c_{𝐣_2}^{𝐛_2}`$. For example, in Figure 2, we see that $`c_{(0,1)}^{(1,1)}c_{(0,2)}^{(1,2)}`$. Abusing the language and notation, we say $`c_𝐣^𝐛`$ is translation-equivalent to 0, and write $`c_𝐣^𝐛0`$, to indicate that the arrow can be translated to a position such that one more step in some direction of decreasing degree would cause the head of the arrow to exit the first orthant, with the tail remaining a monomial outside of $`\beta `$. The arrow $`c_{(0,0)}^{(1,2)}`$ in Figure 2 provides an example: one further step in the decreasing $`x_2`$-direction would cause the head of this arrow to exit the first quadrant.
From the foregoing, it is clear that
(11)
$$\begin{array}{ccc}\hfill c_{𝐣_1}^{𝐛_1}c_{𝐣_2}^{𝐛_2}& & c_{𝐣_1}^{𝐛_1}c_{𝐣_2}^{𝐛_2}(modM^2),\text{ and}\hfill \\ \hfill c_𝐣^𝐛0& & c_𝐣^𝐛0(modM^2).\hfill \end{array}$$
Furthermore, the reasoning in the quoted passage can be adapted to prove
###### Proposition 3.1.1.
Let $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ be any one of the polynomial generators of the kernel of the map $`\delta `$ *(9)*. Then
1. Each term in $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ is, up to sign, either a single indeterminate $`C_𝐣^𝐛`$ or a product of two such indeterminates.
2. The number of linear terms in $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ is equal to $`1`$ or $`2`$.
3. If there are two linear terms in $`\tau _{𝐣_0}^{(𝐝,x_i)}`$, then these terms have the form $`C_{𝐣_0}^𝐝^{}`$ and -$`C_{𝐣_1}^𝐝`$, where $`𝐱^𝐝^{}`$ $`=`$ $`x_i𝐱^𝐝`$ and $`𝐱^{𝐣_0}`$ $`=`$ $`x_i𝐱^{𝐣_1}`$; in particular, the signs differ, and the corresponding arrows $`c_{𝐣_0}^𝐝^{}`$ and $`c_{𝐣_1}^𝐝`$ are translation-equivalent (i.e., $`c_{𝐣_0}^𝐝^{}`$ $``$ $`c_{𝐣_1}^𝐝`$).
4. If there is only one linear term in $`\tau _{𝐣_0}^{(𝐝,x_i)}`$, then it is the term $`C_{𝐣_0}^𝐝^{}`$, and the corresponding arrow $`c_{𝐣_0}^𝐝^{}`$ $``$ $`0`$. ∎
### 3.2. $`𝗄`$-Linearly independent sets of arrows in $`M/M^2`$
Our main goal in this section is to prove the following
###### Theorem 3.2.1.
Let $`S`$ be a finite set of arrows $`c_𝐣^𝐝`$ having $`d`$ members. Then $`S`$ has maximal rank $`(\text{mod }M^2)`$ (that is, the $`𝗄`$-span of $`S`$ in $`M/M^2`$ has dimension $`d`$) if and only if the following conditions hold:
1. $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ for all $`c_𝐣^𝐝`$ $``$ $`S`$.
2. $`c_{𝐣_1}^{𝐝_1}`$ $`\sim ̸`$ $`c_{𝐣_2}^{𝐝_2}`$ for all $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_{𝐣_2}^{𝐝_2}`$ in $`S`$.
Before proceeding with the proof, we make a few preparations. We first extend the notion of translation-equivalence to the set of indeterminates $`C_𝐣^𝐛`$ in the obvious way:
$$\begin{array}{ccc}\hfill C_{𝐣_1}^{𝐛_1}C_{𝐣_2}^{𝐛_2}& & c_{𝐣_1}^{𝐛_1}c_{𝐣_2}^{𝐛_2}\hfill \\ \hfill C_𝐣^𝐛0& & c_𝐣^𝐛0.\hfill \end{array}$$
Then Proposition 3.1.1 immediately yields
###### Lemma 3.2.2.
Let $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ be any one of the polynomial generators of the kernel of the map $`\delta `$ *(9)*. Then
1. If $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ has two linear terms $`C_{𝐣_0}^𝐝^{}`$ and -$`C_{𝐣_1}^𝐝`$, then $`C_{𝐣_0}^𝐝^{}`$ $``$ $`C_{𝐣_1}^𝐝`$.
2. If $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ has just one linear term $`C_{𝐣_0}^𝐝^{}`$, then $`C_{𝐣_0}^𝐝^{}`$ $``$ $`0`$. ∎
Given a polynomial $`P`$ $``$ $`𝗄[(C_𝐣^𝐛)]`$ of degree $`q`$, we write
$$P=P^{(0)}+P^{(1)}+\mathrm{}+P^{(q)},$$
where each $`P^{(j)}`$ is homogeneous of degree $`j`$. We also define the map
(12)
$$E:𝗄[(C_𝐣^𝐛)]𝗄[(C_𝐣^𝐛)],C_𝐣^𝐛C_𝐣^𝐛,$$
where the target is the ring of polynomials in the translation-equivalence classes $`C_𝐣^𝐛`$ of the indeterminates $`C_𝐣^𝐛`$.
###### Lemma 3.2.3.
Let $`a_1,\mathrm{},a_s`$ $``$ $`𝗄`$, and $`\tau _1,\mathrm{},\tau _s`$ be $`s`$ of the polynomials $`\tau _{𝐣_0}^{(𝐝,x_i)}`$. Let $`N`$ denote the set of indices $`n`$ such that $`\tau _n`$ has just one linear term $`C_{𝐣_n}^{𝐛_n}`$. Then
$$E(\underset{j=1}{\overset{s}{}}a_j\tau _j^{(1)})=\underset{nN}{}a_nC_{𝐣_n}^{𝐛_n}.$$
*Proof:* It suffices to observe that for the $`\tau _j`$ containing two linear terms $`C_{𝐣_0}^{𝐛_j^{}}`$ and $`C_{𝐣_1}^{𝐛_j}`$, we have that
$$E(a_j\tau _j^{(1)})=a_j\left(C_{𝐣_0}^{𝐛_j^{}}C_{𝐣_1}^{𝐛_j}\right)=0,$$
where the last equality follows from Lemma 3.2.2. ∎
*Proof of Theorem* 3.2.1: From the implications (11), it is clear that if either of the conditions (a), (b) in the statement of the theorem fails, then the dimension of the $`𝗄`$-span of $`S`$ in $`M/M^2`$ is $`<`$ $`d`$; therefore, if the dimension is $`d`$, then (a) and (b) must hold. To prove the converse, we argue by contradiction: Suppose that (a) and (b) hold, but that the $`𝗄`$-span of $`S`$ in $`M/M^2`$ has dimension $`<`$ $`d`$. Then there exists a nontrivial $`𝗄`$-linear combination of (distinct) elements of $`S`$ that is congruent to 0 modulo $`M^2`$, say
$$\underset{i=1}{\overset{m}{}}\alpha _ic_{𝐣_i}^{𝐝_i}0(modM^2),$$
which implies that
$$\underset{i=1}{\overset{m}{}}\alpha _ic_{𝐣_i}^{𝐝_i}(\text{elt. of }M^2)=0R(R=\text{coord. ring of }U_\beta ).$$
From the description of $`R`$ as a quotient of the polynomial ring $`𝗄[(C_𝐣^𝐝)]`$ provided by Proposition 2.3.1, we see that we have an equation
$$\underset{i=1}{\overset{m}{}}\alpha _iC_{𝐣_i}^{𝐝_i}(\text{terms in }C_𝐣^𝐝\text{ of degree }2)=\underset{j=1}{\overset{s}{}}g_j\tau _j,$$
where the coefficients $`g_j`$ $``$ $`𝗄[(C_𝐣^𝐝)]`$ and each $`\tau _j`$ is one of the polynomials $`\tau _{𝐣_0}^{(𝐝,x_i)}`$ that generate the kernel of the map $`\delta `$ (9). Since the $`\tau _j`$ have only linear and quadratic terms (by Proposition 3.1.1), we see that the terms of the form $`g_j^{(i)}\tau _j`$ for $`i1`$ can be transposed to the left to yield an equation
$$\underset{i=1}{\overset{m}{}}\alpha _iC_{𝐣_i}^{𝐝_i}(\text{terms in }C_𝐣^𝐝\text{ of degree }2)=\underset{j=1}{\overset{s}{}}g_j^{(0)}\tau _j.$$
Equating the degree-1 terms on both sides, we obtain
$$\underset{i=1}{\overset{m}{}}\alpha _iC_{𝐣_i}^{𝐝_i}=\underset{j=1}{\overset{s}{}}g_j^{(0)}\tau _j^{(1)}.$$
As in Lemma 3.2.3, we let $`N`$ denote the set of indices $`n`$ such that $`\tau _n^{(1)}`$ consists of a single term $`C_{𝐣_n}^{𝐛_n}`$. Applying the map $`E`$ (12) to both sides and rewriting the RHS using the lemma, we find that
$$\underset{i=1}{\overset{m}{}}\alpha _iC_{𝐣_i}^{𝐝_i}=\underset{nN}{}g_n^{(0)}C_{𝐣_n}^{𝐛_n}.$$
Recall that conditions (a) and (b) hold, by hypothesis, and that we began with a nontrivial linear combination of the $`c_𝐣^𝐝`$ $``$ $`S`$, so that at least one of the coefficients $`\alpha _i`$ $``$ $`0`$. If the corresponding equivalence class $`C_{𝐣_i}^{𝐝_i}`$ does not appear on the RHS of the last equation, then the term $`\alpha _iC_{𝐣_i}^{𝐝_i}`$ must cancel with one or more other terms $`\alpha _jC_{𝐣_j}^{𝐝_j}`$ on the LHS; whence, $`c_{𝐣_i}^{𝐝_i}`$ $``$ $`c_{𝐣_j}^{𝐝_j}`$, which contradicts condition (b). Therefore, we must have that $`C_{𝐣_i}^{𝐝_i}`$ $`=`$ $`C_{𝐣_n}^{𝐝_n}`$ for some $`n`$ $``$ $`N`$. That is, we have $`C_{𝐣_i}^{𝐝_i}`$ $``$ $`C_{𝐣_n}^{𝐛_n}`$, but since $`C_{𝐣_n}^{𝐛_n}`$ $``$ $`0`$ by Lemma 3.2.2, it follows immediately that $`C_{𝐣_i}^{𝐝_i}`$ $``$ $`0`$ $``$ $`c_{𝐣_i}^{𝐝_i}`$ $``$ $`0`$, which contradicts condition (a), and the proof is complete. ∎
As a corollary, we obtain the following converses to the implications (11):
###### Corollary 3.2.4.
For $`c_𝐣^𝐛,c_{𝐣_1}^{𝐛_1},c_{𝐣_2}^{𝐛_2}`$ any arrows, we have that
1. $`c_𝐣^𝐛0(modM^2)c_𝐣^𝐛0`$.
2. $`c_{𝐣_1}^{𝐛_1}c_{𝐣_2}^{𝐛_2}0(modM^2)c_{𝐣_1}^{𝐛_1}c_{𝐣_2}^{𝐛_2}`$.
*Proof:* Apply the theorem to the sets $`\{c_𝐣^𝐛\}`$ and $`\{c_{𝐣_1}^{𝐛_1},c_{𝐣_2}^{𝐛_2}\}`$, neither of which has maximal rank (mod $`M^2`$), provided that $`c_{𝐣_1}^{𝐛_1}`$ $``$ $`c_{𝐣_2}^{𝐛_2}`$. ∎
## 4. A linearly independent set in $`M/M^2`$ of cardinality $`rn`$
Given a basis set of monomials $`\beta `$ of size $`n`$ in $`r`$ variables (or equivalently the associated monomial ideal $`I_\beta `$), we exhibit a set $`𝒮`$ of arrows $`c_𝐣^𝐝`$ of cardinality $`rn`$ whose $`𝗄`$-span in $`M/M^2`$ has dimension $`rn`$. This gives a second proof that the $`𝗄`$-dimension of the cotangent space of the point $`t_\beta `$ $``$ $`\text{H}^n`$ is at least $`rn`$, as promised in Remark 2.4.2. More importantly, if $`t_\beta `$ $``$ $`\text{H}^n`$ is nonsingular, then $`𝒮`$ must be a basis of the cotangent space. Because of the particular form of the arrows in $`𝒮`$, it is often easy to show in particular cases that $`𝒮`$ does *not* span the cotangent space; we conclude that $`t_\beta `$ is a singular point in such cases. On the other hand, there are several families of basis sets $`\beta `$ for which we can prove that $`𝒮`$ *is* a basis of the cotangent space.
### 4.1. Standard and minimal arrows
We begin by identifying the type of arrow that will belong to our set $`𝒮`$. First we define the vector of the arrow $`c_𝐣^𝐝`$ to be the tuple $`𝐣𝐝`$ (recall that we only speak of an “arrow” when $`𝐝`$ $``$ $`\beta `$ $``$ $`𝐱^𝐝`$ $``$ $`I_\beta `$). Since $`𝐱^𝐝`$ $`|`$ $`𝐱^𝐣`$ is then impossible, the following is immediate:
###### Lemma 4.1.1.
For every arrow $`c_𝐣^𝐝`$, the vector of the arrow has at least one negative component. ∎
We say that $`c_𝐣^𝐝`$ is a standard arrow provided that the vector of the arrow has exactly one negative component; if this component is the $`i`$-th, corresponding to the variable $`x_i`$, then we say that $`c_𝐣^𝐝`$ is standard for $`x_i`$, or $`x_i`$-standard. Concretely, this means that the monomial $`𝐱^𝐝`$ at the tail of the arrow has a strictly larger $`x_i`$-degree than the monomial $`𝐱^𝐣`$ at the head, but the $`x_k`$-degree of $`𝐱^𝐝`$ is $``$ the $`x_k`$-degree of $`𝐱^𝐣`$ for all other variables $`x_k`$ $``$ $`x_i`$. For example, in Figure 2, the arrows $`c_{(0,1)}^{(1,1)}`$, $`c_{(0,2)}^{(1,1)}`$, and $`c_{(0,1)}^{(2,0)}`$ are standard for $`x_1`$, and $`c_{(0,0)}^{(1,2)}`$ is not a standard arrow.
Recall that the monomials in $`x_1,\mathrm{},x_r`$ are partially ordered by divisibility, and $`I_\beta `$ is generated by the minimal monomials in $`I_\beta `$ under this ordering (the minimal generators of $`I_\beta `$). We call an arrow $`c_𝐣^𝐝`$ whose tail $`𝐱^𝐝`$ is a minimal generator a minimal arrow; if the arrow is standard (for $`x_i`$), we call it a minimal standard arrow (for $`x_i`$). In light of the results of Section 3, the next lemma shows that in seeking a set of arrows to span the $`𝗄`$-vector space $`M/M^2`$, it suffices to restrict one’s attention to minimal arrows.
###### Lemma 4.1.2.
Let $`c_𝐣^𝐝`$ be an arbitrary arrow, and $`𝐱^𝐛`$ a minimal generator of $`I_\beta `$ such that $`𝐱^𝐛`$ $`|`$ $`𝐱^𝐝`$ (at least one such minimal generator must clearly exist). Then either $`c_𝐣^𝐝`$ $``$ $`0`$ or $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_1}^𝐛`$ for some $`𝐣_1`$ $``$ $`\beta `$.
*Proof:* We can clearly translate $`𝐱^𝐝`$ to $`𝐱^𝐛`$ by a series of degree-reducing steps in the various variable directions (each step involves dividing the head and tail of the arrow by one of the $`x_k`$). Since dividing a basis monomial (at the head of the arrow) by $`x_k`$ either keeps us inside $`\beta `$ or causes us to exit the first orthant across the hyperplane ($`x_k`$-degree $`=`$ $`0`$), the result follows immediately. ∎
In fact, the set $`𝒮`$ that we are out to construct will consist of minimal *standard* arrows, as our main theorem asserts:
###### Theorem 4.1.3.
Let $`\beta `$ be a basis set of $`n`$ monomials in the variables $`x_1`$, $``$, $`x_r`$, and $`I_\beta `$ $``$ $`𝗄[𝐱]`$ the associated monomial ideal. Then there exists a set $`𝒮`$ of minimal standard arrows that has cardinality $`rn`$ and maximal rank $`(\text{mod }M/M^2)`$; that is, the $`𝗄`$-span of $`S`$ in $`M/M^2`$ has dimension $`rn`$.
The proof will be given in Section 4.5, after the necessary preparations have been made.
### 4.2. Advancement of minimal standard arrows
We begin with a host of definitions. Let $`\beta `$, $`I_\beta `$, etc., be as above. Note first of all that, since $`I_\beta `$ has finite colength, there is for each variable $`x_i`$ a minimal exponent $`w_i`$ $`>`$ $`0`$ such that $`x_i^{w_i}`$ $``$ $`I_\beta `$; we will call $`w_i`$ the $`x_i`$-width of $`\beta `$. It is then clear that
(13)
$$x_i\text{-degree}(𝐱^𝐣)<w_i\text{ for all }𝐱^𝐣\beta ,1ir.$$
For example, the basis set shown in Figure 2 has $`x_1`$-width $`w_1`$ $`=`$ $`2`$ and $`x_2`$-width $`w_2`$ $`=`$ 3.
Given $`𝐱^𝐣`$ $``$ $`\beta `$, we define the $`x_i`$-column of $`𝐱^𝐣`$ to be the set of monomials $`𝐱^𝐣`$, $`𝐱^𝐣/x_i`$, $`𝐱^𝐣/x_i^2`$, …, $`𝐱^𝐣/x_i^q`$, where $`q`$ $`=`$ $`x_i`$-degree$`(𝐱^𝐣)`$. Clearly these monomials all belong to $`\beta `$. We define the $`x_i`$-shadow of an arrow $`c_𝐣^𝐝`$ to be the set of arrows $`c_𝐣^{}^𝐝`$ such that $`𝐱^𝐣^{}`$ is in the $`x_i`$-column of $`𝐱^𝐣`$, and
$$x_i\text{-height}(c_𝐣^𝐝)=x_i\text{-degree}(𝐱^𝐝).$$
If $`c_𝐣^𝐝`$ is a standard arrow for $`x_i`$, we define its $`x_i`$-offset to be the $`(r1)`$-tuple $`v`$ of non-negative integers obtained by deleting the $`i`$-th (negative) component of the vector $`𝐣𝐝`$.
We say that an $`x_i`$-standard arrow $`c_𝐣^𝐝`$ can be advanced if either $`c_𝐣^𝐝`$ $``$ $`0`$ or $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_1}^{𝐝_1}`$ with the latter arrow having strictly smaller $`x_i`$-height. In each case, the idea is that $`c_𝐣^𝐝`$ can be translated so that its tail moves closer to the hyperplane ($`x_i`$-degree $`=`$ $`0`$), provided that we allow the head to exit the first orthant across this hyperplane in the first case. Note that the head of an $`x_i`$-standard arrow can only be translated out of the first orthant across the hyperplane ($`x_i`$-degree $`=`$ $`0`$), since the tail must remain within the first orthant (and outside of $`\beta `$) during translation.
###### Lemma 4.2.1.
Suppose that $`c_𝐣^𝐝`$ is a standard arrow for $`x_i`$. Then:
1. All of the arrows in the $`x_i`$-shadow of $`c_𝐣^𝐝`$ are $`x_i`$-standard and have equal offsets.
2. If $`c_𝐣^𝐝`$ can be advanced, then so can all the arrows in its $`x_i`$-shadow.
3. If $`c_𝐣^𝐝`$ $`\sim ̸0`$ can be advanced, then we can advance this arrow until we reach a minimal standard (for $`x_i`$) arrow $`c_𝐣^{}^𝐝^{}`$ $``$ $`c_𝐣^𝐝`$ that cannot be advanced.
*Proof:* Statement (a) is immediate. Statement (b) follows from the observation that if $`c_𝐣^𝐝`$ can be translated to $`c_{𝐣_1}^𝐛`$, then every arrow in the shadow of $`c_𝐣^𝐝`$ is simultaneously translated “in parallel” either to an arrow in the shadow of $`c_{𝐣_1}^𝐛`$ or out of the first orthant (across the hyperplane ($`x_i`$-degree $`=`$ $`0`$)). Statement (c) follows easily from Lemma 4.1.2: after advancing $`c_𝐣^𝐝`$ to $`c_{𝐣_1}^{𝐝_1}`$ of smaller $`x_i`$-height, we translate by degree-reducing steps to a minimal standard arrow $`c_𝐣^{}^𝐝^{}`$, and repeat as often as necessary. ∎
### 4.3. Shadow promotion
The simultaneous translation and advancement of the arrows in a shadow leads to a process of arrow replacement that is the key to the proof of Theorem 4.1.3; we call this process, which we proceed to describe, *shadow promotion*. Suppose that $`c_𝐣^𝐝`$ is an $`x_i`$-standard arrow that can be advanced, which implies by statement (b) of Lemma 4.2.1 that every arrow in the $`x_i`$-shadow of $`c_𝐣^𝐝`$ can be advanced as well. Consider a sequence of translation steps that advances $`c_𝐣^𝐝`$, and let $`c_{𝐣_1}^{𝐝_1}`$ be the first position in the sequence of steps from which it is possible to move one step in the decreasing $`x_i`$-direction either to reach an arrow of $`x_i`$-height $`<`$ $`x_i`$-height of $`c_𝐣^𝐝`$ or to move the head of the arrow out of the first orthant; put another way, $`c_{𝐣_1}^{𝐝_1}`$ is the first position such that
$$𝐱^{𝐝_1}/x_iI_\beta ,\text{ and }x_i\text{-degree}(𝐱^{𝐝_1})=x_i\text{-degree}(𝐱^𝐝).$$
Note that
(14)
$$x_i\text{-degree}(𝐱^{𝐣_1})=x_i\text{-degree}(𝐱^𝐣).$$
Let $`𝐱^𝐛`$ be a minimal generator of $`I_\beta `$ that divides $`𝐱^{𝐝_1}/x_i`$; it is clear that
$$x_i\text{-degree}(𝐱^𝐛)<x_i\text{-degree}(𝐱^𝐝).$$
We must have that
(15)
$$x_i\text{-degree}(𝐱^𝐛)>x_i\text{-degree}(𝐱^{𝐣_1}),$$
since otherwise $`x_i\text{-degree}(𝐱^𝐛)x_i\text{-degree}(𝐱^{𝐣_1})`$ holds in addition to
(16)
$$x_k\text{-degree}(𝐱^𝐛)x_k\text{-degree}(𝐱^{𝐝_1})x_k\text{-degree}(𝐱^{𝐣_1})\text{ for all }ki,$$
where the last inequality holds because $`c_𝐣^𝐝`$ and its translate $`c_{𝐣_1}^{𝐝_1}`$ are standard arrows for $`x_i`$. In other words, $`𝐱^𝐛`$ $`|`$ $`𝐱^{𝐣_1}`$, which is a contradiction since $`𝐱^{𝐣_1}`$ is a basis monomial and $`𝐱^𝐛`$ is not. The contradiction establishes (15); in particular, we have that $`x_i`$-degree$`(𝐱^𝐛)`$ $`>`$ $`0`$.
Furthermore, (16) shows that we can translate the arrow $`c_{𝐣_1}^{𝐝_1}`$ in degree-reducing directions, excluding the $`x_i`$-th, to reach an arrow $`c_{𝐣_2}^{𝐝_2}`$ such that $`𝐱^{𝐝_2}`$ and $`𝐱^𝐛`$ differ only in $`x_i`$-degree. (Because the arrows involved are $`x_i`$-standard, the head of the arrow can never leave the first orthant during any of these steps.) We then have that
$$x_i\text{-degree}(𝐱^{𝐣_2})=x_i\text{-degree}(𝐱^{𝐣_1})<x_i\text{-degree}(𝐱^𝐛);$$
therefore, the arrow $`c_{𝐣_2}^𝐛`$ is a minimal standard arrow for $`x_i`$ having the same offset as $`c_{𝐣_2}^{𝐝_2}`$, $`c_{𝐣_1}^{𝐝_1}`$, and $`c_𝐣^𝐝`$, and having the same number of arrows in its shadow as $`c_{𝐣_1}^{𝐝_1}`$ and $`c_𝐣^𝐝`$ do (recall (14)). We will call the shadow of $`c_{𝐣_2}^𝐛`$ the promotion image of the shadow of the original standard arrow $`c_𝐣^𝐝`$; the promotion image has the same number of arrows as the original shadow, and the arrows in the promotion image are minimal standard arrows for $`x_i`$ having strictly smaller $`x_i`$-height than the original arrows. Indeed, we can view shadow promotion as the replacement of every arrow $`c_𝐣^{}^𝐝`$ in the shadow of $`c_𝐣^𝐝`$ with the arrow $`c_{𝐣_2^{}}^𝐛`$ in the shadow of $`c_{𝐣_2}^𝐛`$ such that
(17)
$$x_i\text{-degree}(𝐱^{𝐣_2^{}})=x_i\text{-degree}(𝐱^𝐣^{}).$$
Since the promotion image depends on the path chosen to advance $`c_𝐣^𝐝`$, it is not unique in general.
For example, in Figure 2, the arrow $`c_{(0,1)}^{(2,0)}`$ is a standard arrow for $`x_1`$ that can be advanced; its shadow is the singleton set $`\{c_{(0,1)}^{(2,0)}\}`$. The promotion image of this shadow is easily seen to be $`\{c_{(0,2)}^{(1,1)}\}`$.
### 4.4. Iterated shadow promotion
Select one of the variables $`x_i`$, and recall that $`x_i^{w_i}`$ $`=`$ $`𝐱^𝐝`$ is the least power of $`x_i`$ that belongs to the monomial ideal $`I_\beta `$; it is clear that $`𝐱^𝐝`$ is a minimal generator of the ideal, which we will call the $`i`$-th corner monomial of $`\beta `$. Note that *every* arrow $`c_𝐣^𝐝`$ (with tail at the $`i`$-th corner monomial) is minimal and standard for $`x_i`$, since only the $`x_i`$-degree decreases when we move from the tail to the head of the arrow. In this case, the offset of the arrow is the $`(r1)`$-tuple $`v`$ obtained by deleting the $`i`$-th component of $`𝐣`$. Let $`𝒮_0(i,v)`$ denote the set of all (minimal $`x_i`$-standard) arrows having tail $`𝐱^𝐝`$ and offset $`v`$, and let $`c_{𝐣_1}^𝐝`$ be the arrow in $`𝒮_0(i,v)`$ whose head $`𝐱^{𝐣_1}`$ has maximal $`x_i`$-degree $`q(i,v)`$; it follows that $`𝒮_0(i,v)`$ is the $`x_i`$-shadow of the arrow $`c_{𝐣_1}^𝐝`$, and the number of arrows in $`𝒮_0(i,v)`$ is $`q(i,v)`$ $`+`$ $`1`$.
By iterating the process of shadow promotion, we can construct a set of $`q(i,v)+1`$ minimal $`x_i`$-standard arrows that have offset $`v`$ and cannot be advanced. Indeed, if none of the arrows in $`𝒮_0(i,v)`$ can be advanced, then $`𝒮_0(i,v)`$ is itself the desired set, which we will denote $`𝒮(i,v)`$. Otherwise, one or more of the arrows in $`𝒮_0(i,v)`$ can be advanced; in this case we let $`c_{𝐣_2}^𝐝`$ $``$ $`𝒮_0(i,v)`$ be the advanceable arrow whose head has maximal $`x_i`$-degree, so that the set of all advanceable arrows in $`𝒮_0(i,v)`$ is the $`x_i`$-shadow of $`c_{𝐣_2}^𝐝`$, by Lemma 4.2.1. We then promote the $`x_i`$-shadow of $`c_{𝐣_2}^𝐝`$, as described in the previous subsection, and denote the (not necessarily unique) promotion image by $`P_0(i,v)`$; this permits us to form a new set of minimal $`x_i`$-standard arrows of offset $`v`$:
$$𝒮_1(i,v)=\left(𝒮_0(i,v)(x_i\text{-shadow of }c_{𝐣_2}^𝐝)\right)P_0(i,v).$$
Since the number of arrows in the promotion image is equal to the number of arrows in the shadow being promoted, it is clear that the number of arrows in $`𝒮_1(i,v)`$ is equal to the number of arrows in $`𝒮_0(i,v)`$. Furthermore, the $`x_i`$-height of the arrows in $`P_0(i,v)`$ is strictly less than the height of the arrows in $`𝒮_0(i,v)`$, which is $`w_i`$. If none of the arrows in $`𝒮_1(i,v)`$ can be advanced, then $`𝒮_1(i,v)`$ is the desired set $`𝒮(i,v)`$; otherwise, the set of advanceable arrows in $`𝒮_1(i,v)`$ is equal to the shadow of an advanceable arrow $`c_𝐣^{}^𝐝^{}`$ $``$ $`P_0(i,v)`$. Replacing this shadow by its promotion image $`P_1(i,v)`$, we obtain yet another set of minimal $`x_i`$-standard arrows of offset $`v`$:
$$𝒮_2(i,v)=\left(𝒮_1(i,v)(x_i\text{-shadow of }c_𝐣^{}^𝐝^{})\right)P_1(i,v).$$
Continuing in this way, we eventually arrive at the desired set $`𝒮(i,v)`$ of minimal $`x_i`$-standard arrows of offset $`v`$, none of which can be advanced. The process must terminate because at each stage the $`x_i`$-height of the arrows that can be advanced is strictly less than the corresponding height at the previous stage, and this height cannot decrease to 0.
### 4.5. Proof of Theorem 4.1.3
For each of the variables $`x_i`$, $`1`$ $``$ $`i`$ $``$ $`r`$, we will construct a set $`𝒮(i)`$ consisting of $`n`$ minimal standard arrows for $`x_i`$ such that no two of the arrows in $`𝒮(i)`$ are translation-equivalent to each other, and none is $``$ $`0`$. But then the same is true of
$$𝒮=\underset{i=1}{\overset{r}{}}𝒮(i),$$
since $`x_i`$\- and $`x_j`$-standard arrows cannot be translation-equivalent if $`i`$ $``$ $`j`$; in particular, the cardinality of $`𝒮`$ is $`rn`$. Theorem 3.2.1 now yields that the $`𝗄`$-span of $`𝒮`$ in $`M/M^2`$ has dimension $`rn`$, as desired.
It remains to construct the sets $`𝒮_i`$, but this is not difficult. Simply form the sets $`𝒮(i,v)`$ for every possible offset $`v`$ of a minimal $`x_i`$-standard arrow having tail the $`i`$-th corner monomial $`x_i^{w_i}`$, as in Subsection 4.4, and let
$$𝒮(i)=\underset{\{\text{offsets }v\}}{}𝒮(i,v).$$
It is clear that every monomial $`𝐱^𝐣`$ $``$ $`\beta `$ is the head of an arrow in one of the initial sets $`𝒮_0(i,v)`$. Furthermore, the sets $`𝒮(i,v)`$ are pairwise disjoint, since arrows of different offsets cannot be equal. Therefore, counting elements in the various sets, we find that
$$\begin{array}{ccc}\hfill |𝒮(i)|& =& _{\{\text{offsets }v\}}|𝒮(i,v)|\hfill \\ & =& _{\{\text{offsets }v\}}|𝒮_0(i,v)|\hfill \\ & =& |\beta |=n.\hfill \end{array}$$
By construction, the arrows $`c_𝐣^𝐛`$ $``$ $`𝒮(i,v)`$ are minimal $`x_i`$-standard arrows of offset $`v`$ that cannot be advanced; in particular, we have that $`c_𝐣^𝐛`$ $`\sim ̸`$ $`0`$. Therefore, the set $`𝒮(i)`$ consists of $`n`$ minimal $`x_i`$-standard arrows, none of which are $``$ $`0`$. We must now show that no two distinct arrows in $`𝒮(i)`$ are translation-equivalent to one another. To this end, let $`c_{𝐣_1}^{𝐛_1}`$ and $`c_{𝐣_2}^{𝐛_2}`$ be distinct arrows in $`𝒮(i)`$. If these arrows have different offsets, then they cannot possibly be translation-equivalent, so suppose that they each have offset $`v`$; that is, $`c_{𝐣_1}^{𝐛_1}`$, $`c_{𝐣_2}^{𝐛_2}`$ $``$ $`𝒮(i,v)`$. If the two arrows have different heights (i.e., different tails), then, since neither arrow can be advanced, it is again clear that they cannot be translation-equivalent. If the two arrows have the same height, then they have the same tail; therefore, if they were translation-equivalent, they would be equal, which we are assuming is not the case. We conclude that distinct arrows in $`𝒮(i)`$ cannot be translation-equivalent, and the proof is complete. ∎
###### Remark 4.5.1.
It is clear that none of the arrows in the set $`𝒮`$ constructed in Theorem 4.1.3 can be advanced.
###### Remark 4.5.2.
Since the promotion image of a shadow is not necessarily unique, neither are the sets $`𝒮(i)`$ and $`𝒮`$.
### 4.6. Summary and terminology
We will call a set $`𝒮(i)`$ constructed as in Subsection 4.5 a standard $`x_i`$-sub-bunch, and the union
$$𝒮=\underset{i=1}{\overset{r}{}}𝒮(i),$$
a standard bunch, of arrows for $`\beta `$. These sets have the following properties:
contains $`n`$ $`=`$ $`|\beta |`$ minimal $`x_i`$-standard arrows; has maximal rank (mod $`M^2`$); and no arrow in the set can be advanced.
is the union of sets $`𝒮(i)`$, $`1`$ $``$ $`i`$ $``$ $`r`$, and accordingly: contains $`rn`$ minimal standard arrows; has maximal rank (mod $`M^2`$); and no arrow in the set can be advanced.
We will have occasion to consider more general sets of arrows $`𝒮^{}(i)`$ that we call near-standard $`x_i`$-sub-bunches: these are similar to standard $`x_i`$-sub-bunches in that they have cardinality $`n`$, have maximal rank (mod $`M^2`$), and consist of $`x_i`$-standard arrows that cannot be advanced; they differ in that the arrows they contain need not be *minimal*. For example, one way to obtain a near-standard $`x_i`$-sub-bunch is to replace some or all of the arrows $`c_𝐣^𝐝`$ $``$ $`𝒮(i)`$ (a standard $`x_i`$-sub-bunch) with arrows
$$c_𝐣^{}^𝐝^{}c_𝐣^𝐝,\text{ such that }x_i\text{-height}(c_𝐣^{}^𝐝^{})=x_i\text{-height}(c_𝐣^𝐝).$$
We will call the union
$$𝒮^{}=\underset{i=1}{\overset{r}{}}𝒮^{}(i),$$
of near-standard sub-bunches $`𝒮^{}(i)`$ a near-standard bunch of arrows for $`\beta `$; it is clear that $`𝒮^{}`$ consists of $`rn`$ standard arrows that cannot be advanced, and has maximal rank (mod $`M^2`$).
We have the following useful corollaries of the proof of Theorem 4.1.3:
###### Corollary 4.6.1.
Let $`c_𝐣^𝐝`$ be an $`x_i`$-standard arrow for the basis set $`\beta `$. Then every standard $`x_i`$-sub-bunch $`𝒮(i)`$ of arrows for $`\beta `$ contains a unique arrow $`c_𝐣^{}^𝐝^{}`$ such that
$$x_i\text{-degree}(𝐱^𝐣^{})=x_i\text{-degree}(𝐱^𝐣)\text{ and }\text{offset}(c_𝐣^{}^𝐝^{})=\text{offset}(c_𝐣^𝐝).$$
*Proof:* Let
$$p=x_i\text{-degree}(𝐱^𝐝),m=𝐱^𝐝/x_i^p,\text{and }𝐱^{𝐣_0}=𝐱^𝐣/m\beta .$$
Then $`𝐱^{𝐣_0}`$ has offset $`v`$ from the corner monomial $`x_i^{w_i}`$ = $`𝐱^𝐛`$, so $`c_{𝐣_0}^𝐛`$ $``$ $`𝒮_0(i,v)`$, and is the only such arrow whose head has $`x_i`$-degree $`=`$ $`x_i`$-degree$`(𝐱^𝐣)`$. Under iterated shadow promotion, arrows are replaced by other arrows having the same offset and $`x_i`$-degree of the head (17), so we see that any standard $`x_i`$-sub-bunch $`𝒮(i)`$ contains a unique arrow $`c_{𝐣_1}^{𝐝_1}`$ as stated in the lemma, the arrow that ultimately replaces $`c_{𝐣_0}^𝐛`$. ∎
###### Corollary 4.6.2.
If $`c_𝐣^𝐝`$ is an $`x_i`$-standard arrow that cannot be advanced, and whose tail is the corner monomial $`𝐱^𝐝`$ $`=`$ $`x_i^{w_i}`$, then $`c_𝐣^𝐝`$ belongs to every standard $`x_i`$-sub-bunch $`𝒮(i)`$.
*Proof:* Clear from the construction. ∎
## 5. Consequences of Theorem 4.1.3
### 5.1. Necessary and sufficient conditions for $`t_\beta `$ $``$ $`\text{H}^n`$ to be nonsingular
By Proposition 2.4.1, we know that the point $`t_\beta `$ $``$ $`\text{H}^n`$ is nonsingular if and only if the cotangent space $`M/M^2`$ has $`𝗄`$-dimension $`rn`$. Therefore, if $`𝒮^{}`$ is any standard or near-standard bunch of arrows for $`\beta `$ (see Subsection 4.6 for terminology), we have that
$$t_\beta \text{ is nonsingular }𝒮^{}\text{ spans }M/M^2𝒮^{}\text{ is a }𝗄\text{-basis of }M/M^2.$$
We will call $`\beta `$ a smooth basis set whenever any of these equivalent conditions holds.
###### Theorem 5.1.1.
In order that the basis set $`\beta `$ be smooth, it is necessary and sufficient that the following conditions hold:
1. Every non-standard arrow $`c_𝐣^𝐛`$ is translation-equivalent to $`0`$.
2. There exists a near-standard bunch of arrows $`𝒮^{}`$ such that for any standard arrow $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$, there exists an arrow $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮^{}`$ such that $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$.
*Proof:* We first prove the necessity. Let $`𝒮^{}`$ $`=`$ $`𝒮`$ be the standard bunch given by Theorem 4.1.3; since $`\beta `$ is assumed smooth, $`𝒮^{}`$ is a $`𝗄`$-basis of $`M/M^2`$. Therefore, given a non-standard arrow $`c_𝐣^𝐛`$, we have that the $`𝗄`$-span in $`M/M^2`$ of the enlarged set $`𝒮^{}`$ $``$ $`\{c_𝐣^𝐛\}`$ of $`rn`$ $`+`$ $`1`$ elements has dimension $`rn`$. Since no two distinct arrows in $`𝒮^{}`$ are translation-equivalent, none is translation-equivalent to $`0`$, and no standard arrow can be translation-equivalent to a non-standard arrow, Theorem 3.2.1 yields that $`c_𝐣^𝐛`$ $``$ $`0`$. Now let $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ be a standard arrow. If $`c_𝐣^𝐝`$ $``$ $`𝒮^{}`$, then again we must have that the $`𝗄`$-span in $`M/M^2`$ of the enlarged set $`𝒮^{}`$ $``$ $`\{c_𝐣^𝐝\}`$ has dimension $`<`$ $`rn`$ $`+`$ $`1`$; Theorem 3.2.1 now yields that $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$ for some $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮^{}`$.
It remains to prove the sufficiency; that is, assuming that conditions (a) and (b) hold, we must show that $`\beta `$ is smooth, for which it is enough to prove that $`𝒮^{}`$ spans the $`𝗄`$-vector space $`M/M^2`$. However, condition (a) implies that $`M/M^2`$ is spanned by the standard arrows, and condition (b) ensures that every standard arrow $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ is in the $`𝗄`$-linear span of $`𝒮^{}`$, so we are done. ∎
### 5.2. A sufficient condition for $`t_\beta `$ $``$ $`\text{H}^n`$ to be singular
It is often easy to identify basis sets $`\beta `$ for which the conditions (a) and (b) of Theorem 5.1.1 do *not* hold, so that the corresponding point $`t_\beta `$ $``$ $`\text{H}^n`$ is singular; we now present one simple way to do this. (Subsequent sections of the paper will be devoted to identifying basis sets $`\beta `$ for which conditions (a) and (b) *do* hold.)
If $`𝐱^𝐣`$ $``$ $`\beta `$ is maximal among monomials in $`\beta `$ for the divisibility ordering, we call $`𝐱^𝐣`$ a maximal basis monomial. If $`𝐱^𝐝`$ is a minimal generator of $`I_\beta `$ and $`𝐱^𝐣`$ is a maximal basis monomial, we call the arrow $`c_𝐣^𝐝`$ rigid, because we have
###### Lemma 5.2.1.
An arrow $`c_𝐣^𝐝`$ such that $`𝐱^𝐝`$ is a minimal generator of $`I_\beta `$ and $`𝐱^𝐣`$ is a maximal basis monomial cannot be translated except to itself. Such an arrow must belong to any $`𝗄`$-basis $``$ of $`M/M^2`$ consisting of arrows.
*Proof:* Since the tail of the arrow is a minimal generator of $`I_\beta `$, it is clear that the first step in any nontrivial translation of $`c_𝐣^𝐝`$ must be in a degree-increasing direction. However, such a step would cause the head of the arrow to exit $`\beta `$ and enter $`I_\beta `$, which is forbidden. Therefore, $`c_𝐣^𝐝`$ cannot be translated except to itself.
Now suppose that $``$ is a $`𝗄`$-basis of $`M/M^2`$ consisting of arrows, and that $`c_𝐣^𝐝`$ $``$ $``$. Then the $`𝗄`$-span in $`M/M^2`$ of the set of arrows $``$ $``$ $`\{c_𝐣^𝐝\}`$ has dimension $`<`$ $`||`$ $`+`$ $`1`$; moreover, $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$, since it has no nontrivial translations. It then follows from Theorem 3.2.1 that $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_1}^{𝐝_1}`$ for some $`c_{𝐣_1}^{𝐝_1}`$ $``$ $``$, but then $`c_𝐣^𝐝`$ $`=`$ $`c_{𝐣_1}^{𝐝_1}`$ by rigidity, which contradicts $`c_𝐣^𝐝`$ $``$ $``$. We conclude that $`c_𝐣^𝐝`$ $``$ $``$. ∎
###### Corollary 5.2.2.
If the basis set $`\beta `$ has a non-standard rigid arrow $`c_𝐣^𝐝`$, then condition (a) of Theorem 5.1.1 is false, so $`\beta `$ is not smooth.
*Proof:* It is clear that $`c_𝐣^𝐝`$ is a non-standard arrow that is not translation-equivalent to $`0`$. ∎
### 5.3. Example: $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_3\}`$
Figure 3 illustrates the basis set $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_3\}`$, the smallest basis set in three variables that has non-standard rigid arrows, and therefore fails to be smooth by Corollary 5.2.2.
Going further, we recall that by Lemma 4.1.2, the $`𝗄`$-vector space $`M/M^2`$ is spanned by minimal arrows. However, in this case, there are $`64`$ $`=`$ $`24`$ minimal arrows, 6 of which are $`0`$ (the ones with head $`𝐱^{(0,0,0)}=1`$), and the other 18 of which are rigid, and so must belong to any basis of $`M/M^2`$ consisting of arrows, by Lemma 5.2.1. We conclude from this that the $`𝗄`$-dimension of the cotangent space of $`t_\beta `$ is 18. It is known (see, e.g., \[10, Sec. 1, p. 147\], or \[11, Sec. 5.1, p. 443\]) that $`\text{H}^4`$ is irrreducible of dimension $`34`$ $`=`$ $`12`$, so we reconfirm that $`\beta `$ is not smooth.
Finally, note that three of the 18 rigid minimal arrows are non-standard (see Figure 3), leaving 15 rigid minimal standard arrows. It is clear that any near-standard bunch $`𝒮^{}`$ (which contains $`34`$ $`=`$ $`12`$ standard arrows) must exclude at least three of these rigid standard arrows, showing that $`\beta `$ also fails to satisfy condition (b) of Theorem 5.1.1.
## 6. Thickening of basis sets
In this section of the paper, we begin to identify certain families of smooth basis sets $`\beta `$. Recall that a basis set is smooth if it is associated to a nonsingular point $`t_\beta `$ $``$ $`\text{H}^n`$, and therefore satisfies conditions (a) and (b) of Theorem 5.1.1. We first consider a natural way to obtain a smooth basis set in $`r`$ variables by “thickening” a smooth basis set in $`r1`$ variables.
### 6.1. Definition of thickening
Suppose that $`\beta _0`$ is a basis set of $`n_0`$ monomials in the variables $`x_1`$, $`x_2`$, …, $`x_{r1}`$, and $`w_r`$ is a positive integer. We define a basis set in $`r`$ variables as follows:
(18)
$$\beta =\{𝐱^{𝐣_0}x_r^s𝐱^{𝐣_0}\beta _0,0sw_r1\};$$
it is indeed easy to check that $`\beta `$ is a basis set. We say that $`\beta `$ is a thickening of $`\beta _0`$ from $`r1`$ to $`r`$ variables. Note that the number of monomials in $`\beta `$ is
(19)
$$|\beta |=n=n_0w_r.$$
### 6.2. Minimal generators of $`I_\beta `$
###### Lemma 6.2.1.
Let $`\beta _0`$ and $`\beta `$ be as above. Then a minimal generator of the ideal $`I_\beta `$ is either a minimal generator $`𝐱^𝐛`$ of the ideal $`I_{\beta _0}`$ (and therefore not divisible by $`x_r`$) or else is the corner monomial $`x_r^{w_r}`$. Furthermore, every arrow with tail $`x_r^{w_r}`$ is a minimal standard arrow for $`x_r`$, and none of these arrows can be advanced.
*Proof:* Let $`m`$ $`=`$ $`𝐱^𝐝`$ be a minimal generator of $`I_\beta `$. If $`m`$ is not divisible by $`x_r`$, then for any other variable $`x_j`$, $`1`$ $``$ $`j`$ $``$ $`r1`$, we have that either $`m/x_j`$ is a monomial in $`\beta _0`$ or is undefined; it follows that $`m`$ is a minimal generator of $`I_{\beta _0}`$ $``$ $`𝗄[x_1,\mathrm{},x_{r1}]`$. If $`m`$ is divisible by $`x_r`$, then $`m/x_r`$ $``$ $`\beta `$. Therefore, $`m/x_r`$ $`=`$ $`𝐱^{𝐣_0}x_r^s`$, where $`𝐱^{𝐣_0}`$ $``$ $`\beta _0`$ and $`0`$ $``$ $`s`$ $``$ $`w_r1`$. In fact, we must have $`s`$ $`=`$ $`w_r1`$, since otherwise $`m`$ $``$ $`\beta `$, which contradicts $`m`$ $``$ $`I_\beta `$. But then $`x_r^{w_r}`$ $`|`$ $`m`$, and since it is clear that $`x_r^{w_r}`$ is a minimal generator of $`\beta `$, we must have $`m`$ $`=`$ $`x_r^{w_r}`$. As in Section 4.4, we have that every arrow having tail the corner monomial $`x_r^{w_r}`$ is a minimal $`x_r`$-standard arrow. None of these arrows can be advanced, since $`x_r^{w_r}`$ is the only minimal generator that is divisible by $`x_r`$, and hence is the only minimal generator that can serve as the tail of an $`x_r`$-standard arrow; this completes the proof of the lemma. ∎
### 6.3. A standard bunch $`𝒮`$ for $`\beta `$
Let $`\beta `$ be the thickening (18) of $`\beta _0`$. For the proof of our main result on thickenings, it is convenient to have available a standard bunch $`𝒮`$ for $`\beta `$ that is closely related to a standard bunch $`𝒮_0`$ for $`\beta _0`$. We proceed to construct such a set of arrows.
Let $`\pi `$ denote the projection map taking monomials in $`x_1`$, …, $`x_r`$ to monomials in $`x_1`$, …, $`x_{r1}`$ defined by
$$m\stackrel{\pi }{}m/x_r^{x_r\text{-deg}(m)}.$$
If $`c_𝐣^𝐝`$ is an arrow for $`\beta `$ such that $`\pi (𝐱^𝐝)`$ $``$ $`\beta _0`$, then it is clear that
$$\pi (c_𝐣^𝐝)=c_{\pi (𝐣)}^{\pi (𝐝)}$$
is an arrow for $`\beta _0`$ that we will call the projection of $`c_𝐣^𝐝`$. Furthermore, if we can translate $`c_𝐣^𝐝`$ to $`c_{𝐣_1}^{𝐝_1}`$ in such a way that the tails of the arrows in the translation path always project to monomials $``$ $`\beta _0`$, then the projections of the arrows define a translation path from $`\pi (c_𝐣^𝐝)`$ to $`\pi (c_{𝐣_1}^{𝐝_1})`$ that only involves steps in the directions $`x_1`$, …, $`x_{r1}`$. In this case we will say that the translation path projects from $`\beta `$ to $`\beta _0`$.
For each $`c_{𝐣_0}^{𝐝_0}`$ such that
$$x_r\text{-degree}(𝐱^{𝐝_0})=0\text{ and }x_r\text{-degree}(𝐱^{𝐣_0})=0$$
(that is, $`c_{𝐣_0}^{𝐝_0}`$ is an arrow for $`\beta _0`$), define
$$\text{tower}(c_{𝐣_0}^{𝐝_0})=\{c_𝐣^{}^{𝐝_0}𝐱^𝐣^{}=𝐱^{𝐣_0}x_r^s,0sw_r1\}.$$
It is then clear that any translation of $`c_{𝐣_0}^{𝐝_0}`$ in directions other than the $`x_r`$-th can be applied to the entire tower “in parallel”.
###### Lemma 6.3.1.
For $`x_i`$ $``$ $`x_r`$, let $`c_𝐣^𝐝`$ be an $`x_i`$-standard arrow for $`\beta `$ that can be advanced. Then the projection
$$\pi (c_𝐣^𝐝)=c_{𝐣_0}^{𝐝_0}$$
is an $`x_i`$-standard arrow for $`\beta _0`$ that can be advanced; consequently, all the arrows in *tower*$`(c_{𝐣_0}^{𝐝_0})`$ can be advanced in parallel.
*Proof:* We have that
$$x_r\text{-degree}(𝐱^𝐝)x_r\text{-degree}(𝐱^𝐣)<r,$$
so the tail $`𝐱^𝐝`$ must be divisible by a minimal generator of $`I_{\beta _0}`$, by Lemma 6.2.1. Therefore, the projection $`c_{𝐣_0}^{𝐝_0}`$ is defined. Since the vector $`𝐣_0`$ $``$ $`𝐝_0`$ is the same as the original vector $`𝐣`$ $``$ $`𝐝`$ except in the $`x_r`$-component (which is $`0`$ for the former and $``$ $`0`$ for the latter), it follows easily that $`c_{𝐣_0}^{𝐝_0}`$ is $`x_i`$-standard. Furthermore, the translation path that advances $`c_𝐣^𝐝`$ consists entirely of $`x_i`$-standard arrows, so it projects from $`\beta `$ to $`\beta _0`$, implying that the projection $`c_{𝐣_0}^{𝐝_0}`$ can be advanced (as an arrow for $`\beta _0`$). As observed earlier, all the arrows in tower$`(c_{𝐣_0}^{𝐝_0})`$ then advance in parallel. ∎
We now select a particular standard bunch of arrows $`𝒮`$ for $`\beta `$ (Theorem 4.1.3 ensures that at least one such set exists). Recall that $`𝒮`$ can be constructed as the union of standard $`x_i`$-sub-bunches $`𝒮(i)`$, $`1`$ $``$ $`i`$ $``$ $`r`$; each $`𝒮(i)`$ consists of $`n`$ minimal standard arrows for $`x_i`$ that cannot be advanced. By the last statement of Lemma 6.2.1, together with Corollary 4.6.2, we have that $`𝒮(r)`$ is the set of all $`n`$ arrows having tail $`x_r^{w_r}`$. For any $`x_i`$ $``$ $`x_r`$, we can construct $`𝒮(i)`$ as in Subsection 4.5, taking care to translate all towers in parallel, using Lemma 6.3.1. That is, if $`c_𝐣^{𝐝_0}`$ is a *minimal* $`x_i`$-standard arrow that can be advanced, (so the tail $`𝐱^{𝐝_0}`$ is a minimal generator of $`\beta _0`$, by Lemma 6.2.1), then the projection
$$\pi (c_𝐣^{𝐝_0})=c_{𝐣_0}^{𝐝_0}$$
can be advanced (as an arrow for $`\beta _0`$); whence, all the arrows in tower$`(c_{𝐣_0}^{𝐝_0})`$ (which contains the original arrow $`c_𝐣^{𝐝_0}`$) can be advanced in parallel. With this understanding, we see that the construction (for $`\beta `$) of an $`x_i`$-standard sub-bunch $`𝒮(i)`$ can be carried out as follows: First, restricting to the basis set $`\beta _0`$ and the variables $`x_1`$, …, $`x_{r1}`$, choose a standard $`x_i`$-sub-bunch $`𝒮_0(i)`$, as in the proof of Theorem 4.1.3. Then set
$$𝒮(i)=\underset{c_{𝐣_0}^𝐛𝒮_0(i)}{}\text{tower}(c_{𝐣_0}^𝐛).$$
Hence, by setting
$$𝒮_0=_{i=1}^{r1}𝒮_0(i),𝒮=_{i=1}^r𝒮(i),$$
we obtain a standard bunch $`𝒮`$ for $`\beta `$ that “lies over” a standard bunch $`𝒮_0`$ for $`\beta _0`$ in the sense that every arrow in $`𝒮(i)`$ projects to an arrow in $`𝒮_0(i)`$ for $`1`$ $``$ $`i`$ $``$ $`r1`$.
### 6.4. Thickenings of smooth basis sets are smooth
###### Theorem 6.4.1.
If $`\beta _0`$ is a smooth basis set in $`r1`$ variables, then for each integer $`w_r`$ $`>`$ $`0`$, the thickening $`\beta `$ *(18)* is a smooth basis set.
*Proof:* It suffices to show that $`\beta `$ satisfies conditions (a) and (b) of Theorem 5.1.1. First suppose that $`c_𝐣^𝐝`$ is a non-standard arrow for $`\beta `$. If $`x_r^{w_r}`$ $`|`$ $`𝐱^𝐝`$, then we can translate the arrow $`c_𝐣^𝐝`$ in degree-decreasing steps so that its tail approaches $`x_r^{w_r}`$; since no non-standard arrow can have tail $`x_r^{w_r}`$, we see that at some point the head of the arrow must exit the first orthant, demonstrating that $`c_𝐣^𝐝`$ $``$ $`0`$. On the other hand, if $`x_r^{w_r}`$ $``$ $`𝐱^𝐝`$, then we can assume by Lemmas 4.1.2 and 6.2.1 that we have translated our arrow to $`c_𝐣^{}^𝐛`$, with tail a minimal generator $`𝐱^𝐛`$ of $`I_{\beta _0}`$ (we are done if the head of the arrow exits the first orthant during this translation). Let $`c_{𝐣_0}^𝐛`$ = $`\pi (c_𝐣^{}^𝐛)`$, and observe that $`c_{𝐣_0}^𝐛`$ is a non-standard arrow for $`\beta _0`$. Since $`\beta _0`$ is a smooth basis set, by hypothesis, we have that $`c_{𝐣_0}^𝐛`$ $``$ $`0`$, using only translations in the first $`r1`$ variable directions. It is clear that the arrow $`c_𝐣^{}^𝐛`$ translates “in parallel” with $`c_{𝐣_0}^𝐛`$, leading to the conclusion that $`c_𝐣^{}^𝐛`$ $``$ $`0`$. Therefore $`\beta `$ satisfies condition (a) of Theorem 5.1.1.
To prove that $`\beta `$ satisfies condition (b) of Theorem 5.1.1, we use the standard bunches $`𝒮`$ and $`𝒮_0`$ constructed in Subsection 6.3, and we let $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ be an $`x_i`$-standard arrow for $`\beta `$. If $`i`$ $`=`$ $`r`$, then we can translate in degree-reducing steps until we reach an arrow whose tail is the corner monomial $`𝐱^{𝐝_1}`$ $`=`$ $`x_r^{w_r}`$, so that $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`𝒮(r)`$ $``$ $`𝒮`$. If $`i`$ $``$ $`r`$, then we can translate by degree-reducing steps until we reach an arrow $`c_{𝐣_1}^𝐛`$ whose tail is one of the minimal generators $`𝐱^𝐛`$ of $`I_{\beta _0}`$. Let $`c_{𝐣_0}^𝐛`$ $`=`$ $`\pi (c_{𝐣_1}^𝐛)`$, and note that $`c_{𝐣_1}^𝐛`$ $``$ tower$`(c_{𝐣_0}^𝐛)`$. Since $`\beta _0`$ is a smooth basis set, we know that there exists an arrow $`c_{𝐣_0^{}}^𝐛^{}`$ $``$ $`𝒮_0(i)`$ such that $`c_{𝐣_0}^𝐛`$ $``$ $`c_{𝐣_0^{}}^𝐛^{}`$. The translations only involve the variable directions $`x_1`$, …, $`x_{r1}`$, and will carry $`c_{𝐣_1}^𝐛`$ “in parallel” to an arrow $`c_{𝐣_1^{}}^𝐛^{}`$ $``$ tower$`(c_{𝐣_0^{}}^𝐛^{})`$ $``$ $`𝒮(i)`$. We conclude that $`\beta `$ satisfies condition (b) of Theorem 5.1.1. This completes the proof that the thickening $`\beta `$ of $`\beta _0`$ is a smooth basis set. ∎
###### Remark 6.4.2.
The converse of Theorem 6.4.1 also holds; that is, if $`\beta `$ is a smooth basis set that is a thickening (18) of $`\beta _0`$, then $`\beta _0`$ is smooth. We leave this as an exercise for the reader.
### 6.5. Example: “Boxes” are smooth basis sets
Let $`w_1`$, …, $`w_r`$ be positive integers, and let
(20)
$$(w_1,\mathrm{},w_r)=\{x_1^{d_1}x_2^{d_2}\mathrm{}x_r^{d_r}0d_i<w_i,1ir\},$$
which is clearly a basis set containing
$$|(w_1,\mathrm{},w_r)|=n=\underset{i=1}{\overset{r}{}}w_i$$
monomials; for obvious reasons, we call this type of basis set a box. The following useful results are immediate:
###### Lemma 6.5.1.
The minimal generators of $`I_{(w_1,\mathrm{},w_r)}`$ are the corner monomials $`x_i^{w_i}`$, $`1`$ $``$ $`i`$ $``$ $`r`$. ∎
###### Lemma 6.5.2.
If $`m_1`$, $`m_2`$ $``$ $`(w_1,\mathrm{},w_r)`$, then
$$\mathrm{lcm}(m_1,m_2)(w_1,\mathrm{},w_r).\text{}$$
We could give an easy proof of the next proposition using Lemma 6.5.1, but instead we offer a proof based on thickenings.
###### Proposition 6.5.3.
The box $`(w_1,\mathrm{},w_r)`$ is a smooth basis set.
*Proof:* It is clear that $`(w_1,\mathrm{},w_r)`$ is a thickening of $`(w_1,\mathrm{},w_{r1})`$, so if the latter is smooth, then so is the former, by Theorem 6.4.1. The desired result therefore follows by induction, provided that the result holds in the base case: boxes $`\beta `$ = $`(w_1)`$ in one variable. But in this case, the monomial ideal $`I_\beta `$ has just one minimal generator $`x_1^{w_1}`$; therefore, there are exactly $`w_1`$ $`=`$ $`|\beta |`$ minimal arrows, all of which are standard and cannot be advanced; in fact, the set of minimal arrows is equal to the standard bunch $`𝒮`$ $`=`$ $`𝒮(1)`$, in the notation of the proof of Theorem 4.1.3. Since $`M/M^2`$ is spanned by minimal arrows (Lemma 4.1.2), we conclude that $`𝒮`$ is a $`𝗄`$-basis of $`M/M^2`$ $``$ $`\beta `$ is smooth, as desired. ∎
## 7. Truncation of basis sets
Truncation is another natural way to obtain a basis set from a given basis set. In this section we will define the truncation operation and establish one of its key properties: If $`\beta `$ is such that every non-standard arrow $`c_𝐣^𝐝`$ is translation-equivalent to $`0`$ (that is, $`\beta `$ satisfies condition (a) of Theorem 5.1.1), then any truncation of $`\beta `$ also has this property; in particular, this is so if $`\beta `$ is smooth. Under an additional hypothesis, one can further prove that certain truncations of a smooth basis set are smooth (that is, also satisfy condition (b) of Theorem 5.1.1); we discuss this in the next section.
### 7.1. Definition of truncation
Let $`\beta `$ be a basis set of $`n`$ monomials in the variables $`x_1`$, …, $`x_r`$. Choose one of the variables $`x_j`$, and a positive integer $`h`$ such that $`\beta `$ contains at least one monomial that is divisible by $`x_j^h`$. Then we define the $`x_j`$-truncation of $`\beta `$ at height $`h`$ to be the basis set
(21)
$$\beta _t(x_j,h)=\beta _t=\{mx_j^hm\beta \}.$$
That is, $`\beta _t`$ is obtained by discarding the monomials in $`\beta `$ with $`x_j`$-degree $`<`$ $`h`$, and dividing the remaining monomials by $`x_j^h`$. Figure 4 provides an example. We write
$$n=|\beta |,n_t=|\beta _t|$$
for the number of monomials in $`\beta `$, $`\beta _t`$, respectively.
### 7.2. Minimal generators of $`I_{\beta _t}`$
We denote by $`w_i`$ (resp. $`w_i^{(t)}`$) the $`x_i`$-width of $`\beta `$ (resp. $`\beta _t`$); that is, $`x_i^{w_i}`$ (resp. $`x_i^{w_i^{(t)}}`$) is the $`i`$-th corner monomial of $`I_\beta `$ (resp. $`I_{\beta _t}`$). We have the following
###### Lemma 7.2.1.
Let $`\beta `$ be a basis set, and $`\beta _t`$ the $`x_j`$-truncation of $`\beta `$ at height $`h`$. Then
1. If $`m`$ is a minimal generator of $`I_\beta `$ such that $`x_j`$-degree$`(m)`$ $``$ $`h`$, then $`m/x_j^h`$ $`=`$ $`m_t`$ is a minimal generator of $`I_{\beta _t}`$; consequently, for all $`x_k`$ $``$ $`x_j`$, we have that
$$x_k\text{-degree}(m)=x_k\text{-degree}(m_t)w_k^{(t)}.$$
2. The minimal generators $`m`$ of $`I_\beta `$ that have $`x_j`$-degree $`>`$ $`h`$ are in bijective correspondence with the minimal generators $`m_t`$ $`=`$ $`m/x_j^h`$ of $`I_{\beta _t}`$ that are divisible by $`x_j`$.
3. If $`m_t^{}`$ is a minimal generator of $`I_{\beta _t}`$ that is not divisible by $`x_j`$, then $`x_j^hm_t^{}`$ $`=`$ $`m^{}`$ is an $`x_j`$-multiple of a minimal generator of $`I_\beta `$.
4. $`w_j`$ $`=`$ $`w_j^{(t)}+h`$ $`>`$ $`h`$, and for all $`x_k`$ $``$ $`x_j`$, $`w_k`$ $``$ $`w_k^{(t)}`$.
*Proof:* We begin by noting that a monomial $`m`$ $``$ $`I_\beta `$ (resp. $`I_{\beta _t}`$) is a minimal generator if and only if for each variable $`x_i`$, either $`x_i`$ $``$ $`m`$ or $`m/x_i`$ $``$ $`\beta `$ (resp. $`\beta _t`$).
Let $`m`$ and $`m_t`$ be as in assertion (a). Since $`x_j^hm_t`$ $`=`$ $`m`$ $``$ $`\beta `$, we have that $`m_t`$ $``$ $`\beta _t`$ $``$ $`m_t`$ $``$ $`I_{\beta _t}`$. For any variable $`x_k`$ $``$ $`x_j`$, note that
$$x_k|m_tx_k|mm/x_k\beta ;$$
moreover,
$$x_j\text{-degree}(m/x_k)=x_j\text{-degree}(m)h;$$
therefore, $`m_t/x_k`$ $``$ $`\beta _t`$. If $`x_j`$ $`|`$ $`m_t`$, then
$$x_j\text{-degree}(m)=\left(x_j\text{-degree}(m_t)+h\right)>h;$$
whence,
$$x_j\text{-degree}(m/x_j)h,\text{ and as before }m/x_j\beta ,$$
which yields $`m_t/x_j`$ $``$ $`\beta _t`$. This completes the proof that $`m_t`$ is a minimal generator of $`I_{\beta _t}`$, and the stated consequence is immediate.
Since the map $`m`$ $``$ $`m_t`$ is injective, to prove assertion (b), it remains to show that if $`m_t^{}`$ is a minimal generator of $`\beta _t`$ that is divisible by $`x_j`$, then $`x_j^hm_t^{}`$ $`=`$ $`m^{}`$ is a minimal generator of $`\beta `$. It is clear that $`m^{}`$ $``$ $`I_\beta `$, since $`m_t^{}`$ $``$ $`\beta _t`$. Furthermore, for any variable $`x_i`$, we have that
(22)
$$x_i|m^{}x_i|m_t^{}m_t^{}/x_i\beta _tm^{}/x_i\beta $$
(the first implication is trivial for $`x_i`$ $`=`$ $`x_j`$), so $`m^{}`$ is indeed a minimal generator of $`\beta `$.
Let now $`m_t^{}`$, $`m^{}`$ be as in assertion (c). As in the preceding paragraph, $`m^{}`$ $``$ $`I_\beta `$. We can translate $`m^{}`$ to a minimal generator $`m`$ of $`I_\beta `$ by degree-reducing steps, but since the implications (22) hold for all $`x_i`$ $``$ $`x_j`$, we see that we can only move in the direction of decreasing $`x_j`$-degree; whence, $`m^{}`$ is an $`x_j`$-multiple of $`m`$, as desired.
Turning to assertion (d), it is clear that
$$\begin{array}{ccc}\hfill x_{\mathrm{}}|x_j^{w_j^{(t)}+h}& & x_{\mathrm{}}=x_j\hfill \\ & & x_j^{w_j^{(t)}+h}/x_{\mathrm{}}=x_j^{w_j^{(t)}+h1}=x_j^hx_j^{w_j^{(t)}1}x_j^h\beta _t\beta ,\hfill \end{array}$$
which implies that $`x_j^{w_j^{(t)}+h}`$ is a minimal generator of $`I_\beta `$; therefore, $`w_j`$ $`=`$ $`w_j^{(t)}+h`$. Finally, for any $`x_k`$ $``$ $`x_j`$, we have that
$$\begin{array}{ccc}\hfill x_k^{w_k^{(t)}1}\beta _t& & x_j^hx_k^{w_k^{(t)}1}\beta \hfill \\ & & x_k^{w_k}(x_j^hx_k^{w_k^{(t)}1})\hfill \\ & & w_kw_k^{(t)},\hfill \end{array}$$
and the proof is complete. ∎
### 7.3. Lifting arrows from $`\beta _t`$ to $`\beta `$
Let $`\beta `$ and $`\beta _t`$ be as above, and let $`c_𝐣^𝐝`$ be an arrow for $`\beta _t`$. Since $`𝐱^𝐣`$ $``$ $`\beta _t`$ (resp. $`𝐱^𝐝`$ $``$ $`\beta _t`$), we have that $`𝐱^𝐣^{}`$ $`=`$ $`𝐱^𝐣x_j^h`$ $``$ $`\beta `$ (resp. $`𝐱^𝐝^{}`$ $`=`$ $`𝐱^𝐝x_j^h`$ $``$ $`\beta `$); therefore, $`c_𝐣^{}^𝐝^{}`$ is an arrow for $`\beta `$. We will call $`c_𝐣^{}^𝐝^{}`$ the lifting of $`c_𝐣^𝐝`$ from $`\beta _t`$ to $`\beta `$, and say that $`c_𝐣^{}^𝐝^{}`$ descends to $`c_𝐣^𝐝`$. We extend this terminology to sets of arrows in the obvious way. It is clear that the arrows $`c_𝐣^𝐝`$ and $`c_𝐣^{}^𝐝^{}`$ have the same vector
$$𝐣𝐝=𝐣^{}𝐝^{}.$$
### 7.4. Non-standard arrows on $`\beta `$ and $`\beta _t`$
###### Theorem 7.4.1.
Let $`\beta `$ be a basis set with the property that every one of its non-standard arrows is translation-equivalent to $`0`$ (condition *(a)* of Theorem *5.1.1*), and let $`\beta _t`$ be the $`x_j`$-truncation of $`\beta `$ at height $`h`$ *(21)*. Then $`\beta _t`$ also has this property; that is, every non-standard arrow for $`\beta _t`$ is translation-equivalent to $`0`$.
*Proof:* Let $`c_𝐣^𝐝`$ be a non-standard arrow for $`\beta _t`$. We must show that $`c_𝐣^𝐝`$ $``$ $`0`$. To do this, we consider the lifting $`c_𝐣^{}^𝐝^{}`$ of $`c_𝐣^𝐝`$ to $`\beta `$. Suppose first that
$$x_j\text{-degree}(𝐱^𝐝)x_j\text{-degree}(𝐱^𝐣);$$
that is, the $`j`$-th component of the vector $`𝐣`$ $``$ $`𝐝`$ is $``$ $`0`$. By hypothesis, we can translate the arrow $`c_𝐣^{}^𝐝^{}`$ so that the head eventually leaves the first orthant. In fact, we *claim* that we can perform this translation without ever passing through a position $`c_{𝐣_1}^{𝐝_1}`$ for which
$$x_j\text{-height}(c_{𝐣_1}^{𝐝_1})<x_j\text{-height}(c_𝐣^{}^𝐝^{})=x_j\text{-degree}(𝐱^𝐝^{}).$$
*Proof of claim:* We argue by induction on the $`x_j`$-height of the original arrow $`c_𝐣^{}^𝐝^{}`$. The case of height $`0`$ is trivial: no arrows can have $`x_j`$-height $`<`$ $`0`$. So suppose that
$$x_j\text{-height}(c_𝐣^{}^𝐝^{})=p>0,$$
and that the claim holds for all heights $`<`$ $`p`$. Consider a sequence of steps that will translate the arrow’s head out of the first orthant. Let $`c_{𝐣_0}^{𝐝_0}`$ be the first point on this path (if any) from which translation to an $`x_j`$-height $`<`$ $`p`$ is possible; that is,
$$x_j\text{-degree}(𝐱^{𝐝_0})=p,\text{ and }𝐱^{𝐝_0}/x_jI_\beta .$$
Let $`s`$ $``$ $`1`$ be the largest integer such that
$$𝐱^𝐛=𝐱^{𝐝_0}/x_j^sI_\beta .$$
Then the arrow $`c_{𝐣_0}^𝐛`$ is non-standard with the same negative and non-negative components in its vector as the original arrow (the only change in the vector is an increased $`j`$-th component, which was $``$ $`0`$ to begin with). In light of the hypothesis on $`\beta `$, and our induction hypothesis, we can translate the arrow $`c_{𝐣_0}^𝐛`$ (of $`x_j`$-height $`q`$ $`<`$ $`p`$) so that the head exits the first orthant without ever reaching a position of height $`<`$ $`q`$. It is now clear that the arrow $`c_{𝐣_0}^{𝐝_0}`$ can be translated “in parallel” with the arrow $`c_{𝐣_0}^𝐛`$, and the height of the former never becomes $`<`$ $`p`$. This shows that the claim holds for any non-standard arrow at $`x_j`$-height $`p`$, which completes the proof of the claim.
Returning to the proof of Theorem 7.4.1, we now see that the lifted arrow $`c_𝐣^{}^𝐝^{}`$ can be translated so that the head leaves the first orthant and no position on the path has $`x_j`$-height less than the original height (which is $``$ $`h`$). By descending the path to $`\beta _t`$, we obtain that the original arrow $`c_𝐣^𝐝`$ is translation-equivalent to $`0`$ for the basis set $`\beta _t`$.
It remains to show that the same conclusion holds when the $`j`$-th component of the vector of our non-standard arrow $`c_𝐣^𝐝`$ is $`<`$ $`0`$; that is,
$$x_j\text{-degree}(𝐱^𝐝)>x_j\text{-degree}(𝐱^𝐣).$$
The lifted arrow $`c_𝐣^{}^𝐝^{}`$ is by hypothesis translation-equivalent to $`0`$ on $`\beta `$. Note that any step $`c_{𝐣_1}^{𝐝_1}`$ on the corresponding translation path such that
$$x_j\text{-degree}(𝐱^{𝐣_1})h$$
descends to $`\beta _t`$. Therefore, by descending the initial segment of the path, up to the first position $`c_{𝐣_1}^{𝐝_1}`$ for which the $`x_j`$-degree$`(𝐱^{𝐣_1})`$ $`<`$ $`h`$ (if any), we conclude that $`c_𝐣^𝐝`$ is translation-equivalent to $`0`$ for $`\beta _t`$, and we are done. ∎
## 8. Sufficient conditions for a truncation to be smooth
Let $`\beta `$ be a *smooth* basis set in the variables $`x_1,\mathrm{},x_r`$, and $`\beta _t`$ the $`x_j`$-truncation of $`\beta `$ at height $`h`$ (21). In the last section (Theorem 7.4.1) we saw that $`\beta _t`$ necessarily satisfies condition (a) of Theorem 5.1.1. Our main goal in this section is to prove Theorem 8.6.1, which states that $`\beta _t`$ will also satisfy condition (b) of Theorem 5.1.1, and therefore be smooth, if we assume an additional hypothesis. Much of our work involves the construction of a (near-)standard $`x_i`$-sub-bunch of arrows for $`\beta `$ that contains the lifting of a standard $`x_i`$-sub-bunch of arrows for $`\beta _t`$. There are two cases: $`x_i`$ = $`x_j`$ (the easier case), discussed in Subsection 8.2, and $`x_i`$ $`=`$ $`x_k`$ $``$ $`x_j`$, discussed in Subsection 8.5.
### 8.1. The additional hypothesis
By the first assertion of Lemma 7.2.1, we know that if $`m`$ is any minimal generator of $`I_\beta `$, then
$$x_kx_j,x_j\text{-degree}(m)hx_k\text{-degree}(m)w_k^{(t)}.$$
To prove Theorem 8.6.1, we need to control $`x_k`$-degree$`(m)`$ when $`x_j`$-degree($`m`$) $`<`$ $`h`$, using the following
###### Hypothesis 8.1.1.
For every minimal generator $`m`$ of $`I_\beta `$, we have that
$$x_kx_j,x_j\text{-degree}(m)<h\text{and }x_k|mx_k\text{-degree}(m)w_k^{(t)}.$$
### 8.2. $`x_j`$-sub-bunches of arrows for $`\beta `$ and $`\beta _t`$
###### Lemma 8.2.1.
Let $`\beta `$ be an arbitrary (not necessarily smooth) basis set such that $`\beta _t`$, the $`x_j`$-truncation at height $`h`$ (21), is defined. Let $`𝒮(j)`$ be a standard $`x_j`$-sub-bunch of arrows for $`\beta `$, constructed as in the proof of Theorem 4.1.3. Then the set of arrows $`c_𝐣^𝐝`$ $``$ $`𝒮(j)`$ with heads $`𝐱^𝐣`$ divisible by $`x_j^h`$ (so that $`𝐱^𝐣/x_j^h`$ $``$ $`\beta _t`$) is the lifting of a standard $`x_j`$-sub-bunch of arrows $`𝒮^{(t)}(j)`$ for $`\beta _t`$.
*Proof:* Consider a minimal $`x_j`$-standard arrow $`c_𝐣^𝐝`$ for $`\beta `$ such that
$$x_j\text{-degree}(𝐱^𝐣)hx_j\text{-degree}(𝐱^𝐝)>h,$$
and suppose that $`c_𝐣^𝐝`$ can be advanced. We can then promote the entire $`x_j`$-shadow of $`c_𝐣^𝐝`$, as described in Subsection 4.3; the promotion image is the $`x_j`$-shadow of an arrow $`c_{𝐣_2}^𝐛`$, where $`𝐱^𝐛`$ is a minimal generator of $`I_\beta `$,
$$x_j\text{-degree}(𝐱^𝐛)>x_j\text{-degree}(𝐱^{𝐣_2})=x_j\text{-degree}(𝐱^𝐣)h$$
and
$$x_j\text{-degree}(𝐱^𝐛)<x_j\text{-degree}(𝐱^𝐝).$$
Let $`c_{𝐣_t}^{𝐝_t}`$ (resp. $`c_{𝐣_{2,t}}^{𝐛_t}`$) denote the arrow that $`c_𝐣^𝐝`$ (resp. $`c_{𝐣_2}^𝐛`$) descends to (on the truncation $`\beta _t`$). Lemma 7.2.1 implies that $`𝐱^{𝐝_t}`$ and $`𝐱^{𝐛_t}`$ are minimal generators of $`I_{\beta _t}`$; moreover, it is clear that the $`x_j`$-shadow of $`c_{𝐣_{2,t}}^{𝐛_t}`$ is the promotion image of the shadow of $`c_{𝐣_t}^{𝐝_t}`$, and that the lifting of the shadow of $`c_{𝐣_t}^{𝐝_t}`$ (resp. $`c_{𝐣_{2,t}}^{𝐛_t}`$) to $`\beta `$ consists of the arrows in the shadow of $`c_𝐣^𝐝`$ (resp. $`c_{𝐣_2}^𝐛`$) whose heads have $`x_j`$-degree $``$ $`h`$. Applying this observation to the iterated shadow promotion operations used to construct $`𝒮(j)`$, as described in Subsection 4.4, one sees that the desired result follows readily. ∎
### 8.3. $`x_k`$-standard arrows of $`x_j`$-height $``$ $`h`$
Recall from Subsection 4.4 that $`𝒮_0(k,v)`$ denotes the set of all (minimal $`x_k`$-standard) arrows for $`\beta `$ having tail $`x_k^{w_k}`$ and offset $`v`$. Suppose that $`c_𝐣^𝐝`$ $``$ $`𝒮_0(k,v)`$ has head $`𝐱^𝐣`$ $``$ $`\beta _t`$. It is then clear that the entire $`x_k`$-shadow of $`c_𝐣^𝐝`$ has this property; therefore, the subset of all arrows in $`𝒮_0(k,v)`$ with heads in $`\beta _t`$ is equal to the $`x_k`$-shadow of $`c_{𝐣_v}^𝐝`$, the member of the subset whose head has maximal $`x_k`$-degree
(23)
$$r(k,v)=x_k\text{-degree}(𝐱^{𝐣_v}).$$
If no arrow in $`𝒮_0(k,v)`$ has head in $`\beta _t`$, we define $`r(k,v)`$ $`=`$ $`1`$, and $`𝐱^{𝐣_v}`$ is undefined. Recalling that the $`x_j`$-height of an arrow is the $`x_j`$-degree of its tail, we have the following
###### Lemma 8.3.1.
Let $`c_𝐣^𝐛`$ be an $`x_k`$-standard arrow for $`\beta `$ having $`x_j`$-height $``$ $`h`$ and offset $`v`$. Then
$$x_k\text{-degree}(𝐱^𝐣)r(k,v).$$
*Proof:* If not, then, since $`c_𝐣^𝐛`$ is $`x_k`$-standard, we have that
$$\begin{array}{c}x_j\text{-degree}(𝐱^𝐣)x_j\text{-degree}(𝐱^𝐛)h,\text{and}\\ x_i\text{-degree}(𝐱^𝐣)x_i\text{-degree}(𝐱^𝐛),ij,ik.\end{array}$$
Let
$$p=x_k\text{-degree}(𝐱^𝐛),m=𝐱^𝐛/x_k^p,\text{ and }𝐱^{𝐣_1}=𝐱^𝐣/m\beta .$$
Then we have that $`𝐱^{𝐣_1}`$ has offset $`v`$ from the corner monomial $`x_k^{w_k}`$ $`=`$ $`𝐱^𝐝`$, and
$$x_k\text{-degree}(𝐱^{𝐣_1})=x_k\text{-degree}(𝐱^𝐣).$$
We further have that $`𝐱^{𝐣_1}`$ $``$ $`\beta _t`$ because
$$x_j^h|m(x_j^h𝐱^{𝐣_1})|(m𝐱^{𝐣_1}),\text{ and }m𝐱^{𝐣_1}=𝐱^𝐣\beta ,\text{so }x_j^h𝐱^{𝐣_1}\beta .$$
It then follows from the definition (23) of $`r(k,v)`$ that
$$x_k\text{-degree}(𝐱^𝐣)=x_k\text{-degree}(𝐱^{𝐣_1})r(k,v),$$
as desired. ∎
### 8.4. Linear independence of lifts of $`x_k`$-sub-bunches
###### Lemma 8.4.1.
Let $`\beta `$ and $`\beta _t`$ be as above, and let $`x_k`$ $``$ $`x_j`$. Suppose in addition that the specialization of Hypothesis *8.1.1* to $`x_k`$ holds; that is, for every minimal generator $`m`$ of $`I_\beta `$, we have that
(24)
$$x_j\text{-degree}(m)<h\text{and }x_k|mx_k\text{-degree}(m)w_k^{(t)}.$$
Then, given two $`x_k`$-standard arrows (for $`\beta `$) $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ of $`x_j`$-height $``$ $`h`$, there is a translation path from $`c_{𝐣_1}^{𝐝_1}`$ to $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ such that every arrow in the path has $`x_j`$-height $``$ $`h`$.
*Proof:* Consider a translation path from $`c_{𝐣_1}^{𝐝_1}`$ to $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ that at some point involves an arrow of $`x_j`$-height $`<`$ $`h`$, and let $`c_{𝐣_2}^{𝐝_2}`$ be the arrow in the path that immediately precedes the very first such arrow; in particular, since the next step must be in the negative $`x_j`$-direction to reach $`x_j`$-height $`h1`$, we have that
$$𝐱^{𝐝_2}/x_jI_\beta ,\text{and }x_j\text{-degree}(𝐱^{𝐝_2}/x_j)=h1.$$
There exists a minimal generator $`m`$ of $`I_\beta `$ such that
$$m|(𝐱^{𝐝_2}/x_j)x_j\text{-degree}(m)<h.$$
Furthermore, we have that $`x_k`$ $`|`$ $`m`$, since the reasoning leading to the inequality (15) implies that
$$x_k\text{-degree}(m)>x_k\text{-degree}(𝐱^{𝐣_2})0,$$
otherwise we have the contradiction $`m`$ $`|`$ $`𝐱^{𝐣_2}`$ $``$ $`m`$ $``$ $`\beta `$.
The hypothesis (24) now implies that
$$w_k^{(t)}x_k\text{-degree}(m)x_k\text{-degree}(𝐱^{𝐝_2}/x_j)=x_k\text{-degree}(𝐱^{𝐝_2}).$$
Indeed, similar reasoning shows that any arrow in the path having $`x_j`$-height $`<`$ $`h`$ must have $`x_k`$-height $``$ $`w_k^{(t)}`$. Since the path terminates in the arrow $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ of $`x_j`$-height $``$ $`h`$, the path must eventually reach an arrow $`c_{𝐣_3}^{𝐝_3}`$ following $`c_{𝐣_2}^{𝐝_2}`$ such that the $`x_j`$-height of $`c_{𝐣_3}^{𝐝_3}`$ equals $`h`$, all subsequent arrows in the path have $`x_j`$-height $``$ $`h`$, and the arrow preceding $`c_{𝐣_3}^{𝐝_3}`$ has $`x_j`$-height $`=`$ $`h1`$, so that the step to $`c_{𝐣_3}^{𝐝_3}`$ is in the increasing $`x_j`$-direction. An argument similar to that for $`𝐱^{𝐝_2}`$ yields that $`x_k`$-degree$`(𝐱^{𝐝_3})`$ $``$ $`w_k^{(t)}`$. It follows that both $`𝐱^{𝐝_2}`$ and $`𝐱^{𝐝_3}`$ are divisible by $`𝐱^𝐛`$ = $`x_j^hx_k^{w_k^{(t)}}`$ $``$ $`I_\beta `$; therefore, each of the arrows $`c_{𝐣_2}^{𝐝_2}`$ and $`c_{𝐣_3}^{𝐝_3}`$ can be translated by degree-reducing steps to an arrow $`c_{𝐣_4}^𝐛`$. The heads of the arrows cannot leave the first orthant during these translations because, the arrows being $`x_k`$-standard (and of $`x_k`$-height $``$ $`w_k^{(t)}`$), the exit would have to be across the hyperplane ($`x_k`$-degree $`=`$ $`0`$), but then the original arrow $`c_{𝐣_1}^{𝐝_1}`$ of $`x_k`$-height $``$ $`w_k^{(t)}`$ could not have existed. We now replace the original path segment from $`c_{𝐣_2}^{𝐝_2}`$ to $`c_{𝐣_3}^{𝐝_3}`$ by the translation from $`c_{𝐣_2}^{𝐝_2}`$ to $`c_{𝐣_4}^𝐛`$ and the reversal of the translation from $`c_{𝐣_3}^{𝐝_3}`$ to $`c_{𝐣_4}^𝐛`$. Since the latter translations only involve arrows of $`x_j`$-height $`=`$ $`h`$, we have produced a path from $`c_{𝐣_1}^{𝐝_1}`$ to $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ that involves only arrows of $`x_j`$-height $``$ $`h`$, as desired. ∎
###### Corollary 8.4.2.
With the hypotheses of Lemma *8.4.1*, suppose given two $`x_k`$-standard arrows $`c_{𝐣_1}^{𝐝_1}`$ and $`c_{𝐣_2}^{𝐝_2}`$ for $`\beta _t`$, and let $`c_{𝐣_1^{}}^{𝐝_1^{}}`$, $`c_{𝐣_2^{}}^{𝐝_2^{}}`$ be the associated liftings to $`\beta `$. If $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ $``$ $`c_{𝐣_2^{}}^{𝐝_2^{}}`$, then $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_{𝐣_2}^{𝐝_2}`$. Furthermore, if $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ can be advanced, then $`c_{𝐣_1}^{𝐝_1}`$ can be advanced.
*Proof:* Since $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ and $`c_{𝐣_2^{}}^{𝐝_2^{}}`$ are $`x_k`$-standard and have $`x_j`$-height $``$ $`h`$, Lemma 8.4.1 implies that there is a translation path from $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ to $`c_{𝐣_2^{}}^{𝐝_2^{}}`$ consisting entirely of arrows of $`x_j`$-height $``$ $`h`$. But then this path descends to give the translation equivalence of $`c_{𝐣_1}^{𝐝_1}`$ and $`c_{𝐣_2}^{𝐝_2}`$, which proves the first assertion.
Note that the second assertion is trivially true if
$$x_k\text{-degree}(𝐱^{𝐝_1})>w_k^{(t)},$$
for then we can advance $`c_{𝐣_1}^{𝐝_1}`$ by moving its tail in degree-decreasing steps toward the minimal generator $`x_k^{w_k^{(t)}}`$. It therefore remains to show that we can advance $`c_{𝐣_1}^{𝐝_1}`$ provided that its lifting $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ can be advanced and
$$x_k\text{-degree}(𝐱^{𝐝_1^{}})=x_k\text{-degree}(𝐱^{𝐝_1})w_k^{(t)}.$$
Since we can advance $`c_{𝐣_1^{}}^{𝐝_1^{}}`$, we can translate it to an arrow $`c_{𝐣_3^{}}^{𝐝_3^{}}`$ such that
$$𝐱^{𝐝_3^{}}/x_kI_\beta \text{ and }x_k\text{-degree}(𝐱^{𝐝_3^{}})=x_k\text{-degree}(𝐱^{𝐝_1^{}}).$$
Then there is a minimal generator $`m`$ of $`I_\beta `$ such that
$$m|(𝐱^{𝐝_3^{}}/x_k)\text{ and }x_k\text{-degree}(m)>0,$$
where the inequality follows from (15).
If $`x_j`$-degree$`(m)`$ $`<`$ $`h`$, then hypothesis (24) implies that
$$\begin{array}{ccc}\hfill x_k\text{-degree}(m)w_k^{(t)}& & x_k\text{-degree}(𝐱^{𝐝_3^{}})>w_k^{(t)}\hfill \\ & & x_k\text{-degree}(𝐱^{𝐝_1^{}})>w_k^{(t)},\hfill \end{array}$$
which is a contradiction. We therefore have that
$$\begin{array}{ccc}\hfill x_j\text{-degree}(m)h& & x_j\text{-degree}(𝐱^{𝐝_3^{}})h;\hfill \end{array}$$
whence, Lemma 8.4.1 ensures that there is a translation path from $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ to $`c_{𝐣_3^{}}^{𝐝_3^{}}`$ consisting entirely of arrows of $`x_j`$-height $``$ $`h`$, but this path then descends to advance $`c_{𝐣_1}^{𝐝_1}`$, which proves the second assertion. ∎
###### Corollary 8.4.3.
Again with the hypotheses of Lemma *8.4.1*, we have that the lifting, to $`\beta `$, of a standard $`x_k`$-sub-bunch $`𝒮^{(t)}(k)`$ of arrows for $`\beta _t`$, has maximal rank *(mod $`M^2`$)*; furthermore, the lifted arrows cannot be advanced.
*Proof:* It suffices, by Theorem 3.2.1, to show that for any two distinct arrows $`c_{𝐣_1}^{𝐝_1}`$, $`c_{𝐣_2}^{𝐝_2}`$ $``$ $`𝒮^{(t)}(k)`$, with liftings $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ and $`c_{𝐣_2^{}}^{𝐝_2^{}}`$, we have that $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ $`\sim ̸`$ $`0`$ and $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ $`\sim ̸`$ $`c_{𝐣_2^{}}^{𝐝_2^{}}`$. Arguing by contradiction, suppose that $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ $``$ $`c_{𝐣_2^{}}^{𝐝_2^{}}`$. Then Corollary 8.4.2 implies that $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_{𝐣_2}^{𝐝_2}`$, which contradicts the hypothesis that $`𝒮^{(t)}(k)`$ is an $`x_k`$-standard sub-bunch. To prove that $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ $`\sim ̸`$ $`0`$, it suffices to prove, more generally, that $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ cannot be advanced. However, if this were false, then Corollary 8.4.2 would yield that $`c_{𝐣_1}^{𝐝_1}`$ can be advanced, which would again contradict the hypothesis that $`𝒮^{(t)}(k)`$ is an $`x_k`$-standard sub-bunch. ∎
### 8.5. $`x_k`$-sub-bunches of arrows for $`\beta `$ and $`\beta _t`$, $`x_k`$ $``$ $`x_j`$
We are now ready to prove
###### Lemma 8.5.1.
Let $`\beta _t`$ be the $`x_j`$-truncation of $`\beta `$ at height $`h`$, let $`x_k`$ $``$ $`x_j`$, and suppose that the specialization *(24)* of Hypothesis *8.1.1* to $`x_k`$ holds. Then there exists a near-standard $`x_k`$-sub-bunch $`𝒮^{}(k)`$ of arrows for $`\beta `$ that contains the lifting of a standard $`x_k`$-sub-bunch $`𝒮^{(t)}(k)`$ of arrows for $`\beta _t`$.
*Proof:* We begin by constructing a standard $`x_k`$-sub-bunch $`𝒮(k)`$ for $`\beta `$ as in the proof of Theorem 4.1.3. With $`r(k,v)`$ defined as in (23), we let
$$𝒮_1^{}(k)=\{c_{𝐣_1}^{𝐝_1}𝒮(k)x_k\text{-deg}(𝐱^{𝐣_1})>r(k,v),\text{where }v=x_k\text{-offset}(c_{𝐣_1}^{𝐝_1})\};$$
it is evident that $`𝒮_1^{}`$ is a maximal rank (mod $`M^2`$) set of $`x_k`$-standard arrows that cannot be advanced. It follows from Lemma 8.3.1 that the $`x_j`$-height of an arrow $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`𝒮_1^{}(k)`$ is $`<`$ $`h`$. Furthermore, the cardinality of this set is
$$|𝒮_1^{}(k)|=|\beta ||\beta _t|=nn_t,$$
because shadow promotion does not change the $`x_k`$-heights of the heads of the promoted arrows (recall Equation (17)), and the arrows in the original sets $`𝒮_0(k,v)`$ with heads of height $``$ $`r(k,v)`$ are precisely the arrows with heads in $`\beta _t`$.
We next construct a standard $`x_k`$-sub-bunch $`𝒮^{(t)}(k)`$ for $`\beta _t`$, and denote the lifting of this set to $`\beta `$ by $`𝒮_2^{}(k)`$; this set consists of
$$|𝒮_2^{}(k)|=|\beta _t|=n_t$$
$`x_k`$-standard arrows $`c_{𝐣_2}^{𝐝_2}`$ having $`x_j`$-height $``$ $`h`$. Corollary 8.4.3 implies that $`𝒮_2^{}(k)`$ has maximal rank (mod $`M^2`$), and that its arrows cannot be advanced. By Lemma 7.2.1, we know that the tail $`𝐱^{𝐝_2}`$ is a minimal generator of $`I_\beta `$ provided that its $`x_j`$-degree is $`>`$ $`h`$; however, if the $`x_j`$-degree $`=`$ $`h`$, we only know that $`𝐱^{𝐝_2}`$ is an $`x_j`$-multiple of a minimal generator; therefore, some of the arrows in $`𝒮_2^{}(k)`$ may not be minimal.
The desired set is
$$𝒮^{}(k)=𝒮_1^{}(k)𝒮_2^{}(k).$$
By comparing $`x_j`$-heights of arrows, we see that the two sets in the union do not overlap; therefore,
$$|𝒮^{}(k)|=|𝒮_1^{}(k)|+|𝒮_2^{}(k)|=(nn_t)+n_t=n=|\beta |.$$
The arrows in $`𝒮^{}(k)`$ are $`x_k`$-standard arrows that cannot be advanced, but need not be minimal. To show that $`𝒮^{}(k)`$ is a near-standard $`x_k`$-sub-bunch, it remains to show that it has maximal rank (mod $`M^2`$). Since $`𝒮_1^{}(k)`$ and $`𝒮_2^{}(k)`$ each have maximal rank (mod $`M^2`$) and consist of $`x_k`$-unadvanceable arrows, it suffices (by Theorem 3.2.1) to prove that no arrow in $`𝒮_1^{}(k)`$ can be translated to an arrow in $`𝒮_2^{}(k)`$. So let $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`𝒮_1^{}(k)`$ have $`x_k`$-offset $`v`$. By definition we have that $`x_k`$-degree$`(𝐱^{𝐣_1})`$ $`>`$ $`r(k,v)`$. If this arrow could be translated to an arrow $`c_{𝐣_2}^{𝐝_2}`$ of $`x_j`$-height $``$ $`h`$, then Lemma 8.3.1 would yield that $`x_k`$-degree$`(𝐱^{𝐣_2})`$ $``$ $`r(k,v)`$, implying that $`c_{𝐣_1}^{𝐝_1}`$ can be advanced, which is a contradiction. Therefore, $`𝒮^{}(k)`$ is a near-standard $`x_k`$-sub-bunch that contains the lifting of a standard $`x_k`$-sub-bunch for $`\beta _t`$, as desired. ∎
### 8.6. Main theorem on truncations
###### Theorem 8.6.1.
Let $`\beta `$ be a smooth basis set, and $`\beta _t`$ the $`x_j`$-truncation of $`\beta `$ at height $`h`$. If in addition Hypothesis *8.1.1* holds, then $`\beta _t`$ is a smooth basis set.
*Proof:* To show that $`\beta _t`$ is a smooth basis set, it suffices to prove that it satisfies conditions (a) and (b) of Theorem 5.1.1. Theorem 7.4.1 ensures that $`\beta _t`$ satisfies condition (a) — that every non-standard arrow is translation-equivalent to 0 — because $`\beta `$, smooth by hypothesis, has this property; therefore, it remains to show that $`\beta _t`$ satisfies condition (b).
Let $`𝒮(j)`$ be a standard $`x_j`$-sub-bunch of arrows for $`\beta `$, and $`𝒮^{(t)}(j)`$ the associated $`x_j`$-sub-bunch for $`\beta _t`$ whose lifting to $`\beta `$ lies in $`𝒮(j)`$, as in Lemma 8.2.1. For each variable $`x_k`$ $``$ $`x_j`$, let $`𝒮^{}(k)`$ be a near-standard $`x_k`$-sub-bunch of arrows for $`\beta `$ constructed as in Lemma 8.5.1, and $`𝒮^{(t)}(k)`$ the associated $`x_k`$-sub-bunch for $`\beta _t`$ whose lifting to $`\beta `$ lies in $`𝒮^{}(k)`$. Taking unions, we obtain a near-standard bunch $`𝒮^{}`$ for $`\beta `$ that contains the lifting of a standard bunch $`𝒮^{(t)}`$ for $`\beta _t`$. To complete the proof, it suffices to show that if $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ is a standard arrow for $`\beta _t`$, then $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`𝒮^{(t)}`$.
By assertion (c) of Lemma 4.2.1, we may assume that $`c_𝐣^𝐝`$ is a minimal $`x_i`$-standard arrow for $`\beta _t`$ that cannot be advanced. Consider the lifting $`c_𝐣^{}^𝐝^{}`$ of $`c_𝐣^𝐝`$ to $`\beta `$. We have that $`c_𝐣^{}^𝐝^{}`$ cannot be advanced: Indeed, if $`x_i`$ $`=`$ $`x_j`$, then one checks easily that a translation path advancing $`c_𝐣^{}^𝐝^{}`$ would descend to a translation path advancing $`c_𝐣^𝐝`$, a contradiction. The same contradiction arises in case $`x_i`$ $`=`$ $`x_k`$ $``$ $`x_j`$ by Corollary 8.4.2. Since $`\beta `$ is assumed smooth, we know that
$$c_𝐣^{}^𝐝^{}c_{𝐣_1^{}}^{𝐝_1^{}}𝒮^{}.$$
Furthermore, since neither of the arrows $`c_𝐣^{}^𝐝^{}`$, $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ can be advanced, we see that these arrows must have the same $`x_i`$-height. If $`x_i`$ $`=`$ $`x_j`$, we obtain that $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ lies in the lifting of $`𝒮^{(t)}(j)`$, and the translation path from $`c_𝐣^{}^𝐝^{}`$ to $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ descends to yield
$$c_𝐣^𝐝c_{𝐣_1}^{𝐝_1}𝒮^{(t)}(j).$$
If $`x_i`$ $`=`$ $`x_k`$ $``$ $`x_j`$, we obtain a similar conclusion as follows: recalling that $`𝒮_2^{}(k)`$ $``$ $`𝒮^{}(k)`$ denotes the lifting of the $`x_k`$-standard sub-bunch $`𝒮^{(t)}(k)`$ (in the notation of Theorem 8.5.1), we *claim* that
(25)
$$c_{𝐣_1^{}}^{𝐝_1^{}}𝒮_2^{}(k)𝒮^{}(k).$$
*Proof of claim:* The lifted arrow $`c_𝐣^{}^𝐝^{}`$ has $`x_j`$-height $``$ $`h`$; therefore, by Lemma 8.3.1, we have that
$$x_k\text{-degree}(𝐱^𝐣^{})r(k,v),$$
where $`v`$ $`=`$ $`x_k`$-offset$`(c_𝐣^{}^𝐝^{})`$. Since the translation-equivalent arrow $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ $``$ $`𝒮^{}(k)`$ has the same $`x_k`$-height and $`x_k`$-offset, it satisfies
$$x_k\text{-degree}(𝐱^{𝐣_1^{}})=x_k\text{-degree}(𝐱^𝐣^{})r(k,v).$$
The claim follows immediately, because $`𝒮^{}(k)`$ $`=`$ $`𝒮_1^{}(k)`$ $``$ $`𝒮_2^{}(k)`$, and by definition the arrows $`c_𝐣_{}^𝐝_{}`$ $``$ $`𝒮_1^{}(k)`$ of $`x_k`$-offset $`v`$ satisfy
$$x_k\text{-degree}(𝐱^𝐣_{})>r(k,v).$$
The claim (25) implies that the arrow $`c_{𝐣_1^{}}^{𝐝_1^{}}`$ is the lifting of an arrow $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`𝒮^{(t)}(k)`$. Corollary 8.4.2 now yields that $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_1}^{𝐝_1}`$, and the proof is complete. ∎
## 9. Addition of boxes to basis sets
In this section we explore another way to construct a smooth basis set from a given smooth basis set, by “adding a box” in an appropriate way. The undoing of this operation (that is, the removal of the added box) is accomplished by a truncation.
### 9.1. Definition of box addition
Suppose that $`\beta `$ is an arbitrary basis set of monomials in the variables $`x_1,\mathrm{},x_r`$; recall that $`w_i`$ denotes the $`x_i`$-width of $`\beta `$ for $`1ir`$. Choose one of the variables, say $`x_j`$, and an integer $`h`$ $``$ $`1`$, and form the set
$$x_j^h\beta =\{x_j^hmm\beta \},$$
which can be viewed geometrically as the translation of $`\beta `$ a distance of $`h`$ steps in the positive $`x_j`$-direction. Then, for each variable $`x_k`$ $``$ $`x_j`$, choose an integer $`w_k^{}`$ $``$ $`w_k`$, and form the box (20)
(26)
$$=(w_1^{},w_2^{},\mathrm{},w_{j1}^{},h,w_{j+1}^{},\mathrm{},w_r^{}).$$
We then set
(27)
$$\beta ^{}=(x_j^h\beta ),$$
and say that $`\beta ^{}`$ is obtained by adding the box $``$ in the $`x_j`$-direction to $`\beta `$ (see Figure 5).
###### Lemma 9.1.1.
Let $`\beta ^{}`$ be the set of monomials obtained by adding the box $``$ *(26)* to the basis set $`\beta `$ in the $`x_j`$-direction, as in *(27)*. Then
1. $`\beta ^{}`$ is a basis set.
2. $`\beta `$ is the $`x_j`$-truncation of $`\beta ^{}`$ at height $`h`$.
*Proof:* To prove that $`\beta ^{}`$ is a basis set, we must show that if $`m_1`$ and $`m_2`$ are monomials such that $`m_1`$ $``$ $`\beta ^{}`$ and $`m_2`$ $`|`$ $`m_1`$, then $`m_2`$ $``$ $`\beta ^{}`$. If $`m_1`$ $``$ $``$, then it is clear that $`m_2`$ $``$ $``$ $``$ $`\beta ^{}`$. If $`m_1`$ $``$ $``$, then $`m_1/x_j^h`$ $``$ $`\beta `$. Since $`m_2`$ $`|`$ $`m_1`$, it is clear that
$$x_kx_jx_k\text{-degree}(m_2)x_k\text{-degree}(m_1/x_j^h)<w_kw_k^{};$$
therefore,
$$x_j\text{-degree}(m_2)<hm_2\beta ^{},$$
and
$$\begin{array}{ccc}\hfill x_j\text{-deg}(m_2)h& & (m_2/x_j^h)|(m_1/x_j^h)\hfill \\ & & m_2/x_j^h\beta \hfill \\ & & m_2x_j^h\beta \beta ^{}.\hfill \end{array}$$
This completes the proof that $`\beta ^{}`$ is a basis set, and the second statement follows easily from the definitions. ∎
### 9.2. Minimal generators of $`I_\beta ^{}`$
Let $`\beta `$, $`\beta ^{}`$, $``$, etc., be as above. Since $`\beta `$ $`=`$ $`\beta _t^{}`$ is the $`x_j`$-truncation of $`\beta ^{}`$ at height $`h`$, we have that
$$w_j^{}=w_j+h$$
is the $`x_j`$-width of $`\beta ^{}`$, by assertion (d) of Lemma 7.2.1. We then have the following
###### Lemma 9.2.1.
For each $`i`$, $`1`$ $``$ $`i`$ $``$ $`r`$, the $`x_i`$-width of $`\beta ^{}`$ is $`w_i^{}`$; that is, $`x_i^{w_i^{}}`$ is a minimal generator of $`I_\beta ^{}`$. Furthermore, for all minimal generators $`m`$ of $`I_\beta ^{}`$ we have that
$$x_j\text{-degree}(m)<hm=x_k^{w_k^{}}\text{ for some }x_kx_j;$$
in particular, Hypothesis *8.1.1* holds for $`\beta ^{}`$ and its truncation $`\beta `$.
*Proof:* For all variables $`x_i`$ $`=`$ $`x_k`$ $``$ $`x_j`$, we have that
$$x_k^{w_k^{}}I_\beta ^{}I_{},$$
and $`x_k^{w_k^{}}`$ is a minimal generator of $`I_{}`$, by Lemma 6.5.1; therefore, $`x_k^{w_k^{}}`$ is a minimal generator of $`I_\beta ^{}`$. For $`x_i`$ $`=`$ $`x_j`$, we have already observed that $`w_j^{}`$ is the $`x_j`$-width of $`\beta ^{}`$, so the first assertion holds. Suppose now that $`m`$ is a minimal generator of $`I_\beta ^{}`$ such that $`x_j`$-degree$`(m)`$ $`<`$ $`h`$. A moment’s reflection shows that in fact $`m`$ must be a minimal generator of the monomial ideal $`I_{}`$; the second assertion then follows at once from Lemma 6.5.1. ∎
In addition, recall that Lemma 7.2.1 further describes the relationship between the minimal generators of $`I_\beta ^{}`$ and the minimal generators of $`I_\beta `$, since $`\beta `$ is the $`x_j`$-truncation of $`\beta ^{}`$ at height $`h`$.
### 9.3. Main theorem on box additions
Box addition is a convenient tool for building smooth basis sets, as the following result suggests:
###### Theorem 9.3.1.
Let $`\beta `$ be a basis set, and let $`\beta ^{}`$ be obtained by adding the box $``$ *(26)* in the $`x_j`$-direction. Then:
$$\beta \text{ is smooth }\beta ^{}\text{ is smooth.}$$
*Proof:* $`()`$: Since $`\beta ^{}`$ is smooth and Hypothesis 8.1.1 holds, by Lemma 9.2.1, Theorem 8.6.1 implies that the truncation $`\beta _t^{}`$ $`=`$ $`\beta `$ is smooth.
$`()`$: Given that $`\beta `$ is a smooth basis set, we must prove that $`\beta ^{}`$ is a smooth basis set. To do this we will use Theorem 5.1.1: we must show (for $`\beta ^{}`$) that every non-standard arrow is translation-equivalent to $`0`$, and that there exists a near-standard bunch of arrows $`𝒮^{}`$ such that if $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ is an $`x_i`$-standard arrow, then there is an $`x_i`$-standard arrow $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮^{}`$ such that $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$.
First, let $`c_𝐣^𝐝`$ be a non-standard arrow, which we can assume is minimal by Lemma 4.1.2. Since the arrow is non-standard we know that the vector $`𝐣`$ $``$ $`𝐝`$ has negative components in at least two variable directions, say $`x_{i_1}`$ and $`x_{i_2}`$. Note that the tail of the arrow must be a minimal generator of $`I_\beta ^{}`$ that has $`x_j`$-degree $``$ $`h`$, since the minimal generators of $`\beta ^{}`$ with $`x_j`$-degree $`<`$ $`h`$ are the corner monomials $`x_k^{w_k^{}}`$ (Lemma 9.2.1), and a non-standard arrow cannot have a corner monomial as its tail.
Suppose first that the head of the arrow $`𝐱^𝐣`$ $``$ $``$. By translating the arrow in the increasing $`x_{i_2}`$-direction, one eventually reaches an arrow $`c_{𝐣_1}^{𝐝_1}`$ $``$ $`c_𝐣^𝐝`$ such that
$$x_{i_2}\text{-degree}(𝐱^{𝐣_1})=w_{i_2}^{}1x_{i_2}\text{-degree}(𝐱^{𝐝_1})w_{i_2}^{};$$
therefore, $`c_{𝐣_1}^{𝐝_1}`$ can be translated in the direction of decreasing $`x_{i_1}`$-degree until the head exits the first orthant; whence, $`c_𝐣^𝐝`$ $``$ $`0`$.
If $`𝐱^𝐣`$ $``$ $``$, then $`𝐱^𝐣`$ $``$ $`x_j^h\beta `$, and the arrow $`c_𝐣^𝐝`$ descends to $`c_{𝐣_0}^{𝐝_0}`$. Since $`\beta `$ is assumed smooth, we have that $`c_{𝐣_0}^{𝐝_0}`$ $``$ $`0`$. If the associated translation path causes the head to cross the hyperplane ($`x_k`$-degree $`=`$ $`0`$) for any $`x_k`$ $``$ $`x_j`$, then we can lift the translation path to $`\beta ^{}`$ to obtain that $`c_𝐣^𝐝`$ $``$ $`0`$. If the head crosses the hyperplane ($`x_j`$-degree $`=`$ $`0`$), then the lifted path shows that $`c_𝐣^𝐝`$ $``$ $`c_{𝐣_2}^{𝐝_2}`$ with $`𝐱^{𝐣_2}`$ $``$ $``$, and $`c_{𝐣_2}^{𝐝_2}`$ $``$ $`0`$ as before.
Now let $`𝒮^{}`$ be the near-standard bunch of arrows for $`\beta ^{}`$ that was constructed in the proof of Theorem 8.6.1 (based on Lemmas 8.2.1 and 8.5.1); recall that $`𝒮^{}`$ contains the lifting of a standard sub-bunch $`𝒮^{(t)}`$ for $`\beta `$ $`=`$ $`\beta _t^{}`$. Suppose that $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ is an $`x_i`$-standard arrow; we can assume that this arrow is minimal and cannot be advanced by assertion (c) of Lemma 4.2.1. We must show that $`c_𝐣^𝐝`$ is translation-equivalent to an arrow in $`𝒮^{}(i)`$.
Suppose first that $`x_i`$ = $`x_j`$, so that $`𝒮^{}(i)`$ $`=`$ $`𝒮(j)`$ is a standard $`x_j`$-sub-bunch. If $`𝐱^𝐣`$ $``$ $``$, then the arrow $`c_𝐣^𝐝`$ is the lifting of a minimal standard arrow $`c_{𝐣_0}^{𝐝_0}`$ for $`\beta `$. Since $`\beta `$ is smooth, we can find a minimal standard arrow $`c_{𝐣_0^{}}^{𝐝_0^{}}`$ $``$ $`𝒮^{(t)}(j)`$ such that $`c_{𝐣_0}^{𝐝_0}`$ $``$ $`c_{𝐣_0^{}}^{𝐝_0^{}}`$. Lifting the translation path, we find that $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮^{}(j)`$.
If $`𝐱^𝐣`$ $``$ $``$, then let $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮^{}(j)`$ be the unique arrow such that
$$x_j\text{-degree}(𝐱^𝐣^{})=x_j\text{-degree}(𝐱^𝐣)\text{and }x_j\text{-offset}(c_𝐣^{}^𝐝^{})=x_j\text{-offset}(c_𝐣^𝐝),$$
the existence of which is ensured by Corollary 4.6.1. Suppose that the lengths of these two arrows differ; in other words, suppose that
$$x_j\text{-degree}(𝐱^𝐝)x_j\text{-degree}(𝐱^𝐝^{}).$$
Recalling Lemma 6.5.2, we let
$$𝐱^{𝐣_1}=\mathrm{lcm}(𝐱^𝐣,𝐱^𝐣^{})\beta ^{}.$$
It is clear that we can translate both $`c_𝐣^𝐝`$ and $`c_𝐣^{}^𝐝^{}`$ by degree-increasing steps (excluding the $`x_j`$-direction) to arrows $`c_{𝐣_1}^{𝐝_1}`$ and $`c_{𝐣_1}^{𝐝_1^{}}`$, respectively. The tails $`𝐱^{𝐝_1}`$ and $`𝐱^{𝐝_1^{}}`$ differ only in $`x_j`$-degree; it follows that the arrow corresponding to the tail of larger $`x_j`$-degree can be advanced, which is a contradiction, since neither $`c_𝐣^𝐝`$ nor $`c_𝐣^{}^𝐝^{}`$ can be advanced. We therefore have that
$$c_{𝐣_1}^{𝐝_1}=c_{𝐣_1}^{𝐝_1^{}}c_𝐣^𝐝c_𝐣^{}^𝐝^{}𝒮^{}(j),$$
as desired.
Finally, we have to consider the case in which $`c_𝐣^𝐝`$ is a minimal $`x_k`$-standard arrow that cannot be advanced, where $`x_k`$ $``$ $`x_j`$. Let $`v`$ denote the $`x_k`$-offset of $`c_𝐣^𝐝`$. From Lemma 9.2.1, we know that the tail $`𝐱^𝐝`$ is either equal to the corner monomial $`x_k^{w_k^{}}`$ or else $`x_j`$-degree$`(𝐱^𝐝)`$ $``$ $`h`$.
In case $`𝐱^𝐝`$ $`=`$ $`x_k^{w_k^{}}`$, if the head
$$𝐱^𝐣\beta =\beta _t^{},$$
then
$$x_k\text{-degree}(𝐱^𝐣)>r(k,v)c_𝐣^𝐝𝒮_1^{}(k)𝒮^{}(k),$$
in the notation of Lemma 8.5.1. If the head
$$𝐱^𝐣\beta =\beta _t^{},$$
then, since $`c_𝐣^𝐝`$ cannot be advanced, but *can* be translated $`h`$ steps in the direction of increasing $`x_j`$-degree to $`c_𝐣^{}^𝐝^{}`$, we must have that $`w_k^{}`$ $`=`$ $`w_k`$. This in turn implies that $`c_𝐣^𝐝`$, considered as an arrow for the truncation $`\beta `$, lies in the sub-bunch $`𝒮^{(t)}(k)`$; therefore,
$$c_𝐣^𝐝c_𝐣^{}^𝐝^{}𝒮_2^{}(k)𝒮^{}(k).$$
In case $`x_j`$-degree$`(𝐱^𝐝)`$ $``$ $`h`$, we can descend the arrow to $`c_{𝐣_0}^{𝐝_0}`$; then, because $`\beta `$ is assumed smooth, we know that $`c_{𝐣_0}^{𝐝_0}`$ $``$ $`c_{𝐣_0^{}}^{𝐝_0^{}}`$ $``$ $`𝒮^{(t)}(k)`$, and the translation path lifts to yield $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮^{}(k)`$. This completes the proof of the theorem. ∎
### 9.4. Compound boxes
By a compound box, we mean a basis set $`\beta `$ that is constructed by starting with a box, and then performing a finite sequence of box additions in various variable directions; Figure 6 illustrates the idea. Note that a box $``$ is a compound box, since it can be generated by starting with itself and performing a sequence of box additions of length $`0`$.
Since the starting point (a box) is smooth, by Proposition 6.5.3, and adding a box to a smooth basis set yields a smooth basis set, by Theorem 9.3.1, we obtain the following corollary by induction:
###### Corollary 9.4.1.
If $`\beta `$ is a compound box, then $`\beta `$ is a smooth basis set. ∎
### 9.5. Example: Basis sets in two variables
It is easy to verify (see Figure 7) that *every* basis set $`\beta `$ in two variables is a compound box; whence, Corollary 9.4.1 yields the following
###### Corollary 9.5.1.
Every basis set $`\beta `$ in two variables is smooth. ∎
Haiman’s lovely proof of this result (part of the proof of \[6, Proposition 2.4\]) introduced the idea of arrow translation, and inspired the present paper. Note that the $`𝗄`$-basis of $`M/M^2`$ that he obtains is slightly different from ours; his basis arrows are typically not *minimal* standard arrows. From the smoothness of the points $`t_\beta `$, Haiman deduces that $`\text{H}^n`$ is everywhere nonsingular and irreducible (facts first proved by Fogarty ). Corollary 9.4.1 can be viewed as a generalization of the two-variable smoothness phenomenon to higher dimensions.
### 9.6. Example: The lexicographic point
The “lexicographic point” of $`\mathrm{Hilb}_{_𝗄^r}^{p(z)}`$ is the point corresponding to the unique saturated lexicographic ideal $`L`$ such that $`𝗄[X_0,\mathrm{},X_r]/L`$ has Hilbert polynomial $`p(z)`$. A. Reeves and M. Stillman prove in general that the lexicographic point is a smooth point . In the case of a constant Hilbert polynomial $`p(z)`$ $`=`$ $`n`$, one checks that
$$L=(X_0,X_1,\mathrm{},X_{r2},X_{r1}^n).$$
Dehomogenizing with respect to the variable $`X_r`$, we obtain the ideal $`I_\beta `$ $``$ $`𝗄[x_0,\mathrm{},x_{r1}]`$, where
$$\beta =\{1,x_{r1},x_{r1}^2,\mathrm{},x_{r1}^{n1}\}.$$
This is clearly a smooth basis set, since it is a special case of a box.
### 9.7. Example: $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_1x_2,x_3\}`$
In Example 5.3 we considered the basis set
$$\{1,x_1,x_2,x_3\},$$
which is non-smooth. If we add the monomial $`x_1x_2`$, we obtain a compound box $`\beta `$ (see Figure 8); we will verify “by hand” that $`\beta `$ is smooth.
The minimal generators of $`I_\beta `$ are the monomials
$$x_1^2,x_2^2,x_1x_3,x_2x_3,x_3^2;$$
whence, one has $`5`$ $``$ $`5`$ $`=`$ $`25`$ minimal arrows, and these span $`M/M^2`$ by Lemma 4.1.2. However, inspecting the minimal arrows, we find that:
* Four are non-standard arrows, all of which are translation-equivalent to $`0`$: $`c_{(0,0,0)}^{(1,0,1)}`$, $`c_{(0,1,0)}^{(1,0,1)}`$, $`c_{(0,0,0)}^{(0,1,1)}`$, and $`c_{(1,0,0)}^{(0,1,1)}`$.
* Five are standard arrows that are translation-equivalent to $`0`$: $`c_{(0,0,0)}^{(2,0,0)}`$, $`c_{(0,0,0)}^{(0,2,0)}`$, $`c_{(0,0,0)}^{(0,0,2)}`$, $`c_{(1,0,0)}^{(0,0,2)}`$, and $`c_{(0,1,0)}^{(0,0,2)}`$.
* Two are standard arrows that are not translation-equivalent to $`0`$, but *are* translation-equivalent to each other: $`c_{(1,0,0)}^{(1,0,1)}`$ $``$ $`c_{(0,1,0)}^{(0,1,1)}`$.
This means that there are at most 15 non-trivial translation-equivalence classes of arrows available to span $`M/M^2`$, but $`r`$ $``$ $`n`$ $`=`$ $`3`$ $``$ $`5`$ = 15 is a lower bound on the $`𝗄`$-dimension of the cotangent space, by Proposition 2.4.1. It follows that $`M/M^2`$ has $`𝗄`$-dimension 15; whence, $`\beta `$ is smooth.
## 10. Smooth basis sets in three variables are compound boxes
The main goal of this section is to prove Theorem 10.3.1, which states that a basis set $`\beta `$ in three variables is smooth if and only if $`\beta `$ is a compound box. In the next section we will show by example that smooth basis sets in four or more variables need not be compound boxes.
### 10.1. The main lemma
Let $`\beta `$ be a basis set in the variables $`x_1`$, $`x_2`$, $`x_3`$, and $`I_\beta `$ the associated monomial ideal. For
(28)
$$(i,j)\{(1,2),(1,3),(2,3)\},$$
we write $`G(i,j)`$ for the set of minimal generators of $`I_\beta `$ that involve only variables in the set $`\{x_i,x_j\}`$. For example, we always have that
$$x_i^{w_i}G(i,j)\text{ and }x_j^{w_j}G(i,j),$$
so the number of elements
$$|G(i,j)|2.$$
###### Lemma 10.1.1.
If $`|G(i,j)|`$ $`>`$ $`2`$ for all of the ordered pairs $`(i,j)`$ in *(28)*, then there exists a non-standard arrow for $`\beta `$ that is not translation-equivalent to $`0`$; consequently, $`\beta `$ is not smooth by Theorem *5.1.1*.
### 10.2. Proof of Lemma 10.1.1
By the hypothesis, we may choose
(29)
$$\begin{array}{ccccc}\hfill m_{(1,3)}& =& x_1^ax_3^b& & G(1,3),\hfill \\ \hfill m_{(2,3)}& =& x_2^cx_3^d& & G(2,3),\hfill \\ \hfill m_{(1,2)}& =& x_1^ex_2^f& & G(1,2)\hfill \end{array}$$
such that all the exponents are positive and $`m_{(1,3)}`$ (resp. $`m_{(2,3)}`$) has maximal $`x_1`$-degree (resp. maximal $`x_2`$-degree) subject to the stated constraints; note that this implies that $`m_{(1,3)}`$ (resp. $`m_{(2,3)}`$) has minimal $`x_3`$-degree subject to the stated constraints. We proceed to construct the desired non-standard arrow; there are two cases (see Figure 9):
#### 10.2.1. Case 1: one of the monomials in *(29)* is dominant
We say that $`m_{(i,j)}`$ is dominant provided that the degree of each of its constituent variables is greater than or equal to the degree of the same variable in the other monomial in which it appears. For example, $`m_{(1,3)}`$ is dominant provided that (as shown in the left-hand portion of Figure 9)
$$ae\text{ and }bd.$$
In this case, let $`g`$ = $`\mathrm{max}(c,f)`$, and note that $`x_2^g`$ is a basis monomial. Starting at $`x_2^g`$, we can move to a maximal basis monomial $`m^{}`$ $`=`$ $`x_1^px_2^qx_3^r`$ by a sequence of degree-increasing steps. Since $`x_1^ex_2^f`$ is a minimal generator of the ideal, we know that $`p`$ $`<`$ $`e`$. Similarly, since $`x_2^cx_3^d`$ is a minimal generator, we have that $`r`$ $`<`$ $`d`$. Then the rigid arrow $`A`$ with tail $`m_{(1,3)}`$ and head $`m^{}`$ is also non-standard, since its vector
$$(p,q,r)(a,0,b)=(pa,q,rb)$$
has the first and third coordinates negative:
$$\begin{array}{ccccccc}\hfill pa& <& ea& & aa& =& 0,\hfill \\ \hfill rb& <& db& & bb& =& 0.\hfill \end{array}$$
Since $`A`$ is rigid, it is not translation-equivalent to $`0`$, as desired.
#### 10.2.2. Case 2: None of the monomials in *(29)* is dominant
Writing out what this condition says, we find that:
$$\begin{array}{c}(\neg (aebd))(\neg (cfdb))(\neg (eafc))\\ \\ (a<eb<d)(c<fd<b)(e<af<c)\\ \\ (a<ef<cd<b)(b<dc<fe<a).\end{array}$$
Suppose the first alternative in the last line holds; that is, suppose (as shown in the right-hand portion of Figure 9) that
$$a<e,f<c,\text{and }d<b.$$
Note first of all that the monomial
$$m=x_2^{(w_21)}x_3^{(d1)}\beta .$$
If not, there would exist a minimal generator $`m^{}`$ of $`I_\beta `$ such that
$$m^{}|mm^{}G(2,3).$$
It is clear that $`m^{}`$ cannot equal either of the corner monomials $`x_2^{w_2}`$, $`x_3^{w_3}`$, so $`m^{}`$ must involve both $`x_2`$ and $`x_3`$ nontrivially. However, the $`x_3`$-degree of $`m^{}`$ is $``$ $`d1`$, which contradicts our choice of $`m_{(2,3)}`$ as having minimal $`x_3`$-degree $`d`$ among the members of $`G(2,3)`$ that involve two variables nontrivially.
Let $`s`$ $`>`$ $`0`$ denote the minimal exponent such that
$$x_1^sm\beta ;$$
since
$$(x_1^ex_2^f=m_{(1,2)})|(x_1^em)(x_1^em)\beta ,$$
we have that
(30)
$$se<w_1,$$
where the second inequality holds because the minimal generator $`x_1^{w_1}`$ cannot divide the minimal generator $`m_{(1,2)}`$. Let
(31)
$$\alpha =\{\begin{array}{c}s1,\text{ if }sa,\hfill \\ a1,\text{ if }s>a,\hfill \end{array}$$
and form the arrow
$$A=c_{(\alpha ,w_21,d1)}^{(a,0,b)}$$
with tail $`m_{(1,3)}`$ and head
$$m_h=x_1^\alpha x_2^{w_21}x_3^{d1}\beta (\text{ since }\alpha <s).$$
Notice that $`A`$ is non-standard, since its vector
$$(\alpha ,w_21,d1)(a,0,b)=(\alpha a,w_21,(d1)b)$$
has negative first and third components (and non-negative second component $`w_21`$).
We assert that this arrow is not translation-equivalent to $`0`$. If it were, the head would have to exit the first octant across either the hyperplane ($`x_1`$-degree $`=`$ $`0`$) or the hyperplane ($`x_3`$-degree $`=`$ $`0`$). In the former case, we would have to translate $`A`$ to an arrow $`A^{}`$ of the same $`x_1`$-height ($`=`$ $`a`$), but with tail divisible by a minimal generator $`\stackrel{~}{m}`$ of $`x_1`$-degree $`<`$ $`a`$. Since $`\stackrel{~}{m}`$ cannot divide $`m_{(1,2)}`$ and $`m_{(1,3)}`$, we must have that
(32)
$$\text{either }x_2\text{-degree}(\stackrel{~}{m})>f\text{ or }x_3\text{-degree}(\stackrel{~}{m})>b.$$
However, we have that $`A`$ is “rigid” with respect to motion in the $`x_2`$-direction; that is, neither $`A`$ nor any of its translates $`A^{}`$ can be moved in either the increasing or decreasing $`x_2`$-direction. Indeed, an easy induction on the length of the path from $`A`$ to $`A^{}`$ shows that
$$x_2\text{-degree}(\text{head}(A^{}))=w_21\text{ and }x_2\text{-degree}(\text{tail}(A^{}))=0,$$
and clearly such an arrow cannot be translated in either the increasing or decreasing $`x_2`$-degree directions. Furthermore, neither $`A`$ nor any of its translates $`A^{}`$ has $`x_3`$-height greater than the initial value $`b`$, for we have already seen that the $`x_2`$-degree of the head of $`A^{}`$ is invariantly $`w_21`$, and if the $`x_3`$-height of $`A^{}`$ were to exceed $`b`$, then the $`x_3`$-degree of the head would be $``$ $`d`$, implying that the head would be divisible by $`m_{(2,3)}`$. Therefore, we see that it is impossible to translate $`A`$ to $`A^{}`$ with tail divisible by $`\stackrel{~}{m}`$ as in (32), so we cannot move the head of $`A`$ across the hyperplane ($`x_1`$-degree $`=`$ $`0`$).
The only remaining possibility is to translate $`A`$ so that the head exits the first octant across the hyperplane ($`x_3`$-degree $`=`$ $`0`$). This requires us to translate $`A`$ to an arrow $`A^{}`$ of the same $`x_3`$-height ($`=`$ $`b`$), but with tail divisible by a minimal generator $`\widehat{m}`$ of $`x_3`$-degree $`<`$ $`b`$. Since the $`x_2`$-height of $`A^{}`$ is invariantly $`0`$, we see that
$$x_2\text{-degree}(\widehat{m})=0\widehat{m}G(1,3)\widehat{m}=x_1^{w_1},$$
where the last implication follows from our choice of $`m_{(1,3)}`$ as the monomial of minimal $`x_3`$-degree in $`G(1,3)`$ among those involving both $`x_1`$ and $`x_3`$ nontrivially. In view of the constraints on the motion of $`A`$, we see that we would have to be able to translate $`A`$ a distance of $`w_1a`$ units in the positive $`x_1`$-direction, but this motion would move the head to
$$x_1^{\alpha +(w1a)}x_2^{w_21}x_3^{d1},$$
which lies outside of $`\beta `$, because (recalling (30) and (31))
$$\begin{array}{ccc}\hfill sa& & \alpha +(w_1a)=(s1)+(w_1a)s,\text{ and }\hfill \\ \hfill s>a& & \alpha +(w_1a)=(a1)+(w_1a)=w_11s.\hfill \end{array}$$
The required translation is therefore impossible; whence, $`A`$ is not translation-equivalent to $`0`$, and the proof of Lemma 10.1.1 is complete. ∎
### 10.3. The main theorem
###### Theorem 10.3.1.
Let $`\beta `$ be a basis set in the variables $`x_1`$, $`x_2`$, $`x_3`$. Then
$$\beta \text{ is smooth }\beta \text{ is a compound box.}$$
*Proof:* ($``$): Immediate from Corollary 9.4.1.
($``$): We proceed by induction on $`n`$ $`=`$ $`|\beta |`$.
*Base case:* $`n`$ $`=`$ $`1`$. The only basis set of cardinality $`1`$ is the box
$$\{1\}=(1,1,1),$$
which is smooth by Proposition 6.5.3, and trivially a compound box.
*Inductive step:* Suppose that $`|\beta |`$ $`=`$ $`n`$, and that any smooth basis set of cardinality $`<`$ $`n`$ is a compound box. Since $`\beta `$ is by hypothesis smooth, Lemma 10.1.1 implies that for at least one of the pairs $`(i,j)`$ in (28), we have that
$$G(i,j)=\{x_i^{w_i},x_j^{w_j}\};$$
that is, no minimal generator of $`I_\beta `$ exists that involves both $`x_i`$ and $`x_j`$ nontrivially, but does not involve the third variable $`x_k`$. Let $`h`$ be the minimal $`x_k`$-degree among the minimal generators of $`I_\beta `$ that *do* involve $`x_k`$. If $`h`$ $`=`$ $`w_k`$, then
$$\beta =(w_1,w_2,w_3)$$
is a (compound) box, as desired. If $`h`$ $`<`$ $`w_k`$, we let $`\beta _t`$ be the $`x_k`$-truncation of $`\beta `$ at height $`h`$ (21). Then one sees easily that $`\beta `$ is obtained from $`\beta _t`$ by adding a box in the $`x_k`$-direction, as described in Section 9: more precisely, the added box has dimensions $`w_i`$ in the $`x_i`$-direction, $`w_j`$ in the $`x_j`$-direction, and $`h`$ in the $`w_k`$-direction. It follows from Theorem 9.3.1 that $`\beta _t`$ is smooth, since $`\beta `$ is smooth by hypothesis. The induction hypothesis now implies that $`\beta _t`$ is a compound box $``$ $`\beta `$ is a compound box, and we are done. ∎
More can be gleaned from the preceding proof: suppose that $`\beta `$ is a basis set in the variables $`x_1`$, $`x_2`$, $`x_3`$ for which condition (a) of Theorem 5.1.1 holds; that is, every non-standard arrow for $`\beta `$ is translation-equivalent to $`0`$. Then Lemma 10.1.1 implies that for at least one of the pairs $`(i,j)`$ in (28), we have that
$$G(i,j)=\{x_i^{w_i},x_j^{w_j}\};$$
arguing as in the preceding proof, we then see that $`\beta `$ is either a box or is the result of adding a box to a truncation $`\beta _t`$. In the former case, $`\beta `$ is smooth, and in the latter case, $`\beta _t`$ also satisfies condition (a) of Theorem 5.1.1, by Theorem 7.4.1; therefore, induction yields that $`\beta _t`$ is smooth $``$ $`\beta `$ is smooth, by Theorem 9.3.1. Whence:
###### Corollary 10.3.2.
Let $`\beta `$ be a basis set in the variables $`x_1`$, $`x_2`$, $`x_3`$. Then
$$\beta \text{ is smooth}\{\begin{array}{c}\text{condition }\text{(a)}\text{ of Theorem }\text{5.1.1}\text{ holds; that is, every}\hfill \\ \text{non-standard arrow for }\beta \text{ is translation-equivalent to }0\text{.}\text{}\hfill \end{array}$$
As of this writing, I do not know if this result extends to higher dimensions.
## 11. The union of two boxes
To end this paper, we study one more family of smooth basis sets that does not consist entirely of compound boxes (in four or more variables): basis sets that are unions of two boxes.
### 11.1. Notation
We will use the following notation throughout this section. Let
$$\begin{array}{ccc}\hfill _1& =& (w_{1,1},w_{2,1},\mathrm{},w_{r,1}),\hfill \\ \hfill _2& =& (w_{1,2},w_{2,2},\mathrm{},w_{r,2})\hfill \end{array}$$
be two boxes (20) in the variables $`x_1`$, …, $`x_r`$, and let
(33)
$$\beta =_1_2.$$
One checks easily that $`\beta `$ is a basis set. Lacking inspiration, we call $`\beta `$ a two-box union.
### 11.2. Minimal generators of $`I_\beta `$
As usual, we write $`w_i`$ for the $`x_i`$-width of $`\beta `$; that is, $`x_i^{w_i}`$ is the corner (minimal) monomial divisible by $`x_i`$. The following result is clear:
###### Lemma 11.2.1.
The $`x_i`$-width $`w_i`$ of the two-box union $`\beta `$ *(33)* is given by
$$w_i=\mathrm{max}(w_{i,1},w_{i,2}),1ir.\text{}$$
We now write the set of variables as a union
$$\{x_1,x_2,\mathrm{},x_r\}=V_1V_2V_3,$$
where
(34)
$$\begin{array}{ccc}\hfill V_1& =& \{x_jw_{j,1}>w_{j,2}\},\hfill \\ \hfill V_2& =& \{x_kw_{k,1}<w_{k,2}\},\hfill \\ \hfill V_3& =& \{x_{\mathrm{}}w_{\mathrm{},1}=w_{\mathrm{},2}\};\hfill \end{array}$$
henceforth, we will use the subscripts $`j`$, $`k`$, and $`\mathrm{}`$ to denote membership in $`V_1`$, $`V_2`$, and $`V_3`$, respectively. One sees easily that
$$\begin{array}{c}V_1=\mathrm{}_1_2\beta =_2,\hfill \\ V_2=\mathrm{}_2_1\beta =_1,\hfill \end{array}$$
so the most interesting case is when both $`V_1`$ and $`V_2`$ are non-empty.
We have the following
###### Lemma 11.2.2.
Let $`\beta `$ be a two-box union *(33)*. Then a minimal generator $`m`$ of the monomial ideal $`I_\beta `$ is either a corner monomial $`x_i^{w_i}`$ or a two-variable monomial of the form $`x_j^{w_{j,2}}x_k^{w_{k,1}}`$, with $`x_j`$ $``$ $`V_1`$, $`x_k`$ $``$ $`V_2`$ *(34)*.
*Proof:* Let
$$m=x_1^{s_1}x_2^{s_2}\mathrm{}x_r^{s_r}$$
be a minimal generator of $`I_\beta `$, and suppose that three or more of the exponents (say $`s_1`$, $`s_2`$, and $`s_3`$) are positive. Then each of the monomials
$$x_1^{s_11}x_2^{s_2}x_3^{s_3}\mathrm{}x_r^{s_r},x_1^{s_1}x_2^{s_21}x_3^{s_3}\mathrm{}x_r^{s_r},x_1^{s_1}x_2^{s_2}x_3^{s_31}\mathrm{}x_r^{s_r}$$
belong to $`\beta `$, which implies that two of these monomials must belong to the same box ($`_1`$ or $`_2`$). But then the least common multiple of the two monomials must also belong to this box (Lemma 6.5.2); that is, $`m`$ $``$ $`\beta `$, which is a contradiction. We conclude that $``$ $`2`$ of the exponents $`s_i`$ can be positive. If only one of the exponents $`s_i`$ is positive, then $`m`$ is the corner monomial $`x_i^{(s_i=w_i)}`$. If two of the exponents (say $`s_1`$ and $`s_2`$) are positive, so that
$$m=x_1^{s_1}x_2^{s_2},$$
we again have that the monomials
$$x_1^{s_11}x_2^{s_2},x_1^{s_1}x_2^{s_21}$$
belong to $`\beta `$; if both belonged to the same box, then we would arrive once more at the contradiction $`m`$ $``$ $`\beta `$. Therefore, the latter two monomials belong to different boxes, say
$$x_1^{s_11}x_2^{s_2}_2_1,x_1^{s_1}x_2^{s_21}_1_2;$$
consequently,
$$\begin{array}{c}x_1^{s_11}x_2^{s_2}_2s_11<w_{1,2}\text{ and }s_2<w_{2,2},\hfill \\ x_1^{s_1}x_2^{s_21}_1s_1<w_{1,1}\text{ and }s_21<w_{2,1},\hfill \end{array}$$
and
$$\begin{array}{c}x_1^{s_11}x_2^{s_2}_1s_11w_{1,1}\text{ or }s_2w_{2,1},\hfill \\ x_1^{s_1}x_2^{s_21}_2s_1w_{1,2}\text{ or }s_21w_{2,2}.\hfill \end{array}$$
Note that
$$s_2<w_{2,2}s_21<w_{2,2}s_1w_{1,2};$$
whence,
$$(s_1w_{1,2}\text{ and }s_11<w_{1,2})s_1=w_{1,2},$$
and
$$w_{1,2}=s_1<w_{1,1}x_1V_1,$$
as desired. A similar argument yields
$$s_2=w_{2,1}\text{ and }x_2V_2,$$
and the proof is complete. ∎
### 11.3. Two-box unions are smooth basis sets
Retaining the notation of Subsections 11.1 and 11.2, we begin with the following
###### Lemma 11.3.1.
Let $`\beta `$ be a two-box union, and let
$$𝐱^𝐝=x_j^{w_{j,2}}x_k^{w_{k,1}},x_jV_1,x_kV_2,$$
be a two-variable minimal generator of $`I_\beta `$, as in Lemma *11.2.2*. If $`c_𝐣^𝐝`$ is an $`x_j`$\- (resp. $`x_k`$-) standard arrow for $`\beta `$, then the head of the arrow $`𝐱^𝐣`$ $``$ $`_2`$ (resp. $`_1`$).
*Proof:* If $`c_𝐣^𝐝`$ is an $`x_j`$-standard arrow, then we have that
$$x_k\text{-degree}(𝐱^𝐣)x_k\text{-degree}(𝐱^𝐝)=w_{k,1}𝐱^𝐣_1.$$
Since $`𝐱^𝐣`$ $``$ $`\beta `$, we must have that $`𝐱^𝐣`$ $``$ $`_2`$, as asserted. A similar argument applies in the case that $`c_𝐣^𝐝`$ is $`x_k`$-standard. ∎
###### Theorem 11.3.2.
A two-box union $`\beta `$ *(33)* is a smooth basis set.
*Proof:* We will show that the conditions (a) and (b) of Theorem 5.1.1 hold for $`\beta `$.
First suppose that $`c_𝐣^𝐝`$ is a non-standard arrow; we must show that $`c_𝐣^𝐝`$ is translation-equivalent to $`0`$. We may assume, by Lemma 4.1.2, that $`c_𝐣^𝐝`$ is a minimal arrow; that is, its tail $`𝐱^𝐝`$ is a minimal generator of $`I_\beta `$. Since every arrow with tail a corner monomial $`x_i^{w_i}`$ is standard, Lemma 11.2.2 implies that
$$𝐱^𝐝=x_j^{w_{j,2}}x_k^{w_{k,1}},x_jV_1,x_kV_2.$$
Without loss of generality, suppose that the head $`𝐱^𝐣`$ $``$ $`_1`$. Then we can translate the arrow in the direction of increasing $`x_j`$-degree to reach $`c_{𝐣_1}^{𝐝_1}`$, where
$$\begin{array}{c}x_j\text{-degree}(𝐱^{𝐣_1})=w_j1=w_{j,1}1\text{ and }\\ x_i\text{-degree}(𝐱^{𝐣_1})=x_i\text{-degree}(𝐱^𝐣),\text{for }ij.\end{array}$$
Then
$$x_j\text{-degree}(𝐱^{𝐝_1})>x_j\text{-degree}(𝐱^{𝐣_1})x_j\text{-degree}(𝐱^{𝐝_1})w_j,$$
so we may translate $`c_{𝐣_1}^{𝐝_1}`$ by degree-decreasing steps so that its tail approaches the corner monomial $`x_j^{w_j}`$. The head must eventually exit the first orthant, since a non-standard arrow cannot have a corner monomial as its tail. Therefore, $`c_𝐣^𝐝`$ $``$ $`0`$, and condition (a) holds.
Now let $`𝒮`$ be a standard bunch of arrows for $`\beta `$, and $`c_𝐣^𝐝`$ $`\sim ̸`$ $`0`$ be a standard arrow for $`\beta `$. We must show that there exists an arrow $`c_𝐣^{}^𝐝^{}`$ $``$ $`𝒮`$ such that $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$. We may assume that $`c_𝐣^𝐝`$ is a minimal $`x_i`$-standard arrow that cannot be advanced, by assertion (c) of Lemma 4.2.1. If the tail $`𝐱^𝐝`$ is the corner monomial $`x_i^{w_i}`$, we are done, since then
$$c_𝐣^𝐝𝒮(i)𝒮$$
by Corollary 4.6.2. Otherwise, by Lemma 11.2.2, we have that
$$𝐱^𝐝=x_j^{w_{j,2}}x_k^{w_{k,1}},x_jV_1,x_kV_2,$$
and we may assume without loss of generality that $`x_i`$ = $`x_j`$. Let $`v`$ denote the $`x_j`$-offset of $`c_𝐣^𝐝`$, and let $`c_𝐣^{}^𝐝^{}`$ be the unique arrow in $`𝒮(j)`$ such that
$$x_j\text{-degree}(𝐱^𝐣^{})=x_j\text{-degree}(𝐱^𝐣)\text{ and }x_j\text{-offset}(c_𝐣^{}^𝐝^{})=v;$$
the existence of $`c_𝐣^{}^𝐝^{}`$ is guaranteed by Corollary 4.6.1. The tail $`𝐱^𝐝^{}`$ is a minimal generator of $`I_\beta `$ that is divisible by $`x_j`$. We proceed to show that $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$.
If $`𝐱^𝐝^{}`$ $`=`$ $`x_j^{w_j}`$, then one sees easily that
$$𝐱^𝐣^{}=𝐱^𝐣/x_k^{w_{k,1}};$$
it is then apparent that $`c_𝐣^{}^𝐝^{}`$ can be translated $`w_{k,1}`$ steps in the direction of increasing $`x_k`$-degree to reach an arrow $`c_𝐣^{𝐝_1^{}}`$, whose tail $`𝐱^{𝐝_1^{}}`$ is divisible by $`𝐱^𝐝`$; therefore, $`c_𝐣^{𝐝_1^{}}`$ and $`c_𝐣^{}^𝐝^{}`$ can be advanced, which is a contradiction. Hence, Lemma 11.2.1 yields that
$$𝐱^𝐝^{}=x_j^{w_{j,2}}x_k^{}^{w_{k^{},1}},x_k^{}V_2.$$
We now know that $`c_𝐣^𝐝`$ and $`c_𝐣^{}^𝐝^{}`$ have the same length, $`x_j`$-offset, and $`x_j`$-height $`=`$ $`w_{j,2}`$. By Lemma 11.3.1, we know that
$$𝐱^𝐣,𝐱^𝐣^{}_2\mathrm{lcm}(𝐱^𝐣,𝐱^𝐣^{})=𝐱^𝐣^{}_2\beta ,$$
where the implication uses Lemma 6.5.2. It is now clear that $`c_𝐣^𝐝`$ and $`c_𝐣^{}^𝐝^{}`$ can each be translated by degree-increasing steps to the same arrow $`c_𝐣^{}^𝐝^{}`$; whence, $`c_𝐣^𝐝`$ $``$ $`c_𝐣^{}^𝐝^{}`$, as desired. Therefore, condition (b) holds, and the proof is complete. ∎
###### Remark 11.3.3.
Note that *three*-box unions need not be smooth; for example, the left-hand basis set illustrated in Figure 9 is a three-box union, but is not smooth, since a non-standard rigid arrow exists.
### 11.4. Example: $`\beta `$ $`=`$ $`\{1,x_1,x_2,x_1x_2,x_3,x_4,x_3x_4\}`$
This basis set $`\beta `$ is the two-box union
$$\beta =(2,2,1,1)(1,1,2,2)𝗄[x_1,x_2,x_3,x_4].$$
Therefore, by Theorem 11.3.2, $`\beta `$ is a smooth basis set. However, it is clear that $`\beta `$ is not a compound box, so the characterization of smooth basis sets in three variables given by Theorem 10.3.1 does not extend to higher dimensions.
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# Estimating the Porosity of the Interstellar Medium from Three-Dimensional Photoionization Modeling of H ii Regions
## 1 Introduction
The interstellar medium (ISM) is observed to be clumpy on a very wide range of size scales from parsec scale molecular clouds to AU sizes inferred from interstellar scintillation measurements (Hill et al. 2004; Scalo & Elmegreen 2005). This clumpiness is observed to be self-similar over a range of scales with hierarchical or fractal dimensions in the range 2.2 to 2.7 (Sanchez, Alfaro, & Perez 2005). The complex structures are thought to result from mechanical and radiative energy input into the ISM from many processes including stellar winds, supernovae, shocks, photoionization, and cosmic rays. This is reflected in the multi-phase nature of the ISM, where the hot ionized, warm ionized, warm neutral, cold neutral, and molecular gas co-exist. There is considerable debate among theorists regarding the volume filling factors of each phase and the overall structure of the ISM. Over the last decade large multiwavelength surveys of the ISM are providing new high spatial and spectral resolution observations against which observational signatures of theoretical models may be tested. These velocity resolved surveys include neutral gas (Hartmann & Burton 1997), molecular gas (Dame et al. 2001), and diffuse ionized gas (Haffner et al. 2003). The combination of the new datasets and advances in 3D dynamical and radiation transfer simulations now allows for critical testing of global models of the structure and dynamics of the ISM. In this paper we use three-dimensional photoionization models to study the structure of a low density H ii region in the Galaxy on scales of tens of parsecs. This is a first step towards understanding the structure, clumpiness, and ionization of the large scale diffuse ionized gas in the ISM.
Widespread diffuse ionized gas is a significant component of the ISM in the Milky Way (Kulkarni & Heiles 1988; Reynolds 1995) and other galaxies (Rand 1997; Hoopes & Walterbos 2003; Thilker et al. 2002; Collins et al. 2000). The diffuse ionized gas, hereafter referred to as the warm ionized medium (WIM), in the Milky Way is revealed primarily through faint diffuse H$`\alpha `$ emission and pulsar dispersion measures. The physical characteristics of this gas, including its temperature structure and ionization state, as well as the spectrum of its ionizing sources are probed through analysis of optical forbidden emission lines. The H$`\alpha `$ data combined with the pulsar data suggest that the WIM has a volume filling factor of about 20% (Reynolds 1991), and is one of the principal components of the ISM. However, there is considerable uncertainty regarding the structure, dynamics, source of ionization, and heating of the WIM (e.g., see recent review by Reynolds, Haffner, & Madsen 2002). The Wisconsin H$`\alpha `$ Mapper (WHAM) has completed an H$`\alpha `$ survey of the Northern sky (Haffner et al. 2003) revealing the distribution and kinematics of the WIM in the Galaxy for the first time. Several large regions are also being mapped with WHAM in other lines such as H$`\beta `$, He i $`\lambda `$5876, \[N ii$`\lambda `$6583, \[S ii$`\lambda `$6716, \[O i$`\lambda `$6300, and \[O iii$`\lambda `$5007. The data have revealed the clumpy and filamentary structure of the WIM on large scales, such as the bipolar loop structure in the Perseus Arm which extends $`\pm 1`$ kpc above and below the Galactic plane (Haffner at al. 1999; Madsen 2004). In addition to the large scale structures, WHAM has mapped several nearby low density H ii regions with sizes of tens of parsecs (e.g., Haffner et al. 1999). Three-dimensional photoionization modeling of these H ii regions may allow us to estimate the structure and porosity of the ISM on these scales. Such models of well-defined H ii regions with known ionizing sources are a first step in modeling the WIM on larger scales with multiple and unidentified ionizing sources.
In this paper we study the three-dimensional structure of the H ii region surrounding the O9.5V star $`\zeta `$ Oph, incorporating recent WHAM observations in the lines of H$`\alpha `$, \[N ii$`\lambda `$6583, and \[S ii$`\lambda `$6716 (Haffner et al. 1999; Baker et al. 2005). We use the three-dimensional photoionization code of Wood, Mathis, & Ercolano (2004) to model the observed H$`\alpha `$ intensity and line ratios in order to estimate the clumpiness of the ISM around $`\zeta `$ Oph. Our models enable us to investigate the penetration and escape of ionizing radiation from clumpy H ii regions into the more widespread WIM. In addition, we study the observational signatures of ionized/neutral interfaces in smooth and clumpy models. Section 2 presents the H$`\alpha `$ intensity and line ratio maps of the $`\zeta `$ Oph region, §3 summarizes model predictions for line ratios in H ii regions and at ionized/neutral interfaces, §4 and §5 present smooth and clumpy models for the observations, and we summarize our findings in §6.
## 2 The $`\zeta `$ Oph H ii Region
Figure 1 shows the structure of the H ii region around $`\zeta `$ Oph. The upper left panel shows the H$`\alpha `$ intensity from the SHASSA survey (Gaustad et al. 2001; Finkbeiner 2003) with an angular resolution of $`0.8\mathrm{}`$. The upper right panel shows the same region at $`1^{}`$ resolution from the WHAM Northern Sky Survey (Haffner et al. 2003). We have interpolated the irregularly-spaced WHAM observations onto a regular $`0\stackrel{}{\mathrm{.}}5\times 0\stackrel{}{\mathrm{.}}5`$ grid. The location of $`\zeta `$ Oph is shown with a cross. The WHAM image shows that the overall structure of the H$`\alpha `$-emitting gas is remarkably circularly symmetric, although there are strong fluctuations on smaller angular scales as shown in the SHASSA image. Diffuse H$`\alpha `$ emission associated with the $`\zeta `$ Oph H ii region is traced out to around $`6^{}`$ from the star, corresponding to a radius of about 15 pc, before trailing off into more diffuse foreground/background emission toward higher Galactic longitudes (on the left). On the other side of the H ii region the H$`\alpha `$ intensity becomes confused with that from the $`\delta `$ Sco H ii region (see Haffner et al. 2003). We do not consider data from this region in our analysis.
The sightline towards $`\zeta `$ Oph is well studied with many absorption line determinations of interstellar abundances (e.g., Morton 1975; Cardelli et al. 1993, 1994; Howk & Savage 1999). The star itself is catalogued as O9.5V, with an estimated temperature in the range $`\mathrm{32\hspace{0.17em}000}<T_{}<\mathrm{33\hspace{0.17em}000}`$ K (Code et al. 1976; Markova et al. 2004), and a distance of 140 pc (Perryman et al. 1997; Howk & Savage 1999). $`\zeta `$ Oph has a heliocentric radial velocity of $`10.7`$ km s<sup>-1</sup>, which is $`+3.3`$ km s<sup>-1</sup> relative to the local standard of rest. The WHAM data shows that the ionized gas in the H ii region lies at 0 km s<sup>-1</sup>. Its proper motion from HIPPARCOS is $`13\pm 1`$ mas yr<sup>-1</sup> in right ascension and $`25.5\pm 0.5`$ mas yr<sup>-1</sup> in declination, resulting in a space velocity of $`19.3\pm 0.4`$ km s<sup>-1</sup> relative to the LSR for a distance of 140 pc. With such a velocity, $`\zeta `$ Oph travels more than 19 pc in $`10^6`$ years, larger than the 15 pc radius of its own H ii region. Therefore over its lifetime, it appears that $`\zeta `$ Oph has moved from its birthplace and is now ionizing a different region of the ISM from where it was born. We will return to this point when discussing our photoionization models below.
The H$`\alpha `$ maps have not been corrected for foreground extinction by interstellar dust that may lie between us and the H ii region. In particular, the dust lane cutting across the lower left quadrant in the SHASSA image is likely in the foreground and not and integral part of the H ii region. Future H$`\beta `$ observations with WHAM will allow for dust correction and the determination of which features are internal to the H ii region and which are due to foreground dust. However, the \[N ii\]/H$`\alpha `$ and \[S ii\]/H$`\alpha `$ line ratio maps will not be affected by dust since H$`\alpha `$, \[N ii\], and \[S ii\] are very close in wavelength.
In addition to the H$`\alpha `$ maps, Fig. 1 shows WHAM line ratio maps of \[N ii\]/H$`\alpha `$ and \[S ii\]/H$`\alpha `$. These maps show these line ratios increase with distance away from $`\zeta `$ Oph. Figure 2 shows a scatter plot of H$`\alpha `$, \[N ii\]/H$`\alpha `$ and \[S ii\]/H$`\alpha `$ against radial distance from $`\zeta `$ Oph. The line ratios increase fairly smoothly with increasing distance and do not show the rapid increase expected towards the edge of a uniform density H ii region (see §3). For comparison with models, in the rest of the paper when showing the WHAM data we show the mean radial surface brightness and line ratios and the dispersions about the mean. The mean radial surface brightness is the average value for the H$`\alpha `$ intensity (or line ratios) in circular annuli centered on $`\zeta `$ Oph.
## 3 Interface Emission: Predictions and Observations
One dimensional photoionization models predict rapid increases in the ratios of the projected intensities of \[O i\], \[N ii\], and \[S ii\] relative to H$`\alpha `$ towards the edge of uniform density H ii regions (e.g., Henney et al. 2005). The increased line ratios occur in the transition region at the edge of the Strömgren sphere where the fraction of neutral gas is increasing and the temperature is rising due to hardening of the radiation field. The wavelength dependence of H<sup>0</sup> opacity allows the highest energy photons to penetrate the largest distances, resulting in a rising temperature in the interface region (e.g., Osterbrock 1989).
Observations of the diffuse ISM near the Galactic plane show low \[O i\]/H$`\alpha `$ line ratios (Reynolds et al. 1998), suggesting that interfaces between ionized and neutral gas have either not been detected or do not contribute significantly to the emission. Photoionization models can explain the low \[O i\]/H$`\alpha `$ observations if the WIM is highly ionized ($`\mathrm{H}^0/\mathrm{H}<0.1`$ throughout), implying that ionizing radiation escapes from fully ionized (i.e., density bounded) H ii regions having few H i – H ii interfaces (e.g., Mathis 1986; Domgoergen & Mathis 1994; Mathis 2000; Sembach et al. 2000). Also, recent work by Giammanco et al. (2004) showed that if a clumpy H ii region is comprised of a mixture of fully and partially ionized spherical clouds, the line ratios increase more slowly with radius than in uniform density, ionization bounded models (e.g., see the radial dependence of \[O i\]/H$`\beta `$ in their Fig. 6). This is because the combination of many spherical clouds along a given sightline through the H ii region sample a wide range of ionization parameters and hence the line ratios do not show the rapid increase of uniform, ionization bounded models.
While the predicted rapid rise of temperatures and line ratios at the edges of H ii regions has not been observed (e.g., Pauls & Wilson 1977), gradual rises in the line ratios have been observed with height above the midplane ionizing sources in the Milky Way (Haffner et al. 1999) and several other edge-on galaxies (Rand 1998; Otte et al. 2001, 2002). The elevated \[N ii\]/H$`\alpha `$ line ratios observed at large distances from the midplane in the Perseus Arm (Haffner et al. 1999) have been interpreted as due to an extra heat source in addition to photoionization heating (Reynolds et al. 1999). This additional heat source must have the property that it dominates over photoionization heating at low densities. A constant energy source or one that is proportional to $`n_\mathrm{e}`$ would suffice, since photoionization heating is proportional to $`n_\mathrm{e}^2`$. Candidates for the additional heating include photoelectric heating from grains (Reynolds & Cox 1992), magnetic reconnection (Raymond 1992), dissipation of turbulence (Minter & Spangler 1997), shocks and cooling hot gas (Slavin, Shull, & Begelman 1993; Collins & Rand 2001). However, hardening of the radiation field also increases the temperature of photoionized gas and may help to explain in part some of the observed elevated line ratios (e.g., Bland-Hawthorn, Freeman, & Quinn 1997; Elwert 2003; Wood & Mathis 2004).
In this paper we present models for the H ii region associated with $`\zeta `$ Oph. Additional heating will not be important for this source since the gas density is larger and the line ratios are much lower than observed at high latitude in the Perseus arm. The WHAM observations presented in Figures 1 and 2 show that \[N ii\]/H$`\alpha `$ increases away from the source, but does not appear to show the rapid rise at the edge of the H ii region predicted by spherical models with uniform density (see §4.1 below). To understand this behavior we have investigated photoionization models that incorporate various smooth and clumpy density structures for the H ii region and the effects of diffuse foreground/background emission. We discuss the plausibility of the smooth models and the insight three-dimensional models can provide into the structure of the ISM around $`\zeta `$ Oph and the escape of ionizing photons from H ii regions. Our 3D models produce similar results to those of Giammanco et al. (2004) for radial line ratio gradients, but our 3D radiation transfer allows for shadowing of clumps and multiple clumps along any sightline from the star.
## 4 Photoionization Models
To model the WHAM observations of the $`\zeta `$ Oph H ii region, we use the three-dimensional photoionization code described in Wood, Mathis, & Ercolano (2004). This code calculates the 3D ionization and temperature structure for arbitrary geometries and illuminations, keeping track of the ionization structure of H, He, C, N, O, Ne, and S. The opacity is from to H<sup>0</sup> and He<sup>0</sup> and we only consider photon energies in the range 13.6 eV to 54 eV. Heating is from photoionization of H and He and cooling is from H and He recombination, free-free emission, and collisionally excited line emission from C, N, O, Ne, and S. The output of our code is the ionization and temperature structure from which we calculate emissivities and intensity maps for the various emission lines we wish to study. A complete description of the code is presented in Wood et al. (2004).
We perform the radiation transfer for stellar and diffuse recombination ionizing radiation on a 65<sup>3</sup> linear Cartesian grid. Due to the resolution of our grid, we do not consider emission from cells that are more neutral than $`\mathrm{H}^0/\mathrm{H}0.25`$ (see discussion in Wood et al. 2004). These cells lie towards the edge of the ionized volume and account for less than 20% of the ionized cells in our ionization bounded models. We have constructed intensity and line ratio maps including these cells and find that our results for the \[N ii\]/H$`\alpha `$ line ratio maps do not change at the $`5`$% level. Moreover, in the cells at the ionization boundary our code is incomplete, since we ignore the effects of dynamics and shocks at the ionization front (Henney et al. 2005). Higher resolution simulations ($`128^3`$ grids) do not appreciably change our results, again at around the $`5`$% level.
Our code performs well compared to other photoionization codes in predicting ionization fractions, temperatures, and line strengths. Here we focus on modeling the observed H$`\alpha `$ intensity and \[N ii\]/H$`\alpha `$ maps and do not consider the \[S ii\]/H$`\alpha `$ observations. Modeling emission from sulfur is problematic due to the unknown dielectronic recombination rates for S (Ali et al. 1990). In addition we do not consider the effects of dust within the H ii region and note that our temperatures may be slightly higher than calculations that include cooling from more elements (e.g., see Sembach et al. 2000). Since the goal is to investigate the 3D structure of the H ii region by modeling variations in the line ratios towards the edge of the H ii region, ignoring these effects will not change our conclusions on the H ii region density structure. Inclusion of these additional effects will be required for more comprehensive models of future observations of H$`\beta `$, \[O i\], \[O ii\], \[O iii\], and He i. In particular, we may be able to model data from several H ii regions to place constraints on the S dielectronic recombination rates.
Relative to H, the adopted abundances for He, C, N, O, Ne, and S in our models are 0.1, 140 ppm, 75 ppm, 319 ppm, 117 ppm, and 18 ppm. These are interstellar gas phase abundances from the compilation of Sembach et al. (2000). The ionizing spectrum is taken to be that of a 32 000 K WM-basic model atmosphere from the library of Sternberg, Hoffmann, & Pauldrach (2003). This temperature is within the range estimated for an O9.5V star (see compilation of temperature/spectral type scales and discussion in Harries, Hilditch, & Howarth 2003, Table 4). The WM-basic model atmospheres include the effects of line blanketing and stellar winds (Pauldrach et al. 2001). The ionizing luminosity and density are varied in the models to reproduce the observations of H$`\alpha `$ and \[N ii\]/H$`\alpha `$. The following sections present results for smooth models (with and without density gradients) and three-dimensional, hierarchically clumped models.
### 4.1 Smooth Spherically Symmetric Models
Our spherically symmetric models are centered on a source with an ionizing photon luminosity $`Q`$ (s<sup>-1</sup>), and the density within the cloud (cm<sup>-3</sup>) is given by
$$n(r)=n_0\left(\frac{r}{1\mathrm{pc}}\right)^p.$$
(1)
where the parameter $`p`$ controls the radial density gradient within the cloud. In all models (smooth and clumpy) we have left the inner 10% of the grid empty, which prevents the clumping algorithm randomly placing the star in the center of a dense clump. The grid is 60 pc on a side and contains $`65^3`$ cells, while the observed $`\zeta `$ Oph H ii region is about 40 pc in diameter but has some extensions to greater distances. We have investigated many ($`p`$, $`n_0`$, $`Q`$) combinations for ionization and density bounded models and show those that reproduce the overall H$`\alpha `$ intensity level observed by WHAM. Table 1 summarizes the range of adopted model parameters and the results of the models are compared to the observations in figure 3.
The radial variation of H$`\alpha `$ and \[N ii\]/H$`\alpha `$ may be influenced by the presence of foreground and background emission. This is illustrated in Figure 3 which shows models without (left panels) and with (right panels) the addition of a constant intensity representing foreground/background emission. Both the smooth (shown here) and clumpy models (shown in §4.2) have been convolved with a Gaussian beam to simulate the $`1^{}`$ angular resolution of WHAM. The diffuse background is taken to be 1.5 R at H$`\alpha `$ and 0.6 R for the \[N ii\] simulation. These values provide a good fit to the observations (see Fig. 2) and also are typical for the diffuse interstellar medium near the Galactic latitude ($`b=24^{}`$) of $`\zeta `$ Oph (see also Haffner et al. 2003; Madsen 2004). The addition of diffuse foreground/background emission (which is an unavoidable contaminant in real observations) clearly suppresses the very large line ratios that are predicted in the faintest outermost regions ($`20`$ pc) of the H ii region. However, even with the addition of foreground/background emission, the uniform density, ionization bounded model (i.e., a classic Strömgren sphere) exhibits a rise in the \[N ii\]/H$`\alpha `$ at the edge of the H ii region that is steeper than the observations. The H$`\alpha `$ intensity gradient in such models is also much steeper at the edge of the ionized region than observed (see also the giant H ii region models of Giammanco et al. 2004). Density bounded or “leaky H ii region” models (Fig. 3) do not exhibit the rapid rise in line ratios because there is no ionized/neutral interface in the simulations (e.g., Mathis 2000; Sembach et al. 2000). The density bounded models shown in Fig. 3 allow about 40% of the ionizing photons to escape. Changing the radial density gradient from uniform ($`p=0`$) to steeper laws ($`p=1`$, 2) changes the H$`\alpha `$ radial surface brightness profile. Models with gradients steeper than $`p=1`$ do not match the WHAM data. Models with $`p=1`$ appear to give a reasonable match to the H$`\alpha `$ intensity, but all the density bounded models underpredict the \[N ii\]/H$`\alpha `$ observations.
As mentioned previously, $`\zeta `$ Oph is at a different radial velocity from its associated H ii region, so models that place the star at the centre of a spherical cloud with radial density gradients seem contrived. A more realistic scenario is that $`\zeta `$ Oph is ionizing a clumpy medium and we explore 3D cloud geometries and the effects on the intensity and line ratio maps in the next section.
### 4.2 Hierarchically Clumped Models
There is much evidence for hierarchical structure in the ISM, with surveys of clouds revealing fractal structures (e.g., Elmegreen & Falgarone 1996). We therefore investigate 3D H ii regions which have hierarchically clumped density structures. A 3D hierarchical density structure is created using the algorithm presented by Elmegreen (1997). We use five hierarchical levels and the density in the grid is set as follows. At the first level we randomly cast $`N`$ points with $`x`$, $`y`$, and $`z`$ coordinates in the range (0,1). At each subsequent hierarchical level we cast $`N`$ random points around each of the points cast at the previous level. The casting length for the $`x`$, $`y`$, and $`z`$ points at each subsequent level is in the range $`\pm \mathrm{\Delta }^{(1H)}/2`$, where $`H`$ is the hierarchical level being cast, and $`f=\mathrm{log}N/\mathrm{log}\mathrm{\Delta }`$ is the fractal dimension of the hierarchical structure. The casting length, $`\mathrm{\Delta }`$, gets smaller for subsequent hierarchical levels. Note that in this description $`\mathrm{\Delta }`$ is dimensionless and we construct the hierarchical grid within a cube with side of unit length. Actual dimensions are constructed by mapping the hierarchical grid onto our 3D density grid, so that the unit cube corresponds to the physical size of our density grid. We use a fractal dimension $`f=2.6`$, appropriate for interstellar clouds (e.g., Sanchez, Alfaro, & Perez 2005). Using this algorithm, the density of a cell in our grid is proportional to the number of points cast at the final level that lie in the cell. If points are cast beyond the $`x`$, $`y`$, or $`z`$ boundaries of the grid, they are added into the corresponding cells on the opposite side of the grid, e.g. if a point is cast at $`x=x_{\mathrm{max}}+ϵ`$, we change its coordinate to be $`x=ϵ`$ and similarly for $`y`$ and $`z`$ coordinates. Figure 4 illustrates this algorithm in 2D for a three tier hierarchical scheme, with four random points cast at each level. The filled in squares are the first castings, diamonds the second level, and crosses at the third level. The dotted lines show overlaid grid cells. We assume that a fraction, $`f_{\mathrm{smooth}}`$, of the average density is present in every cell. The remaining mass is distributed in proportion to the number of crosses that are in each cell.
The smooth component represents unresolved fractal structure that cannot be resolved within our grid. In real nebulae the hot, shocked stellar winds may create very low density regions between the clumps of gas. The resulting clumpy density grid is renormalized so that it has the same mean density structure as the smooth grid. We generally adopt $`f_{\mathrm{smooth}}=1/3`$, but also consider the case where there is zero smooth density component. Our adopted value of $`f_{\mathrm{smooth}}`$ comes from radiation transfer models of the penetration of radiation into turbulent cloud models from hydrodynamic simulations. A value of $`f_{\mathrm{smooth}}=1/3`$ provides a good match for the internal intensity levels between the hydrodynamical simulations and density structures generated with the fractal algorithm (Bethell et al. 2004).
To investigate the clumpiness of the ISM around $`\zeta `$ Oph we varied $`N`$ at the first hierarchical level, using $`N_1=32`$, 64, 128, 256, and 512, but keeping $`N=32`$ at all four subsequent hierarchical levels. With this modification to the clumping algorithm the relation between casting length, $`\mathrm{\Delta }`$, and fractal dimension, $`f=\mathrm{log}N/\mathrm{log}\mathrm{\Delta }`$, is determined assuming $`N=32`$. Increasing $`N_1`$ at the first level gives a progressively smoother density structure at large spatial scales, but maintains the observed fractal dimension of interstellar clouds on smaller scales. We have not investigated the effects on the ionization structure and line ratios of changing the value $`N=32`$ at the second and higher hierarchical levels. Increasing $`N`$ will give a larger dynamic range and finer structure in the hierarchical density grid, because the total number of points used to determine the density is given by the number of castings at the final level, $`N_1\times N^4`$. The choice of $`N=32`$ is not observationally motivated, but is determined by the resolution of our density grid. Using very large values of $`N`$ is not justified because the fine structure will not be resolved by our grid.
The algorithm produces maximum density contrasts between the densest clumps and the smooth component of 65, 35, 20, 10, for $`N_1=32`$, 64, 128, and 512 respectively. A small $`N_1`$ allows ionizing radiation to penetrate further between the clumps than in a uniform medium (Elmegreen 1997). The fraction of cells that contain only the smooth component of density, without any of the points that were cast randomly, is 65%, 45%, 20%, and 0.1% for these models. If the density in the smooth component is almost zero, ionizing radiation can escape along the empty sightlines into the larger scale ISM.
Having a smooth component of zero density did not significantly change the morphology of our H$`\alpha `$ maps compared to the $`f_{\mathrm{smooth}}=1/3`$ that we adopt because the emission is dominated by the higher density ionized clumps. For $`N_1=32`$ and $`f_{\mathrm{smooth}}=1/3`$, the smooth component accounts for about 1/3 of the H$`\alpha `$ intensity level. The smooth contribution to the H$`\alpha `$ intensity decreases with increasing $`N_1`$. Elmegreen (1997) speculates that much of the volume of a fractal ISM is at very low density, like our models with no smooth component.
We investigated many different clumpy models for each $`N_1`$ value by changing the random number seed used for setting up the hierarchical density grid. For a given $`N_1`$, each model has a different density structure and hence slightly different H$`\alpha `$ and \[N ii\] intensity profile due to the random casting of clumps around the star, but the overall features in the intensity maps are similar. For each $`N_1`$ of our clumpy models, Figures 5 through 8 show one pixel thick slices through one of the density grids and slices showing the hydrogen ionization fraction $`\mathrm{H}^0/\mathrm{H}`$. The figures also show the H$`\alpha `$ intensity and \[N ii\]/H$`\alpha `$ line ratio maps, without the addition of uniform foreground/background emission. Figures 9 to 12 show the radial variation of H$`\alpha `$ and \[N ii\]/H$`\alpha `$ with the addition of uniform foreground/background emission. For each $`N_1`$ model, the figures show one model with the standard deviations about the mean and also the radial profiles for simulations of five different density grids.
## 5 Results
Our models are constrained by the emission line intensities observed at various angles from the star (or distances projected onto the sky) and their dispersions. For all clumpy models the ionizing luminosity is $`Q=8\times 10^{47}\mathrm{s}^1`$ and the mean density of the H ii region is $`n=2\mathrm{cm}^3`$. Both the luminosity and density are about 50% smaller than estimated by Elmergreen (1975) from analysis of older H$`\alpha `$ data from Reynolds et al. (1974). Our estimates for the luminosity and density are well constrained by the intensity and angular extent of the H$`\alpha `$ data. We find that $`N_164`$, followed by coarser clumping on the smaller scales, provides the best fit to the WHAM data. This implies that the real ISM in this direction is not hierarchical, but is smoother on a scale of $`10`$ pc than at smaller scales. This result reflects the general circular symmetry in the WHAM H$`\alpha `$ image in Fig. 1, but the small scale irregularities are not revealed in the WHAM data because of its relatively coarse angular resolution ($`1^{}`$).
As expected, Figures 5 through 12 show that the clumpy model that most resembles a Strömgren sphere has $`N_1=512`$, since distributing many random points approximates a uniform distribution. This model has a somewhat irregular boundary, but there are no sightlines from the star that do not intersect dense clumps, so overall the model has a very smooth structure and exhibits the large increase in the \[N ii\]/H$`\alpha `$ line ratio at the edge of the H ii region. As $`N_1`$ is reduced, the medium becomes less smooth and there are sightlines from the star through which ionizing radiation can traverse large distances before intersecting a clump, or not intersect any clumps. Models with $`N_1=32`$ show a gradual radial decline of the H$`\alpha `$ intensity. As $`N_1`$ is increased, the fraction of low density sightlines decreases and the H$`\alpha `$ exhibits a steeper radial gradient. The concavity of the radial H$`\alpha `$ surface brightness profiles change for larger $`N_1`$ and approach that of the uniform density Strömgren sphere in Figure 3. The trends in our models are similar to the radial surface brightness profiles and line ratios found by Giammanco et al. (2004) in their H ii region models comprising spherical blobs surrounding an ionizing source.
Within the context of our clumping algorithm, we find that models with $`N_1=32`$ produce too much H$`\alpha `$ at large radii and variations about the mean radial intensity are larger than observed. Models with $`N_1256`$ resemble Strömgren spheres, producing too steep a gradient of the H$`\alpha `$ and a rapid rise of \[N ii\]/H$`\alpha `$ at the edge of the H ii region and small variations about the mean values. Models with $`N_1=64`$ and 128 provide better matches to the WHAM data in reproducing the observed levels of H$`\alpha `$ and \[N ii\]/H$`\alpha `$, their radial variations, and standard deviations.
While the intensity maps for this low density H ii region model are not very sensitive to the smooth component, the smooth component has a large influence on the ionizing radiation that may escape the H ii region and ionize the larger scale diffuse ionized gas. Models with $`N_1=32`$, 64, and 128 and zero smooth density component ($`f_{\mathrm{smooth}}=0`$) typically allow around 10%–25%, 3%–15%, and 1%–3% respectively of the stellar ionizing photons to escape the H ii region directly along the zero density paths from the star. The large variations in the escape fractions for each $`N_1`$ value are due to the random placing of clumps around the star. Our 3D photoionization models therefore confirm the analysis of Elmegreen (1997) that a fractal ISM will allow for ionizing photons to traverse much larger distances than a smooth density structure. Further clearing of a fractal ISM by feedback from photoionization (Dale et al. 2005) and stellar winds will likely provide an even larger escape fraction from H ii regions than those determined from our models.
Since the widespread H$`\alpha `$ emission from the Galactic WIM requires about 15% of the ionizing photons from Galactic O stars (Reynolds 1995), it appears that a 3D fractal ISM structure could allow for such leakage from H ii regions. Our hierarchical density models for the ISM around $`\zeta `$ Oph with $`N_1=64`$ and 128 typically provide escape fractions of less than 15%, and even lower if the density in the smooth component is not zero. Therefore, the ionizing sources responsible for powering the WIM in the Galaxy must reside in regions of the Galaxy that are more porous than the region around $`\zeta `$ Oph. Perhaps their environment has been shaped more by turbulent motions or are regions cleared out by stellar winds, thus enabling the escape of ionizing photons. Future models of WHAM observations of other H ii regions will enable us to determine whether the porosity levels for the ISM around $`\zeta `$ Oph are typical.
Our models placed the star at the center of a fractal cloud, so models where the star is at the edge of the cloud will allow for larger leakage, since the cloud no longer covers the whole sky as seen from the star. The major question from this work relating to ionization of the diffuse ionized gas is what is the density of the smooth component in fractal models? We cannot determine how much H$`\alpha `$ emission from the WIM comes from an extensive, relatively smooth component and how much is from ionized clumps and the ionized faces of clouds. Further progress to answer this question may be made by conducting photoionization simulations on the 3D density grids being produced by global models of the ISM (e.g., de Avillez & Berry 2001). Such models predict the ISM density structure and the locations of ionizing sources and their validity may be tested against observations through 3D photoionization simulations.
## 6 Summary
We have presented smooth and clumpy 3D photoionization models for the $`\zeta `$ Oph H ii region. We find that a star with $`T=\mathrm{32\hspace{0.17em}000}`$ K and ionizing photon luminosity $`Q=8\times 10^{47}\mathrm{s}^1`$ and mean density of the H ii region of $`n2\mathrm{cm}^3`$ reproduce the observations. These numbers are well constrained by the WHAM data. The stellar temperature and ionizing luminosity are remarkably close to recent determinations of stellar properties for O9.5V stars (Martins et al. 2005, Table 4). Larger densities require a larger ionizing luminosity to match the extent of the H ii region, but then the H$`\alpha `$ values are larger than observed. Smaller densities have the opposite effect, requiring a lower ionizing luminosity to match the extent of the H ii region, but then the simulated H$`\alpha `$ intensity is much smaller than observed.
Smooth models that best match the observations are density bounded and have a radial density gradient that falls off no faster than $`r^1`$. However, the velocity offset between $`\zeta `$ Oph and its H ii region suggest that these uniform models are unrealistic as they require the chance alignment of the star at the center of the cloud. The inclusion of foreground/background emission suppresses the rapid rise in the \[N ii\]/H$`\alpha `$ line ratio predicted at the edge of uniform Strömgren spheres.
Hierarchically clumped models, which reproduce the observed fractal structure in the ISM, also reproduce the H$`\alpha `$ surface brightness and the shallow gradient of \[N ii\]/H$`\alpha `$ away from $`\zeta `$ Oph. The best models suggest that around 20% ($`N_1=128`$) to 50% ($`N_1=64`$) of the H ii region volume is occupied by gas distributed in a smooth component or at very low density, with the remainder in clumps. Three-dimensional models that are more uniform do not match the observations.
The large, low density regions in 3D models could allow ionizing photons to escape at levels that sustain the ionization of the large scale diffuse ionized gas in the Galaxy. However, the escape fractions from our preferred clumpy models for $`\zeta `$ Oph are less than the 15% escape fraction required to ionize the WIM. Therefore we conclude that the stars responsible for ionizing the WIM must reside in regions of the ISM that are more porous than that surrounding $`\zeta `$ Oph, possibly evacuated by supernovae and/or stellar winds. A more detailed investigation into the leakage of ionizing photons to the WIM requires consideration of dynamical clearing and the time evolution of the ionizing luminosity of O stars.
Future work will include more comprehensive modeling of the planned mapping observations of $`\zeta `$ Oph in additional emission lines. The WHAM survey has mapped out several other low emission measure H ii regions and 3D modeling will allow us to test whether the porosity values we estimate for the $`\zeta `$ Oph region are common in other regions. Modeling the \[S ii\] emission from many such regions may help to place constraints on the as yet unknown dielectronic recombination rates for sulfur. From a theoretical perspective, global dynamical models of the ISM can be tested by applying 3D photoionization codes and comparing their observational signatures with the high resolution observations now available of diffuse ionized gas in our own and other galaxies.
We thank Lynn Matthews and Jane Greaves for comments on an early version of this paper. We acknowledge funding from a PPARC Advanced Fellowship (KW) and a PPARC Visiting Fellowship to the University of St Andrews (JSM, RJR, GM). WHAM is funded by grants from the NSF (AST 0204973) and the University of Wisconsin Graduate School.
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# Dynamical quintessence fields Press-Schechter mass function: detectability and effect on dark haloes
## 1 Introduction
Recent evidence is piling up in favour of a dark energy component in the dynamics of the universe forming up to 70% of the energy density: from SNIa magnitudes \[Riess et al. 1998, Perlmutter et al. 1998\]; and more recently \[Wang & Mukherjee 2004, Nesseris & Perivolaropoulos 2004, Riess et al. 2004, Daly & Djorgovski 2004, Biesiada et al. 2005\]; from the CMB (Cosmic Microwave Background) measured by WMAP (Wilkinson Microwave Anisotropy Probe) coupled with complementary inputs \[Bennett et al. 2003, Page et al. 2003\] (for reviews, see \[Padmanabhan 2003, Sahni 2004\]). The nature of this dark energy is still unknown but could lie within two major families: it is either coming from a static cosmological constant or developing from a dynamical scalar field just emerging from subdominance called quintessence \[Ratra & Peebles 1988\]. The interest of quintessence lies in the possibility it holds to solve the four major problems raised by the cosmological constant model: the problem of a need for fine-tuning the cosmological constant to unnatural values of vacuum energy, the problem of coincidence between its present energy density and the order of magnitude of the critical density, the problem of the equation of state of the cosmological fluid, which can only take one value for a cosmological constant, and the problem of building a model that would naturally come up with solutions to these problems while coming from fundamental physics.
The drive that lead to quintessence models was based on the attempt to address those four problems of the observed behaviour of dark energy. The first problem requires a theory which yields the present value for the dark energy without the need to tune initial conditions at unnaturally small values compared with the natural energy scales of the early universe. This addresses the huge discrepancy (typically $`M_{Pl}^4/\rho _{DE}10^{120}`$; $`M_{Pl}`$: Planck mass; $`\rho _{DE}`$: Dark Energy density) commonly found between classical quantum vacuum energy calculations and its cosmological measurements. The coincidence problem deals with the current transition status of the universe between dominations by matter and dark energy. It requires a theory that yields such a tuned peculiar equilibrium without strong constraints on its initial conditions. it deals with the value of the dark energy component $`P/\rho `$ ratio, which is constrained to lie within $`\left[\begin{array}{ccc}1\hfill & ;\hfill & 0.6\hfill \end{array}\right[`$ \[Wang et al. 2000\] or even within $`\left[\begin{array}{ccc}1\hfill & ;\hfill & 0.8\hfill \end{array}\right[`$ \[Efstathiou 2000, Hannestad & Mörtsell 2004\] (for dark energy equation of state supposed to respect the weak energy condition, thus $`P/\rho >1`$). Finally the model building problem express the need to have a theoretically motivated dark energy potential.
The fine tuning problem has been the main drive for quintessence models, hence it is solved by almost all proposed quintessence potentials, involving high energy physics natural energy scales. Various models have been proposed but few are deeply grounded in high energy physics. The main examples of motivated potentials are the original Ratra-Peebles’ potentials \[Ratra & Peebles 1988\] that have been linked with the context of global SUSY \[Binetruy 1999\] and the so-called SUGRA potentials derived from supergravity arguments \[Brax & Martin 1999, Brax & Martin 2000\]. The tracking property is attracting much interest because it solves the coincidence problem naturally. It defines classes of potentials which Klein-Gordon equation admits a time varying solution that acts as an attractor for wide ranges of field initial conditions. In \[Steinhardt et al. 1999\], general conditions to obtain a tracking potential are given.
In front of the wealth of scalar fields available from physics beyond the standard model, the need for discrimination between different quintessential models proposed in the literature is a strong drive to confront their predictions with observable features. Since the advent of the COBE satellite results and the pursuit of refined Cosmic Microwave Background Radiation (CMBR) measurements, the attention of cosmologists dealing with quintessence had been focused on CMBR anisotropy measurement as primary probes for the various quintessence models (e.g. \[Brax et al. 2000\]). That has foreshadowed the possibility offered by quintessential cosmic dynamics to alter in turn the formation of large scale structures, that are more readily available to observations. Some of the earlier attempts in this direction that have been started out have restricted themselves to linear or perturbation theory of structure formation \[Benabed & Bernardeau 2001\], to pseudo-quintessence models (approximation of a constant equation of state parameter different than the $`\mathrm{\Lambda }`$ term; \[Lokas & Hoffman 2001, Lokas et al. 2004, Kuhlen et al. 2005\]) or have not pushed the envelope beyond the study of the spherical collapse model \[Mainini et al. 2003b, Mota & van de Bruck 2004, Nunes & Mota 2004\] or the impact on halo concentrations \[Dolag et al. 2004, Kuhlen et al. 2005\].
The aim of this paper is to compute the non-linear mass function of collapsed structures in the presence of a quintessence field. One approach is to use direct numerical simulation \[Klypin et al. 2003, Linder & Jenkins 2003, Dolag et al. 2004, Macciò et al. 2004, Kuhlen et al. 2005, Solevi et al. 2005\], but this is only practical for a few models with well defined parameters and gives little physical insight into the influence of the quintessence potential. Our approach takes into account the fully dynamical nature of the field using the methodology developed by Press & Schechter \[Press & Schechter 1974, hereafter PS\]. This method uses a spherically symmetric dynamical model to relate the collapse of massive structures to a density threshold in the linearly extrapolated density field. In this way, it is possible to apply Gaussian statistics to the initial density field in order to count the numbers of collapsed structures above a given mass threshold at a particular epoch. Because the method is nearly analytic, it can be applied to rapidly expore the parameter space of possible quintessence models and inverted to constrain quintessence models using observational data. Some approaches using semi-analytical methods have been published in the course of the present work with a restricted range of quintessence potentials \[Mainini et al. 2003a, Mainini et al. 2003b, Solevi et al. 2004, Nunes & Mota 2004\], some restricting to very indirect observables \[Mainini et al. 2003b, Nunes & Mota 2004\]. None have yet looked at the same time at observables like the mass function with many different potentials and an approach aimed at understanding the dominant physics involved, using the full potential of semi-analytical methods.
The paper is laid out as follows. In section (2), homogeneous universe models with quintessence are explored and familiarity with the physics of quintessence is developed. Some results on the behaviour of the universe’s radial scale in various such models are presented, showing the possibility offered by large scale structures to probe the quintessence potentials. In section (3), the collapse of a top hat spherical model for a single primordial inhomogeneity is computed. This provides a simple way of linking perturbations in the initial density field with collapsed structures. In section (4), we combined this model with information from the statistics of the initial inhomogeneities. Using a PS type scheme we obtain predictions for the mass function in the presence of various quintessential models, and compare the evolution of different models. Finally, in section 5 we summarise the results, and explore the constraints that future astronomical surveys will be able to set on the form of the quintessence potential.
## 2 Homogeneous quintessence evolution
In this section we will review the features in the homogeneous models of dynamical dark energy and emphasize that which points toward an effect on structure formation. This roadmap will later be useful for interpreting our results. The first step is to establish the behaviour of various quintessence models within FLRW-type solutions.
### 2.1 The dynamical system
In a homogeneous model, the system evolution is entirely defined by the evolution of its scale factor and its quintessence field with its time derivative. The Einstein’s Field Equations governing the Friedman Lemaitre Robertson Walker universe and the Klein-Gordon Equation for the quintessence homogeneous, time dependent, scalar field can be written as evolutions for the scale factor $`a`$ and the scalar field $`Q`$ .
The first Friedman equation of our model can be expressed as
$`\dot{a}^2`$ $`=`$ $`a^1\left(\mathrm{\Omega }_{m_0}+\mathrm{\Omega }_{r_0}a^1+{\displaystyle \frac{\rho _Q}{\rho _{c_0}}}a^3+\mathrm{\Omega }_{\mathrm{\Lambda }_0}a^3+\mathrm{\Omega }_{k_0}a\right),`$ (2.1)
the dots refer to cosmic time derivatives and will refer (e.g. in Eq.2.2 below) to the total derivative with respect to the field Q. Conventionally, subscript 0 refers to the present epoch. $`\mathrm{\Omega }_X=\rho _X/\rho _c`$ represents the density parameter of species X, $`\rho _c`$ is the critical total density required by a flat universe model. We introduce here the density factor $`\mathrm{\Omega }_X^{}=\rho _X/\rho _{c_0}`$; note $`\mathrm{\Omega }_{X_0}^{}=\mathrm{\Omega }_{X_0}`$.
Each term corresponds respectively to the contribution in matter (dark and luminous: $`\mathrm{\Omega }_{m_0}`$), radiation (neutrinos and photons: $`\mathrm{\Omega }_{r_0}`$), cosmological constant ($`\mathrm{\Omega }_{\mathrm{\Lambda }_0}`$), either time dependent ($`\mathrm{\Omega }_Q^{}=\frac{\rho _Q}{\rho _{c_0}}`$) or present epoch ($`\mathrm{\Omega }_{Q_0}=\rho _{Q_0}/\rho _{c_0}`$, as follows) quintessence density factor and the corresponding FLRW curvature $`\mathrm{\Omega }_{k_0}`$, with the definition
$`\mathrm{\Omega }_{k_0}`$ $`=`$ $`1\left(\mathrm{\Omega }_{m_0}+\mathrm{\Omega }_{r_0}+\mathrm{\Omega }_{Q_0}+\mathrm{\Omega }_{\mathrm{\Lambda }_0}\right).`$
In this work we have assumed throughout a globally flat geometry of the universe ($`\mathrm{\Omega }_{k_0}=0`$) and a fully dynamical dark energy ($`\mathrm{\Omega }_{\mathrm{\Lambda }_0}=0`$).
It should be noted that the very non-linear Klein-Gordon equation requires forward integration to reach the tracking solution while the boundary conditions in the density parameters are set for a backward time integration.
The natural units for a primordial quantum field like quintessence are Planck units (Planck time: $`t_{Pl}`$). In those units, expressing the characteristic time, which is the Hubble time: (\[Bahcall et al. 1999, Freedman 2000\] $`t_{H_0}^1=H_065kms^1Mpc^12.10610^{18}s^1=(8.65\times 10^{60}t_{Pl})^1`$), shows a discrepancy between the two timescales of the problem of more than 60 orders of magnitudes. Since the Planck time should govern the evolution of the field, this could have revealed problematic if we were not focusing on cosmic evolution of the quintessence field: neglecting the effects very of rapid fluctuations, we can observe the evolution of the field as following a regular evolution of its energy density if we adopt the Hubble time as characterizing our system.
We also assumed spatial homogeneity of the field. \[Ma et al. 1999, Lokas & Hoffman 2001\] argue for homogeneous Q, given that they found the smallest scale of Q fluctuations to be larger than the clusters scale. Nevertheless it has to be mentioned that their results hold their validity from their studies in the mildly non-linear regime. \[Mota & van de Bruck 2004, Nunes & Mota 2004, Macciò et al. 2004\] argue that the highly non-linear regime involved might require some quintessence clustering or coupling to the dark matter. We will concentrate here on effects in non-clustering minimally coupled models as a threshold model to ascertain the impact of quintessence on non-linear clustering.
In Eq.(2.1), where we also assume no curvature nor cosmological constant contribution, the effect of quintessence is focused in the density factor, involving the energy density of the field $`Q`$ defined with its potential: in our units it reads $`\mathrm{\Omega }_Q^{}=\frac{8\pi }{3}\left(\frac{\dot{Q}^2}{2}+\frac{V(Q)}{_0^2}\right)`$ following notations defined above and with $`\dot{Q}`$, the quintessence field (Hubble) time derivative, $`V(Q)`$ its potential energy and $`_0`$ the Hubble constant in Planck mass units. Its dynamics is governed by the Klein-Gordon equation in the case of homogeneity:
$`\ddot{Q}`$ $`=`$ $`3{\displaystyle \frac{\dot{a}}{a}}\dot{Q}_0^2V^{}(Q).`$ (2.2)
Though the pressure of the field is not involved in its homogeneous evolution, it is crucial to the effect of quintessence on non-linear collapse: for a scalar field recall that (with cosmic time) $`P_Q=\left(\frac{1}{2}\dot{Q}^2V(Q)\right),`$ and together with the density, they define the acceleration state of the universe. One can characterise it using the equation of state
$$\omega _Q=P_Q/\rho _Q$$
(2.3)
(see figure 1’s upper panels). In terms of the energy momentum tensor of the field, recall that for scalar fields with these definitions of density and pressure, we have $`\rho _QT_{00};T_j^iP_Q\delta _j^i.`$
Now, we will discuss a variety of potentials proposed for quintessence models and their previously known homogeneous properties, emphasizing their impact on matter domination.
### 2.2 Explorations with several tracking potentials
In this section we restrict our choice upon a set of potentials and recall their previously studied equation of state (Eq.2.3) and density parameter ($`\mathrm{\Omega }_Q`$) homogeneous evolution, emphasizing behaviours that can affect the formation of large scale structures.
We have narrowed our study on such potentials from the literature that we found to agree (at least marginally) with the \[Steinhardt et al. 1999\] general conditions to obtain a tracking potential. We therefore selected several forms of potential (the list is not exhaustive) \[Ratra & Peebles 1988, Brax & Martin 1999, Brax & Martin 2000, Ferreira & Joyce 1998, Steinhardt et al. 1999\] for the rest of this study. The Ratra-Peebles potential \[Ratra & Peebles 1988, hereafter RP or Ratra-Peebles\], first discussed potential in the literature has retained its interest with its more recent discussion within the context of global SUSY \[Binetruy 1999\]. The \[Brax & Martin 2000, hereafter SUGRA\] potential has been motivated in the framework of supergravity low energy approximation. The simple exponential \[Ferreira & Joyce 1998, hereafter FJ or Ferreira & Joyce\] potential displays a generic form for moduli fields from extradimensional theories flat directions while the \[Steinhardt et al. 1999, hereafter Steinhardt et al.\] potential has been proposed as an infinite sum of Ratra-Peebles-type potentials. This set of potential has been chosen as a well motivated starting point. It also spans the main types of potentials, power laws, gaussian and exponentials and display very different behaviours.
Using those potentials, spelled out in table 1, we can explore the impact of different models on the timescale of the quintessence dominated epoch and the apparent strength of matter domination between radiation and quintessence eras.
We also can compare the observable effects of various potentials at the homogeneous level through the equation of state evolution. The equation of state evolution should be attained by SNIa measurements which constrains directly the integrated luminosity distance evolution, although it relies on SNIa to be good standard candles without systematic errors. In this paper, we are even more interested in the density parameters: since, on the rough, dark energy with its negative pressure has a freezing effect on clustering of dark matter, the differences between potentials in the equivalence epoch for matter-quintessence and in the strength of the matter dominance phase are signs respectively of differing inhibiting times for structure formation and of overall matter clustering activity.
We thus emphasize those features in the otherwise known homogeneous evolution with Ratra-Peebles potentials for the slope $`\alpha _Q=6`$ and $`11,`$with the SUGRA potentials for the slope $`\alpha _Q=6`$ and $`11,`$ Ferreira & Joyce potential for $`\lambda =10`$ <sup>1</sup><sup>1</sup>1This choice is historical and does not allow the marginal tracking behaviour of FJ to be reached. It allows nevertheless to explore a very different behaviour of the equation of state that yields crucial insights (see section 4.3).and the \[Steinhardt et al. 1999\] potential (the choices for parameter values follow the authors). Comparison of the density parameters evolutions allows to conclude on the fact that we expect a stronger inhibition of structure formation with the Ratra-Peebles potential than with the \[Steinhardt et al. 1999\], than with the SUGRA, than with the Ferreira & Joyce potentials. Within models (i.e. for Ratra-Peebles and SUGRA), variations of their respective evolution points towards the degree of freedom inherent to each potential and thus towards the falsifiability of structure formation tests on each model. Another remark drawn from the density parameters concerns the scales of the structures affected: in the hierarchical CDM scenario, clusters and superclusters scales being formed last, we expect them to be most affected by the changes in inhibition epochs from the various potentials because those epochs occur during the most recent periods and inhibition is expected to act most on them. Eventually it should be stressed that the observational constraints being applied nowadays, differences between models are expected to increase as we look back in time. It should be noted that we do not expect subdominant quintessence to alter the matter radiation equivalence, and in relation, the recombination, thus the power spectrum to remain essentially unchanged for the purpose of structure formation.
The lower panels of figure 1 show some changes for the Ratra-Peebles model that are much less pronounced in the case of the SUGRA model, and also illustrate, together with the FJ/Steinhardt *et al.* panel, the variations in equivalence epochs and strength of matter domination. The upper panels all illustrate the fact that the equation of state cannot be considered constant during the epoch of structure formation (contrary to \[Lokas & Hoffman 2001\]).
The solving of Eqs.(2.1, 2.2) was effected with a second order Runge-Kutta integration method. As previously mentioned, giving initial conditions for each potential shown in table 1 is not constrained with input observations. The search for tracking solutions excludes reverse integration methods: the selection of the tracker solution which leads to observed quintessence density was obtained with a simple iteration on the characteristic energy scale $`\mathrm{\Lambda }_Q`$ with respect to the target final quintessential density.
The cosmic time numerical increment is chosen so as to keep a constant logarithmic scale factor increment: $`dt=\alpha \frac{a}{\dot{a}},`$ where $`\alpha `$ is fixed (we usually take $`\alpha =10^2`$, a value of $`10^3`$changing the outcome by less than 1%).
Given the roadmap that we now have on variations between models in the influence of quintessence on large scale structure formation. We are going to effectively test them in the mass functions of cosmological haloes within quintessential context as computed by the PS scheme \[Press & Schechter 1974\]. First, in order to get the mass function, we need to study the spherical collapse model. This is done in the next section.
## 3 Spherical collapse in the presence of quintessence
The influence of quintessence models on structure formation comes from the repulsive gravitational effect of its negative pressure on the bulk of spacetime during its phase of less than -1/3 equation of state. Collapsing structures are then slowed down in their build up by this relative aggravated expansion. An elementary model of this phenomenon can be used to get simple quantitative results that can then be inputted in a PS-type scheme that will be discussed in section 4: the spherical collapse model (pioneered in \[Larson 1969, Penston 1969, Gunn & Gott 1972, Fillmore & Goldreich 1984, Bertschinger 1985\] and summarised in \[Peebles 1980\]) in the presence of quintessence holds the key to this exploration from the beginning of the collapse phase after recombination down to shell crossing. After shell crossing, mass conservation does not allow to follow the system with just its outer shell but a prescription can be used for the virialization of the model. We will now describe the dynamics of the cosmological spherical collapse.
### 3.1 The dynamical system
The cosmological spherical collapse is embedded into an FLRW-type universe for which all the component of energy density are supplemented with a spherical overdense region. Birkhoff’s theorem holds the key to the spherical non linear collapse model: any spherical region embedded in a spherically symmetric universe behaves shell by shell as a patch of FLRW universe with each characteristics modified to match the corresponding average inner ones. In this case, any given shell behaves according to the average of its inner density. With a flat background this leads to positive curvature inside the overdensity. However, because this inner curvature may change with time, one has to be cautious in using the Friedmann’s equations \[Mota & van de Bruck 2004\].
For simplicity reasons and because it is a building block of the original PS scheme, we will use the spherical top hat model (sphere of constant overdensity). In this case, the averages are equal to the local values and every shell reach the center of the overdensity at the same time, when shell crossing occurs. Problematic for the virialization, this feature allows one to follow only the evolution of the outermost shell representing the whole system. Virialization is then assumed to take place soon after shell crossing and a halo of mass given by the extent of the initial overdensity is formed.
Thus the radius $`r_{od}=r`$ of the overdensity can be written as a function of its initial value $`r_{i_{od}}=r_i`$ as $`r_{od}(t)=r_{i_{od}}a_{k>0}(t)=r_is(t),`$ where we note the scale factor of the positive curvature patch $`a_k>0`$ as a fiducial or rescaled radius of the overdensity region $`s=r/r_i.`$ Thus the radius of the overdensity follows Friedman’s equation for the modified overdensity patch of universe. For this shifted FLRW model, the Einstein’s Equations yield the evolution (recall $`\kappa =8\pi G`$, the gravitational coupling):
$`\dot{s}^2`$ $`=`$ $`{\displaystyle \frac{\kappa s^2}{3}}{\displaystyle \rho }k(t),`$ (3.1)
$`\ddot{s}`$ $`=`$ $`s\left[{\displaystyle \frac{\kappa }{6}}\left({\displaystyle \left(3P+\rho \right)}\right)\right]`$ (3.2)
where the sums on energy density and pressure concern each cosmic species inside the patch radius. Again note that the curvature $`k`$ inside the overdensity is not constant in general, so that we prefer to use Eq.(3.2). Since we assumed homogeneous behaviour for the scalar field (i.e. no clustering), that means that there is no conservation inside the spherical patch for the quintessence field \[Mota & van de Bruck 2004\] and its pressure is that for the background universe.
The acceleration follows the shifted second Friedman equation (Eq. 3.2), which contains as well the pressure terms. The density term for the matter follows the same pattern as for Eq. (3.1):
$$\frac{\kappa r}{6}\rho =\frac{4\pi Gr}{3}\rho =\frac{GM}{r^2},$$
with the total mass M conserved and contained initially inside the spherical patch. The pressure terms depend on the respective state equations. Hence the acceleration equation reads, with our time and radius units,
$`\ddot{s}`$ $`=`$ $`\left[\lambda _0\left(\mathrm{\Omega }_{r_0}a^4+{\displaystyle \frac{4\pi G}{3_0^2}}(3P_Q+\rho _Q)\right)\right]s{\displaystyle \frac{\mathrm{\Omega }_{m_0}a_i^3\left(1+\mathrm{\Delta }_i\right)}{2s^2}}.`$ (3.3)
Initial conditions of the homogeneous evolution are taken after inflation ($`a=10^{30}`$), the field is taken in a reasonable range allowing for the tracking solution to establish. For the collapse evolution, initial time is chosen in the relevant overdensity linear regime (we usually take the arbitrary cut $`a_i=10^5`$), the overdensity is set in a Hubble flow, that is following the general expansion of the universe at the initial onset of the overdensity, (the definition of $`s`$ sets its initial condition to be $`s_i=1`$ at initial time) so
$`\dot{s}_i`$ $`=`$ $`{\displaystyle \frac{\dot{a}_i}{a_i}}.`$ (3.4)
In the model of the spherical collapse, the non-linear density at the border of the overdense region can be monitored using the Lagrangian mass conservation. For the calculation of the overdensity, one just follows the canonical definition $`\delta =\frac{\rho _m\rho _b}{\rho _b}`$, then
$$\delta =\left(1+\mathrm{\Delta }_i\right)\left(\frac{a}{a_is}\right)^31.$$
(3.5)
The corresponding initial overdensity value thus proceeds from Eq.(3.5) at $`a=a_i`$ and $`s=1`$ and using Eq.(3.4):
$`\delta _i`$ $`=`$ $`\mathrm{\Delta }_i,`$ (3.6)
$`\dot{\delta _i}`$ $`=`$ $`3(\delta _i+1)\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{s}}{s}}\right)_i=0.`$ (3.7)
Here we have used a fiducial initial time from which the collapse was evolved. This choice is of course arbitrary and in the PS scheme, this arbitrariness is resolved by the extrapolating to present time (redshift $`z=0`$) the initial overdensity $`\delta _i=\mathrm{\Delta }_i`$ using the linear evolution theory
$$\ddot{\delta _L}+2\frac{\dot{a}}{a}\dot{\delta _L}=\frac{3}{2}H_0^2\mathrm{\Omega }_{m_0}a_i^3\delta _L.$$
(3.8)
This theory is only valid for overdensities in the linear regime ($`\delta 1`$) but the non-linear collapse leads to infinite density at a finite time given by $`a=a_c`$ while the linear theory allows the density to remain finite up to $`z=0`$.
The immediate interest of solving this model lies in the possible comparison between different collapse times. The relevant quantity extracted was the function $`\delta _{c_0}=f(a_c),`$ the linearly extrapolated overdensity as a function of collapse epoch expressed in terms of the scale factor. Given a potential, we construct its characteristic extrapolated density contrast function of collapse scale: each initial overdensity yields both the collapse scale and extrapolated linear scale using the non linear and linear collapse respectively (see figure 2). The construction of the function involved scanning a whole range of initial overdensities as illustrated in figure 2. The critical density contrast function of collapse scale factor $`a_c`$ is constructed by the linear extrapolation to the present from the arbitrarily chosen initial epoch’s starting overdensity. Therefore $`\delta _{c_0}`$ yields an object that collapses at an epoch given by the value of scale factor $`a_c`$. Thus $`\delta _{c_0}`$ is found by evolving our initial conditions with Eq.(3.8)’s linear theory, and the corresponding collapse scale is found by the non-linear collapse of Eq.(3.3) – given by the vertical asymptote at $`a_c`$, in the homogeneous background provided by Eqs.(2.1, 2.2).
For implementing this scheme, we first determine the correct QCDM potential energy scale, as described in section 2.2, then we bring the model to the field’s tracking regime and produce initial parameters for the spherical collapse. Eventually the coupled evolution of the background FLRW, the quintessence field, the non-linear spherical collapse and the linear overdensity evolution obtains the values of $`\delta _{c_0}(a_c).`$
A fourth order Runge-Kutta integration method was implemented over the whole Eqs.(2.1, 2.2, 3.3, 3.8) system. We still used the logarithmic increment $`dt=\alpha \frac{s}{\dot{s}},`$ where we usually take $`\alpha =10^2`$, as described in section 2.2, but we had to limit the lower increment value of $`dt`$ to avoid inflation of numerical expenses when integration approaches the overdensity turnaround point.
We are now ready to use the Top Hat Spherical collapse model to decide when objects are considered to have collapsed, that is when their non-linear overdensity diverges. The following section discuss those first results.
### 3.2 Critical densities
Once we have computed a series of $`\delta _{c_0}(a_c),`$ for a range of initial conditions, set in the linear regime, at our arbitrary starting epoch, we can display for each model their extrapolated critical overdensity as a function of collapse scale factor and compare them among models, which is done in figures 3 (comparing pseudo-quintessence models) and 4 (for dynamical quintessence).
Since we are dealing with density contrasts over the background cosmological matter density, the main evolution effects are expected to come from the lower homogeneous matter density yielded by the quintessence models during their matter dominated era. Also the observation setting the models corresponding to z=0 ($`a=\frac{1}{1+z}=1`$; present epoch), we expect the differentiation between models to increase backwards in time.
A comparison between our various potentials, the $`\mathrm{\Lambda }\mathrm{CDM}`$ model<sup>2</sup><sup>2</sup>2It should be noted that for a standard CDM ($`\mathrm{\Omega }_m=1`$) model, the spherical collapse yields a linear evolution $`\delta =\delta _ia/a_i`$\[Peebles 1980\] so, extrapolating the collapse value nowadays, $`\delta _{c_0}=\delta _c1/a_c`$ which gives the straight line asymptote to each models in these log-log coordinates and the well known $`\delta _c=1.686`$. and three pseudo-quintessential models (i.e. with constant equation of state) have been performed. The overdensity as a function of collapse shows clear distinctions between models are possible. The results obtained are plotted on figures 3, for the pseudo quintessences and two common models (Ratra-Peebles and SUGRA), and 4, for all our studied models.
The pseudo-quintessence models are shown for reference and comparison to other works even though our homogeneous study has pointed that such an assumption is highly unlikely to hold compared to real quintessence models (e.g. figure 1’s upper panels). However figure 3 reveals that even though the region scanned by the pseudo-quintessence models is covering most of the values taken by our sample of models (except for Ratra-Peebles), there is still a strong distinction if measurements are confronted at different epoch since no pseudo-quintessence model is matching quintessence evolution; for instance, the SUGRA11 (notation defined in the figure caption) model cannot be fitted with only one wQCDM model.
The main point is that there are distinctions in this type of representations between the various models when examined at earlier epochs and between successive era. Whilst evaluation of variability for the SUGRA model yields very little separation, that of the Ratra-Peebles are so large that it allows for an observational selection of the free power index. Whereas the FJ model crosses the curves from various values of the free power index in the Ratra-Peebles model, the \[Steinhardt et al. 1999\] model is not very distinguishable from the SUGRA6. However we will see that the sensitivity of the mass function to the overdensity allows for some distinction.
To summarise, it is the departure, in quintessential spherical collapse, from the constant value of the collapse critical overdensity of the standard CDM collapse that will affect the mass function. On figures 3 and 4, this is translated into the departure from the simple linear log-log relation of the critical overdensity with the scale factor. Thus the linear density contrast indexed by the non-linear scale of collapse sets the redshift evolution of the PS mass function observable as we will see in the next section.
## 4 The mass function in presence of a quintessence field
In this section we will lay down our assumptions for constructing the mass dispersion of our models which imprints initial conditions in the PS scheme, we will then recall it before discussing the integrated mass functions for the quintessential models studied.
### 4.1 Mass dispersion
Initial conditions of large scale structure formation are condensed in the PS scheme into their mass dispersion. To do this for gaussian fluctuations the basic information lies in the density field power spectrum.
In general, the Fourier power spectrum P(k) for the density fluctuations of scale $`\frac{1}{k}`$ can be decomposed into a primordial part and a linear evolution part, written as follows (\[Cole & Lucchin 1995\]pp 267,282, \[Peebles 1980\]p169, see discussion in section 4.1.2):
$$P(k,a)=Bk^nT_Q(k,a)^2=Bk^nT_\mathrm{\Lambda }(k,a)^2T_{Q/\mathrm{\Lambda }}(k,a)^2$$
(4.1)
where B is a normalization factor and $`T_Q`$ the quintessence transfer function written as $`T_Q=T_\mathrm{\Lambda }`$. $`T_{Q/\mathrm{\Lambda }}`$, $`T_\mathrm{\Lambda }`$ being the $`\mathrm{\Lambda }`$ dark energy transfer function. The transfer function is itself written as: T$`(k,a)=`$T$`(k,a=a_i).D(a)/D(a_i)`$ where the linear growth D gives the time evolution.
In the following section, we briefly decompose the construction of the structure formation initial spectrum between a primordial part and a linear evolution, we discuss our assumed spectrum and show the mass dispersion and its normalisation to observations.
#### 4.1.1 primordial power spectrum
In the context of inflation, for commonly used potentials, one gets a gaussian primordial spectrum P(k,$`a_{init.}`$)=$`\left|\delta (k,a_{init.})\right|^2`$=A$`\lambda _{infl.}k`$ (where $`a_{init.}`$ is the epoch at which a perturbation of size l=1/k is reentering inside the horizon): i.e the scale invariant Harrison-Zeldovich spectrum, with $`\lambda _{infl.}`$ and A constants depending the inflationary potential and expansion period.
#### 4.1.2 Linear power spectrum
To describe the growth of the fluctuations and the evolution of the initial power spectrum P(k, $`a_{init.}`$) one usually evolves linearly fluctuations up to a given epoch $`a`$ (typically after recombination when $`\frac{1}{a}1=`$z$`<`$10<sup>3</sup>) as P(k,t)=P(k, $`a_{init.}`$).T$`{}_{}{}^{2}(k,a),`$ where T is the above transfer function.
Previous studies have shown that P(k) depends only weakly on the equation of state $`w_Q(a)=P(a)/\rho (a)`$ \[Lokas & Hoffman 2001\] since the Q field is sub-dominant at epochs where z is large (e.g., at recombination, z$``$ 10<sup>3</sup> ). In the case of a universe dominated by cold dark matter (CDM) the initial T(k) function is known to be well fitted \[Bardeen et al. 1986, Sugiyama 1995\]. The $`\mathrm{\Lambda }`$CDM (case $`w_Q`$=-1) transfer function T<sub>Λ</sub> fit is very close to the CDM function with a slightly smaller cutoff scale \[Efstathiou et al. 1992\] and identical power laws. Flat cosmologies transfer functions fits from N-body simulations of pseudoQuintessence models (!real bf!models with constant equation of state $`w_Q`$) have been proposed in the form $`T_Q=T_\mathrm{\Lambda }`$. $`T_{Q/\mathrm{\Lambda }}`$ \[Ma et al. 1999\]. They show only modifications to $`T_\mathrm{\Lambda }`$ of order unity in the observationally relevant range of $`w_Q`$.
Following previous work, we therefore apply the same initial density fluctuation power spectrum to all of the models. This ensures that the differences that we see result from the dynamical evolution of the density fluctuations rather than differences in the initial power spectrum. In practice, CMBR experiments will establish the fluctuation power spectrum independently of the mass functions that we consider here and a successful model will have to successfully reproduce both datasets. We thus consider that the T<sub>Λ</sub> function given for a model with a cosmological constant can describe suitably the fluctuations and do not allow the Q field to influence the initial power spectrum of density fluctuations. We therefore adopt the \[Bond & Efstathiou 1984\] form for the power spectrum, following our use of the \[Jenkins et al 2001\]<sup>3</sup><sup>3</sup>3note their code was built to encompass $`\mathrm{\Lambda }`$CDM dark energy only mass dispertion and mass function code, modified with the present work’s quintessential non-linear collapse program.
#### 4.1.3 mass dispersion
The mass dispersion $`\sigma ^2(M,a)`$ is computed using a top-hat filter in real space with a radius of filtering R, corresponding to $`M=4\pi \rho _b(a)R^3/3,`$ $`\rho _b(a)`$ being the density of mass of the Universe at a given epoch marked with the scale factor $`a`$ (thus we can express $`R\left(M\right)=\left(3M/4\pi \rho _b(a)\right)^{1/3}`$.
One has, for the filtered mass dispersion, the relation:
$$\sigma _Q^2(R,a)=D^2(a)d^3k\left|W(R(a),k)\right|^2P(k)/8\pi ^3$$
(4.2)
with the filter Fourier transform W (which is here the top hat filter), indices in $`Q`$ denotes the model dependence. Thus we get the relation:
$$\sigma _Q^2(R,a)=\sigma _Q^2(f(a)M^{1/3},a)=D^2(a)𝑑k\left|W(R,k)\right|^2P_{BondEfsthathiou}(k)/8\pi ^3$$
(4.3)
with $`f(a)=(3/4\pi \rho _b(a))^{1/3}`$ and D, the linear growing mode for fluctuations.
#### 4.1.4 spectrum normalization
To get the value of the constant in P(k,a) one usually normalizes $`\sigma `$ to the value $`\sigma _{8h^1}`$ which is observed today for a typical radius R= 8h<sup>-1</sup>Mpc (here we take h=0.65) and to reproduce the present amplitude of the density contrast one gets that $`\sigma _{8h^1}`$ is of order 1 (see for example reference \[Eke et al. 1996\] that yields the product $`\sigma _{8h^{1.}}`$ $`\mathrm{\Omega }_0^{1/2}=0.6,`$ or $`\sigma _{8h^1}=`$1.1 for $`\mathrm{\Omega }_0=0.3`$).
### 4.2 The mass function evaluation
To compute the mass function in the non-linear regime without using fits to N-body simulations we restrict here to the popular PS prescription.
This semi-analytical tool turns out to be the simplest way to construct a mass function from the assumed gaussian statistics of the initial density. It neglects a large number of effects that tend to compensate each other yielding often accurate fits to the numerical simulations.
It is yet an empirical recipe yielding the number of (collapsed) objects with a mass greater than a given one \[Press & Schechter 1974\].
Hence, we obtain the fraction of collapsed mass linearly extrapolated at a selected present time t<sub>0</sub> defined by $`\delta _{c_0}(t_c)`$ as in section (3.8), and account for Eq.(4.3) for the variance $`\sigma `$ itself linearly extrapolated (and normalised with $`\sigma _8`$) at the same time t$`_0.`$ The mass fraction thus writes:
$$F(m>M,a)=1Erf(\delta _{c_0}(a)/2\sigma (M,a)),$$
(4.4)
and its derivative with respect to the mass leads to the density of collapsed objects
$`n_{PS}(m,a)`$ $`=`$ $`\rho _b(a)\left(dF(m>M,a)/dM\right)/M`$ (4.6)
$`=`$ $`(2/\pi )^{1/2}\left(\rho _b(a)\delta _{c_0}(a)\right)\left(d\sigma (M,a)/dM\right)/\left(M\sigma ^2(M,a)\right)`$
$`.\mathrm{exp}(\delta _{c_0}^2/(2\sigma ^2(M,a)))`$
$`=`$ $`\left(\rho _b/\sigma M\right)\left(d\sigma (M,a)/dM\right)F_{PS}(M)`$ (4.7)
with the characteristic for the PS scheme condensed into the function
$$F_{PS}(M)=(2/\pi )^{1/2}\left(\delta _{c_0}(a)/\sigma (M,a)\right)\mathrm{exp}\left(\delta _{c_0}^2(a)/(2\sigma ^2(M,a))\right)$$
(4.8)
that can be replaced in later studies with more complex schemes.
We now have set in place the machinery for computing mass functions for non-linear structures in the presence of non-clustering dynamical quintessence. We will then present comparisons of integrated mass functions that can be obtained for the potentials we selected in the following section.
### 4.3 Integrated mass functions
One of the main interests of this study is to have shown the sensitivity of $`\delta _{c_0}`$(a) as function of the chosen Q field potential V(Q). This can be seen above (section 3.2) with figures 3 and 4 when dealing with the collapse of an initial overdensity. Placing these results within the PS scheme, we have found that some differences among the various models can emerge from the mass functions. In figure 5, the upper panel shows as a reference the mass functions for all the selected models at present time (z=0). The models there are all very close to each others, which reflects our z=0 normalisation and confirms the usual hypothesis taken of a normalization of the mass functions to present days observations. This is why figure 5’s lower panel only recalls the LCDM (i.e. $`\mathrm{\Lambda }\mathrm{CDM}`$) mass function as a reference.
When we examine the mass function at larger z, as shown on figure 5’s upper panel for z=0.5 and lower panel for z=1, discrimination can occur between the various models. All models’ mass functions are lying above the LCDM (i.e. $`\mathrm{\Lambda }\mathrm{CDM}`$) curve for the same epoch. The distinctions remain nevertheless mild at z=0.5, as shown by the relative proximity between the curves; they, however, become more pronounced at z=1, with the following observations: the FJ model displays the most structure suppression, that is the least evolution of mass function, i.e. the least structure formation during recent times. Then, in the same decreasing hierarchy, one encounters the Ratra-Peebles models, in decreasing order of their power law, the Steinhardt *et al.* model, the SUGRA models, also in decreasing order of their power law, before reaching the $`\mathrm{\Lambda }\mathrm{CDM}`$ model (the Steinhardt *et al.* and the SUGRA11 models’ mass functions are almost on top of each other). Surprisingly, those results are not completely agreeing with the naive interpretation given in section 2.2, that would have yielded the same hierarchy of structure inhibition and matter dominance except for the FJ model, that would be expected to have a lower mass function than all other quintessence models (see figure 1’s lower panels). Confronted with their lookback effects on the number of structures present at a certain epoch, these are both agreeing with the mass function hierarchy except for the ordering of the FJ model. We are thus compelled to consider the only qualitative difference that springs to the eye between the FJ model and the rest of them: the evolution of the equation of state (figure 1’s upper panels). The FJ model appears as the only one with $`\omega _Q`$ evolving towards higher values at its latest stages, during structure formation. It seems thus that the effect of its negative pressure in the acceleration equation of the spherical collapse (Eq. 3.3) is strong enough to counter the strength of matter domination earlier than expected in the natural first approach. This can be seen from the ratio of quintessence- and matter-induced accelerations of the overdensity, from Eq.(3.3), so we get $`(3\left|\omega _Q\right|1)\frac{\mathrm{\Omega }_Q}{\mathrm{\Omega }_m}>\frac{\left(1+\mathrm{\Delta }_i\right)}{2(a_i.s/a)^3}`$ despite having $`\frac{\mathrm{\Omega }_Q}{\mathrm{\Omega }_m}`$ small at a given stage of collapse and cosmic evolution. These results cast very interesting light into the intimate mechanisms of quintessence that would escape to non semi-analytic approaches.
The evaluation of model’s variability is expected after section 3.2’s more or less pronounced differences among each class of model and their respective parameter variations (Ratra-Peebles and SUGRA). In the light of integrated mass functions, each model evaluated displays little variability so discrimination among power law values may prove difficult with this method.
## 5 Conclusions
In this paper we have extended the evidence that several different recent models of quintessence predict significant difference in the evolution of the halo mass function, in general agreement with previously more restricted works \[Mainini et al. 2003a, Mainini et al. 2003b, Nunes & Mota 2004, Solevi et al. 2004\], and used our wider range of potentials to emphasize the impact of the equation of state in the formation of dark haloes. Although the best evidence seems to come from the largest mass structures, as expected from the homogeneous exploration (section 1), the models studied do not behave simply according to matter and quintessence dominance epochs. Nevertheless, the sensitivity of mass functions to the value of $`\delta _{c_0}`$ at the limits of accuracy cautions us against strong assertions.
The models can be distinguished if observational measurements of the cluster mass function at $`z1`$ achieve a precision of better than $``$10% at $`10^{14}h^1M_{}`$. Because the models are normalised to give the same mass function at $`z=0`$, a higher level of precision, approaching 1%, is required to distinguish between the models at $`z=0.5`$. Several experiments are planned that will exceed this level of accuracy and will have good control of the systematic errors.
One of the most promising is the South Pole Telescope (SPT) survey \- a cosmic microwave background experiment aimed at detecting the Sunyaev-Zeldovich effect (SZE) due to clusters of galaxies \[Ruhl et al. 2004\]. The effect arises because inverse Compton scattering by the hot intra-cluster plasma distorts the spectrum of the cosmic microwave background. Importantly, the amplitude of the decrement is independent of the cluster redshift, and the integrated signal decreases only with the clusters angular diameter. The aim of the South Pole Telescope experiment is to map a region of 4000 deg.<sup>2</sup>, detecting essentially all clusters with masses greater than $`2\times 10^{14}h^1M_{}`$ and with redshift less than 2 \[Carlstrom et al. 2002\]. The survey is expected to detect around 30,000 clusters in the $`\mathrm{\Lambda }`$CDM cosmology, with 30% expected to have redshifts greater than 0.8. The amplitude of the SZ decrement is determined by the integrated pressure of the intra-cluster plasma, and can thus be used to determine the cluster mass. Random errors in the normalisation of the mass function will be $`1\%`$ at both $`z1`$ and $`z0.5`$. Systematic errors in the determination of mass and the scatter between cluster mass and the observed SZ decrement are a significant source of concern, but these can be controlled using the “self calibration” techniques discussed in \[Hu 2003\] and \[Majumdar & Mohr 2003\]. An other technique based on the optical detection of clusters may also yield promising results when combined with similar self-calibration methods \[Gladder & Yee 2005\].
These observational programmes will map the development of the mass function from $`z=1`$ to $`z=0`$ at the level precision that is required to distinguish between the quintessence models considered in this paper. The techniques we have developed allow this map to be directly related to the form and parameterisation of the quintessence potential.
Although this work focuses on the actual mass functions of studied models, that is on the real volume number densities predicted, some authors have argued that detection of such density should be hindered by the global geometric effect of dark energy \[Solevi et al. 2004, Solevi et al. 2005\]. However, their treatment of observed cluster samples with a unique mass cut-off forgets about the bias-geometry dependence, discussed by \[Kaiser 1984\] and retained as mass-observable relation in dark energy studies \[Mohr 2004\], that should certainly call for further scrutiny. Such a study is in progress that should test whether a break of degeneracy is induced by this bias-geometry dependence and that will be at the core of a companion paper.
More potentials have attracted attention \[Barreiro et al. 2000, Albrecht & Skordis 2000\] and although they combine our studied exponential and power laws, they should be confronted using our method. Confirmation could be sought using parallel semi-analytical methods like the \[Jenkins et al 2001\] one in the context of quintessence \[Linder & Jenkins 2003\]. Eventually, the homogeneity and minimal coupling of the field lacks proof in the highly non-linear regime \[Mota & van de Bruck 2004, Nunes & Mota 2004\] and calls for some wider explorations in the line of \[Mota & van de Bruck 2004, Nunes & Mota 2004, Macciò et al. 2004\]. This will be the subject of a follow up paper.
MLeD wishes to thank J.-M. Alimi and the LUTH for their hospitality and in the initiation of this work, H.Courtois and the CRAL, A. Falvard, E. Giraud and the GAM-LPTA and J.P. Mimoso and the CFTC, for their hospitality and support, Tom Theuns for fruitfull discussions and advice, A. Jenkins for letting me modify his mass function code with my quintessence non-linear collapse program, and T. Lehner and D. Steer and also R. Bower for their encouragements, efforts and many discussions.
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# Magnetoresistance Anisotropy of Polycrystalline Cobalt Films: Geometrical-Size- and Domain-Effects
## I Introduction
In applied magnetism, the coupling of the magnetic moment to spatial degrees of freedom plays a key role, and this especially applies to modern magneto- or spin-electronicsThompson et al. (1975). Basically, this coupling is provided by the spin-orbit interaction, which in the example of magnetotransport causes the scattering rate of the conduction electrons, $`\tau ^1`$, to depend on the direction of the local magnetization $`\stackrel{}{M}`$ with respect to the current. In the archetypal bulk ferromagnets iron, cobalt, nickel, and their alloys, the resistance difference for orientations of $`\stackrel{}{M}`$ parallel and perpendicular to the current, i.e. the socalled anisotropic magnetoresistance ratio (AMR), $`\mathrm{\Delta }\rho /\rho \left(\rho _{}\rho _{}\right)/\rho `$, amounts to some percent. In nanostructured devices like magnetic multilayers, wires, or constrictions in the ballistic regimeVelev et al. (2005), this ratio may be enhanced to several ten percent.
Basically, the determination of the scattering rate $`\tau ^1(\stackrel{}{M})`$ and of the AMR requires the knowledge of the scattering potential and also of the spin-orbit split bands at the Fermi-surface $`_F`$. Some special aspects of the AMR have been evaluated by SmitSmit (1951), BergerBerger (1964), PotterPotter (1974), and Fert and CampbellFert and Campbell (1976), however, the evaluation of $`\tau ^1(\stackrel{}{M})`$ for a realistic case is still lacking, at least to the best of our knowledge. In this context, we note a recent ab initio calculation of the intrinsic anomalous Hall-effect which, in contrast to the AMR, depends only to first order on the spin-orbit interaction and *not* on a scattering potential. This quantity was obtained by integrating the k-space Berry-phase over the occupied spin-orbit split states of ironYao et al. (2004) and was found to agree up to some 30 percent with data on iron whiskersDheer (1967).
The present work is intended to a fairly systematic study of the AMR, which is of second-order in the spin-orbit interaction, also in an elemental 3d-ferromagnet. By selecting hcp cobalt with a rather well known band-structurePapconstantopoulos (1986), some deeper insight into the AMR may be facilitated. By choosing polycrystalline films, we are closer to devices which invariably use polycrystalline materials. We will vary the structural disorder and the temperature in the films to probe the role of different scattering mechnisms. These basic properties of the films under study are examined in Sect.II. Section III is devoted to the AMR in the technically saturated state with main emphasis to a still unexplained phenomenon of the AMR, i.e. the socalled geometrical-size effect (GSE), previously observed in thin NiChen and Marsocci (1972) and PermalloyRijks et al. (1997) films. Another point of interest will be the absolute value of $`\mathrm{\Delta }\rho `$ at low temperatures: for Ni-alloys, already McGuire and PotterMcGuire and Potter (1975) pointed out the unsensitivity of $`\mathrm{\Delta }\rho `$ against significant variations of the residual resistivity $`\rho (0)`$. The influence of different domain structures, depending on the film thicknesses, on the magnetoresistance, is investigated in Section IV and will be discussed by using the results on the GSE. This low-field regime, where the in-plane AMR switches at rather small coercive fields $`(H_c10\text{Oe})`$, may be of interest for applications despite the fact that $`\mathrm{\Delta }\rho /\rho `$ lies in the range of some percent. The summary and conclusions are contained in Section V.
## II Characterization of the films
By means of DC-magnetron sputtering at an Ar-pressure of 2$``$10<sup>-9</sup> bar, cobalt films of thicknesses 10 nm, 20 nm, and 188 nm were deposited on Synsil-glass and oxidized Si(100) surfaces and capped by 3 nm thick Al-layers. The thicknesses were measured by a profilometer to an accuracy of 0.6 nm and controlled by high-resolution SQUID-magnetometry. X-ray diffraction diagrams (XRD), as shown in Fig.1, revealed a polycrystalline hcp-structure with a slight texture of the hexagonal axis normal to the plane. Surface images recorded by an atomic force microscope (AFM, Q-Scope$`^{\text{TM}}`$250, Quesant Instruments Co.) yielded surface roughnesses between $`(1.5\pm 0.3)`$ nm for 10 nm and $`(3.8\pm 0.5)`$ nm for 188 nm and indicated the grain sizes to increase from $`(25\pm 5)`$ nm to $`(80\pm 5)`$ nm. Within the error margins, these results turned out to be the same for both substrates. It is interesting to note that the grain sizes and their increase with thickness are consistent with a recent report for polycrystalline Co on glass and Si(100) substrates Kharmouche et al. (2004).
The magnetic properties of all films have been investigated by ferromagnetic resonance (FMR), hysteresis loops, and magnetic force microscopy (MFM). Using a home-made FMR spectrometer operating at 9.1 GHz, the directions and magnitudes of small uniaxial anisotropy fields, $`\stackrel{}{H}_u`$, in the film planes were determined . On 20 nm Co:Si and 188 nm Co:glass, for example, $`H_u=22.3`$ Oe and 15.3 Oe, respectively, was obtained and the orientation of $`\stackrel{}{H}_u`$ could be related to the direction of the deposition process. Magnetization isotherms were measured by a SQUID-magnetometer (Quantum Design MPMS2) along three orthogonal directions of the applied field $`\stackrel{}{H}`$ at temperatures, which were of interest for the analyses of the magnetoresistances (MR’s) in Section IV. There also MFM images are presented in order to visualize the domain structure underlying the magnetization processes, see Fig.7 below. For this purpose, the Q-scope was equipped with a commercial tip, coated by a 40 nm thin hard Co-alloy (Nanosensors$`^{\text{TM}}`$), and magnetized perpendicularly to the scanning directions. The directions of the in-plane magnetization were determined by monitoring the domain wall motion induced by a small magnetic field produced by external Helmholtz coils.
The resistances have been measured by an array of four in-line contacts prepared parallel to $`\stackrel{}{H}_u`$ by ultrasonic bonding. The driving currents were kept small enough to produce linear responses and the resulting $`U/I`$-ratios were corrected for the sample geometry Smit (1958) to determine the resistivities of the films. The sample chip was mounted to the end of an cold-finger extending from the cold-plate of a pulse-tube cooler (PRK, Transmit Co. Giessen, Germany) to the center of a warm-bore superconducting magnet (130 kOe, Oxford Instruments). A PID controller and a heater allowed stable sample temperatures between 70 K and 350 K. Measurements of the magnetoresistance in the domain states, i.e. at low magnetic fields, were performed by means of an electromagnet, by which also the angle between current and field could be varied. More experimental details are given in Ref.Gil, 2004. We should mention here, that the structural, magnetic, and transport properties proved to be largely independent on the substrate, i.e. synsil-glass or oxidized Si(100)Gil (2004). This feature indicates a dominant effect of the polycrystallinity of the films, i.e. of the deposition process. For some practical reasons, we selected three films with thicknesses between 10 nm and 188 nm for the present study.
The temperature dependence of the zero-field resistivities is depicted in Fig.2 for these three films. The data can be well described by a sum of three contributions
$$\rho (T)=\rho (0)+\rho _{ph}(T)+\rho _m(T).$$
(1)
According to the inset, the residual resistivity increases linearly with the inverse thickness,
$$\rho (0)=\rho _b(0)[1+d_c/d].$$
The characteristic thickness, $`d_c=(18\pm 1)\text{nm}`$ cannot be related to an extra scattering by the film surfacesSondheimer (1952) or grain boundariesMayadas and Shatzkes (1970) since, the theories predict the 1/d-behavior only for small deviations from the bulk limit, $`\rho (0)\rho _b(0)`$. Hence, the observed increase of $`\rho (0)`$ indicates scattering by an additional, yet unidentified disorder in the thinner films. Using the extrapolated bulk value, $`\rho _b(0)=(11+1)\mu \mathrm{\Omega }\text{cm}`$, the carrier density $`5.810^{22}\text{cm}^3`$ from Hall-data for these filmsGil (2004), and the free electron model for the conduction electrons in CoGurney et al. (1993), we find an upper limit for the mean free path, $`\mathrm{}_e\left(0\right)=\mathrm{}k_F/n_ee^2\rho _b\left(0\right)11\text{nm}`$. Since this length is significantly smaller than the mean grain sizes observed by AFM, it may be associated with point-defect scattering within the otherwise crystalline grains.
Since the electron-magnon scattering in Co, $`\rho _m(T)=1.510^5(\mu \mathrm{\Omega }\text{cm K}^2)T^2`$ Raquet et al. (2002) is small, the temperature variation of the resistivities should be dominated by phonons. Indeed, by fitting $`\rho (T)`$ to the Bloch-Grueneisen form
$$\rho _{ph}(T)=\rho _{ph}\left(\frac{T}{\mathrm{\Theta }_D}\right)^n\underset{0}{\overset{\mathrm{\Theta }_D/2T}{}}\frac{x^n}{\left(\mathrm{sinh}x\right)^2}𝑑x$$
and taking for the Debye temperature $`\mathrm{\Theta }_D`$=445 K, we find an excellent agreement by setting for the exponent n=3, valid for phonon-mediated sd-scatteringWilson (1938); Goodings (1963). The strength of this scattering, $`\rho _{ph}`$, becomes smaller in the thinner, more disordered films, however, due to coupling of the phonons to the complicated structure of the d-states, it is difficult to estimate, $`\rho _{ph}`$, even for single crystalsGoodings (1963).
Finally, it may be interesting to note that the present resistivities of the 188 nm film are almost identical to those obtained by Freitas et al. Freitas et al. (1990) on a 300 nm Co film deposited by magnetron sputtering on glass. This applies to the residual resistivity, $`\rho (0)=14\mu \mathrm{\Omega }\text{cm}`$, as well as to $`\rho (T)`$ at room temperature. Significantly smaller $`\rho (0)`$-values have been detected on diode sputteredFreitas et al. (1990) and epitaxialRüdiger et al. (1999) films of similar thickness.
## III High-Field Magnetoresistance
The MR of all films has been studied for three principal directions of $`\stackrel{}{H}`$, defined by the directions of the current $`(\stackrel{}{I}||\stackrel{}{H}_u)`$ and the film plane, see inset to Fig.3a. To give an example, Fig.3a shows the three MR’s of the 20 nm film at room-temperature. Starting from a common value at low fields, a negative MR is found for all directions of $`\stackrel{}{H}`$. While the longitudinal MR, $`\rho (H_{\mathrm{}})`$ decreases linearly with field, the transverse and polar MR’s contain additional contributions. Above the saturation fields $`H_s`$, where the films become homogeneously magnetized, $`\stackrel{}{M}(H>H_s)=M_s\stackrel{}{H}/H`$, these additional contributions to the MR also saturate at values $`\mathrm{\Delta }\rho _t=\rho _{\mathrm{}}\rho _t`$ and $`\mathrm{\Delta }\rho _p=\rho _{\mathrm{}}\rho _p`$, both indicated in Fig.3a. This contribution results from the spin-orbit induced AMR, since upon rotation of the magnetization either to the transverse or to the polar direction, we realize the angular dependence
$$\rho (\phi )=\rho (0)\mathrm{\Delta }\rho sin^2\phi ,$$
(2)
where $`\phi `$ is the angle between current and the direction of the magnetization $`\stackrel{}{M}`$. Such behavior is characteristic of the AMR of polycrystalline samples of cubic or hexagonal ferromagnetsMcGuire and Potter (1975), and is illustrated by Fig.3b for the in-plane rotation of $`\stackrel{}{M}`$ in a field $`H=0.6\text{kOe}>H_s`$. Details of the MR during saturation by (weak) in-plane fields will be discussed in Section IV. Here we look at the polar MR by increasing $`H_p`$ in Fig.3a. SQUID magnetization dataGil (2004) reveal $`M_p(H_p<M_s)=H_p`$ due to a rotation of $`\stackrel{}{M}`$ against the in-plane demagnetizing field $`H_N=M_s`$, so that the angle between $`\stackrel{}{M}`$ and current $`\stackrel{}{I}(||\stackrel{}{H}_u)`$ is determined by $`sin\phi =M_p/M_s=H_p/M_s`$. For this case, Equ.2 predicts a parabolic decrease, $`\rho (H_p)\rho (0)=\mathrm{\Delta }\rho _p(H_p/M_s)^2`$, which is depicted in Fig.3a by the dashed curve in full agreement with the data.
### III.1 Spin-wave contribution
It is evident from Fig.3a that the linear MR, $`d\rho /dH`$, is the same in all directions of $`\stackrel{}{H}`$. No signature of the classical Lorentz-MR, which is positive and proportional to $`(M+H)^2`$, is realized for $`H>H_s`$, even not at room temperature. Due to the small mean free path the absence of this effect is plausible, while in epitaxial films it becomes visibleRüdiger et al. (1999). The linear MR has been realized before in $`\rho (H_{\mathrm{}})`$ on epitaxially grown iron, cobalt, and nickel films on $`MgO`$ and $`Al_2O_3`$Raquet et al. (2002) and was quantitatively discussed in terms of elastic scattering by thermally excited magnons. Roughly spoken, the negative MR can be ascribed to the suppression of low energy magnons, which results from the increase of the magnon gap proportional to H. The strong thermal increase of $`d\rho /dH`$ is illustrated by Fig.4a for the longitudinal MR to which the AMR does not contribute. In Fig.4b, their temperature dependence is shown for the three films under study and compared to the result for a 7 nm thin Co-film obtained by Raquet et al.Raquet et al. (2002). These authors fitted their data to a simplified model for sd-scattering by magnonsGoodings (1963),
$$\frac{d\rho }{dH}=AT\left(1+2d_1T^2\right)\mathrm{ln}\left(\frac{T}{T_0}\right),$$
(3)
where the amplitude A changed only little from 3 to 4 $`p\mathrm{\Omega }\text{cm/K kOe}`$. Since A depends on the sd-exchange, numerical estimates are rather difficult. The coefficient $`d_1=D_1/D_0`$ is determined by the ratio of the mass renormalization coefficient $`D_1`$ and the zero-temperature stiffness of the spin-waves $`D_0`$. Independent experimental data for Co yield $`d_1=1.5710^6\text{K}^2`$ in good agreement with calculations, and it was arguedRaquet et al. (2002) that $`d_1`$ might be rather insensitive to microstructural details of the films. Consequently, we fitted our data to Equ.3 admitting (plausible) variations in the amplitude A and found a larger value, $`d_1=310^6\text{K}^2`$. We believe that the difference is related to the rather strong disorder in the present films with a residual resistance ratio (RRR) near 2 (see Fig.2), which contrasts to RRR=27 reported by Raquet et al.Raquet et al. (2002) for their thickest films. Hence, one may suspect that the granular disorder in our films gives rise to a stronger thermal renormalization of ’the spin-wave energies’.
### III.2 Anisotropic Magnetoresistance
At low temperatures, where the spin-wave contribution vanishes, the AMR effect should prevail. This is demonstrated in Fig.5 by the MR curves of the 20 nm and 188 nm films measured at T=78 K along the three principal directions of the field. The significant difference between the MR’s of both films at smaller fields is related to the domain structure and will be discussed in Section IV. Here we focus on the saturated transverse and polar AMR’s, $`\mathrm{\Delta }\rho _t`$ and $`\mathrm{\Delta }\rho _p`$, which differ significantly from each other, but do not change very much with thickness (essentially the same observation is made on the 10 nm film). This phenomenon is one of our main results: for all thicknesses, the polar AMR turns out to be about twice as large as the transverse AMR.
At first, a sizable difference between both MR’s, $`\mathrm{\Delta }\rho _p>\mathrm{\Delta }\rho _t`$, has been reported by Chen and MarsocciChen and Marsocci (1972) for single- and poly-crystalline nickel films. They coined this feature as ’geometrical size effect’ (GSE) and believed that it may arise from the electronic structure inside the film material. More recently, this size-effect has also been detected on sputter-deposited 4.5 nm to 100 nm thin Permalloy filmsRijks et al. (1997) at a low temperature, T=5 K. This study revealed that by raising the degree of (111)-texture in the film, $`\mathrm{\Delta }\rho _p`$ was increased so that the ratio $`\mathrm{\Delta }\rho _p/\mathrm{\Delta }\rho _t`$ tended towards two. An attempt to explain this GSE by assuming an anisotropic scattering rate due to diffuse scattering at the film boundaries, however, did not provide conclusive resultsRijks et al. (1997).
In order to explore the AMR and the GSE of our Co-films in some more detail, the absolute values and the thermal behavior of both $`\mathrm{\Delta }\rho ^{}s`$ are summarized in Fig.6. Two remarkable features should be emphasized: (i) despite different temperature variations, the MR’s of all films can be extrapolated to the same value at T=0, as shown in Fig.6a for the polar direction; (ii) Fig.6b demonstrates that the GSE, i.e. the ratio of the polar and transverse AMR’s, remains almost independent of temperature.
At first, we address to the AMR postponing the discussion of the GSE to the following subsection. A thickness-independence of $`\mathrm{\Delta }\rho `$ itself rather than of the ratio $`\mathrm{\Delta }\rho /\rho `$ has been pointed out earlier for $`Ni_{0.7}Co_{0.3}`$ and $`Ni_{0.8}Fe_{0.2}`$ alloys (see Fig.17 of Ref.McGuire and Potter, 1975). For all present Co-films, $`\mathrm{\Delta }\rho _p(0)=0.19\mu \mathrm{\Omega }\text{cm}`$ follows from Fig.6a, and we suspect that the origin of this AMR resides in the crystalline regions, to which we tentatively assigned already the bulk residual resistivity, $`\rho _b(0)=11\mu \mathrm{\Omega }\text{cm}`$, in Sect.II. There we determined the mean-free path, $`\mathrm{}_e=11\text{nm}`$, which turned out to be much smaller than the grain sizes estimated from AFM imagesGil (2004). Therefore, we relate the low-temperature AMR $`\mathrm{\Delta }\rho (T0)`$ also to the scattering within the crystalline grains and believe that the extra scattering, which enhances $`\rho (0)`$ in the thinner films (see inset to Fig.2), produces a negligible AMR. In fact, a weak AMR is expected for scattering potentials with reduced symmetry, e.g. associated with phononsMcGuire and Potter (1975) or correlated structural disorder (grain boundaries, dislocations), because in these cases the directional symmetry-breaking effect by the magnetization $`\stackrel{}{M}`$ via the spin-orbit interaction becomes less effective.
Analyzing the effect of temperature, i.e. of phonon-scattering, we employ the widely used Parker-plot Parker (1951), based on the relation for the AMR ratio,
$$\frac{\mathrm{\Delta }\rho \left(T\right)}{\rho \left(T\right)}=\left[\left(\frac{\mathrm{\Delta }\rho }{\rho }\right)_d\left(\frac{\mathrm{\Delta }\rho }{\rho }\right)_T\right]\frac{\rho \left(0\right)}{\rho \left(T\right)}+\left(\frac{\mathrm{\Delta }\rho }{\rho }\right)_T.$$
(4)
This equation is valid under the two premises: (i) the electric transport is dominated by one spin-channel, i.e. the majority channel in Co Gurney et al. (1993), and (ii) Matthiessen’s rule applies for the thermal and defect scattering. The validity of the latter has been demonstrated by the fits of $`\rho (T)/\rho (0)`$ to Equ.1 indicated in Fig.2. Then a plot of $`\mathrm{\Delta }\rho (T)/\rho (T)`$ vs. $`\rho (0)/\rho (T)`$ allows to separate the thermal contribution to the AMR, $`(\mathrm{\Delta }\rho /\rho )_T`$, from the defect one, $`(\mathrm{\Delta }\rho /\rho )_d`$. In fact, the extrapolation of the ’high-temperature’ data, shown by the inset to Fig.6a, is consistent with a common intercept at $`(\mathrm{\Delta }\rho /\rho )_T`$ = -0.40 $`\%`$. Such negative contribution to the AMR has been realized early on crystalline PermalloyBerger and Friedberg (1968) and, more recently, also on polycrystalline Co-filmsFreitas et al. (1990). It was associated with phonon-scattering rather than with magnon contributions. At lower temperatures, our data break away from the straight lines, which in the thickest film occurs at a rather high temperature, where $`\rho (0)/\rho (T)0.6`$. This feature indicates a change of the dominant defect type for scattering and has also been observed by Freitas et al.Freitas et al. (1990) on various Co-films with different $`\rho (0)^{}s`$, i.e. different degrees of disorder.
### III.3 Geometrical Size Effect
As a guide for discussing the GSE, we refer to Potter’s workPotter (1974), who calculated the AMR’s produced by majority and minority spins for polycrystalline cubic ferromagnets. He assumed an isotropic scattering potential, as it may be present in the grains of our films. Calculating the sd-scattering rates, Potter considered the effect of the spin-orbit interaction on localized 3d-states, but ignored the influence on the band-structure. Therefore, we expect only a more or less qualitatively correct guidance by infering the AMR-ratio from Ref.Potter, 1974:
$$\frac{\mathrm{\Delta }\rho }{\rho }=\frac{1+r}{2+r}\left\{\frac{3\sqrt{3}}{64}\left(\frac{K_{SO}}{_d}\right)^2\frac{r}{560}\left(\frac{K_{SO}}{2_{ex}}\right)^2\right\}.$$
(5)
Here $`K_{SO}`$0.1 eV measures the spin-orbit coupling energy $`_{SO}=K_{SO}\stackrel{}{L}\stackrel{}{S}`$. The positive contribution to Equ.5 arises from the longitudinal part of $`_{SO}`$ mixing two 3d-orbitals of the minority bands separated by $`_d`$ near the Fermi-surface $`_F`$ . The negative term is due to the nondiagonal part of $``$<sub>SO</sub>, which admixes some of the exchange-split majority states to the minority band. The parameter $`r=\tau _{sd}^1/\tau _{ss}^1`$ accounts for the different scattering rates of the conduction electrons into the 4s- and 3d-states and is mainly determined by the density of states of the 3d-bands at $`_F`$. Because the exchange splitting $`2_{ex}`$ is significantly larger than $`_d`$, the negative majority spin contribution to the AMR may be small relative to the positive one. Taking $`r10`$ from a recent experiment on Co-filmsGurney et al. (1993), $`\mathrm{\Delta }\rho /\rho =(3\sqrt{3}/64)(K_{SO}/_d)^2`$ follows from Equ.5. Comparing this estimate with our result for the transverse AMR at low temperatures, $`\mathrm{\Delta }\rho _t(0)/\rho _b(0)10^2`$, we obtain for the effective splitting of the two unperturbed 3d-levels, $`_d3.0K_{SO}0.3\text{eV}`$. This finding for $`_d`$ becomes smaller if a finite contribution by the majority spins would be considered in Equ.5, but it seems to be reasonable regarding the other simplifying assumptions of the theoryPotter (1974). Here we mention the neglect of the effects of the lattice potential and the spin-orbit interaction on the Fermi-surface and on the density of states at $`_F`$Potter (1974), and also of possible hybridizations between the s- and d-orbitalsFert and Campbell (1976).
Nevertheless, we will extend Potter’s results derived for an ’isotropic’, i.e. polycrystalline cubic ferromagnet, to films with polar texture. Let us recall that the AMR originates from a symmetry breaking of the 3d-orbitals by the magnetization $`\stackrel{}{M}`$ via the spin-orbit coupling. The resulting anisotropic charge distribution gives rise to the scattering asymmetry of the conduction electrons into these 3d-states. On general grounds one may expect that a reduction of the symmetry of the ferromagnet structure weakens the AMRMcGuire and Potter (1975), loosely spoken, because then the magnetization induced axial anisotropy of the orbitals becomes less effective. Since the texture in the permalloyRijks et al. (1997) and in our cobalt films, both perpendicular to the plane, appear to be strongly correlated with the GSE, we assume the mixing parameter in Equ.5, $`k_\alpha ^2=3\sqrt{3}(K_{SO}/4_\alpha )^2`$, to be different for the in-plane $`(\alpha =i)`$ and the polar $`(\alpha =p)`$ directions of $`\stackrel{}{M}`$. Then Equ.5 remains still valid for the in-plane orientations of $`\stackrel{}{M}`$ and ignoring again the small contribution by the majority spins, we have
$$\frac{\rho _{\mathrm{}}\rho _t}{\rho }=\frac{1+r}{2+r}\frac{3\sqrt{3}}{4}k_i^2.$$
(6)
In order to determine the effect of the film anisotropy on the polar MR, we introduce $`k_\alpha ^2`$ directly into PottersPotter (1974) result for the perpendicular conductivity of the minority spins, $`\sigma _{}^\alpha /\sigma _0=(3\sqrt{3}/2r)\mathrm{ln}\left[r/(1+\frac{1}{2}rk_\alpha ^2)\right]`$. For small spin-orbit perturbations, $`rk_\alpha ^21`$, the difference between the transverse and polar resistivities becomes:
$$\frac{\rho _t\rho _p}{\rho }=\frac{1+r}{2+r}\frac{3\sqrt{3}}{4}(k_i^2k_p{}_{}{}^{2}).$$
(7)
By some trivial algebra we obtain for the GSE from Equs.6, 7:
$$\frac{\mathrm{\Delta }\rho _p}{\mathrm{\Delta }\rho _t}=2\frac{k_p^2}{k_i^2}.$$
(8)
Hence, this simple model can explain the upper limit of two of the GSE, which emerges from our data in Fig.6b and also from Fig.6 in Ref.Rijks et al., 1997 for Permalloy films. Moreover, this model ascribes the GSE to the electronic structure, as it was suspected by Chen and MarsocciChen and Marsocci (1972). Consequently, the GSE should not depend on the temperature which is fully consistent with our results depicted in Fig.6b.
Equation 8 also predicts that the upper limit of two is reached, when the mixing effect due to the polar oriented magnetization is small compared to mixing by the in-plane $`\stackrel{}{M}`$, i.e. $`k_p^2<<k_t^2`$. This case seems to be realized in our films, see Fig.6b, and also in the Permalloy films with increased $`111`$-epitaxy (Fig.6 of Ref.Rijks et al., 1997). These observations indicate that the spin-orbit induced anisotropy in the 3d-orbitals near $`_F`$ is smallest, if $`\stackrel{}{M}`$ is aligned parallel to the existing axial perturbation built in by hcp- or $`111`$-epitaxy. In this case, the 3d-orbitals have already the axial symmetry so that an induced magnetization along the epitaxial (polar) direction may have only a moderate effect on the scattering probability into these states. This is in distinct contrast to the in-plane orientation of $`\stackrel{}{M}`$ which breaks the symmetry of these orbitals. Therefore, the mechanism proposed here for the GSE qualifies the film anisotropy of the AMR more precisely as structural, rather than as a geometrical effect.
## IV Low-field magnetoresistance
### IV.1 Domain Structures
The formation of domains affects the MR’s of the 188 nm thick film and of the thinner films, d$``$ 20 nm, in quite different ways. The interesting features can already be realized on the large field scale of Fig.5: (i) for d=20 nm Co (and also for 10 nm, not shown) both, the polar and the transverse MR’s approach the field-independent longitudinal MR, $`\rho _{\mathrm{}}`$, whereas the polar and the longitudinal MR’s of the 188 nm film tend to the field-independent transverse resistance $`\rho _t`$. In order to provide some solid basis for a detailed discussion of these characteristic features of the domain MR’s, we examine the domain structures by magnetic force microscopy (MFM).
The essential difference between the thick (188 nm) and the thinner films can be infered from MFM images of the demagnetized states, shown in Fig.7a. The images have been recorded in the dynamic mode of the Q-scope which is sensitive to the polar gradient of the polar force, i.e. $`\delta F_p/\delta x_p=M_p\delta ^2H_p/\delta x_p^2`$. The 20 nm film consists of large, some 10 $`\mu m`$ wide domains with in-plane magnetizations separated by 180 Neel walls. The domain magnetizations are oriented parallel to the uniaxial anisotropy field $`\stackrel{}{H}_u`$ as determined by FMR. A slight longitudinal ripple of $`\stackrel{}{M}`$ about $`\stackrel{}{H}_u`$ is visible, which most likely arises from the polycrystallinity of the film.
In constrast, the 188 nm film exhibits a maze configuration of stripe domains with sizable polar components of the domain magnetizations. The Fourier transform of the image in Fig.7a yields a mean width of the stripes, $`d_D=(205\pm 15)`$ nm, being rather close to the thickness as expected for weak stripes by magnetostatic reasons Hubert and Schäfer (1998). Recently, the same observations were reported for a 195 nm thin polycrystalline Co-film on glass and related to a hexagonal texture perpendicular to film planeKharmouche et al. (2004). MFM images depicted on epitaxial Co-films revealed a reorientation of the domain magnetization from in-plane to polar between 10 nm and 50 nm Hehn et al. (1996); Demand et al. (2002) which was explained in terms of the perpendicular magnetocrystalline anisotropy of Co. These results suggest that also in our case the hcp texture, realized by the XRD (Fig.1), generates such a crystalline anisotropy, which in the 188 nm film becomes large enough to produce a significant polar component of $`\stackrel{}{M}`$. Let us also recall that we supposed this texture already in the discussion of the GSE.
The other interesting property of these weak stripes is seen in Fig.7b. In the remanent states, stripe patterns are found aligned with the direction of the previously applied fields $`H_{\mathrm{}}`$ or $`H_b`$. This socalled rotable anisotropy can be attributed to the stiffness of the domain walls against deformations Hubert and Schäfer (1998) and is probably supported by a pinning of the walls by local anisotropies in the granular structure. The rotatable anisotropy suggests also an ’isotropic’ hysteresis loop, the shape of which should be independent on the direction of the in-plane field. In fact, we do observe this feature on the 188 nm film, see Fig.9 below, and will refer to it when discussing the MR in the domain state.
### IV.2 Anisotropic Magnetoresistance
We begin with the low-field resistance of the thin films, exemplified by Figs.3a, 5a for d=20 nm: both the transverse and the polar MR’s, $`\rho (H_t0)`$ and $`\rho (H_p0)`$, tend to the longitudinal one, $`\rho (H_{\mathrm{}})`$. This behavior is readily explained by the fact that the resistance is measured along $`\stackrel{}{H}_u`$, and that at low fields the domain magnetization is also directed parallel to $`H_u`$ evidenced by MFM (Fig.7a). The parabolic decrease of $`\rho `$ in larger polar fields, $`\mathrm{\Delta }\rho (H_p<M_s)H_p^2`$, was already attributed to the AMR resulting from the rotation of $`\stackrel{}{M}`$ from an in-plane to the polar direction. Also the detailed variation of the in-plane MR’s, shown in Fig.8a, can be explained by the AMR. Using the hysteresis loops $`M(H_i)`$ in Fig.8b, and assuming the relations for the angle $`\phi `$ in Equ.2,
$$\mathrm{cos}\phi \left(H_{\mathrm{}}\right)=M\left(H_{\mathrm{}}\right)/M_s,$$
(9a)
$$\mathrm{sin}\phi \left(H_t\right)=M\left(H_t\right)/M_s,$$
(9b)
the in-plane MR can be described rather nicely. The physical arguments for these agreements are: (i) the longitudinal magnetization process, $`M(H_{\mathrm{}})`$, is due to the nucleation of 180 Neel walls (see Fig.7a) at the coercive field $`H_c=H_u`$ (determined by FMR), which then rapidly cross the film leaving the resistance unchanged; (ii) upon reduction of the transverse field, on the other hand, a longitudinal ripple of $`\stackrel{}{M}`$ about $`\stackrel{}{H}_t`$ appears which originates from $`\stackrel{}{H}_u`$ (see e.g. Ch.5.5 of Ref.Hubert and Schäfer, 1998). Accordingly, the components of $`\stackrel{}{M}`$ parallel and antiparallel to the current $`\stackrel{}{I}||\stackrel{}{H}_u`$ are growing continuously so that $`\rho (H_t)`$ increases until the transverse coercive field $`H_c^t<H_u`$ is reached. There the magnetization component along $`\stackrel{}{H}_u`$ changes sign and increases at the expense of the ripple, so that $`\rho (H_t)`$ back to $`\rho _t`$ at larger negative fields.
A rather different behavior is displayed by the 188 nm thick film. Already in Fig.5b we noticed that at low fields the polar and the longitudinal MR’s tended to the transverse MR. As a rather unexpected feature, the transverse MR turned out to be nearly independent of the field also in the domain regime, $`\rho (H_t)=\rho _t`$. The detailed variation of the in-plane MR’s at low fields is shown in Fig.9a revealing just the opposite to the behavior of the thin films (see Fig.8a): the longitudinal MR displays a strong field dependence, while the transverse MR remains very small. These results are explained also by the AMR effect. The in-plane MR, shown in Fig.9b, is rather nicely reproduced by the solid curves which have also been calculated from Equ.2. Again, the mean angle $`\phi `$ between current and magnetizations $`\stackrel{}{M}(\stackrel{}{H})`$ has been determined from Equ.9 and the hysteresis loops, Fig.9b. As a matter of fact, we emphasize that the shape of these loops does not depend on the direction of the in-plane field (’rotatable loops’). This is consistent with the corresponding behavior of the weak stripe domains depicted by MFM in Fig.7b. In contrast to the thin films, d$``$20 nm, no effect by the uniaxial in-plane anisotropy field, $`H_u=15`$ Oe, determined by FMR Gil (2004), is realized. The much larger coercive field, $`H_c200`$ Oe, stems most likely from the pinning of the stripe domain walls by the random polycrystalline anisotropy in the films.
### IV.3 Effective Medium Approach
Aiming at a more detailed description of the MR in the 188 nm film, again the domain structure has to be taken into account. For this purpose, we use an effective medium model, by which Rüdiger et al.Rüdiger et al. (1999) successfully interpreted the AMR of epitaxial Co-films. Introducing the volume fractions $`v_i`$ for different domain species, the AMR is approximated by
$$\mathrm{\Delta }\rho (\stackrel{}{H})=\underset{i=1}{\overset{3}{}}\upsilon _i(\stackrel{}{H})\mathrm{\Delta }\rho _i.$$
(10)
Here the $`\mathrm{\Delta }\rho _i`$ denote the AMR’s of the corresponding domain with polar, transverse or longitudinal orientations of $`\stackrel{}{M}`$ relative to current and film plane. By definition is $`\mathrm{\Delta }\rho _{\mathrm{}}=0`$ and if, for convenience, $`\mathrm{\Delta }\rho (\stackrel{}{H})`$ is normalize to the transverse MR, Equ.10 takes the form
$$\frac{\mathrm{\Delta }\rho (\stackrel{}{H})}{\mathrm{\Delta }\rho _t}=\upsilon _t(\stackrel{}{H})+g_s\upsilon _p(\stackrel{}{H}),$$
(11)
where $`g_s=\mathrm{\Delta }\rho _p/\mathrm{\Delta }\rho _t`$ denotes the GSE-ratio. The simplest case, $`\upsilon _t=\upsilon _p=0`$ and, hence, $`\mathrm{\Delta }\rho =0`$ has been realized on the thin films at low fields.
The most interesting example is the 188 nm film, where (i) the low-field MR appears to be inverted relative to the thin films and, moreover, (ii) the transverse resistivity remains at the saturation value $`\mathrm{\Delta }\rho (H_t)=\mathrm{\Delta }\rho _t`$, even in the domain state. We will now attempt to relate these striking features displayed by Figs.5b, 9 to the domain structure observed by MFM, see Fig.7. Observation (ii) in connection with Equ.11 implies for the concentration of polar oriented domains,
$$\upsilon _p(H_t)=\frac{1}{g_s}\left[1\upsilon _t(H_t)\right].$$
(12)
Below the saturation field, the magnetization $`M(H_t)`$ and therefore, $`\upsilon _t(H_t)`$, starts to decrease at the expense of a finite polar component $`\upsilon _p`$, which leads to the nucleation of stripe domains. Upon further reduction, $`H_t0`$, the hysteresis loops display a normalized remanent magnetization $`M(H_t0)/M_s=0.66(2)`$, i.e. volume fraction $`\upsilon _t(0)2/3`$. For an estimate, we take the maximum GSE, $`g_s=2`$, to find from Equ.11 $`\upsilon _p(0)=1/6`$ and by using $`\underset{i=1}{\overset{3}{}}\upsilon _i=1`$,the same longitudinal volume fraction $`\upsilon _{\mathrm{}}(0)=1/6=\upsilon _p(0)`$. The agreement of both volumes implies that the nucleation of polar domains is accompanied by the creation of an equal amount of longitudinally oriented domain. Considering the square-like cross-section of the stripes following from Fig.7b, this result indicates that the flux extending from the polar phase is closed by the longitudinal volume $`\upsilon _{\mathrm{}}(H_c)`$. The rotatable symmetry of the hysteresis loops implies for the longitudinal direction also $`\upsilon _p(H_{\mathrm{}}0)=\upsilon _t(H_{\mathrm{}}0)=1/6`$. For the longitudinal MR Equ.10 predicts then $`\mathrm{\Delta }\rho (H_{\mathrm{}}0)/\mathrm{\Delta }\rho _t=1/6+21/6=1/2`$, which is in close agreement with the measured value, see Fig.9a.
Finally, upon reduction of $`H_t`$ to the coercive field, $`\upsilon _p(H_t)`$ increases further. The volume fraction of the polar domains at $`H_c`$ can be estimated from the stripe maze of the demagnetized state, Fig.7a, which suggests $`\upsilon _t(H_c)=\upsilon _{\mathrm{}}(H_c)`$. Then, from Equ.11 and simple algebra we obtain $`\upsilon _p(H_c)=1/3=\upsilon _{\mathrm{}}(H_c)=\upsilon _{\mathrm{}}(H_c)`$. Thus the demagnetized state consist of equal volumes for all six possible magnetization directions, which by considering the symmetry of the stripe structure is again a plausible result.
## V Summary and Conclusions
The magnetoresistance of polycrystalline Co-films, which were characterized by XRD, FMR, SQUID-magnetometry, AFM, MFM and temperature variable resistivity, has been investigated in fields up to 100 kOe directed along three principal directions. In the saturated state, the MR displayed the socalled geometrical size-effect (GSE), according to which the MR for the polar orientation of $`\stackrel{}{M}`$ is up to twice as large as for the in-plane $`\stackrel{}{M}`$. The determination of the GSE was facilitated by the facts that the spin-wave contributions could be easily subtracted and that in the present disordered films the classical Lorentz MR proved to be negligible. Basing on a correlation between the GSE and a texture detected previously on Permalloy filmsRijks et al. (1997) and also on our Co-films by XRD, we proposed to attribute the GSE to an anisotropic mixing of the 3d-levels near $`_F`$ by the longitudinal part of the spin-orbit interaction. By extending Potter’sPotter (1974) prediction for the AMR of the minority spin channel, we obtained a result which is consistent with the observed upper limit of two for the GSE and also with the temperature independence of the GSE. A relation of the GSE to the electronic structure has already been conjectured in the literatureChen and Marsocci (1972), but not yet been worked out. Of course, regarding the simplicity of the proposed extension and the assumption of a simple spherical Fermi-surface for the final 3d-states in Ref.Potter, 1974, which considers only the local aspect of the spin-orbit interaction, these consistencies may be fortuitous. However, we believe that the central argument for the appearance of the GSE, i.e. the presence of an additional uniaxial symmetry in polycrystalline films through a texture in thin films remains valid. Hence, a more detailed reasoning for the GSE considering also the effect of the spin-orbit interaction on the band-structure is indicated in order to check the present rough model.
The MR’s in the domain state were interpreted using the saturated AMR’s and the GSE, the hysteresis loops, and MFM images of the domain structure. For thin Co-films, $`d20\text{nm}`$, where the magnetization remained the film-plane and became for $`\stackrel{}{H}0`$ parallel to the weak uniaxial anisotropy field, the MR attained the maximum (longitudinal) value at zero magnetization. The MR in the domain state of the thickest film, d=188 nm, on the other hand, displayed a rather different behavior. As a function of the transverse field in the film-plane, the resistance turned out to be almost constant, whereas upon reducing longitudinal and polar fields the resistances decreased and increased, respectively, from their different saturation values to the transverse MR. MFM images and hysteresis loops revealed the formation of rotatable stripe domains with square cross-section due to the hcp texture. By means of an effective medium modelRüdiger et al. (1999), the MR’s could be quantitatively explained in terms of a flux closure configuration of the magnetization components about the directions of the stripes. Approaching the coercive field, the stripes terminated in a maze configuration, and the fractional volumes of all three magnetization components proved to be equal. In this model, the surprising field independence of the transverse MR results from the squared corss-section of the stripes with transverse flux-closure and from a GSE ratio of two. We should note that this discussion did not invoke (possible small) contributions to the MR by the Neel- and Bloch-walls in the thin and thick films, respectively. Such effects have been reported before in epitaxial Co-filmsGregg et al. (1996); Viret et al. (1996); Rüdiger et al. (1999) with strong hcp crystalline anisotropy and quantitatively different domain dimensions.
The authors are indebted to the late Prof. J. Appel (Hamburg) for his encouraging interest during the early steps of this work. Discussions with Proff. P. Böni (München), D. Grundler (Hamburg) and M. Morgenstern (Aachen) are gratefully acknowledged.
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# Does Standard Cosmology Express Cosmological Principle Faithfully?
## I A 1+1 Dimensional Model
Consider a one dimensional infinitely long system consisting of uniformly placed galaxies, see FIG.1. Suppose the system is expanding uniformly, i.e., from any galaxies (such as $`O`$), we will see that the two galaxies ($`A`$ and $`B`$) mostly nearest to us are recessing away from us at equal speeds, and the distances between us and this two neighbors are equal.
In standard cosmology, the scale factor is scale independent, i.e., if on the scale of $`|OB|`$, the scale factor of the system is $`a(t)`$, then on the scale of $`|OC|`$, the scale factor is also $`a(t)`$. So the physical length of $`|OB|`$ is half of $`|OC|`$ and the metric describing the system is
$`ds^2=dt^2+a^2(t)dx^2`$ (1)
However, for a one-dimensional gravitation system, Einstein equation cannot give us anything about its dynamical evolutions, note $`G_{\mu \nu }=R_{\mu \nu }\frac{1}{2}Rg_{\mu \nu }==0`$. So if we insists eq(1) is the only correct metric ansaltz of the system illustrated in FIG.1, then we have no way to determine the function form of $`a(t)`$.
The reason that standard cosmology insists eq(1) is, the density of the system is not function of the position co-ordinate $`x`$. We doubt this statement faithfully expresses the requirements of cosmological principle. Let us consider the following series
$`v_B=v;`$
$`v_C={\displaystyle \frac{v+v}{1+v^2}};`$
$`v_D={\displaystyle \frac{v+v_C}{1+vv_C}};`$
… …
$`v_X={\displaystyle \frac{v+v_{X1}}{1+vv_{X1}}};`$ (2)
$`|AB|=2a`$
$`|OC|=2a\sqrt{1v_B^2}`$
$`|BD|=2a\sqrt{1v_C^2}`$
… …
$`|X^{}X^+|=2a\sqrt{1v_X^2},`$ (3)
From which we get
$`v_X={\displaystyle \frac{(1+v)^X(1v)^X}{(1+v)^X+(1v)^X}};`$ (4)
$`|OX|`$ $`a{\displaystyle \underset{N=0}{\overset{X}{}}}\sqrt{1v_N^2}`$ (5)
$`=l{\displaystyle _0^X}𝑑x\sqrt{1v_x^2}`$
$`={\displaystyle \frac{4a}{\text{ln}\frac{1+v}{1v}}}\left[\text{arctg}[({\displaystyle \frac{1+v}{1v}})^{\frac{x}{2}}]{\displaystyle \frac{\pi }{4}}\right],`$
In eqs(2)-(5), $`v`$ is the relative recessing speed between two nearest galaxies, $`a`$ is the physical distance between them, it can also be considered as the scale factor on the smallest scales,
$`a=vt,`$ (6)
where $`t`$ is understood as the time from the system being created (the distance between any two galaxies is zero) to the epochs we observe it.
In our models, we will not consider dark energies. But we assume that the relative recessing velocity between any two nearest galaxies is a time-independent constance. (One reason for our assumption is, if the system is at rest initially, it will not collapse at self-gravitations from symmetry analysis, what matters here is parity symmetry. So when the system is expanding but cosmological principle is always kept, it will not decelerate because the parity symmetry is not broken by expansions)
Let
$`\sigma ={\displaystyle \frac{1}{2}}\text{ln}{\displaystyle \frac{1+v}{1v}}`$ (7)
we can write down the metric of our one dimensional cosmology in FIG. 1 as
$`ds^2=dt^2+{\displaystyle \frac{4v^2t^2}{(e^{\sigma x}+e^{\sigma x})^2}}dx^2`$ (8)
In this metric space, physical co-ordinate $`x_{ph}`$ is related to $`x`$ by
$`x_{ph}={\displaystyle \frac{2vt}{\sigma }}(\text{arctg}[e^{\sigma x}]{\displaystyle \frac{\pi }{4}}).`$ (9)
Note in eqs(8) and (9), $`x`$ is a pure number of no-dimensions. Before a length unit is assigned, the difference of it has no meaning of any distance lengths. But if we let
$`dx_{pr}=dxvt`$ (10)
then $`dx_{pr}`$ can be naturally interpreted as the proper length of line element between points $`(t,x)(t,x+dx)`$ at time $`t`$. Using co-ordinate $`\{t,x_{pr}\}`$, the metric eq(8) can be written in the following form
$`ds^2=dt^2+{\displaystyle \frac{4}{(e^{\frac{\sigma x_{pr}}{vt}}+e^{\frac{\sigma x_{pr}}{vt}})^2}}dx_{pr}^2`$ (11)
Because we are so deeply affected by Friedmann-Robertson-Walker metric and only familiar with only-time-dependent (or although both time- and position-dependent in the non-flat universes but the two are separated) scale factors, we will mostly use eq(8) in this paper. It is worth noting that the $`g_{00}`$ component of eq(8) has different dimension from $`g_{11}`$. Let us put this in mind so that when dimension problems appear, correct interpretation can be given.
Some people may ask, why not redefine the co-ordinate $`x`$ so that the $`x_{ph}`$ has simple linear dependence on it? Yes, we can do that way, but we should note after the re-definition, the relation between $`x_{pr}`$ and $`x`$ will change correspondingly, which will introduce corresponding complexities, so we will not re-define the co-ordinate $`x`$ at this time. Physically, if we re-define $`x`$ to $`\stackrel{~}{x}`$ so that the metric eq(8) has the form $`ds^2=dt^2+t^2d\stackrel{~}{x}^2,`$ then we should note in equal length of line elements $`(t,\stackrel{~}{x}\frac{1}{2}d\stackrel{~}{x})`$ $``$ $`(t,\stackrel{~}{x}+\frac{1}{2}d\stackrel{~}{x})`$ and $`(t,\stackrel{~}{x}^{}\frac{1}{2}d\stackrel{~}{x})`$ $``$ $`(t,\stackrel{~}{x}^{}+\frac{1}{2}d\stackrel{~}{x})`$, we will not find equal number of galaxies, as long as $`|x||\stackrel{~}{x}|`$.
Obviously, eqs(8) or (11) contradicts the standard cosmological results eq(1) remarkablly. Standard cosmology insists that cosmological principle requires the density of the system is not function of the position co-ordinate $`x`$, so get its metric ansaltz eq(1). While we insists that cosmological does not require so, it only requires that on any galaxies, observers will measure that his two nearest neighbors are equally far away from him and recessing at equal velocities. The density of the system can be functions of the position co-ordinate, as long as whereever the origin is chosen, the metric function has the same form. Although we do not think the generalization from 1+1 dimension to 1+3 dimension is trivial, we think this is an evidence that standard cosmology may not express cosmological principle faithfully. We will put aside debates at this moment and focus on the fact if we generalize the metric eq(8) or (11) into three dimensions, the prediction is consistent with observations or not. We will answer the main criticisms from standard cosmologists in CosmoSDSFext .
## II Generalization To 1+3 Dimensions
In generalizing eq(8) to three dimensional case we have two methods, i.e.,
$`ds^2=dt^2+{\displaystyle \frac{v^2t^2}{\mathrm{cosh}^2\sigma r}}(dr^2+r^2d\theta ^2+r^2\text{sin}^2\theta d\varphi ^2)`$ (12)
or
$`ds^2=`$
$`dt^2+{\displaystyle \frac{v^2t^2}{\mathrm{cosh}^2\sigma x}}dx^2+{\displaystyle \frac{v^2t^2}{\mathrm{cosh}^2\sigma y}}dy^2+{\displaystyle \frac{v^2t^2}{\mathrm{cosh}^2\sigma z}}dz^2`$
(13)
It is impossible to get eq(12) from eq(13), because eq(13) describes a cubic lattice system, while eq(12) describes a spherical symmetric system. The most remarkable difference between eq(12) and the usual Friedmann-Robertson-Walker metric are, (i) the maximum symmetric subspace of metric space eq(12) is 2-spheres, while that of standard cosmology is homogeneous 3-balls; (ii) eq(12) contains no unknown function such as standard cosmology’s scale factor, the evolution of the universe are completely prescribed by one parameter $`v`$. We will explain these differences in CosmoSDSFext
As the first step, let us verify that the metric eq(12) indeed describing a homogeneous, isotropic system. Using Einstein equation $`G_{\mu \nu }=8\pi GT_{\mu \nu }`$ we calculate the energy momentum tensor of our cosmology in the following
$`8\pi GT_{\mu \nu }=\text{diag}`$
$`\{{\displaystyle \frac{4\sigma (1+e^{4\sigma r})+\sigma ^2(110e^{2\sigma r}+e^{4\sigma r})r12e^{2\sigma r}rv^2}{4e^{2\sigma r}rt^2v^2}}`$
$`,{\displaystyle \frac{2\sigma (1+e^{4\sigma r})\sigma ^2(1+e^{2\sigma r})^2r+4e^{2\sigma r}rv^2}{(1+e^{2\sigma r})^2r}},`$
$`{\displaystyle \frac{r(\sigma [1+e^{4\sigma r}]+4\sigma ^2e^{2\sigma r}r+4e^{2\sigma r}rv^2)}{(1+e^{2\sigma r})^2}},`$
$`{\displaystyle \frac{r(\sigma [1+e^{4\sigma r}]+4\sigma ^2e^{2\sigma r}r+4e^{2\sigma r}rv^2)\mathrm{sin}^2\theta }{(1+e^{2\sigma r})^2}}\}`$
(14)
Note in our frame-works, the dimension of $`T_{00}`$ is different from that of $`T_{ii}`$, because the component of our metric are of different dimensions. The same problem will occur on the four velocity of an observer, see eq(15). If we use the metric of form eq(11), this will not be a problem. Superficially, our energy momentum tensor is position-dependent, which seems to violate cosmological principles. This is not the case. Let us calculate the energy density and pressures measured by an observer at position $`(t,r,\theta ,\varphi )`$, whose four velocity can be written as
$`u^{(t,r,\theta ,\varphi )}={\displaystyle \frac{1}{\sqrt{N}}}\{1,{\displaystyle \frac{v_r}{vt}},0,0\}`$
$`\text{where }v_r={\displaystyle \frac{e^{\sigma r}e^{\sigma r}}{e^{\sigma r}+e^{\sigma r}}},`$
$`N=1{\displaystyle \frac{4v_r^2}{(e^{\sigma r}+e^{\sigma r})^2}}.`$ (15)
It is very easy to verify $`g_{\mu \nu }u^\mu u^\nu =1`$. Measured by this observer, the energy density and pressure are respectively
$`8\pi G\rho =\left[T_{\mu \nu }u^\mu u^\nu \right]_{v0}=9t^2,`$
$`8\pi Gp=\left[T_{\mu \nu }(g^{\mu \nu }+u^\mu u^\nu )\right]_{v0}=9t^2.`$ (16)
Some people may not understand the limit procedure in eq(16), please see the notations under eq(26). From eqs(14)+(16) we can see that, viewing from any point, we can see an isotropic but not in-homogeneous universe. The inhomogeneity originates from Lorentz contraction, it is just a kinematical effects instead a dynamical one. Obviously, if we can take photos of the universe from different places at the same time, we get the same pictures. We think this is a faithful expression of cosmological principle. While standard cosmology’s statements, the energy momentum tensor should not depends on the position co-ordinate of the universe, does not express cosmological principle faithfully.
Now let us consider super-novaes in the metric space eq(12). We want to calculate their luminosity-distance v.s. red-shift relations. Let us follow the same procedures from SWeinberg , section 14.4, eqs(14.4.11-14). Consider a super-novae at position $`(t,r,\theta ,\varphi )`$, its recessing velocity relative to us is
$`v_r={\displaystyle \frac{e^{\sigma r}e^{\sigma r}}{e^{\sigma r}+e^{\sigma r}}};`$ (17)
so, if the proper frequency of a photon emitted from this super-novae is $`\omega _0`$, the frequency measured by us is $`\omega `$, then the red-shift $`z`$ of this photon satisfy
$`(1+z){\displaystyle \frac{\omega ^1}{\omega _0^1}}=e^{\sigma r};`$ (18)
considering Lorentz dilating, the photons emitted in period $`\delta t_1`$ can only reach us in period $`\delta t_1e^{\sigma r}`$. So we get the luminosity distance v.s. red-shift relation as
$`d_l=(1+z){\displaystyle \frac{2vH_0^1}{\sigma }}[\text{arctg}(1+z){\displaystyle \frac{\pi }{4}}]`$ (19)
The relation between $`\sigma `$ and $`v`$ is given by eq(7). From best fitting observational results of Riess04 , we get $`v=0.79/3000`$, $`H_0=60\text{km/(s}\text{Mpc)}`$, $`\chi ^2=303`$ (186Golden+Silver sample) or $`v=0.899/3000`$, $`H_0=`$$`60\text{km/(s}\text{Mpc)}`$, $`\chi ^2=237`$ (157Golden sample).
We illustrate the numerical results of eq(19) in FIG.2. As comparisons, we also depict the predictions of $`\mathrm{\Lambda }`$CDM and $`s`$CDM cosmologies. From the figure we can easily see that the prediction of our eq(19) is very similar to that of $`\mathrm{\Lambda }`$CDM cosmology. Because we consider special relativity effects on the definition of homogeneity in our theory, we call our results eqs(8), (12) and (19) in this paper as Relativity Cosmology.
## III Comparing Our Models With Standard Cosmology
Now let us return to standard cosmology where we are taught that homogeneity and isotropy of the observed universe directly leads to the conclusion that Freimann-Robertson-Walker metric is the unique metric describing our real universe (we will put perturbations and structure formation problems in the future works),
$`ds^2=dt^2+{\displaystyle \frac{a^2(t)}{1kr^2}}dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),`$
$`\text{where,}k=+1,0,1.`$ (20)
While in the co-moving co-ordinate the energy momentum tensor of the cosmological fluid has the form
$`T_{\mu \nu }=\text{diag}(\rho ,p,p,p)`$
if no radiation and/or
no dark energy is involved
$`=\text{dial}(\rho (t),0,0,0)`$ (21)
For a stable universe, special relativity has no effects on the homogeneity definition of it, please see the left part of FIG. 3. But for an expanding universe, its homogeneity definition is affected by special relativity strongly. If we neglect special relativity effects, then if we were put on a given galaxy such as $`O`$ and were asked to measure the distances between us and the mostly nearest neighbors, see the right part of FIG. 3, we get results, say $`r`$; if we were asked to measure distances between us and the next-nearest neighbors (should be on the same line with the previous galaxies) we get results $`2r`$, and so on. So the system has translation symmetry at a given time, i.e., the maximum symmetric subspace of the whole space-time is homogeneous 3-ball.
If special relativity is considered, then when we were put on a given galaxy $`O`$ and were asked to measure the distances between us and the nearest neighbors ($`B`$), we got results, say $`r`$. But if we were asked to measure the distances between us and the next-nearest but on the same line neighbor ($`C`$), we did not get $`2r`$, we got a number less than $`2r`$ because of Lorentz contraction. In this case the maximum symmetric subspace of our physical universe is only a 2-sphere instead of a homogeneous 3-ball.
According to the results of SWeinberg , section 13.5, the general metric describing a space-time with maximum symmetric subspace of $`S_2`$ can only be reduced to
$`ds^2=dt^2+U(t,r)dr^2+V(t,r)(d\theta ^2+\text{sin}^2\theta d\varphi ^2)`$ (22)
instead of eq(20). Just for the same reason, we can only write the energy momentum tensor describing our real universe as
$`T_{\mu \nu }=\rho u_\mu u_\nu +p(g_{\mu \nu }+u_\mu u_\nu )`$
$`\text{where }u^{(t,r,\theta ,\varphi )}(1,v_r,0,0)`$ (23)
instead of eq(21). If generalize our metric eq(11) into three dimensions it will just has the form of eq(22), while the appropriate energy momentum tensor eq(14), will also have the form of eq(23) correspondingly.
On the contrary, standard cosmology does not consider special relativity when define homogeneities, which will introduce problems to it. We provide one in the following. Starting from metric (20), let $`k=0`$ for the moment, using Einstein equation we calculate the energy momentum tensor
$`8\pi GT_{\mu \nu }=H^2(t)\text{Dial.}\{3,A(t),A(t)r^2,A(t)r^2\mathrm{sin}^2\theta \},`$
$`\text{where }A(t)=a^2(t)+{\displaystyle \frac{a^3(t)a^{\prime \prime }}{a^2}}.`$ (24)
Note $`T_{00}`$ only depending on $`t`$ does not mean observers on different places will get the equal energy densities in measures. It only means that energy density measured by observers on the origin of the co-ordinate is position independent. To calculate the energy density measured by observers on different places, we have to consider observers on general positions $`(t,r,\theta ,\varphi )`$, whose four velocity are
$`u^\mu ={\displaystyle \frac{1}{\sqrt{1[a(t)\dot{a}(t)r]^2}}}\{1,\dot{a}(t)r,0,0\}`$ (25)
The energy density and pressure measured by these observers are respectively
$`8\pi G\rho `$ $`=8\pi GT_{\mu \nu }u^\mu u^\nu `$
$`={\displaystyle \frac{H^2(t)}{1[a(t)\dot{a}(t)r]^2}}(3+A(t)\dot{a}^2(t)r^2)`$
$`8\pi Gp`$ $`=8\pi GT_{\mu \nu }(g^{\mu \nu }+u^\mu u^\nu )`$ (26)
$`={\displaystyle \frac{5r^2a^2\dot{a}^46a\ddot{a}+\dot{a}^2(3+4r^2a^3\ddot{a})}{a^2(1+r^2a^2\dot{a}^2)}}`$
Obviously, without a limiting procedure like that in eq(16), eq(26) will tell us that both the energy density and pressure measured by observers at different places are not the same. However, if we take the limiting $`\dot{a}0`$, the energy density of the cosmological fluid will become zero. We think this is a problem of standard cosmology. But our cosmological models in eqs(12)+(14)+(16) does not have this problem.
## IV Conclusions
We express our suspicions that standard cosmology expresses cosmological principle faithfully. In 1+1 dimension case, we prove that the background metric of the universe is not Friedmann-Robertson-Walker type. We then generalize the 1+1 dimensional results into 1+3 dimensional case and explain the observed luminosity-distance v.s. red-shift relations of super-novaes naturally without introducing any concepts of dark energies. We will answer the criticisms from standard cosmologists in another extended version of this paper, CosmoSDSFext .
Of course, the observed luminosity-distance v.s. red-shift relations of super-novaes is not the only evidence of dark energies, seeRiess98 ; Perlmutter98 ; Knop03 ; Tonry03 ; Riess04 ; WMAP03 ; Tegmark04 for experimental works and Quintessence1 ; Quintessence2 ; Phantom ; Phantom2 ; backReaction for theoretical ones. We will study the perturbations of eqs(12)+(14) and structure formation problems in the future. The original ideal of this paper is also expressed in Assumption .
Acknowledgements
When we finish the first version of this paper, we send it to professor E. Witten, G. ’t Hooft, P. J. Steinhardt, A. Linde, E. Kolb and S. Dodelson and ask them to give us some comments or criticisms, some of them accepted our asking and give us serious comments, we thank them for their comments or criticisms or encouragements very much.
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# Mesoscopic d-Wave Qubits: Can High-Tc Cuprates Play a Role in Quantum Computing?†
## I Time-reversal symmetry breaking on surfaces and interfaces of high-T<sub>c</sub> superconductors
### I.1 Description of transport in high-T<sub>c</sub> structures.
The lack of accepted microscopic theory of superconductivity in high-T<sub>c</sub> cuprates did not prevent successful research in this field since 1986. We have in several respects the repetition of the situation ca 1938, but with a clear advantage of already having BCS theory to provide insight and language for phenomenological treatment.
On this level, high-T<sub>c</sub> superconductors are successfully described by Gor’kov equations for normal and anomalous Green’s functionsMineev99 , which in Matsubara representation are defined in the usual way:
$`G_{\alpha \beta }(𝐤,\tau ;𝐤^{},\tau ^{})=𝒯_\tau a_\alpha (\tau )a_\beta ^{}(\tau ^{}),`$
$`F_{\alpha \beta }=(𝐤,\tau ;𝐤^{},\tau ^{})=𝒯_\tau a_{𝐤\alpha }(\tau )a_{𝐤^{}\beta }^{}(\tau ^{}),`$ (1)
$`F_{\alpha \beta }^+=(𝐤,\tau ;𝐤^{},\tau ^{})=𝒯_\tau a_{𝐤\alpha }(\tau )a_{𝐤^{}\beta }^{}(\tau ^{}).`$
The only difference from conventional superconductivity is in the nontrivial symmetry of the pairing potential,
$$H_{\mathrm{int}}=\frac{1}{2}\underset{𝐤,𝐤^{},𝐪}{}V_{\alpha \beta ,\lambda \mu }(𝐤,𝐤^{})a_{𝐤+𝐪/2,\alpha }^{}a_{𝐤+𝐪/2,\beta }^{}a_{𝐤^{}+𝐪/2,\lambda }a_{𝐤^{}+𝐪/2,\mu }.$$
(2)
As usual, we write equations of motion for each of these functions and use the ”anomalous mean field” recipe to decouple the four-operator products, $`a^{}a^{}aaa^{}a^{}aa`$. This brings out modified self-consistency relations for the order parameter,
$`\mathrm{\Delta }_{\alpha \beta }(𝐤,𝐪)={\displaystyle \underset{𝐤^{}}{}}V_{\beta \alpha ,\lambda \mu }(𝐤,𝐤^{})F_{\lambda \mu }(𝐤^{}+𝐪/2,\tau ;𝐤^{}𝐪/2,\tau ),`$
$`\mathrm{\Delta }_{\lambda \mu }^+(𝐤,𝐪)={\displaystyle \underset{𝐤^{}}{}}V_{\alpha \beta ,\mu \lambda }(𝐤^{},𝐤)F_{\alpha \beta }^+(𝐤^{}𝐪/2,\tau ;𝐤^{}+𝐪/2,\tau ),`$ (3)
and Gor’kov equations, which in the spatially uniform, stationary system read:
$`(i\omega _n\xi _k)G_{\alpha \beta }(𝐤,\omega _n)+\mathrm{\Delta }_{\alpha \gamma }F_{\gamma \beta }^+(𝐤,\omega _n)=\delta _{\alpha \beta };`$
$`(i\omega _n+\xi _k)F_{\alpha \beta }^+(𝐤,\omega _n)+\mathrm{\Delta }_{\alpha \gamma }^+G_{\gamma \beta }(𝐤,\omega _n)=0;`$ (4)
$`(i\omega _n\xi _k)F_{\alpha \beta }(𝐤,\omega _n)\mathrm{\Delta }_{\alpha \gamma }G_{\beta \gamma }(𝐤,\omega _n)=0.`$
Here $`\xi _k`$ is the fourier transform of the kinetic energy operator, $`\widehat{\xi }=(2m)^1\widehat{p}^2\mu `$, $`\mu `$ being the chemical potential.
It was established, that in high-T<sub>c</sub> cuprates, like YBCO, the order parameter is a spin singlet with $`d`$-wave orbital symmetry,
$$\mathrm{\Delta }_{\alpha \beta }(𝐤)=\delta _{\alpha \beta }\mathrm{\Delta }(𝐤),\mathrm{\Delta }(𝐤)\mathrm{cos}^2(k_x)\mathrm{cos}^2(k_y)$$
(5)
(with axes chosen along crystallographic directions (1,0,0), (0,1,0) in the cuprate layer)Sigrist95 ; VanHarlingen95 ; Tsuei00 . For the following, the most important consequence of this symmetry is the sign change of $`\mathrm{\Delta }(𝐤)`$ for certain directions. This means, first, that there exists an intrinsic phase shift of $`\pi `$ between different directions in the crystal; second, that in certain (nodal) directions the order parameter is zero, and therefore the quasiparticle excitation spectrum is not gapped.
On the spatial scale exceeding the coherence length, $`\xi _0`$, it is more convenient to use the Eilenberger equationsEilenberger68 , which follow from (I.1) in the quasiclassical limit. This is certainly justified in high-T<sub>c</sub> cuprates with their small $`\xi _0`$.
The Eilenberger equations are conveniently written in matrix form,
$`𝐯_F\widehat{G}(\omega _n)+[\omega _n\widehat{\tau }_3+\widehat{\mathrm{\Delta }},\widehat{G}(\omega _n)]=0.`$ (6)
Here the matrix Green’s function and order parameter,
$`\widehat{G}(𝐯_F,𝐫;\omega _n)=\left(\begin{array}{cc}g_{\omega _n}(𝐯_F,𝐫)& f_{\omega _n}(𝐯_F,𝐫)\\ f_{\omega _n}^+(𝐯_F,𝐫)& g_{\omega _n}(𝐯_F,𝐫)\end{array}\right);\widehat{\mathrm{\Delta }}(𝐯_F,𝐫;\omega _n)=\left(\begin{array}{cc}0& \mathrm{\Delta }(𝐯_F,𝐫)\\ \mathrm{\Delta }^+(𝐯_F,𝐫)& 0\end{array}\right),`$ (11)
depend both on position $`𝐫`$ and on (direction of) Fermi velocity $`𝐯_F`$. (The layered structure of high-T<sub>c</sub> cuprates allows us to reduce the problem to two dimensions to a good accuracy, unless we have to consider, for example, a twist junctionTafuri00 , or tunneling in the $`c`$-direction.) The components of $`\widehat{G}`$, obtained from Gor’kov’s functions by integration over energies, satisfy the normalization condition, $`g_{\omega _n}=\sqrt{1f_{\omega _n}^+f_{\omega _n}},`$ and the self-consistency relation (I.1) becomes
$$\mathrm{\Delta }(𝐯_F,𝐫)=2\pi N(0)T\underset{\omega _n>0}{}V(𝐯_F,𝐯_F^{})f_{\omega _n}(𝐯_F,𝐫)_\theta .$$
(12)
Here the angle averaging $`_\theta =_0^{2\pi }\frac{d\theta }{2\pi }`$.
In a little different language, the same results are obtained with the Andreev approximation in the Bogoliubov-de Gennes equations for the components of the single-bogolon wave function, $`(u(𝐫),v(𝐫))^{}`$. The original Bogoliubov-de Gennes equations are obtained in the process of diagonalization of pairing BCS Hamiltonian Zagoskin98 :
$`\left(\begin{array}{cc}\frac{1}{2m}^2\mu & \mathrm{\Delta }_k(𝐫)\\ \mathrm{\Delta }_k^{}(𝐫)& \frac{1}{2m}^2+\mu \end{array}\right)\left(\begin{array}{c}u_k(𝐫)\\ v_k(𝐫)\end{array}\right)\left(\begin{array}{cc}𝐯_F& \mathrm{\Delta }_k(𝐫)\\ \mathrm{\Delta }_k^{}(𝐫)& 𝐯_F\end{array}\right)\left(\begin{array}{c}u_k(𝐫)\\ v_k(𝐫)\end{array}\right)=E_k\left(\begin{array}{c}u_k(𝐫)\\ v_k(𝐫)\end{array}\right).`$ (23)
Here the quasimomentum $`𝐤`$ labels the bogolon state, $`E_k`$ is the excitation energy, and the self-consistency relation reads
$$\mathrm{\Delta }_k^{}(𝐫,T)=\underset{k^{}}{}V(𝐤,𝐤^{})u_k^{}^{}(𝐫)v_k^{}(𝐫)\mathrm{tanh}\frac{E_k^{}[\mathrm{\Delta }_k^{}(𝐫)]}{2T}.$$
(24)
Unlike the Gor’kov equations (I.1) (or the initial Bogoliubov-de Gennes equations), the equations (6,23) are of the first order in gradients, which allows us to introduce quasiclassical trajectories (characteristics) along $`𝐯_F`$ and solve the corresponding equations by integration along these trajectories, with proper boundary conditions.
It is known that a quasiparticle (electron or hole, described, in terms of Bogoliubov-de Gennes equations, by a vector $`\stackrel{}{\psi }_k=(1,0)^{}\mathrm{exp}(i\mathrm{𝐤𝐫})`$ ($`(0,1)^{}\mathrm{exp}(i\mathrm{𝐤𝐫})`$) respectively) impinging on the superconductor from the normal metal can undergo an Andreev reflection, switching the branch of the excitation spectrum, acquiring an additional phase, and almost exactly reversing the direction of its group velocity (this happens because for an electron and a hole with the same momentum $`𝐤`$ group velocities are opposite):
$$\stackrel{}{\psi }\widehat{}_A\stackrel{}{\psi },$$
(25)
where
$`\widehat{}_A=\left(\begin{array}{cc}0& e^{i\pi /2+i\chi }\\ e^{i\pi /2i\chi }& 0\end{array}\right).`$ (28)
The phase $`\chi `$ is the phase of the superconducting order parameter; the ($`\pi /2`$)-shift is exact in the limit when the quasiparticle energy is much less than the superconducting gap (generally it is some energy-dependent function $`\delta (E)`$).
Now consider a slab of normal conductor sandwiched between two superconductors (SNS junction). If we neglect the spatial dependence of the order parameter in superconductors, we don’t need to solve the self-consistency equations (I.1,24). Therefore the problem reduces to a single-particle one and is most naturally solved in Bogoliubov-de Gennes language (Eq.(23) becomes a Schrödinger equation for a two-component wave function).
Solutions of this equation with the boundary condition (28) are standing waves. Obviously in the act of Andreev reflection a charge of $`\pm 2e`$ is transferred to the superconductor, therefore every standing wave (Andreev level) carries supercurrent.
Quasiclassically, in order to find Andreev levels in a normal layer of thickness $`L`$, sandwiched between ”left” and ”right” superconductors, with phase difference $`\chi `$, we write the Bohr-Sommerfeld quantization condition,
$$p(E)𝑑q\pm \chi +\delta _l(E)+\delta _r(E)=2\pi n.$$
(29)
Here the kinematic phase gain of the quasiparticle along the closed trajectory, $`p(E)𝑑q=_l^rp_e(E)𝑑q+_r^lp_h(E)𝑑q=_l^r(p_e(E)p_h(E))𝑑q`$, takes into account electron-hole (or vice versa) conversion.
The positions of levels, and therefore the supercurrent, depend on the phase difference $`\chi `$ between the superconducting banks, and we arrive at Josephson effect in SNS structures Kulik70 ; Ishii70 ; Bardeen72 . Actually, the language of Andreev levels can be successfully used to describe the Josephson effect in general (for a review see Furusaki99 ).
### I.2 $`\pi `$-junctions and time-reversal symmetry breaking
The crucial experiments(Wollman93 ; Wollman95 ; Tsuei94 ; see also reviewTsuei00 ) which confirmed $`d`$-wave pairing symmetry in high-T<sub>c</sub> cuprates were directed at catching the intrinsic phase $`\pi `$-shift. The general idea of the experiment follows from the fluxoid quantization condition in a superconductor: if a superconducting contour $`𝒞`$ is penetrated by the magnetic flux $`\mathrm{\Phi }`$, then Tinkham96
$$2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}+_𝒞𝑑𝐬\varphi =2\pi n,\mathrm{\Phi }_0\frac{hc}{2e}$$
(30)
(in CGS units; $`\mathrm{\Phi }_0210^{15}`$Wb in SI).
In the case of a massive superconducting ring the contour can be chosen well inside the superconductor, where there is no current and therefore no superconducting phase gradient. Then the magnetic flux is quantized in units of $`\mathrm{\Phi }_0`$.
If there is a Josephson junction in the ring, the phase change will concentrate there, yielding
$$2\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}+\chi =2\pi n.$$
(31)
Here $`\chi `$ is the phase difference across the junction.
The equilibrium value of $`\chi `$ is determined by the interplay between Josephson and magnetic energy of the system. The Josephson energy is related to the Josephson current via
$$I(\chi )=2e\frac{E(\chi )}{\chi },$$
(32)
and in the simplest case of a tunneling junction, $`I(\chi )=I_c\mathrm{sin}\chi `$, and $`E(\chi )=(I_c/2e)\mathrm{cos}\chi .`$ For a conventional Josephson junction $`E(\chi )`$ is at a minimum when $`\chi =0`$; therefore in the absence of an external field the total energy of the ring,
$`U(\chi )=E(\chi )+E_{\mathrm{magn}.}(\mathrm{\Phi })={\displaystyle \frac{I_c}{2e}}\mathrm{cos}\chi +\left({\displaystyle \frac{\chi }{2\pi }}\right)^2{\displaystyle \frac{\mathrm{\Phi }_0^2}{2L}},`$ (33)
has a single global minimum at $`\chi =0`$ for any value of $`2\pi LI_c/\mathrm{\Phi }_0`$ (here $`L`$ is the self-inductance of the ring).
Not so if the pairing symmetry is $`d`$-wave: if we choose the configuration of the loop in such a way, that the opposite sign lobes contact across the junction (so called $`\pi `$-junction), the current-phase relation switches to $`I(\chi )=I_c\mathrm{sin}(\chi +\pi )=I_c\mathrm{sin}\chi `$, and $`E(\chi )=+(I_c/2e)\mathrm{cos}\chi .`$ Therefore if $`2\pi LI_c/\mathrm{\Phi }_0>1,`$ the system has two degenerate minima; if this ratio is so big that the magnetic energy can be neglected compared to the Josephson energy, the equilibrium phase difference is $`\chi _0=\pi `$, and we obtain from (31), that
$$\mathrm{\Phi }=\left(n+\frac{1}{2}\right)\mathrm{\Phi }_0.$$
(34)
Thus in equilibrium there is a spontaneous flux $`\mathrm{\Phi }_0/2`$ in the ring. Its direction (up or down) allows us to distinguish the two ground states of the system, where the *time-reversal symmetry is thus broken.*
Such behaviour was indeed observed in an experiment Tsuei94 : a tri-crystal ring of YBCO generated a half-flux quantum, which was detected by SQUID microscopy.
A natural question is whether $`\pi `$-junctions are the only possibility provided by $`d`$-wave symmetry, and whether only $`\mathrm{\Phi }_0/2`$ fluxes can be spontaneously generated. The answer is no: in principle, any equilibrium phase difference can be realized in a $`d`$-wave junction, and time-reversal symmetry breaking can be accompanied by generation of an arbitrary magnetic flux, or none at all.
### I.3 Josephson effect and T-breaking in SND and DND junctions
Let us return to an SNS junction (assuming a rectangular normal part, $`L\times W`$). Each quasiclassical trajectory connecting two superconductors acts as a conduit of supercurrent between them, an ”Andreev tube” of diameter $`\lambda _F`$, which carries supercurrent, determined by the phase difference between its ends. If we neglect the normal scattering at NS interfaces, there is no mixing between different trajectories, since Andreev reflections simply reverse the velocity, sending the reflected particle along the same path. It can be shown, that if $`L\xi _0`$, every such trajectory passing through the point $`𝐫`$ contributes a partial supercurrent densityZagoskin98 ; Barzykin99
$$𝐣(𝐫,𝐯_F)=\frac{2e𝐯_F}{\lambda _FW}\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n+1}\frac{(𝐯_F)}{l_T}\frac{\mathrm{sin}(n(\varphi _1\varphi _2))}{\mathrm{sinh}(n(𝐯_F)/l_T)}e^{2n(𝐯_F)/v_F\tau _i}.$$
(35)
In Eq.(35), $``$($`𝐯_F`$) is the length of the trajectory, and $`l_T=v_F/2\pi k_BT`$ is the so called normal metal coherence length. (We also included effects of weak impurity scattering with scattering time $`\tau _i`$.) The physical meaning of $`l_T`$ is that in a clean normal metal an electron and Andreev-reflected hole (or vice versa) with energy $`k_BT`$ maintain phase coherence across distance $`l_T`$ simply because they travel along the same trajectory. Indeed, the momentum of an electron (hole) with energy $`k_BTE_F`$ is $`p_{e,h}p_F\pm (_Ep_F)k_BT=p_F\pm k_BT/v_F.`$ At a distance $`l`$ from the point of Andreev reflection (measured along the trajectory) they would gain phase difference $`(p_ep_h)l=2k_BTl/v_F.`$ If the phase difference is of order $`\pi `$, the coherence is effectively lost, so we get $`l_{T,\mathrm{ballistic}}v_F/k_BT`$. The factor of $`2\pi `$ appears in accurate treatment, like in Eq.(35).
As an aside, in case of very strong scattering in the normal part of the system, when the motion of electrons/holes is diffusive with diffusion coefficient $`D`$, we can still use the same argument. Now the observable length scale is given by the displacement of quasiparticle, $`l^2=Dt,`$ while the phase difference between the electron and the hole is gained along the crooked path of length $`l^{}=v_Ft=v_Fl^2/D`$ they take. So the condition $`2k_BTl^{}/v_F\pi `$ yields $`l_{T,\mathrm{diffusive}}\sqrt{D/k_BT}.`$
Such coherence in a normal metal is a purely kinematic effect, since there is no interaction in the normal (non-magnetic) metal, which would either support or suppress superconductivity. (”Normal metal is neutral with respect to superconducting correlations.”-C.W.J. Beenakker.) Nevertheless its effects are quite real, e.g. the very possibility of coherent supercurrent flow (Josephson effect) in SNS junctions.
Let us apply the approach of (35) to calculation of Josephson current in a clean SNS contact. In every point of NS boundary (e.g. $`x=L`$), we integrate (35) over directions of $`𝐯_F`$ (such that $`v_{Fx}>0`$). As it should be expected, the result reproduces Ishii’s sawtooth Ishii70 : in the limit of zero temperature and no scattering inside the normal layer ($`\chi \varphi _1\varphi _2`$),
$$I(\chi )=𝑑y\underset{v_{Fx}>0}{}j_x(y,𝐯_F)\frac{2}{\pi }\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n+1}\frac{\mathrm{sin}(n\chi )}{n}=\frac{\chi }{\pi },|\chi |<\pi $$
(36)
(periodically extended).
Consider now an SND junction, that is, an SNS contact, where one of the superconductors has $`d`$-wave symmetry Zagoskin97 ; Zagoskin98 . Now, when integrating over the directions of $`𝐯_F`$ we encounter two kinds of trajectories: ”zero”- and ”$`\pi `$”-trajectories, which link the conventional superconductor with the lobes of the $`d`$-wave order parameter of opposite sign (intrinsic phase $`\pi `$-shift). The Josephson current is therefore a sum of two contributions like (36), one of which has the phase argument $`\chi +\pi `$. The relative weight of these contributions depends on the orientation of the SD boundary with respect to the crystal axes of the cuprate (to which the order parameter is nailed). In the case of a $`45^{}`$ orientation, when the boundary is along the (110) plane and is therefore normal to the nodal direction of the $`d`$-wave order parameter, the two groups contribute equally, yielding ($`\mathrm{cos}\theta =v_{Fx}/v_F`$)
$$I(\chi )=\underset{v_{Fx}>0}{\overset{\mathrm{zero}\mathrm{levels}}{}}\frac{ev_F\mathrm{cos}\theta }{L}\frac{2}{\pi }\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n+1}e^{2nL/l_i\mathrm{cos}\theta }\frac{L}{l_T\mathrm{cos}\theta }\frac{\mathrm{sin}(n\chi )+\mathrm{sin}(n(\chi +\pi ))}{\mathrm{sinh}(nL/l_T\mathrm{cos}\theta )}.$$
(37)
All odd harmonics cancel, and we obtain a $`\pi `$-periodic sawtooth: at $`T=0`$, $`l_i\mathrm{}`$
$$I(\chi )\frac{\chi \pi /2}{\pi /2},\mathrm{\hspace{0.25em}\hspace{0.25em}0}<\chi <\pi ,$$
(38)
(periodically extended). (This is simply a sum of two identical, $`2\pi `$-periodic sawtooth functions, shifted by $`\pi `$, see Fig.2.) There are two stable equilibrium phase differences across the junction: $`\chi =\pm \pi /2.`$ This means, that the time-reversal symmetry is broken. (In addition, the frequency of ac Josephson effect in the system will be doubled Zagoskin97 ; Hurd99 ).
The SND junction may seem simply a ”superposition” of two SNS junctions, with phase differences shifted by $`\pi `$, but the situation is more interesting. In equilibrium, spontaneous currents flow in the normal layer, parallel to the boundaryHuck97 . This is clear from our ”Andreev tube” treatment. Take, for example, two trajectories, with $`\theta =+\alpha `$ and $`\theta =\alpha `$, carrying equal partial currents. Their contribution to the Josephson current (normal to the boundary) is zero, since their currents’ projections on this direction cancel each other. On the contrary, the projections on the direction, parallel to the boundary, add. The phase dependence of this current, $`I_s(\chi )`$, can be obtained in the same way as for the Josephson current. In the limit of zero temperature and clean normal layer it reduces to a difference of two sawtooth functions, yielding
$$I_s(\chi )\mathrm{sgn}\chi ,|\chi |<\pi ,$$
(39)
again periodically extended. States with $`\pi /2`$ and $`\pi /2`$ thus carry spontaneous currents in opposite directions. If an SND junction is closed on itself (annular geometry), these currents translate into spontaneous magnetic fluxes, normal to the plane of the system Zagoskin98a .
Note that all of the above is only possible if higher harmonics of current-phase dependence are not all zero, since all the odd harmonics (including the standard Josephson term, $`\mathrm{sin}\chi `$) cancel. (The expected ”total depairing” in the (110) plane.)
In the case of arbitrary orientation of the ND plane , the weights of zero- and $`\pi `$-sawtooth functions will be different (Fig.2c). Then the current-phase dependence regains $`2\pi `$-periodicity. At zero temperature and in the absence of scattering, the T-breaking is generally still present. The equilibrium phase difference is no longer $`\pm \pi /2`$ and depends on the geometry. (The exception is when the boundary is along (100) or (010). Then the contribution of only one group survives, and we have either standard, or $`\pi `$-junction. Of course, it is impossible to tell, whether a single junction is a $`\pi `$-junction - it is necessary to look at the contour, in which it is includedVanHarlingen95 .) Finite temperature and scattering suppress higher harmonics first, and therefore the transition to non-T-breaking state is possible Huck97 .
Essentially the same interplay of zero- and $`\pi `$-levels takes place in DND junctions. There we have a somewhat richer picture. For example, the time-reversal symmetry can be broken without producing spontaneous currents (Fig.3c).
One of the reliable methods of fabrication of Josephson devices in high-T<sub>c</sub> cuprates is based on forming grain boundary junctions (see e.g. Tafuri99 ). The order parameter is suppressed within $`\xi _0`$ around the boundary, and this region can be considered ”normal”. Therefore the SND/DND model applies, but only qualitatively, since the above equations are derived in the limit $`L\xi _0`$. A more accurate approach, based on Eilenberger equations (6), confirms the qualitative similarity between DND and DD junctionsAmin01 .
Integrating the current-phase dependences of Fig.2, we obtain the Josephson energy of the junction, Eq.(32), which has two minima, corresponding to its degenerate ground states. This bistability plays the crucial role in qubit applications of high-T<sub>c</sub> superconductors. This prediction was confirmed in grain boundary YBCO junctionsIlichev01 (Fig.6).
So far we did not take into account normal scattering on NS interfaces, and therefore missed an important point. Let us first consider a boundary of a $`d`$-wave superconductor with a vacuum or an insulator. The order parameter near the surface is suppressed, and a qualitative understanding of the situation can be obtained by ”inserting” a normal layer of thickness $`\xi _0`$ between the insulator and the bulk superconductor, similar to the DND modelLofwander . Energy levels in such a layer can be found from (29) for every quasiparticle trajectory (assuming specular reflection), Fig.7a:
$$p(E)𝑑q+2\delta (E)+\pi s=2\pi n.$$
(40)
Here $`s=0`$ or $`s=1`$ depending on whether the trajectory connects the lobes of the same or opposite sign. Note that $`p(0)𝑑q=0`$ and $`2\delta (0)=\pi `$. Therefore solution $`E=0`$ exists if and only if $`s=1`$. This is the zero-energy, or midgap, state (ZES, MGS), which obviously cannot exist in conventional superconductors, where $`s`$ always equalis zero (see review Lofwander ).
Now consider a DD junction with finite transparency, $`𝒟<1`$ (instead of the case of ideal transparency, $`𝒟=1`$, which we dealt with earlier). We can use the DND model, inserting in the middle of the normal layer an infinitely thin barrier with transparency $`𝒟(\theta )`$, dependent on the incidence angle of the quasiparticle trajectoryLofwander .
Bohr-Sommerfeld quantization conditions for quasiparticle trajectories shown in Fig.7b yield the energy levels of the bound states in the junction. In the limit of low transparency, $`𝒟1`$, the critical current is much larger ($`O(\sqrt{𝒟})`$, not $`O(𝒟)`$), if the orientations of the $`d`$-wave order parameter allow formation of zero energy states on both sides of the barrier.
Even when transparency is not small, the presence of ZES may be qualitatively important. For example, the junction can become a $`\pi `$-junction at low enough temperature (see review Lofwander ); you can see such behaviour in Fig.6b, where below 11 K the potential wells are centered around $`\pi `$ rather than zero.
The appearance of spontaneous currents in SD and DD junctions can be considered as due to a time-reversal symmetry breaking order parameter with $`s+e^{i\chi _0}d`$ or $`d+e^{i\chi _0}d^{}`$ symmetry, which is formed in the junction area due to the proximity effect (here $`\chi _00,\pi `$ is the equilibrium phase difference across such junction). In certain conditions, such combinations could appear near the surface of a $`d`$-wave superconductor, leading to spontaneous currents and magnetic moments. So far there is no conclusive evidence for such currents (see Amin02 and references therein).
## II Mesosopic $`d`$-wave qubits
### II.1 Flux qubits with conventional superconductors
Qubits are the basic building blocks of future quantum computers. Essentially they are two-state quantum systems which can be put in an arbitrary superposition of states (”initialized”), coupled to each other, undergo desired quantum evolution and measured (”read out”) before losing quantum coherence. Here we concentrate on superconducting phase qubits.
The simplest example of such a qubit is an RF SQUID, that is, a loop with a single Josephson junction (like in Section I.2). We saw, that in certain circumstances the system has two degenerate minima, corresponding to a flux of $`\pm \mathrm{\Phi }_0/2`$ through the loop.
The Hamiltonian of such a system is
$$H=U_J(\chi )+U_C(\widehat{Q}),$$
(41)
where $`U_C`$ is the Coulomb energy of charge $`Q`$ on the junction (which has some finite capacitance $`C`$). The charge operator $`\widehat{Q}`$ can be expressed in terms of the phase difference across the junction as $`\widehat{Q}=i_\chi `$ (see e.g. Tinkham96 ; Zagoskin98 ). Due to the presence of the Coulomb term, phase is no longer a ”good” quantum variable, and the eigenstates of the Hamiltonian (41) become linear combinations of ”up” and ”down”, or ”left” and ”right”, states (with spontaneous flux $`\pm \mathrm{\Phi }_0/2`$). In other words, the qubit can now tunnel between the wells of the Josephson potential, corresponding to certain phases (Fig.6). (This description is appropriate as long as the Coulomb energy does not exceed the Josephson energy, otherwise the natural starting point would be the states with definite charge on the junction; we do not consider such systems (charge qubits) here, but they were successfully implemented experimentallyNakamura:1999 .)
Coherent quantum tunneling was indeed observed in an RF SQUIDFreidman:2000 . Simultaneously, this effect was obtained in a different system, consisting of a small inductance loop with three Josephson junctions (the so called persistent current qubit, vdWal:2000 ).
The advantage of the latter design is as follows. As we have seen in Section I.2, the degenerate states appear in the RF SQUID only if the self-inductance of the loop is large enough, and they carry comparatively large spontaneous fluxes, $`\pm \mathrm{\Phi }_0/2`$. Therefore they will couple to the external degrees of freedom and reduce the time $`\tau _d`$ during which the system maintain its quantum coherence. Even more important is the fact that the resulting potential barrier is comparatively high, which may prohibit the tunneling we are after. In an experiment Freidman:2000 the coherent tunneling was indeed observed not between the lowest, but between the excited states in the wells.
In case of the three-junction loop of negligible self-inductance, the fluxoid quantization condition (31) leaves two independent Josephson phase differences in the circuit:
$$\chi _1+\chi _2+\chi _2+2\pi \frac{\mathrm{\Phi }_{\mathrm{external}}}{\mathrm{\Phi }_0}=2\pi n.$$
(42)
The Josephson energy,
$$U_J=E_{J1}\mathrm{cos}(\chi _1)E_{J2}\mathrm{cos}(\chi _2)E_{J3}\mathrm{cos}(\chi _3)$$
(43)
of the system forms a 2D potential profile (e.g. as a function of $`\chi _1`$, $`\chi _2`$), which depends on the external flux $`\mathrm{\Phi }_{\mathrm{external}}f\times \mathrm{\Phi }_0`$ as a parameter.
If $`f=0.5`$, and the $`E_J`$’s are comparable, this potential has two degenerate minima; unlike the case of the RF SQUID, the potential barrier can be small, as are the spontaneous fluxes corresponding to the ”left” and ”right” states.
### II.2 Rationale and proposed designs for qubits with $`d`$-wave superconductors
One of the main problems with the designs of the previous section is the necessity to artificially break the T-symmetry of the system by putting a half flux quantum through it. Estimates show Blatter01 that the required relative accuracy is $`10^510^6`$. The micron-size qubits must be positioned close enough to each other to make possible their coupling; the dispersion of qubit parameters means that applied fields must be locally calibrated; this is a formidable task given such sources of field fluctuations as fields generated by persistent currents in qubits themselves, which depend on the state of the qubit; field creep in the shielding; captured fluxes and magnetic impurities. Moreover, the circuitry which produces and tunes the bias fields is an additional source of decoherence in the system.
These problems are avoided if the qubit is intrinsically bistable. The most straightforward way to achieve this is to substitute the external flux by a static phase shifter, a Josephson junction with unconventional superconductors with nonzero equilibrium phase shift $`\chi _0`$. For example, three-junction (persistent current, Mooij) qubits would require an extra $`\pi `$-junction ($`\chi _0=\pi `$) Blatter01 . The only difference compared to the case of an external magnetic field bias is in the prevalent decoherence sources: instead of noise from field-generating circuits we will have to take into account intrinsic decoherence from nodal quasiparticles and interface bound states (see below).
Another suggestion Newns02 is based on the same tricrystal high-T<sub>c</sub> ring geometry as in exp erimentsTsuei94 ; Tsuei00 (Fig8). The spontaneous flux $`\pm \mathrm{\Phi }_0`$ generated by such a structure labels the qubit states $`|0`$, $`|1`$. Tunnelling between them, necessary for quantum operations, is made possible by applying a magnetic field in the plane of the system. Indeed, then the states $`|0`$, $`|1`$ are no longer the eigenstates of the Hamiltonian.
The tricrystal ring D is surrounded by an s-wave superconductor ring S, aimed at screening the spontaneous flux from the environment (including other qubits). Indeed, due to the fluxoid quantization condition (30) in the ring S, the total flux through it must be an integer multiple of $`\mathrm{\Phi }_0`$, and therefore the states of the rings D and S become entangled (e.g. $`\alpha |0_D|1_S+\beta |1_D|0_S`$).
This entanglement puts forward an interesting problem. In order to perform two-qubit operations, as well as initialization and readout, it is necessary to make the qubits ”visible” to each other and the outside world. To do this, it is suggested in Newns02 to locally destroy superconductivity in the screening ring by using a superconducting field effect transistor (SUFET) (not shown) (i.e. by applying a gate voltage to the part of the screening ring). Will this transition collapse the wave function of the qubit, destroying quantum coherence between $`|0_D`$ and $`|1_D`$?
Now let us consider the case when the bistable $`d`$-wave system is employed $`dynamically`$, that is, when its phase is allowed to tunnel between the degenerate values. In a so called ”quiet” qubit Ioffe99 an SDS’ junction (effectively two SD junctions in the (110) direction) put in a small-inductance SQUID loop in parallel with a conventional Josephson junction and a large capacitor, Fig.9. One of the SD junctions plays the role of a $`\pi /2`$-phase shifter. The other junction’s capacitance $`C`$ is small enough to make possible tunneling between the $`\pi /2`$ and $`\pi /2`$-states due to the Coulomb energy $`Q^2/2C`$. Two consecutive SD junctions are effectively a single junction with equilibrium phases $`0`$ and $`\pi `$ (which are chosen as working states of the qubit, $`|0`$ and $`|1`$). The proposed control mechanisms are based on switches $`c,s`$. Switch $`c`$ connects the small S’D junction to a large capacitor, thus suppressing the tunneling. Connecting $`s`$ for the duration $`\mathrm{\Delta }t`$ creates an energy difference $`\mathrm{\Delta }E`$ between $`|0`$ and $`|1`$, because in the latter case we have a frustrated SQUID with $`0`$\- and $`\pi `$\- junctions, which generates $`\mathrm{\Phi }_0/2`$ spontaneous flux. This is a generalization of applying a $`\sigma _z`$ operation to the qubit. Finally, if switch $`c`$ is open, the phase of the small junction can tunnel between $`0`$ and $`\pi `$. Entanglement between different qubits is realized by connecting them through another Josephson junction in a bigger SQUID loop.(The switches in question are in no way trivial, since they must operate without destroying quantum coherence of the system. One possible solution is to use a frustrated dc-SQUIDIoffe99 , that is, insert in the wire a small inductance loop with two equivalent Josephson junctions in parallel. The total supercurrent through the switch, $`I=I_1(\chi _1)+I_2(\chi _2)`$, goes to zero if the phase difference (tunable by external magnetic flux) $`\chi _1\chi _2=\pi `$. Modifications of this scheme are discussed in Blatter01 . Another possibility is to use superconducting single-electron transistors (SSETs, ”parity keys”)Zagoskin99 .)
Due to the absence of currents through the loop during tunneling between $`|0`$ and $`|1`$ the authors called it ”quiet”, though as we have seen small currents and fluxes are still generated near SD boundaries.
Another design based on the same bistability Zagoskin99 ; Blais00 only requires one SD or DD boundary (Fig.10). Here a small island contacts a massive superconductor, and the angle between the orientation of the $`d`$-wave order parameter and the direction of the boundary can be arbitrary (as long as it is compatible with bistability Amin01 ; Amin02b , Fig.5). The advantage of such a design is, that the potential barrier can be to a certain extent controlled and suppressed; moreover, in general the two ”working” minima ($`\varphi _0,\varphi _0`$; the phase of the bulk superconductor across the boundary is zero) will be separated from each other by a smaller barrier, than from the equivalent states differing by $`2\pi n`$. This allows us to disregard the ”leakage” of the qubit state from the working space spanned by ($`|0`$,$`|1`$), which cannot be done in the ”quiet” design with the exact $`\pi `$-periodicity of the potential profile.
Qubit operations in this system are realized by connecting qubits to each other and to the ground electrode (normal or superconducting) through SSETs (or other kind of switches). When isolated, a qubit undergoes natural evolution between $`|0`$ (state with phase $`\varphi _0`$) and $`|1`$ (phase $`\varphi _0`$), which realizes the $`\sigma _x`$ operation. The $`\sigma _z`$ operation (that is, adding a controllable phase shift to one of the states with respect to the other) can be realized by e.g. connecting the island through a SSET to the massive superconductor (”bus”), the phase of which $`\chi 0`$. The same operation repeated periodically can be used to block unwanted tunneling (so called ”bang-bang” technique) Blais00 : physically, if we keep shifting the levels in the right and left well with respect to each other, the tunneling becomes suppressed, since they are practically never in resonance.
A better design was suggested in Amin03 (”silent qubit”). Here two small bistable $`d`$-wave grain boundary junctions with a small superconducting island between them are set in a SQUID loop (Fig. 11). As usual, ”small” means that the total capacitance of the system allows phase tunneling: the Coulomb energy term in (41) is not negligible.
In the limit of negligible self-inductance of the loop, the quantization condition (30) fixes the sum of the phases to the external flux, $`\chi _1+\chi _2\varphi =2\pi \mathrm{\Phi }/\mathrm{\Phi }_0`$. This leaves only one independent phase combination, the superconducting phase of the island, $`\theta =(\chi _1\chi _2)/2`$.
Keeping for simplicity only the first two harmonics in the current-phase relation of the junctions,
$`I_i=I_i^{(I)}\mathrm{sin}\chi _1I_i^{(II)}\mathrm{sin}2\chi `$
$`I_{0i}\left[\mathrm{sin}\chi _i\mathrm{sin}\gamma _i\mathrm{sin}2\chi _i\mathrm{cos}\gamma _i\right],`$ (44)
we find for the Josephson potential of the qubit the expression (Fig. 12)
$`U_J(\theta ,\varphi )=(I_{01}/2e)\left[f(\varphi /2+\theta ,\gamma _1)+\eta f(\varphi _2\theta ,\gamma _2)\right].`$ (45)
Here $`f(\phi ,\gamma )=\mathrm{cos}(\phi )\mathrm{sin}(\gamma )(1/2)cos(2\phi )\mathrm{cos}(\gamma )`$, and $`\eta =I_{02}/I_{01}`$. Parameter $`\gamma [0,\pi /2]`$ provides a convenient parametrization for the current-phase dependence in a DD junction.
In the absence of the external flux ($`\varphi =0`$) the qubit potential $`U_J(\theta ,0)`$ has two degenerate minima. Moreover, if the junctions only differ in the amplitude of critical current, but have the same $`\gamma `$, there is no spontaneous current in the loop in either minimum, which means that the qubit is decoupled from the external magnetic fields (if we disregard spontaneous currents in the junctions themselves, which can be very small Amin01 ; Amin02b ; Amin02 ). This justifies the moniker ”silent”.
Both the barrier and the bias between the wells can be controlled by the external flux. It is noteworthy that the corrections are of at least second order in $`\varphi `$, which drastically reduces the influence of fluctuations in the external circuitryAmin03 . The mechanism of noise reduction is similar to that of the ”quantronium” qubitVion , but the ”sweet spot” (an extremal point on the energy surface) appears already on the classical level, and at zero external field.
Finally, let us briefly mention two more proposals.
A ”no tunneling” designAmin03a combines the ideas of CBJJ (current-based Josephson junction) qubits Yu02 ; Martinis02 and the intrinsic bistability of d-wave junctions. In a single bistable Josephson junction, the potential barrier between the degenerate levels corresponding to phase difference $`\pm \chi _0`$ is too high to allow tunneling. The transitions between the states are realized through Rabi transitions via an auxiliary energy level, situated above the top of the barrier. Rabi transitions are induced by applying an external high-frequency field.
”Dot/antidot” proposals are based on the spontaneous flux (less than a flux quantum) generated in a high-T<sub>c</sub> island or around a hole in bulk high-T<sub>c</sub> due to the presence of a subdominant order parameterAmin02 ; Zagoskin02 (see the end of Section I.3), but so far lack experimental support.
### II.3 Fabrication and experiment
Due to fabrication difficulties, as well as expected problems with decoherence from e.g. nodal quasiparticles (see Section II.4), experimental research on $`d`$-wave qubits is not as far advanced as on the devices with conventional superconductors. The experimental confirmation of quantum behaviour in these systems is still missing. Nevertheless several recent successes should be noted.
Arrays of half-flux quanta were realized and manipulated (classically) in YBCO-Nb zigzag junctions Hilgenkamp03 . Each facet of the junction effectively constitutes a $`\pi `$-ring, which supports spontaneous flux of $`\pm \mathrm{\Phi }_0`$. Interaction between the fluxes is mainly due to the superconducting connection, which leads to robust antiferromagnetic ordering. In its absence, a weaker, magnetic interaction establishes ferromagnetic flux ordering. The authors consider the possibility of using their structures in qubit design.
Good quality submicron grain boundary YBCO junctions were fabricated and bistable energy vs. phase dependence was demonstrated Ilichev01 ; Tzalenchuk03 . Dc SQUIDs YBCO (15$`\times `$15 $`\mu \mathrm{m}^2`$ square loops with nominally $`2\mu `$m wide grain boundary junctions) were fabricated and tested, and their classical behaviour is very well described by the existing theory Lindstrom03 (Fig. 13). Like in (44), only two harmonics of current-phase dependence were considered. From fitting the experimental data, we had to conclude that the junctions in the same SQUID have not only different critical current amplitudes, but different ratios of first to second harmonic ($`\gamma _1\gamma _2`$), probably due to the variation in the grain boundary properties over a distance of $``$ 15 $`\mu `$m. This is one reason for putting the junctions in the silent qubit, Fig. 11, close to each other, as it was done when fabricating its prototypeAmin03 .
### II.4 Decoherence in $`d`$-wave qubits
Decoherence is the major concern for any qubit realization, especially for solid state qubits, due to abundance of low-energy degrees of freedom. In superconductors, this problem is mitigated by the exclusion of quasiparticle excitations due to the superconducting gap. This explains also why the very fact of existence of gapless excitations in high-T<sub>c</sub> superconductors long served as a deterrent against serious search for macroscopic quantum coherence in these systems. An additional source of trouble may be zero-energy states in DD junctions.
Nevertheless, recent theoretical analysis of DD junctions Fominov03a ; Amin03b ; Fominov03 , all using quasiclassical Eilenberger equations, shows that the detrimental role of nodal quasiparticles and ZES could be exaggerated.
Before turning to these results, let us first do a simple estimate of dissipation due to nodal quasiparticles in bulk $`d`$-wave superconductors.
Consider, for example, a three-junction (”Mooij”) qubit with d-wave phase shifters. The $`|0`$ and $`|1`$ states support, respectively, clockwise and counterclockwise persistent currents around the loop, with superfluid velocity $`𝐯_s`$. Tunnelling between these states leads to nonzero average $`\dot{𝐯}_s^2`$ in the bulk of the superconducting loop.
Time-dependent superfluid velocity produces a local electric field
$$𝐄=\frac{1}{c}\dot{𝐀}=\frac{m}{e}\dot{𝐯}_s,$$
(46)
and quasiparticle current $`𝐣_{qp}=\sigma 𝐄.`$ The resulting average energy dissipation rate per unit volume is
$$\dot{}=\sigma E^2m\tau _{qp}\overline{n}(v_s)\dot{v}_s^2.$$
(47)
Here $`\tau _{qp}`$ is the quasiparticle lifetime, and
$$\overline{n}(v_s)=_0^{\mathrm{}}𝑑ϵ\overline{N}(ϵ)\left[n_F(ϵp_Fv_s)+n_F(ϵ+p_Fv_s)\right]$$
(48)
is the effective quasiparticle density. The angle-averaged density of states inside the d-wave gap is Xu95
$$\overline{N}(ϵ)N(0)\frac{2ϵ}{\mu \mathrm{\Delta }_0},$$
(49)
where $`\mu =\frac{1}{\mathrm{\Delta }_0}\frac{d|\mathrm{\Delta }(\theta )|}{d\theta }`$, and $`\mathrm{\Delta }_0`$ is the maximal value of the superconducting order parameter. Substituting (49) in (48), we obtain
$$\overline{n}(v_s)N(0)\frac{2}{\mu \mathrm{\Delta }_0}\left(T^2\right)\left(\mathrm{Li}_2\left(e^{\frac{p_Fv_s}{T}}\right)+\mathrm{Li}_2\left(e^{\frac{p_Fv_s}{T}}\right)\right),$$
(50)
where Li$`{}_{2}{}^{}(z)=_z^0dt\frac{\mathrm{ln}(1t)}{t}`$ is the dilogarithm. Expanding for small $`p_Fv_sT`$, and taking into account that Li$`{}_{2}{}^{}(1)=\pi ^2/12`$, Li$`{}_{2}{}^{}(1)=\mathrm{ln}2`$, and Li$`{}_{2}{}^{\prime \prime }(1)=1/2+\mathrm{ln}2,`$ we obtain
$$\overline{n}(v_s)\frac{N(0)}{\mu \mathrm{\Delta }_0}\left(\frac{\pi ^2T^2}{3}+\left(p_Fv_s\right)^2\right).$$
(51)
The two terms in parentheses correspond to thermal activation of quasiparticles and their ”Cherenkov” generation by current-carrying state. Note that finite quasiparticle density does not lead by itself to any dissipation.
In the opposite limit ($`Tp_Fv_s`$) only the ”Cherenkov” contribution remains,
$$\overline{n}(v_s)\frac{N(0)}{\mu \mathrm{\Delta }_0}\left(p_Fv_s\right)^2$$
(52)
(since Li$`{}_{2}{}^{}(z)(1/2)(\mathrm{ln}z)^2`$ for large negative $`z`$, and $`z`$ for small $`z`$).
The energy dissipation rate provides the upper limit $`\tau _ϵ`$ for the decoherence time (since dissipation is sufficient, but not necessary condition for decoherence). Denoting by $`I_c`$ the amplitude of the persistent current in the loop, by $`L`$ the inductance of the loop, and by $`\mathrm{\Omega }`$ the effective volume of d-wave superconductor, where it flows, we can write
$$\tau _ϵ^1=\frac{2\dot{}\mathrm{\Omega }}{LI_c^2}\frac{2m\tau _{qp}N(0)\mathrm{\Omega }\left(\frac{\pi ^2T^2}{3}\dot{v}_s^2+p_F^2v_s^2\dot{v}_s^2\right)}{\mu \mathrm{\Delta }_0LI_c^2}.$$
(53)
Note that the thermal contribution to $`\tau _ϵ^1`$ is independent on the absolute value of the supercurrent in the loop ($`v_s`$), while the other term scales as $`I_c^2`$. Both contributions are proportional to $`\mathrm{\Omega }`$ and (via $`\dot{v}_s`$) to $`\omega _t`$, the characteristic frequency of current oscillations (i.e. tunneling rate between clockwise and counterclockwise current states).
It follows from the above analysis that the intrinsic decoherence in a $`d`$-wave superconductor due to nodal quasiparticles can be minimized by decreasing the amplitude of the supercurrent through it, and the volume of the material where time-dependent supercurrents flow. From this point of view the designs with phase shifters, as well as the Newns-Tsuei and ”no tunneling” designs are at a disadvantage (the latter, because the microwave field necessary to produce Rabi transitions will affect the whole sample).
Now let us estimate dissipation in a DD junction. First, following Zagoskin98a ; Zagoskin99 , consider a DND model with ideally transmissive ND boundaries. Due to tunneling, the phase will fluctuate, creating a finite voltage on the junction, $`V=(1/2e)\dot{\chi }`$, and normal current $`I_n=GV`$. The corresponding dissipative function and decay decrement are
$`={\displaystyle \frac{1}{2}}\dot{}={\displaystyle \frac{1}{2}}GV^2={\displaystyle \frac{G\dot{\chi }^2}{2}}\left({\displaystyle \frac{1}{2e}}\right)^2;`$ (54)
$`\gamma ={\displaystyle \frac{2}{M_Q\dot{\chi }}}{\displaystyle \frac{}{\dot{\chi }}}={\displaystyle \frac{G}{4e^2M_Q}}={\displaystyle \frac{4N_{}E_Q}{\pi }}.`$ (55)
Here $`E_Q=e^2/2C`$, $`M_Q=C/16e^2=1/32E_Q`$, $`N_{}`$ are the Coulomb energy, effective ”mass” and number of quantum channels in the junction respectively. The latter is related to the critical Josephson current $`I_0`$ and spacing between Andreev levels in the normal part of the system $`\overline{ϵ}=v_F/2L`$ via
$$I_0=N_{}e\overline{ϵ}.$$
(56)
We require, that $`\gamma /\omega _01,`$ where $`\omega _0=\sqrt{32N_{}E_Q\overline{ϵ}}/\pi `$ is the frequency of small phase oscillations near a local minimum. This means,
$$N_{}\frac{\overline{ϵ}}{E_Q}.$$
(57)
The above condition allows a straightforward physical interpretation. In the absence of thermal excitations, the only dissipation mechanism in the normal part of the system is through the transitions between Andreev levels, induced by fluctuation voltage. These transitions become possible, if $`\overline{ϵ}<2e\overline{V}\sqrt{\overline{\dot{\chi }^2}}\omega _0,`$ which brings us back to (57). Another interpretation of this criterion arises if we rewrite it as $`\omega _0^1(v_F/L)^1`$ Zagoskin99 . On the right-hand side we see time for a quasiparticle to traverse the normal part of the junction. If it exceeds the period of phase oscillations (on the left-hand side), Andreev levels simply don’t have time to form. Since they provide the only mechanism for coherent transport through the system, the latter is impossible, unless our ”no dissipation” criterion holds.
For the thickness of the normal layer $`L1000\AA `$ and $`v_F10^7`$ cm/s this criterion limits $`\omega _0<10^{12}\mathrm{s}^1`$, which is a comfortable two orders of magnitude above the usually obtained tunneling splitting in such qubits ($`1`$ GHz) and can be accommodated in the above designs. Nevertheless, while presenting a useful qualitative picture, the DND model is not adequate for the task of extracting quantitative predictions. For example, the coherence length $`l_T`$ in the normal metal can be very large, while in the high-T<sub>c</sub> compound it is short. Therefore the estimates for crucial parameters (like $`\overline{ϵ}`$) based on the assumption $`l_TL`$ can be wrong. Moreover, the assumption of ideal ND boundaries is not realistic.
A calculationNewns02 , which used a model of a DD junction interacting with a bosonic thermal bath, gave an optimistic estimate for the quality of the tricrystal qubit $`Q>10^8`$.
The role of size quantization of quasiparticles in small DD and SND structures was suggested in Zagoskin99 ; Ioffe99 . The importance of the effect is that it would exponentially suppress the quasiparticle density and therefore the dissipation below temperature of the quantization gap, estimated as $`110`$ K. Recently this problem was investigatedFominov03 for a finite width DD junction. Contrary to the expectations, the size quantization as such turned out to be effectively absent on the scale exceeding $`\xi _0`$ (that is, practically irrelevant). On the other hand, the finite transverse size of the system imposed an effective band structure. Namely, due to the direction dependence of the order parameter, a quasiparticle travelling along the trajectory, bouncing between the sides of the junction, goes through a periodic 1D off-diagonal potential. The influence of this band structure on dissipation in the system is not straightforward. Ironically, from the practical point of view this is a moot point, since the decoherence time from the quasiparticles in the junction, estimated in Fominov03 ; Fominov03a , already corresponds to the quality factor $`\tau _\phi /\tau _g10^6`$, which exceeds by two orders of magnitude the theoretical threshold allowing to run a quantum computer indefinitely.
The expression for the decoherence time obtained in Fominov03 ; Fominov03a ,
$$\tau _\phi =\frac{4e}{\delta \varphi _2I(\mathrm{\Delta }_t/e)},$$
(58)
where $`\delta \varphi `$ is the difference between equilibrium phases in degenerate minima of the junction (i.e. $`\delta \varphi =2\chi _0`$ in other notation), contains the expression for quasiparticle current in the junction at finite voltage $`\mathrm{\Delta }_t/e`$ (where $`\mathrm{\Delta }_t`$ is the tunneling rate between the minima). This agrees with our back-of-the-envelope analysis: phase tunneling leads to finite voltage in the system through the second Josephson relation and with finite voltage comes quasiparticle current and decoherence. The aforementioned quality factor is defined as $`Q=\tau _\phi \mathrm{\Delta }_t/2\mathrm{}`$, that is, we compare the decoherence time with the tunneling time. This is a usual optimistic estimate, since it will certainly take several tunneling cycles to perform a quantum gate operation.
It turns out that a much bigger threat is posed by the contribution from zero energy bound states, which can be at least two orders of magnitude larger. We can see this qualitatively from (58): a large density of quasiparticle states close to zero energy (i.e., on Fermi level) means that even small voltages create large quasiparticle currents, which sit in the denominator of the expression for $`\tau _\phi .`$ Fortunately, this contribution is suppressed in the case of ZES splitting, and such splitting is always present due to, e.g., finite equilibrium phase difference across the junctionFominov03 .
A similar picture follows from the analysis presented in Amin03b . A specific question addressed there is especially important: it is known that $`RC`$-constant measured in DD junctions is consistently 1 ps over a wide range of junction sizes Tzalenchuk03 , and it is tempting to accept this value as the dissipation rate in the system. It would be a death knell for any quantum computing application of high-T<sub>c</sub> structures, and nearly that for any hope to see there some quantum effects. Nevertheless, it is not quite that bad. Indeed, we saw that ZES play a major role in dissipation in a DD junction, but are sensitive to phase differences across itFN1 . Measurements of the $`RC`$ constant are done in the resistive regime, when a finite voltage exists across the junction, so that the phase difference grows monotonously in time, forcing ZES to approach the Fermi surface repeatedly. Therefore $`\tau _{RC}`$ reflects some averaged dissipation rate. On the other hand, in a free junction with not too high a tunneling rate phase differences obviously tend to oscillate around $`\chi _0`$ or $`\chi _0`$, its equilibrium values, and do not spend much time near zero or $`\pi `$; therefore ZES are usually shifted from the Fermi level, and their contribution to dissipation is suppressed.
This qualitative picture is confirmed by a detailed calculationAmin03b (Fig. 14). The decoherence time is related to the phase-dependent conductance via
$$\tau _\phi =\frac{1}{\alpha F(\chi _0)^2\delta E}\mathrm{tanh}\frac{\delta E}{2T}.$$
(59)
Here $`\alpha `$ is the dissipation coefficient, $`\delta E`$ is interlevel spacing in the well, and
$$G(\chi )=4e^2\alpha \left(_\chi F(\chi )\right)^2.$$
(60)
For a realistic choice of parameters Eq.(59) gives a conservative estimate $`\tau _\phi =1100`$ ns, and quality factor $`Q1100`$. This is, of course, too little for quantum computing, but quite enough for observation of quantum tunneling and coherence in such junctions.
## III Conclusions
Since the $`d`$-wave character of superconducting pairing in high-T<sub>c</sub> cuprates was established, our understanding of high-T<sub>c</sub> Josephson structures and ability to fabricate them progressed significantly. Now it enables us to ask the question in the title of this paper: Can high-T<sub>c</sub> cuprates play a role in quantum computing?- and to tentatively answer: Yes.
We have seen that submicron junctions of sufficiently high quality are now fabricated. Several designs, which take advantage of the intrinsic bistability of $`d`$-wave structures, were developed, fabricated, and tested in classical regime. If the decoherence time turns out to be large enough, their better scalability can be a decisive advantage for quantum computing applications.
This is, of course, a big if. Still, theoretical estimates tend to be rather optimistic. Even though they widely differ, all of them predict $`\tau _\phi `$ long enough to encourage experimental search for quantum coherence. There are several compelling reasons to do that. First, it would be a spectacular result. Second, it would clarify why different models give different answers. Third, it could indeed lead to practical application of high-T<sub>c</sub> devices in quantum computing.
## P.S.
Since this paper was published (2003), macroscopic quantum tunneling was observed in YBCO MQT\_YBCO and BiSCCO MQT\_BiSCCO junctions. This made the ”big if” of the previous paragraph somewhat smaller.
## Acknowledgements
I am grateful to M.H.S. Amin, A. Blais, A. Golubov, E. Il’ichev, A.N. Omelyanchouk, A. Smirnov, and A. Tzalenchuk for many illuminating discussions, and to B. Wilson for the careful reading and helpful comments on the manuscript.
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# Lorentz invariant supersymmetric mechanism for non(anti)commutative deformations of space-time geometry
## 1 Introduction
The unification of physics and mathematics in the development of noncommutative quantum geometry and field theories resulted in new ideas and approaches (see reviews , and additional references therein). One of them has come from string theory, where the noncommutativity of the bosonic string coordinates $`x_m`$ in the presence of the constant antisymmetric field $`B_{mn}`$ was observed . More recently, the noncommutativity between the components of the odd spinor coordinate $`\theta _a`$ in the presence of a constant graviphoton field $`C_{ab}`$ was considered in . The constant gravitino background $`\mathrm{\Psi }_m^a`$ resulted in the noncommutativity between the $`x_m`$ and $`\theta _a`$ coordinates . These results have focused attention on the role of constant background fields in superspace deformations. Studying field/string theories and supersymmetry preservation in the superspaces deformed by the graviphoton background , was further advanced in and . A general approach to the construction of superspace deformations in a constant background based on the Moyal-Weyl quantization of the Poisson brackets was developed in . The presence of constant background fields in the much discussed deformed (anti)commutation relations for the (super)coordinate operators leads to the well-recognized problem of Lorentz symmetry breaking. The idea to overcome this problem by using twisted Hopf algebra was recently proposed in and its supersymmetric generalization was realized in and further developed in , . Another way was observed in , where the Hamiltonian structure of free twistor-like model of super p-brane in $`N=1`$ superspace extended by tensor central charge coordinates was studied and the Dirac bracket-non(anti)commutativity of the brane (super)coordinates was established. The r.h.s. of these D.B’s. have been constructed from the components of auxiliary twistor-like dynamical variables which are Lorentz covariant and supersymmetric. It gives a hint that a hidden spinor structure, associated with the Penrose twistor picture might be an alternative source for the non(anti)commutativity of the quantum space-time (super)coordinates. Accepting such a possibility we start here from the above mentioned spinor extension of the $`N=1D=4`$ superspace $`(x_m,\theta _a)`$ by one commuting Majorana spinor coordinate $`\lambda _a`$ and construct Lorentz invariant and supersymmetric Poisson and Moyal brackets generating non(anti)commutative relations of the (super)coordinates. An interesting feature of these brackets is the presence of a real (or complex) Grassmannian vector $`\psi _m`$, which is well known from the theory of spinning strings and particles , in the r.h.s. of the brackets of $`x_m`$ with $`x_n`$ and $`\theta _a`$. The odd vector $`\psi _m`$ appears there in the form of an effective variable $`\psi _m=\frac{1}{2}(\overline{\theta }\gamma _m\lambda )`$ composed from the two Majorana spinors $`\lambda _a,\theta _a`$ and encoding primordial degrees of freedom presented by $`\theta _a`$. In the simplest case there is a correspondence between the Lorentz invariant brackets in question and the known brackets including the constant background fields. That correspondence may be schematically illustrated as the map transforming the field dependent brackets into the new brackets and vice versa :
$$B_{mn}i\psi _m\psi _n,C_{ab}\lambda _a\lambda _b,\mathrm{\Psi }_m^a\psi _m\lambda ^a.$$
(modulo the change $`B_{mn}\theta _{mn}B_{mn}^1`$ etc). The schematical correspondence is preserved in the more sophisticated cases considered below and points to a deep correlation between the spin structure, non(anti)commutativity and supergravity. We find also Lorentz invariant and supersymmetric brackets, where nonanticommutativity occurs only for the components of $`\theta ^a`$ with opposite chirality. The generalizations to the higher dimensions $`D=2,3,4(mod8),N>1`$ and several additional spinors are discussed.
## 2 Supersymmetry algebra in the presence of a spinor coordinate
Using the agreements of we accept here the $`D=4N=1`$ supersymmetry transformation law in the presence of the twistor-like Majorana spinor coordinates $`(\nu _\alpha ,\overline{\nu }_{\dot{\alpha }})`$ in the form
$$\begin{array}{c}\delta \theta _\alpha =\epsilon _\alpha ,\delta x_{\alpha \dot{\alpha }}=2i(\epsilon _\alpha \overline{\theta }_{\dot{\alpha }}\theta _\alpha \overline{\epsilon }_{\dot{\alpha }}),\delta \nu _\alpha =0,\end{array}$$
(1)
The supercharges $`Q_\alpha `$ and $`\overline{Q}_{\dot{\alpha }}`$ of the superalgebra (1) are given by the differential operators
$$\begin{array}{c}Q^\alpha =\frac{}{\theta _\alpha }+2i\overline{\theta }_{\dot{\alpha }}^{\alpha \dot{\alpha }},\overline{Q}^{\dot{\alpha }}(Q^\alpha )^{}=\frac{}{\overline{\theta }_{\dot{\alpha }}}+2i\theta _\alpha ^{\alpha \dot{\alpha }},[Q^\alpha ,\overline{Q}^{\dot{\alpha }}]_+=4i^{\alpha \dot{\alpha }},\end{array}$$
(2)
where $`^{\alpha \dot{\alpha }}\frac{}{x_{\alpha \dot{\alpha }}}`$ and the correspondent supersymmetric covariant derivatives $`D,\overline{D}`$ are
$$\begin{array}{c}D^\alpha =\frac{}{\theta _\alpha }2i\overline{\theta }_{\dot{\alpha }}^{\alpha \dot{\alpha }},\overline{D}^{\dot{\alpha }}(D^\alpha )^{}=\frac{}{\overline{\theta }_{\dot{\alpha }}}2i\theta _\alpha ^{\alpha \dot{\alpha }},[D^\alpha ,\overline{D}^{\dot{\beta }}]=4i^{\alpha \dot{\alpha }},\\ [Q^\alpha ,D^\beta ]_+=[Q^\alpha ,\overline{D}^{\dot{\beta }}]_+=[\overline{Q}^{\dot{\alpha }},D^\beta ]_+=[\overline{Q}^{\dot{\alpha }},\overline{D}^{\dot{\beta }}]_+=0.\end{array}$$
(3)
The spinor coordinates $`(\nu _\alpha ,\overline{\nu }_{\dot{\alpha }})`$ and the light-like real vector $`\phi _{\alpha \dot{\alpha }}`$ composed from them
$$\phi _{\alpha \dot{\alpha }}\nu _\alpha \overline{\nu }_{\dot{\alpha }},\phi _{\alpha \dot{\alpha }}\nu ^\alpha =\phi _{\alpha \dot{\alpha }}\overline{\nu }^{\dot{\alpha }}=0,\delta \phi _{\alpha \dot{\beta }}=0$$
(4)
may be used to construct the Lorentz invariant differential operators $`D,\overline{D},`$
$$\begin{array}{c}D=\nu _\alpha D^\alpha ,\overline{D}=\overline{\nu }_{\dot{\alpha }}\overline{D^{\dot{\alpha }}},=\phi _{\alpha \dot{\alpha }}^{\alpha \dot{\alpha }}\end{array}$$
(5)
which form a supersymmetric subalgebra of the algebra of derivatives
$$\begin{array}{c}[D,\overline{D}]_+=4i,[D,D]_+=[\overline{D},\overline{D}]_+=0,[D,]=[\overline{D},]=[,]=0.\end{array}$$
(6)
Introducing $`D_\pm `$ combinations of the invariant derivatives $`D,\overline{D}`$
$$\begin{array}{c}D_\pm D\pm \overline{D}\end{array}$$
(7)
one can split the Lorentz invariant complex subalgebra (6) into two invariant and (anti)commuting subalgebras formed by the generators $`(D_{},)`$ and $`(D_+,)`$
$$\begin{array}{c}[D_\pm ,D_\pm ]_+=8i,[D_+,D_{}]_+=0,[D_\pm ,]=[,]=0.\end{array}$$
(8)
A twistor-like character of the Majorana spinor $`(\nu _\alpha ,\overline{\nu }_{\dot{\alpha }})`$ means that its dilatations, generated by the differential operator $`\mathrm{\Delta }`$
$$\begin{array}{c}\mathrm{\Delta }=\nu _\alpha \frac{}{\nu _\alpha }+\overline{\nu }_{\dot{\alpha }}\frac{}{\overline{\nu }_{\dot{\alpha }}},\end{array}$$
(9)
have to be an additional symmetry of physical theories of massless fields. Taking into account this dilaton symmetry assumes extension of the superalgebra (6) by the dilaton generator $`\mathrm{\Delta }`$
$$\begin{array}{c}[D,\overline{D}]_+=4i,[D,D]_+=[\overline{D},\overline{D}]_+=0,\\ [\mathrm{\Delta },D]=D,[\mathrm{\Delta },\overline{D}]=\overline{D},[\mathrm{\Delta },]=2,\\ [,D]=[,\overline{D}]=[,]=[\mathrm{\Delta },\mathrm{\Delta }]=0,\end{array}$$
(10)
which has two real anticommutative subalgebras formed by the generators $`(D_\pm ,,\mathrm{\Delta })`$
$$\begin{array}{c}[D_\pm ,D_\pm ]_+=8i,[\mathrm{\Delta },D_\pm ]=D_\pm ,[\mathrm{\Delta },]=2,\\ [D_+,D_{}]_+=[D_\pm ,]=[,]=[\mathrm{\Delta },\mathrm{\Delta }]=0.\end{array}$$
(11)
The Lorentz invariant supersymmetric differential operators forming the superalgebras (10), (11) may be used as building blocks for the construction of Lorentz invariant and supersymmetric non(anti)commutative relations among quantum operators of (super)coordinates.
## 3 Lorentz invariant supersymmetric Poisson brackets: non(anti)commutativity of space-time coordinates
To clarify the role of $`(\nu _\alpha ,\overline{\nu }_{\dot{\alpha }})`$ in the formation of Lorentz invariant non(anti)commutative relations among $`x_{\alpha \dot{\alpha }},\theta _\alpha ,\overline{\theta }_{\dot{\alpha }}`$ we consider the Poisson bracket constructed from the three differential operators $`(D_{},,\mathrm{\Delta })`$ forming the simplest superalgebra (11)
$$\begin{array}{c}\{F,G\}=F[\frac{i}{4}\stackrel{}{D}_{}\stackrel{}{D}_{}+(\stackrel{}{}\stackrel{}{\mathrm{\Delta }}\stackrel{}{\mathrm{\Delta }}\stackrel{}{})]G,\end{array}$$
(12)
where $`\{,\}_{P.B.}\{,\}`$ and $`F(x,\theta ,\overline{\theta },\nu ,\overline{\nu }),G(x,\theta ,\overline{\theta },\nu ,\overline{\nu })`$ are generalized superfields depending on both the superspace coordinates $`(x,\theta ,\overline{\theta })`$ and on the spinor coordinates $`(\nu ,\overline{\nu })`$. The Lorentz invariant and supersymmetric differential operators $`\stackrel{}{D}_{},\stackrel{}{},\stackrel{}{\mathrm{\Delta }}`$ define derivatives acting from the right hand side. Conversely, the differential operators (11) act from the l.h.s and coincide with the left derivatives $`\stackrel{}{D}_{},\stackrel{}{},\stackrel{}{\mathrm{\Delta }}`$
$$\begin{array}{c}\stackrel{}{D}_{}GD_{}G,\stackrel{}{}GG,\stackrel{}{\mathrm{\Delta }}G\mathrm{\Delta }G\end{array}$$
(13)
The left and right invariant derivatives in (12) are connected by the relations
$$\begin{array}{c}F\stackrel{}{D}=(1)^f\stackrel{}{D}F,F\stackrel{}{\overline{D}}=(1)^f\stackrel{}{\overline{D}}F,F\stackrel{}{}=\stackrel{}{}F,F\stackrel{}{\mathrm{\Delta }}=\stackrel{}{\mathrm{\Delta }}F,\end{array}$$
(14)
where $`f=0,1`$ is the Grassmannian grading of the superfield $`F`$.
The action of $`D,\overline{D},D_\pm `$ on the composite coordinates $`\psi `$ and $`\phi `$ is given by the relations
$$\begin{array}{c}\stackrel{}{D}\psi _{\alpha \dot{\alpha }}=i\phi _{\alpha \dot{\alpha }},\stackrel{}{\overline{D}}\psi _{\alpha \dot{\alpha }}=i\phi _{\alpha \dot{\alpha }},\psi _{\alpha \dot{\alpha }}i(\nu _\alpha \overline{\theta }_{\dot{\alpha }}\theta _\alpha \overline{\nu }_{\dot{\alpha }}),\\ \stackrel{}{D_{}}\psi _{\alpha \dot{\alpha }}=2i\phi _{\alpha \dot{\alpha }},\stackrel{}{D_+}\psi _{\alpha \dot{\alpha }}=\stackrel{}{D_\pm }\phi _{\alpha \dot{\alpha }}=0.\end{array}$$
(15)
After the substitions of the (super)coordinates under discussion in the P.B. (12) we find non(anti)commutative P.B’s. for them.
The twistor-like coordinates have zero P.B’s. among themselves
$$\{\nu _\alpha ,\nu _\beta \}=\{\nu _\alpha ,\overline{\nu }_{\dot{\beta }}\}=\{\overline{\nu }_\alpha ,\overline{\nu }_{\dot{\beta }}\}=0$$
(16)
and with the Grassmannian spinors $`(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }})`$
$$\begin{array}{c}\{\nu _\alpha ,\theta _\beta \}=\{\nu _\alpha ,\overline{\theta }_{\dot{\beta }}\}=\{\overline{\nu }_{\dot{\alpha }},\theta _\beta \}=\{\overline{\nu }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=0.\end{array}$$
(17)
However, they have non zero P.B’s. with the space-time coordinates $`x_{\alpha \dot{\alpha }}`$
$$\{x_{\alpha \dot{\alpha }},\nu _\beta \}=\phi _{\alpha \dot{\alpha }}\nu _\beta ,\{x_{\alpha \dot{\alpha }},\overline{\nu }_{\dot{\beta }}\}=\phi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }},$$
(18)
The remaining non zero P.B’s. define the P.B’s. among the space-time coordinates $`x_{\alpha \dot{\alpha }}`$ and spinors $`(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }})`$
$$\begin{array}{c}\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\}=i\psi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }},\\ \{x_{\alpha \dot{\alpha }},\theta _\beta \}=\frac{i}{2}\psi _{\alpha \dot{\alpha }}\nu _\beta ,\{x_{\alpha \dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\frac{i}{2}\psi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }},\\ \{\theta _\alpha ,\theta _\beta \}=\frac{i}{4}\phi _{\alpha \beta },\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\}=\frac{i}{4}\phi _{\alpha \dot{\beta }},\{\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\frac{i}{4}\overline{\phi }_{\dot{\alpha }\dot{\beta }},\end{array}$$
(19)
where $`\phi _{\alpha \beta },\overline{\phi }_{\dot{\alpha }\dot{\beta }}`$ are composite symmetric spin-tensors
$$\phi _{\alpha \beta }\nu _\alpha \nu _\beta ,\overline{\phi }_{\dot{\alpha }\dot{\beta }}\overline{\nu }_{\dot{\alpha }}\overline{\nu }_{\dot{\beta }},\delta \phi _{\alpha \beta }=\delta \overline{\phi }_{\dot{\alpha }\dot{\beta }}=0$$
(20)
orthogonal to the vector $`\phi _{\alpha \dot{\alpha }}`$ (4) and to the composite Grassmannian vector $`\psi _{\alpha \dot{\alpha }}`$
$$\psi _{\alpha \dot{\alpha }}i(\nu _\alpha \overline{\theta }_{\dot{\alpha }}\theta _\alpha \overline{\nu }_{\dot{\alpha }}),\phi ^{\alpha \dot{\alpha }}\psi _{\alpha \dot{\alpha }}=\phi ^{\alpha \beta }\psi _{\alpha \dot{\alpha }}=\overline{\phi }^{\dot{\alpha }\dot{\beta }}\psi _{\alpha \dot{\alpha }}=0,\delta \psi _{\alpha \dot{\alpha }}=i(\epsilon _\alpha \overline{\nu }_{\dot{\alpha }}\overline{\epsilon }_{\dot{\alpha }}\nu _\alpha )$$
(21)
The appearance of the odd vector $`\psi _{\alpha \dot{\alpha }}`$ (21) associated with the description of the spin degrees of freedom of fermions in the r.h.s. of P.B’s. (19) hints on a spin structure of superspaces in back of the coordinate’s non(anti)commutativity<sup>1</sup><sup>1</sup>1Let us remind that composite character of the anticommuting vector $`\psi _{\alpha \dot{\alpha }}`$ (21) was revealed in , where the spinor representation (21) was found to be the general solution of Dirac constraints $`p^{\alpha \dot{\alpha }}\psi _{\alpha \dot{\alpha }}=0=p^{\alpha \dot{\alpha }}p_{\alpha \dot{\alpha }}`$ characterising massless spinning particle ,. This spinor representation was important to find equivalence between spinning and Brink-Schwarz superparticles.. The Lorentz covariance of the Poisson brackets (16)-(19) is provided by the spinor, vector and spin-tensor representations of the Lorentz group involved in the r.h.s. of the Poisson brackets. These P.B’s. are also supersymmetric by the construction.
In the next Section we prove the Jacobi identities for the P.B’s. (16)-(19).
## 4 Proof of the Jacobi identities
The graded Jacobi identities for the considered P.B. algebra have the standard form
$$\{\{A,B\},C\}+(1)^{(b+c)a}\{\{B,C\},A\}+(1)^{c(a+b)}\{\{C,A\},B\}=0,$$
(22)
where $`a=0,1`$ is the Grassmannian grading of $`A`$. To prove these identities for the P.B’s. (16)-(19) one needs to study the Poisson brackets of the composite vectors $`\phi _{\alpha \dot{\beta }}`$ (4), $`\psi _{\alpha \dot{\alpha }}`$ (21) and spin-tensors $`\phi _{\alpha \beta },\overline{\phi }_{\dot{\alpha }\dot{\beta }}`$ (20) between themselves and with $`x_{\alpha \dot{\alpha }},\theta _\alpha ,\overline{\theta }_{\dot{\alpha }}`$. The P.B’s. (16), (17) together with the definitions (20), (21) show the P.B.-commutativity of $`\phi _{\alpha \beta },\phi _{\alpha \dot{\beta }},\overline{\phi }_{\dot{\alpha }\dot{\beta }}`$ between themselves and with $`(\nu _\alpha ,\overline{\nu }_{\dot{\alpha }})`$, $`(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }})`$ and $`\psi _{\alpha \dot{\alpha }}`$
$$\begin{array}{c}\{\phi _{},\nu _\alpha \}=\{\phi _{},\overline{\nu }_{\dot{\alpha }}\}=\{\psi _{\alpha \dot{\alpha }},\nu _\beta \}=\{\psi _{\alpha \dot{\alpha }},\overline{\nu }_{\dot{\beta }}\}=0,\\ \{\phi _{},\phi _{}\}=\{\phi _{},\theta _\alpha \}=\{\phi _{},\overline{\theta }_{\dot{\alpha }}\}=\{\phi _{},\psi _{\gamma \dot{\gamma }}\}=0,\end{array}$$
(23)
where $`\phi _{}(\phi _{\alpha \beta },\phi _{\alpha \dot{\beta }},\overline{\phi }_{\dot{\alpha }\dot{\beta }})`$ is a condenced symbol for the composite coordinates (4) and (20). However, the P.B. of the spin-tensors $`\phi _{\alpha \beta },\phi _{\alpha \dot{\beta }},\overline{\phi }_{\dot{\alpha }\dot{\beta }}`$ with $`x_{\gamma \dot{\gamma }}`$ are different from zero
$$\{x_{\alpha \dot{\alpha }},\phi _{\beta \gamma }\}=2\phi _{\alpha \dot{\alpha }}\phi _{\beta \gamma },\{x_{\alpha \dot{\alpha }},\phi _{\beta \dot{\gamma }}\}=2\phi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\gamma }},\{x_{\alpha \dot{\alpha }},\overline{\phi }_{\dot{\beta }\dot{\gamma }}\}=2\phi _{\alpha \dot{\alpha }}\overline{\phi }_{\dot{\beta }\dot{\gamma }},$$
(24)
as well as, the P.B. between $`x_{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}`$ and $`(\theta _\beta ,\overline{\theta }_{\dot{\beta }})`$
$$\begin{array}{c}\{x_{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\}=\phi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }}+\phi _{\beta \dot{\beta }}\psi _{\alpha \dot{\alpha }},\\ \{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\}=i\phi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }},\\ \{\psi _{\alpha \dot{\alpha }},\theta _\beta \}=\frac{1}{2}\phi _{\alpha \dot{\alpha }}\nu _\beta ,\{\psi _{\alpha \dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\frac{1}{2}\phi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }}.\end{array}$$
(25)
A combination of the P.B. relations (23) together with ones (25) results in the relation
$$\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\psi _{\gamma \dot{\gamma }}\}=0$$
(26)
which proves the graded Jacobi identity (12) for the case $`A=B=C=\psi `$
$$Cycle\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\psi _{\gamma \dot{\gamma }}\}=0$$
(27)
The same result occurs for the Jacoby cycles cubic in $`(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }})`$
$$Cycle\{\{\theta _\alpha ,\theta _\beta \},\theta _\gamma \}=\mathrm{}=Cycle\{\{\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\},\overline{\theta }_{\dot{\gamma }}\}=0,$$
(28)
as well as, for the cycles quadratic in $`\theta _\alpha `$ or $`\psi _{\alpha \dot{\alpha }}`$ and linear in $`(\nu _\gamma ,\overline{\nu }_{\dot{\gamma }})`$
$$\begin{array}{c}Cycle\{\{\theta _\alpha ,\theta _\beta \},\nu _\gamma \}=\mathrm{}.=Cycle\{\{\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\},\overline{\nu }_{\dot{\gamma }}\}=0,\\ Cycle\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\nu _\gamma \}=Cycle\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\overline{\nu }_{\dot{\gamma }}\}=0\end{array}$$
(29)
and for other trivial Jacobi cycles cubic or quadratic in $`(\nu _\gamma ,\overline{\nu }_{\dot{\gamma }})`$ and linear in $`(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }})`$ or $`\psi _{\alpha \dot{\alpha }}`$.
To calculate the Jacobi cycle cubic in $`x_{\alpha \dot{\alpha }}`$ we use the relation
$$\{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},x_{\gamma \dot{\gamma }}\}=2i(\psi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }})\phi _{\gamma \dot{\gamma }}+i(\psi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\psi _{\beta \dot{\beta }}\phi _{\alpha \dot{\alpha }})\psi _{\gamma \dot{\gamma }}$$
(30)
arisen from the P.B’s. (19) and (25) and resulting in zero Jacobi cycle
$$Cycle\{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},x_{\gamma \dot{\gamma }}\}=0.$$
(31)
It follows from the mutual cancellation between the contributions of first and last summands in the r.h.s. of the cyclic sum generated by Eq. (30). Next one can see that the Jacobi cycles quadratic in $`x_{\alpha \dot{\alpha }}`$ and linear in $`\psi _{\alpha \dot{\alpha }}`$ or $`(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }})`$ are equal to zero
$$\begin{array}{c}Cycle\{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},\psi _{\gamma \dot{\gamma }}\}=0,\\ Cycle\{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},\theta _\gamma \}=Cycle\{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},\overline{\theta }_{\dot{\gamma }}\}=0,\end{array}$$
(32)
because of the relations
$$\begin{array}{c}\{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},\psi _{\gamma \dot{\gamma }}\}=(\psi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\phi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }})\phi _{\gamma \dot{\gamma }},\\ \{\{x_{\beta \dot{\beta }},\psi _{\gamma \dot{\gamma }}\},x_{\alpha \dot{\alpha }}\}(\alpha \beta )=(\psi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\phi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }})\phi _{\gamma \dot{\gamma }},\\ \{\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\},\theta _\gamma \}=\frac{i}{2}(\psi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\phi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }})\nu _\gamma ,\\ \{\{x_{\beta \dot{\beta }},\theta _\gamma \},x_{\alpha \dot{\alpha }}\}(\alpha \beta )=\frac{i}{2}(\psi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\phi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }})\nu _\gamma \end{array}$$
(33)
and their complex conjugate following from the P.B’s. (18), (19) and (25). A similar cancellation takes place in the Jacobi cycles quadratic in $`\psi _{\alpha \dot{\alpha }}`$ and linear in $`x_{\alpha \dot{\alpha }}`$ or $`\theta _\alpha ,\overline{\theta }_{\dot{\alpha }}`$
$$\begin{array}{c}Cycle\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},x_{\gamma \dot{\gamma }}\}=0,\\ Cycle\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\theta _\gamma \}=Cycle\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\overline{\theta }_{\dot{\gamma }}\}=0,\end{array}$$
(34)
as it follows from the P.B. relations
$$\begin{array}{c}\{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},x_{\gamma \dot{\gamma }}\}=4i\phi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\phi _{\gamma \dot{\gamma }},\\ \{\{x_{\gamma \dot{\gamma }},\psi _{\beta \dot{\beta }}\},\psi _{\alpha \dot{\alpha }}\}+(\alpha \beta )=4i\phi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\phi _{\gamma \dot{\gamma }},\\ \{\{\psi _{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\theta _\gamma \}=\{\{\psi _{\alpha \dot{\alpha }},\theta _\gamma \},\psi _{\beta \dot{\beta }}\}=0.\end{array}$$
(35)
Next we prove the Jacobi identites for cycles quadratic in $`\theta _\alpha ,\overline{\theta }_{\dot{\alpha }}`$ and linear in $`x_{\alpha \dot{\alpha }}`$ or $`\psi _{\alpha \dot{\alpha }}`$
$$\begin{array}{c}Cycle\{\{\theta _\alpha ,\theta _\beta \},x_{\gamma \dot{\gamma }}\}=\{\{\theta _\alpha ,\theta _\beta \},\psi _{\gamma \dot{\gamma }}\}=0,\\ Cycle\{\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\},x_{\gamma \dot{\gamma }}\}=\{\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\},\psi _{\gamma \dot{\gamma }}\}=0\end{array}$$
(36)
and for their complex conjugate using the relations
$$\begin{array}{c}\{\{\theta _\alpha ,\theta _\beta \},x_{\gamma \dot{\gamma }}\}=\frac{i}{2}\phi _{\alpha \beta }\phi _{\gamma \dot{\gamma }},\{\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\},x_{\gamma \dot{\gamma }}\}=\frac{i}{2}\phi _{\alpha \dot{\beta }}\phi _{\gamma \dot{\gamma }},\\ \{\{x_{\gamma \dot{\gamma }},\theta _\beta \},\theta _\alpha \}+(\alpha \beta )=\frac{i}{2}\phi _{\alpha \beta }\phi _{\gamma \dot{\gamma }},\\ \{\{x_{\gamma \dot{\gamma }},\overline{\theta }_{\dot{\beta }}\},\theta _\alpha \}+\{\{x_{\gamma \dot{\gamma }},\theta _\alpha \},\overline{\theta }_{\dot{\beta }}\}=\frac{i}{2}\phi _{\alpha \dot{\beta }}\phi _{\gamma \dot{\gamma }}.\end{array}$$
(37)
together with the relations
$$\begin{array}{c}\{\{\theta _\alpha ,\theta _\beta \},\psi _{\gamma \dot{\gamma }}\}=\{\{\theta _\alpha ,\psi _{\gamma \dot{\gamma }}\},\theta _\beta \}=0,\\ \{\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\},\psi _{\gamma \dot{\gamma }}\}=\{\{\theta _\alpha ,\psi _{\gamma \dot{\gamma }}\},\overline{\theta }_{\dot{\beta }}\}=0\end{array}$$
(38)
and their complex conjugate. The remaining nontrivail and also vanishing Jacobi cycles are formed by any three coordinates from the set $`[x_{\alpha \dot{\alpha }},\psi _{\alpha \dot{\alpha }},(\theta _\alpha ,\overline{\theta }_{\dot{\alpha }}),(\nu _\alpha ,\overline{\nu }_{\dot{\alpha }})]`$
$$\begin{array}{c}Cycle\{\{x_{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\theta _\gamma \}=Cycle\{\{x_{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\nu _\gamma \}=0,\\ Cycle\{\{x_{\alpha \dot{\alpha }},\theta _\beta \},\nu _\gamma \}=Cycle\{\{\psi _{\alpha \dot{\alpha }},\theta _\beta \},\nu _\gamma \}=0.\end{array}$$
(39)
Their complex conjugate cycles equal zero too.. The proof of first and second Jacobi identities in (39) is based on the P.B. relations
$$\begin{array}{c}\{\{x_{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\theta _\gamma \}=\frac{2}{3}\{\{\psi _{\beta \dot{\beta }},\theta _\gamma \},x_{\alpha \dot{\alpha }}\}=2\{\{x_{\alpha \dot{\alpha }},\theta _\gamma \},\psi _{\beta \dot{\beta }}\}=\phi _{\alpha \dot{\alpha }}\phi _{\beta \dot{\beta }}\nu _\gamma ,\\ \{\{x_{\alpha \dot{\alpha }},\psi _{\beta \dot{\beta }}\},\nu _\gamma \}=\{\{\psi _{\beta \dot{\beta }},\nu _\gamma \},x_{\alpha \dot{\alpha }}\}=\{\{\nu _\gamma ,x_{\alpha \dot{\alpha }}\},\psi _{\beta \dot{\beta }}\}=0.\end{array}$$
(40)
The proof of third and fourth Jacobi identitities in (39) uses the P.B. relations
$$\begin{array}{c}\{\{x_{\alpha \dot{\alpha }},\theta _\beta \},\nu _\gamma \}=\{\{\theta _\beta ,\nu _\gamma \},x_{\alpha \dot{\alpha }}\}=\{\{\nu _\gamma ,x_{\alpha \dot{\alpha }}\},\theta _\beta \}=0,\\ \{\{\psi _{\alpha \dot{\alpha }},\theta _\beta \},\nu _\gamma \}=\{\{\theta _\beta ,\nu _\gamma \},\psi _{\alpha \dot{\alpha }}\}=\{\{\nu _\gamma ,\psi _{\alpha \dot{\alpha }}\},\theta _\beta \}=0\end{array}$$
(41)
which follow from the P.B. relations (17), (19), (23), (25). It complets the proof of the Jacobi identities for the above introduced Lorentz invariant Poisson brackets. The next step is to use them for the construction of the Moyal brackets.
## 5 Lorentz invariant and supersymmetric star product
A transition to the quantum picture based on the P.B. (12) may be done by using the well known Weyl-Moyal correspondence which establishes one to one correspondence between quantum field operators and their symbols acting on commutative space-time. Then the quantum information is encoded in the change of the usual product by the Moyal $``$-product of their Weyl symbols. To realise this prescription here we note that the P.B. (12) may be presented as
$$\begin{array}{c}\{F,G\}=F\stackrel{}{𝒟_\mathrm{\Lambda }}C^{\mathrm{\Lambda }\mathrm{\Sigma }}\stackrel{}{𝒟_\mathrm{\Sigma }}G=\\ F(\stackrel{}{},\stackrel{}{\mathrm{\Delta }},\stackrel{}{D_{}})\left(\begin{array}{ccc}0\hfill & \hfill 1& 0\\ 1\hfill & \hfill 0& 0\\ 0\hfill & \hfill 0& \frac{i}{4}\end{array}\right)\left(\begin{array}{cc}\stackrel{}{}& \\ \stackrel{}{\mathrm{\Delta }}& \\ \stackrel{}{D_{}}& \end{array}\right)G,\end{array}$$
(42)
where the condenced notation $`𝒟_\mathrm{\Lambda }=(D_{},,\mathrm{\Delta })`$ was used for the invariant derivatives of the (-)-superalgebra (11) numerated by the index $`\mathrm{\Lambda }`$ running over the even and odd variables. As a result, the superalgebra (11) is presented in a condenced form
$$\begin{array}{c}[𝒟_\mathrm{\Lambda },𝒟_\mathrm{\Sigma }\}=𝒞_{\mathrm{\Lambda }\mathrm{\Sigma }}{}_{}{}^{\mathrm{\Xi }}𝒟_{\mathrm{\Xi }}^{},\end{array}$$
(43)
where $`𝒞_{\mathrm{\Lambda }\mathrm{\Sigma }}^\mathrm{\Xi }`$ are the structural constants defined by the explicit (anti)commutation relations (11) and $`C^{\mathrm{\Lambda }\mathrm{\Sigma }}`$ is represented by the $`3\times 3`$ matrice
$$C^{\mathrm{\Lambda }\mathrm{\Sigma }}=\left(\begin{array}{ccc}0\hfill & \hfill 1& 0\\ 1\hfill & \hfill 0& 0\\ 0\hfill & \hfill 0& \frac{i}{4}\end{array}\right).$$
(44)
The representation (42) defines the Moyal $``$-product of the superfields $`F`$ and $`G`$
$$FG=Fe^{\frac{1}{2}\stackrel{}{𝒟_\mathrm{\Lambda }}C^{\mathrm{\Lambda }\mathrm{\Sigma }}\stackrel{}{𝒟_\mathrm{\Sigma }}}G,$$
(45)
where the Planck constant and the velocity of light are chosen to be equal to unit. The definition (45) together with (12) yield the Moyal products of the (super)coordinates
$$\begin{array}{c}x_{\alpha \dot{\alpha }}x_{\beta \dot{\beta }}=x_{\alpha \dot{\alpha }}x_{\beta \dot{\beta }}\frac{i}{2}\psi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }},\\ x_{\alpha \dot{\alpha }}\theta _\beta =x_{\alpha \dot{\alpha }}\theta _\beta +\frac{i}{4}\psi _{\alpha \dot{\alpha }}\nu _\beta ,x_{\alpha \dot{\alpha }}\overline{\theta }_{\dot{\beta }}=x_{\alpha \dot{\alpha }}\overline{\theta }_{\dot{\beta }}\frac{i}{4}\psi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }},\\ \theta _\alpha \theta _\beta =\theta _\alpha \theta _\beta +\frac{i}{8}\phi _{\alpha \beta },\theta _\alpha \overline{\theta }_{\dot{\beta }}=\theta _\alpha \overline{\theta }_{\dot{\beta }}\frac{i}{8}\phi _{\alpha \dot{\beta }},\\ \overline{\theta }_{\dot{\alpha }}\overline{\theta }_{\dot{\beta }}=\overline{\theta }_{\dot{\alpha }}\overline{\theta }_{\dot{\beta }}+\frac{i}{8}\overline{\phi }_{\dot{\alpha }\dot{\beta }}.\end{array}$$
(46)
Consequently, the (anti)commutators of the coordinate operators are replaced by the following Lorentz invariant and supersymmetric Moyal brackets
$$\begin{array}{c}[x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}]_{}x_{\alpha \dot{\alpha }}x_{\beta \dot{\beta }}x_{\beta \dot{\beta }}x_{\alpha \dot{\alpha }}=i\psi _{\alpha \dot{\alpha }}\psi _{\beta \dot{\beta }},\\ [x_{\alpha \dot{\alpha }},\theta _\beta ]_{}=\frac{i}{2}\psi _{\alpha \dot{\alpha }}\nu _\beta ,[x_{\alpha \dot{\alpha }},\overline{\theta }_{\dot{\beta }}]_{}=\frac{i}{2}\psi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }},\\ [\theta _\alpha ,\theta _\beta ]_\mathrm{?}=\frac{i}{4}\phi _{\alpha \beta },[\theta _\alpha ,\overline{\theta }_{\dot{\beta }}]_+=\frac{i}{4}\phi _{\alpha \dot{\beta }},[\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}]_+=\frac{i}{4}\overline{\phi }_{\dot{\alpha }\dot{\beta }},\end{array}$$
(47)
which in turn are directly restored from the invariant P.B’s. (16)-(19). The change $`\{,\}[,]_{}`$ restores the remaining Moyal brackets originated from the above considered P.B’s that together with the brackets (47) may be used for the studying Lorentz invariant and supersymmetric quantum field models in non(anti)commutative superspace.
## 6 Noncommutativity of the twistor components
The twistor associated with $`\nu _\alpha `$ and $`x_{\alpha \dot{\alpha }}`$ is formed by the pair $`Z^𝒜=(\omega ^\alpha ,\overline{\nu }_{\dot{\alpha }})`$, where the first twistor element $`\omega ^\alpha `$ is composed from $`\nu _\alpha `$ and $`x_{\alpha \dot{\alpha }}`$
$$\omega _\alpha =ix_{\alpha \dot{\alpha }}\overline{\nu }^{\dot{\alpha }}.$$
(48)
The considered Poisson and Moyal brackets result in the commutativity between the twistor components $`\omega _\alpha `$, $`\nu _\beta `$ and their complex conjugate
$$\{\omega _\alpha ,\nu _\beta \}=\{\omega _\alpha ,\overline{\nu }_{\dot{\beta }}\}=\{\overline{\omega }_{\dot{\alpha }},\nu _\beta \}=\{\overline{\omega }_{\dot{\alpha }},\overline{\nu }_{\dot{\beta }}\}=0,$$
(49)
because of the P.B’s. (16), (18) and the orthogonality relations
$$\phi _{\alpha \dot{\alpha }}\nu ^\alpha =\phi _{\alpha \dot{\alpha }}\overline{\nu }^{\dot{\alpha }}=0.$$
(50)
However, $`\omega _\alpha `$ and $`\overline{\omega }_{\dot{\alpha }}`$ have non zero brackets with $`x_{\alpha \dot{\alpha }}`$
$$\{x_{\alpha \dot{\alpha }},\omega _\beta \}=\phi _{\alpha \dot{\alpha }}\omega _\beta i\overline{\eta }\psi _{\alpha \dot{\alpha }}\nu _\beta ,\{x_{\alpha \dot{\alpha }},\overline{\omega }_{\dot{\beta }}\}=\phi _{\alpha \dot{\alpha }}\overline{\omega }_{\dot{\beta }}i\eta \psi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }},$$
(51)
as well as, with $`\theta _\alpha `$ and $`\overline{\theta }_{\dot{\alpha }}`$
$$\begin{array}{c}\{\omega _\alpha ,\theta _\beta \}=\frac{i}{2}\overline{\eta }\phi _{\alpha \beta },\{\omega _\alpha ,\overline{\theta }_{\dot{\beta }}\}=\frac{i}{2}\overline{\eta }\phi _{\alpha \dot{\beta }},\eta \theta _\alpha \nu ^\alpha ,\end{array}$$
(52)
because of the P.B’s. (17)-(19), (49). The Grassmannian scalar $`\eta `$ has zero P.B’s. with $`\nu ,\omega ,\theta `$
$$\{\eta ,\nu _\alpha \}=\{\eta ,\omega _\alpha \}=\{\eta ,\theta _\alpha \}=0$$
(53)
and their complex conjugate. The multiplication of the relations (51) by $`(i\overline{\nu }^{\dot{\alpha }})`$ together with using (49) and (50) yield zero brackets between the components $`\omega _\alpha `$ and $`\omega _\beta `$ of the same chirality
$$\{\omega _\alpha ,\omega _\beta \}=i\overline{\eta }^2\phi _{\alpha \beta }0,\{\overline{\omega }_{\dot{\alpha }},\overline{\omega }_{\dot{\beta }}\}=0,$$
(54)
but yields zero brackets with $`\overline{\omega }_{\dot{\beta }}`$
$$\{\omega _\alpha ,\overline{\omega }_{\dot{\beta }}\}=i\eta \overline{\eta }\phi _{\alpha \dot{\beta }}.$$
(55)
Comparing the latter bracket with the bracket (19) for $`\theta _\alpha `$ and $`\overline{\theta }_{\dot{\beta }}`$ we get the relation
$$\begin{array}{c}\{\omega _\alpha ,\overline{\omega }_{\dot{\beta }}\}=4\eta \overline{\eta }\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\}\end{array}$$
(56)
pointing out a connection of the $`(\omega ,\overline{\omega })`$noncommutativity with the $`(\theta ,\overline{\theta })`$-nonanti-
commutativity. On the other side, it shows the connection of the twistor complex structure with supersymmetry. Therefore, the choice of $`(\theta ,\overline{\theta })`$-nonanticommutative bracket induces the $`(\omega ,\overline{\omega })`$noncommutative bracket. Such a correlation of the spin complex structure with supersymmetry and non(anti)commutativity deserves more carefull studying.
## 7 Lorentz invariant brackets in higher dimensions
The brackets (17)-(19) get more compact form in the Majorana representation
$$\nu _a=\left(\genfrac{}{}{0pt}{}{\nu _\alpha }{\overline{\nu }^{\dot{\alpha }}}\right),\theta _a=\left(\genfrac{}{}{0pt}{}{\theta _\alpha }{\overline{\theta }^{\dot{\alpha }}}\right),C^{ab}=\left(\begin{array}{cc}\epsilon ^{\alpha \beta }& 0\\ 0& \overline{\epsilon }_{\dot{\alpha }\dot{\beta }}\end{array}\right),\chi ^a=C^{ab}\chi _b$$
(57)
for the considered Weyl spinors $`\nu _\alpha ,\theta _\alpha `$, their c.c. and the charge conjugation matrix $`C^{ab}`$. Then the P.B’s. (16)-(19) are presented in a form suitable for generalizations. The P.B’s. (16)-(18) take the form
$$\{\nu _a,\nu _b\}=0,\{\theta _a,\nu _b\}=0,\{x_m,\nu _a\}=\phi _m\nu _a,$$
(58)
where the real vectors $`x_m`$ and $`\phi _m`$ are defined by the relations
$$\begin{array}{c}x_m=\frac{1}{2}(\stackrel{~}{\sigma }_m)^{\dot{\alpha }\beta }x_{\beta \dot{\alpha }},x_{\alpha \dot{\beta }}=(\sigma ^m)_{\alpha \dot{\beta }}x_m,\\ \phi _m=\frac{1}{2}(\stackrel{~}{\sigma }_m)^{\dot{\alpha }\beta }\phi _{\beta \dot{\alpha }}\frac{1}{4}(\overline{\nu }\gamma _m\nu )\end{array}$$
(59)
and $`\gamma _m`$ are the Dirac matrices in the Majorana representation.
To rewrite the rest of the P.B’s. in the Majorana representation it is convenient to change the Majorana spinor $`\nu _a`$ by other Majorana spinor $`\lambda _a`$
$$\lambda _a=\left(\genfrac{}{}{0pt}{}{\lambda _\alpha }{\overline{\lambda }^{\dot{\alpha }}}\right)(\gamma _5\nu )_a,(\gamma _5)_a{}_{}{}^{b}=\left(\begin{array}{cc}i\delta _\alpha ^\beta & 0\\ 0& i\delta _{\dot{\beta }}^{\dot{\alpha }}\end{array}\right)$$
(60)
preserving the form of the P.B’s. (58). In terms of the real Majorana spinor $`\lambda _a`$ and the composed vectors $`\phi _m`$ and $`\psi _m`$
$$\phi _m=\frac{1}{4}(\overline{\lambda }\gamma _m\lambda ),\psi _m=\frac{1}{2}(\stackrel{~}{\sigma }_m)^{\dot{\alpha }\alpha }\psi _{\alpha \dot{\alpha }}\frac{1}{2}(\overline{\theta }\gamma _m\lambda )$$
(61)
the P.B’s. (16)-(19) of the primordial coordinates $`x_m,\theta _a,\lambda _a`$ are presented as follow
$$\begin{array}{c}\{\lambda _a,\lambda _b\}=0,\{\theta _a,\lambda _b\}=0,\{x_m,\lambda _a\}=\phi _m\lambda _a,\\ \{x_m,x_n\}=i\psi _n\psi _m,\{x_m,\theta _a\}=\frac{1}{2}\psi _m\lambda _a,\{\theta _a,\theta _b\}=\frac{i}{4}\lambda _a\lambda _b.\end{array}$$
(62)
The P.B’s. of the secondary composite vectors $`\psi _m`$ and $`\phi _m`$ (61) between themselves and with the primordial coordinates are presented in the form
$$\begin{array}{c}\{x_m,\psi _n\}=\phi _m\psi _n+\phi _n\psi _m,\{\psi _m,\theta _b\}=\frac{i}{2}\phi _m\lambda _b,\{\psi _m,\lambda _a\}=0,\\ \{\psi _m,\psi _n\}=i\phi _m\phi _n,\{\psi _m,\phi _n\}=0\end{array}$$
(63)
and respectively
$$\begin{array}{c}\{x_m,\phi _n\}=2\phi _m\phi _n,\{\theta _a,\phi _m\}=\{\lambda _a,\phi _m\}=\{\phi _m,\phi _n\}=0.\end{array}$$
(64)
The Poisson brackets (62)-(64) originally derived for $`D=4`$ remain to be valid in $`D`$dimensional spaces with $`D=2,3,4(mod8)`$, where the Majorana spinors exist. Using the arguments given in the Section 5 one can restore the Moyal brackets originated from the P.B’s. (62)-(64) in the higher dimensions by the simple change $`\{,\}[,]_{}`$.
## 8 Other Lorentz invariant brackets with one spinor
Using the Majorana spinor $`\nu _a`$ one can constuct other simple supersymmetric and Lorentz invariant brackets. One of the possible invariant Poisson bracket might be
$$\begin{array}{c}\{F,G\}=F[\frac{i}{4}(\stackrel{}{D}\stackrel{}{D}+\stackrel{}{\overline{D}}\stackrel{}{\overline{D}})+c(\stackrel{}{}\stackrel{}{\mathrm{\Delta }}\stackrel{}{\mathrm{\Delta }}\stackrel{}{})]G\end{array}$$
(65)
which changes the brackets (19) to the brackets
$$\begin{array}{c}\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\}=i(\phi _{\alpha \beta }\overline{\theta }_{\dot{\alpha }}\overline{\theta }_{\dot{\beta }}+\overline{\phi }_{\dot{\alpha }\dot{\beta }}\theta _\alpha \theta _\beta ),\\ \{x_{\alpha \dot{\alpha }},\theta _\beta \}=\frac{1}{2}\phi _{\alpha \beta }\overline{\theta }_{\dot{\alpha }}\{x_{\alpha \dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\frac{1}{2}\overline{\phi }_{\dot{\alpha }\dot{\beta }}\theta _\alpha ,\\ \{\theta _\alpha ,\theta _\beta \}=\frac{i}{4}\phi _{\alpha \beta },\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\}=0,\{\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\frac{i}{4}\overline{\phi }_{\dot{\alpha }\dot{\beta }}.\end{array}$$
(66)
We see that the bracket of the spinor components $`\theta _\alpha `$, $`\overline{\theta }_{\dot{\alpha }}`$ having opposite chiralities is not deformed and remains equal zero. Moreover, the brackets $`x_{\alpha \dot{\alpha }}`$ with $`\theta _\alpha `$ and $`\overline{\theta }_{\dot{\alpha }}`$ don’t preserve the chiralities of $`\theta `$ and $`\overline{\theta }`$ spinors. As a result, one can find breaking of the Jacoby identity for the $`(x,\theta ,\overline{\theta })`$-cycle, because of the relation
$$Cycle\{\{x_{\alpha \dot{\alpha }},\theta _\beta ,\overline{\theta }_{\dot{\gamma }}\}\}=\{\{x_{\alpha \dot{\alpha }},\theta _\beta ,\overline{\theta }_{\dot{\gamma }}\}+\{\{x_{\alpha \dot{\alpha }},\overline{\theta }_\gamma \},\overline{\theta }_\beta \}=\frac{i}{4}\phi _{\alpha \beta }\overline{\phi }_{\dot{\alpha }\dot{\beta }}.$$
(67)
So, we conclude that the Lorentz invariant P.B. (65) has to be excluded. But, the next supersymmetric and Lorentz invariant Poisson bracket
$$\begin{array}{c}\{F,G\}=F[\frac{i}{4}(\stackrel{}{D}\stackrel{}{\overline{D}}+\stackrel{}{\overline{D}}\stackrel{}{D})+\frac{1}{2}(\stackrel{}{}\stackrel{}{\mathrm{\Delta }}\stackrel{}{\mathrm{\Delta }}\stackrel{}{})]G\end{array}$$
(68)
is proved to be selfconsistent and yields the following invariant Poisson brackets
$$\begin{array}{c}\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\}=i(\phi _{\alpha \dot{\beta }}\overline{\theta }_{\dot{\alpha }}\theta _\beta \phi _{\beta \dot{\alpha }}\overline{\theta }_{\dot{\beta }}\theta _\alpha ),\\ \{x_{\alpha \dot{\alpha }},\theta _\beta \}=\frac{1}{2}\phi _{\beta \dot{\alpha }}\theta _\alpha ,\{x_{\alpha \dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=\frac{1}{2}\phi _{\alpha \dot{\beta }}\overline{\theta }_{\dot{\alpha }},\\ \{\theta _\alpha ,\theta _\beta \}=\{\overline{\theta }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=0,\{\theta _\alpha ,\overline{\theta }_{\dot{\beta }}\}=\frac{i}{4}\phi _{\alpha \dot{\beta }}\end{array}$$
(69)
added by the brackets
$$\begin{array}{c}\{\nu _\alpha ,\nu _\beta \}=\{\nu _\alpha ,\overline{\nu }_{\dot{\beta }}\}=\{\overline{\nu }_\alpha ,\overline{\nu }_{\dot{\beta }}\}=0,\\ \{\nu _\alpha ,\theta _\beta \}=\{\nu _\alpha ,\overline{\theta }_{\dot{\beta }}\}=\{\overline{\nu }_{\dot{\alpha }},\theta _\beta \}=\{\overline{\nu }_{\dot{\alpha }},\overline{\theta }_{\dot{\beta }}\}=0,\\ \{x_{\alpha \dot{\alpha }},\nu _\beta \}=\frac{1}{2}\phi _{\alpha \dot{\alpha }}\nu _\beta ,\{x_{\alpha \dot{\alpha }},\overline{\nu }_{\dot{\beta }}\}=\frac{1}{2}\phi _{\alpha \dot{\alpha }}\overline{\nu }_{\dot{\beta }},\end{array}$$
(70)
One can see that in contrast to the deformation (65) the new deformation (68) generates zero P.B’s. for the $`\theta _a`$ components with the same chiralities. The P.B’s. (69),(70) satisfy the Jacobi identities and deserve to be studied in physical applications. The proof of the Jacobi identities for the P.B’s. (69) and (70) is analogous to the proof presented in the Section 4.
The P.B’s. (69) may be presented in the vector form as follows
$$\begin{array}{c}\{x_m,x_n\}=\frac{i}{4}(\chi _m\overline{\chi }_n\chi _n\overline{\chi }_m),\\ \{x_m,\theta _\beta \}=\frac{1}{4}\overline{\chi }_m\nu _\beta ,\{x_m,\overline{\theta }_{\dot{\beta }}\}=\frac{1}{4}\chi _m\overline{\nu }_{\dot{\beta }},\\ \{\theta _a,\theta _b\}=\frac{i}{8}(\nu _a^{(+)}\nu _b^{()}+\nu _b^{(+)}\nu _a^{()}),\end{array}$$
(71)
where we introduced the complex Grasssmannian vector $`\chi _m`$ with the real and imaginary parts presented by $`\psi _{1m},\psi _{2m}`$ and the chiral components $`\theta ^{(\pm )}`$ and $`\nu ^{(\pm )}`$
$$\begin{array}{c}\chi _m(\nu \sigma _m\overline{\theta })\overline{\nu }\gamma _m\frac{1+i\gamma _5}{2}\theta \psi _{1m}+i\psi _{2m},\\ \overline{\chi }_m(\chi _m)^{}=\overline{\nu }\gamma _m\frac{1i\gamma _5}{2}\theta ,\psi _{1m}\frac{1}{2}(\overline{\theta }\gamma _m\nu ),\psi _{2m}\frac{1}{2}(\overline{\theta }\gamma _m\gamma _5\nu ),\\ \theta ^{(\pm )}\frac{1}{2}(1\pm i\gamma _5)\theta ,\nu ^{(\pm )}\frac{1}{2}(1\pm i\gamma _5)\nu .\end{array}$$
(72)
Then the P.B’s. (71) are presented in the form directly generalizing the P.B’s. (62)
$$\begin{array}{c}\{x_m,x_n\}=\frac{i}{2}(\psi _{1m}\psi _{1n}+\psi _{2m}\psi _{2n}),\\ \{x_m,\theta _a\}=\frac{1}{4}(\psi _{1m}\nu _a+\psi _{2m}\lambda _a),\\ \{\theta _a,\theta _b\}=\frac{i}{8}(\nu _a\nu _b+\lambda _a\lambda _b),\end{array}$$
(73)
where $`\lambda _a(\gamma _5\nu )_a`$ (60). Comparing (73) with (62) we observe that the change of the P.B. (12) by (68) is equivalent to the complexification of the real Grassmannian vector $`\psi _m`$ (61) accompanied by the appearance of the spinors $`\nu _a`$ and $`(\gamma _5\nu )_a)`$ in the r.h.s. of (73).
The P.B’s. (68) may be generalized to the case of extended supersymmetries with $`N>1`$
$$\begin{array}{c}\{F,G\}=F[\frac{i}{4}(\stackrel{}{D_i}\stackrel{}{\overline{D}^i}+\stackrel{}{\overline{D}^i}\stackrel{}{D_i})+\frac{1}{2}(\stackrel{}{}\stackrel{}{\mathrm{\Delta }}\stackrel{}{\mathrm{\Delta }}\stackrel{}{})]G,\end{array}$$
(74)
where $`D_i=\nu _\alpha D_i^\alpha `$ and $`\overline{D}^i=\overline{\nu }_{\dot{\alpha }}\overline{D}^{\dot{\alpha }i},(i=1,2,..,N)`$. The P.B’s. (74) generate the following brackets for the primordial space-time (super)coordinates
$$\begin{array}{c}\{x_{\alpha \dot{\alpha }},x_{\beta \dot{\beta }}\}=i(\phi _{\alpha \dot{\beta }}\overline{\theta }_{\dot{\alpha }i}\theta _\beta ^i\phi _{\dot{\alpha }\beta }\overline{\theta }_{\dot{\beta }i}\theta _\alpha ^i),\\ \{x_{\alpha \dot{\alpha }},\theta _\beta ^i\}=\frac{1}{2}\phi _{\dot{\alpha }\beta }\theta _\alpha ^i,\{x_{\alpha \dot{\alpha }},\overline{\theta }_{\dot{\beta }i}\}=\frac{1}{2}\phi _{\alpha \dot{\beta }}\overline{\theta }_{\dot{\alpha }i},\\ \{\theta _\alpha ^i,\theta _\beta ^k\}=\{\overline{\theta }_{\dot{\alpha }i},\overline{\theta }_{\dot{\beta }k}\}=0,\{\theta _\alpha ^i,\overline{\theta }_{\dot{\beta }k}\}=\frac{i}{4}\phi _{\alpha \dot{\beta }}.\delta _{}^{i}{}_{k}{}^{}\end{array}$$
(75)
The rest of the P.B’s for $`x_{\alpha \dot{\alpha }},\nu _a,\theta _\alpha ^i`$ coincides with the P.B’s. (70).
## 9 More spinors - more Lorentz invariant brackets
The above consideration to construct the Lorentz covariant Poisson and Moyal brackets was restricted by the simplest case of one additional spinor coordinate which resulted in the appearance of the supersymmetric derivatives $`D^\alpha ,\overline{D}^{\dot{\alpha }}`$ and $`^{\alpha \dot{\alpha }}`$ (3) in the considered P.B’s. only in the form of the scalars (5). Using only these scalars for the construction of the invariant P.B. (12) restricts the class of admissible motions in superspace. To extend this class still preserving the Lorentz invariance and supersymmetry one can introduce additional independent spinor coordinates. In the case $`D=4`$ it is enough to add only one new spinor coordinate $`\mu _\alpha `$, because $`\mu _\alpha `$ and $`\nu _\alpha `$ form the complete spinorial basis and may be identified with the Newman-Penrose dyad
$$\mu ^\alpha \nu _\alpha \mu ^\alpha \epsilon _{\alpha \beta }\nu ^\beta =1,\mu _\alpha \nu _\beta \mu _\beta \nu _\alpha =\epsilon _{\alpha \beta }.$$
(76)
Then one can form four independent Lorentz invariant supersymmetric differential operators
$$\begin{array}{c}D^{(\nu )}=\nu _\alpha D^\alpha ,\overline{D}^{(\nu )}=\overline{\nu }_{\dot{\alpha }}\overline{D^{\dot{\alpha }}},D^{(\mu )}=\nu _\alpha D^\alpha ,\overline{D}^{(\mu )}=\overline{\mu }_{\dot{\alpha }}\overline{D^{\dot{\alpha }}},\end{array}$$
(77)
two of which $`D^{(\nu )},\overline{D}^{(\nu )}`$ coincide with the operators $`D,\overline{D}`$ (5). Their linear combinations
$$\begin{array}{c}D_\pm ^{(\nu )}D^{(\nu )}\pm \overline{D}^{(\nu )},D_\pm ^{(\mu )}D^{(\mu )}\pm \overline{D}^{(\mu )},\end{array}$$
(78)
form four Lorentz invariant and supersymmetric supersubalgebras
$$\begin{array}{c}[D_\pm ^{(\nu )},D_\pm ^{(\nu )}]_+=8i^{(\nu )},[D_\pm ^{(\nu )},^{(\nu )})]=[^{(\nu )},^{(\nu )}]=0,^{(\nu )}(\nu _\alpha \overline{\nu }_{\dot{\alpha }}^{\alpha \dot{\alpha }}),\\ [D_\pm ^{(\mu )},D_\pm ^{(\mu )}]_+=8i^{(\nu )},[D_\pm ^{(\mu )},^{(\mu )})]=[^{(\mu )},^{(\mu )}]=0,^{(\mu )}(\mu _\alpha \overline{\mu }_{\dot{\alpha }}^{\alpha \dot{\alpha }}),\end{array}$$
(79)
which are connected by the P.B. relations
$$\begin{array}{c}[D_\pm ^{(\nu )},D_\pm ^{(\mu )}]_+=4i^{(+)},^{(+)}(\nu _\alpha \overline{\mu }_{\dot{\alpha }}+\mu _\alpha \overline{\nu }_{\dot{\alpha }})^{\alpha \dot{\alpha }},\\ [D_\pm ^{(\nu )},D_{}^{(\mu )}]_+=\pm 4i^{()},^{()}(\nu _\alpha \overline{\mu }_{\dot{\alpha }}\mu _\alpha \overline{\nu }_{\dot{\alpha }})^{\alpha \dot{\alpha }}.\end{array}$$
(80)
It is easy to see that the Lorentz invariant and supersymmetric differential operators $`D_\pm ^{(\nu )},D_\pm ^{(\mu )},^{(\nu )},^{(\mu )},^{()}`$ describe whole class of admissible motions in the superspace and together with the extended dilatation operator $`\mathrm{\Delta }^{}`$
$$\begin{array}{c}\mathrm{\Delta }^{}=(\nu _\alpha \frac{}{\nu _\alpha }+\overline{\nu }_{\dot{\alpha }}\frac{}{\overline{\nu }_{\dot{\alpha }}})(\mu _\alpha \frac{}{\mu _\alpha }+\overline{\mu }_{\dot{\alpha }}\frac{}{\overline{\mu }_{\dot{\alpha }}})\end{array}$$
(81)
preserving the condition (76) may be used as invariant bilding blocks for the construction of more general Lorentz invariant supersymmetric Poisson and Moyal brackets.
Then the Lorentz invariant and supersymmetric Poisson bracket
$$\begin{array}{c}\{F,G\}=F[\frac{i}{4}(\stackrel{}{D}_{}^{(\nu )}\stackrel{}{D}_{}^{(\nu )}+\stackrel{}{D}_{}^{(\mu )}\stackrel{}{D}_{}^{(\mu )})+c(\stackrel{}{^{(\nu )}}+\stackrel{}{^{(\mu )}})\stackrel{}{\mathrm{\Delta }^{}}\stackrel{}{\mathrm{\Delta }^{}}(\stackrel{}{^{(\nu )}}+\stackrel{}{^{(\mu )}})]G.\end{array}$$
(82)
might be considered as a candidate for the generalizations of (12). The P.B. (82) yields the following coordinate P.B’s.
$$\begin{array}{c}\{x_m,x_n\}=i(\psi _n^{(\nu )}\psi _m^{(\nu )}+\psi _n^{(\mu )}\psi _m^{(\mu )}),\\ \{x_m,\theta _a\}=\frac{1}{2}(\psi _m^{(\nu )}\lambda _a^{(\nu )}+\psi _m^{(\mu )}\lambda _a^{(\mu )}),\\ \{\theta _a,\theta _b\}=\frac{i}{4}(\lambda _a^{(\nu )}\lambda _b^{(\nu )}+\lambda _a^{(\mu )}\lambda _b^{(\mu )})\end{array}$$
(83)
for the primordial coordinates $`x_m`$ and $`\theta _a`$, where the additional Majorana spinor $`\lambda _a^{(\mu )}`$ and the Grassmannian vector $`\psi _n^{(\mu )}`$ are defined by the relations
$$\begin{array}{c}\psi _n^{(\nu )}\psi _n,\lambda _a^{(\nu )}\lambda _a,\psi _n^{(\mu )}\frac{1}{2}(\overline{\theta }\gamma _n\lambda ^{(\mu )}),\lambda _a^{(\mu )}(\gamma _5\mu )_a\end{array}$$
(84)
The primordial Majorana spinors $`\lambda _a^{(\nu )}`$ and $`\lambda _a^{(\mu )}`$ have zero P.B’s. between themselves and with $`\theta _a,\psi _m^{(\nu )},\psi _n^{(\mu )}`$, but non zero P.B’s. with $`x_m`$
$$\begin{array}{c}\{x_m,\lambda _a^{(\nu )}\}=c(\phi _m^{(\nu )}+\phi _m^{(\mu )})\lambda _a^{(\nu )}),\{x_m,\lambda _a^{(\mu )}\}=c(\phi _m^{(\nu )}+\phi _m^{(\mu )})\lambda ^{(\mu )}_a,\\ \phi _m^{(\nu )}\phi _m,\phi _m^{(\mu )}\frac{1}{4}(\overline{\lambda }^{(\mu )}\gamma _m\lambda ^{(\mu )}),\end{array}$$
(85)
where the real constant $`c`$ has to be defined from the solution of the Jacobi identities. We intend to give back to the studying this P.B. and other possible generalizations of the P.B’s (62) in other place.
## 10 Conclusion
It was shown here that the extension of the $`N=1`$ superspace ($`x_m,\theta _a`$) by commuting Majorana spinors may be used for the construction of supersymmetric and Lorentz invariant Poisson and Moyal brackets generating deformed non(anti)commutative relations for space-time (super)coordinates. To make clear the proposal we elaborate the case of one additional spinor $`\lambda _a`$ extending the standard $`N=1`$ superspace to the non(anti)commutative superspaces free of background fields. The corresponding Lorentz invariant and supersymmetric coordinate brackets were presented. It was established that noncommutativity of $`x_m`$ with $`x_n`$ and $`\theta _a`$ is measured by the real or complex Grassmannian vectors $`\psi _m`$ composed from $`\theta _a`$ and $`\lambda _a`$, which are known as dynamical variables describing the spin degrees of freedom of spinning string or particle. At the same time, the nonanticommutativity of the $`\theta _a`$ componets between themselves is measured by only additional spinor or its chiral components. These results hint on a hidden spinorial structure of space-time encoded in the Penrose twistor picture and its supersymmetric extensions as an alternative source of (super)coordinate non(anti)commutativity. In the simplest case corresponding to the P.B’s. (62) a correlation between the spinorial structure and supergravity fields may be schematically illustrated by the correspondence:
$$\begin{array}{c}\theta _{mn}i\psi _m\psi _n,C_{ab}\lambda _a\lambda _b,\mathrm{\Psi }_m^a\psi _m\lambda ^a,\end{array}$$
(86)
where $`B_{mn}=\theta _{mn}^1`$, $`C_{ab}`$ and $`\mathrm{\Psi }_m^a`$ are constant antisymmetric field, the graviphoton and the gravitino respectively. The map (86) transforms the well-known field dependent bracket relations into the brackets (62) and vice versa. On the other hand, such a correspondence hints on a connection with the known Feynman-Wheeler picture and its superymmeric generalization , where the Maxwell supermultiplet fields arise as secondary objects constructed from the superspace coordinates. We outlined also a way to construction of more general supersymmetric Lorentz invariant brackets for the cases of $`N`$ extended supersymmetries and additional spinor coordinates using Lorentz invariant supersymmetric derivatives generalizing (5). Studying that generalizations and the corresponding deformations of quantum field models are under consideration.
## 11 Acknowledgements
The author thanks Fysikum at the Stockholm University for kind hospitality and I. Bengtsson and B. Sundborg for useful discussions. The work was partially supported by the grant of the Royal Swedish Academy of Sciences and by the SFFR of Ukraine under Project 02.07/276.
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# SDSS J0246-0825: A New Gravitationally Lensed Quasar from the Sloan Digital Sky Survey1footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS 5-26555. These observations are associated with HST program 9744.
## 1 Introduction
Since the discovery of the first lensed quasar, Q0957+561 (Walsh, Carswell, & Weymann, 1979), approximately 80 lensed quasars have been discovered to date. Gravitationally lensed quasars are now known to be an useful tool for cosmology and astrophysics, such as measurements of dark energy (Turner, 1990; Fukugita, Futamase, & Kasai, 1990; Chiba & Yoshii, 1999; Chae et al., 2002), direct probes of the Hubble constant (Refsdal, 1964; Bernstein & Fischer, 1999; Koopmans et al., 2003), and the study of lensing galaxies (Kochanek, 1991; Keeton, 2001a; Oguri, 2002; Koopmans & Treu, 2003; Rusin, Kochanek, & Keeton, 2003; Ofek, Rix, & Maoz, 2003). In order to use lensed quasars as these cosmological and astrophysical tests, large and homogeneous samples of gravitational lenses are necessary, either to ensure good statistical precisions and to provide sufficient number of lensing systems which are suitable for measurements of time delays and/or for studies of the lensing galaxies. Although several previous large and homogeneous lensed quasar surveys, e.g., Jodrell/VLA Astrometric Survey and the Cosmic Lens All-Sky Survey (Myers et al., 2003; Browne et al., 2003) and the Hubble Space Telescope Snapshot Survey (Maoz et al., 1992, 1993a, 1993b; Bahcall et al., 1992), are of considerable value, discoveries of more lensed quasars in larger surveys, such as the Sloan Digital Sky Survey (SDSS; York et al., 2000; Stoughton et al., 2002; Abazajian et al., 2003, 2004) will greatly improve the accuracy of these cosmological and astrophysical tests. The final number of spectroscopically confirmed quasars in the SDSS is expected to be about $`10^5`$, thus, if we assume the lensing rate of 0.1% (Turner, Ostriker, & Gott, 1984), the expected number of gravitationally lensed quasars in the SDSS will be approximately $`10^2`$; the SDSS could double the number of lensed quasars. Indeed, several newly discovered lensed quasars in the SDSS data have already been reported (Inada et al., 2003a, b, c; Morgan, Snyder, & Reens, 2003; Johnston et al., 2003; Pindor et al., 2004; Oguri et al., 2004a, b).
In this paper, we report the discovery of another lensed quasar, SDSS J024634.11$``$082536.3 (SDSS J0246$``$0825). This object was selected as a lensed quasar candidate (an “extended” quasar) from about 50,000 SDSS quasars using the algorithm described in Inada et al. (2003b), Pindor et al. (2004), Oguri et al. (2004a) and Oguri et al. (2004b). Since the SDSS imaging and spectroscopic data of SDSS J0246$``$0825 are unresolved, we conducted additional observations using the W. M. Keck Observatory’s Keck II telescope, the Magellan Consortium’s Walter Baade 6.5-m telescope (WB6.5m)<sup>17</sup><sup>17</sup>17The first telescope of the Magellan Project; a collaboration between the Observatories of the Carnegie Institution of Washington (OCIW), University of Arizona, Harvard University, University of Michigan, and Massachusetts Institute of Technology (MIT) to construct two 6.5 Meter optical telescopes in the southern hemisphere., and the Hubble Space Telescope (HST). These observations demonstrate that SDSS J0246$``$0825 is a two-image lens (the image separation is 1$`\stackrel{}{\mathrm{.}}`$04) of a source quasar at $`z=1.68`$, and that the likely lensing galaxy has been detected.
The structure of this paper is as follows. Section 2 describes both the photometric and spectroscopic follow-up observations. In §3, we show mass modeling of this lensing system. Finally, we present a summary and give a conclusion in §4. We assume a Lambda-dominate cosmology (Spergel et al. 2003; matter density $`\mathrm{\Omega }_M=0.27`$, cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }=0.73`$, and the Hubble constant $`H_0=70`$ km sec<sup>-1</sup> Mpc<sup>-1</sup>), and we use AB magnitude system (Oke & Gunn, 1983; Fukugita et al., 1996) throughout the paper.
## 2 Follow-up Observation and Data Analysis
### 2.1 SDSS Observations
The SDSS is a project to conduct in parallel a photometric survey (Gunn et al., 1998; Lupton, Gunn, & Szalay, 1999) and a spectroscopic survey (Blanton et al., 2003) of 10,000 square degrees of the sky centered approximately on the North Galactic Pole, using dedicated wide-field ($`3^{}`$ field of view) 2.5-m telescope at Apache Point Observatory in New Mexico, USA. The astrometric positions are accurate to better than about $`0\stackrel{}{\mathrm{.}}1`$ rms per coordinate (Pier et al., 2003) and the photometric errors are less than about 0.03 magnitude (Hogg et al., 2001; Smith et al., 2002; Ivezić et al., 2004). Photometric observations are conducted using five optical broad bands (Fukugita et al., 1996). After automated data processing by the photometric pipeline (Lupton et al., 2001), quasar, galaxy and luminous red galaxy candidates are selected by the target selection pipelines (Eisenstein et al., 2001; Richards et al., 2002; Strauss et al., 2002). Spectra of these candidates are obtained with a multi-fiber spectrograph covering 3800 Å to 9200 Å (resolution R $``$ 1800).
The $`ugriz`$ SDSS images of SDSS J0246$``$0825 (sky and bias subtracted and flat-field corrected) are shown in Figure 1. Two stellar components are marginally-resolved/unresolved in the SDSS images; total magnitudes (in $``$2$`\stackrel{}{\mathrm{.}}`$0 aperture radius) of the two quasar components and the lensing galaxy (and an unknown extended object, see below) are 18.21$`\pm `$0.02, 17.99$`\pm `$0.01, 17.95$`\pm `$0.01, 17.60$`\pm `$0.01, and 17.58$`\pm `$0.02 in $`u`$, $`g`$, $`r`$, $`i`$, and $`z`$, respectively. These errors in the magnitudes are statistical errors. We show the field image (SDSS $`i`$ band image) around SDSS J0246$``$0825, including a nearby star (star S) which is seen in both the WB6.5m images and the SDSS images, in Figure 2. We use star S to estimate a Point Spread Function (PSF) and to calibrate magnitudes. SDSS J0246$``$0825 was targeted as a quasar candidate by the target selection pipelines, and the spectrum was taken by the SDSS spectrograph. The spectrum data are also unresolved due to the spatial resolution of the SDSS spectroscopy, but it clearly shows the quasar emission lines at $`z=1.686\pm 0.001`$.
### 2.2 Follow-up Observation using the Walter Baade Telescope
High-resolution images of SDSS J0246$``$0825 in the $`u`$, $`g`$, $`r`$, and $`i`$ filters were obtained with the Magellan Instant Camera (MagIC, a 2048$`\times `$2048 CCD camera; pixel scale is 0$`\stackrel{}{\mathrm{.}}`$069) on the WB6.5m telescope, on 2002 February 13. The seeing (FWHM) was $``$0$`\stackrel{}{\mathrm{.}}`$5–0$`\stackrel{}{\mathrm{.}}`$6, and the exposure time was 120 sec in each band. The $`u`$, $`g`$, $`r`$, and $`i`$ MagIC images of SDSS J0246$``$0825 are shown in Figure 3. Two stellar components are clearly seen in all these images; we named the two stellar components A and B, with A being the brighter component. The flux ratios between components B and A (estimated by the single Gaussian fit) are 0.32, 0.31, 0.33, and 0.34 in $`u`$, $`g`$, $`r`$, and $`i`$, respectively.
The PSFs derived from star S were subtracted from the original $`i`$ band MagIC image. The peak flux and the center coordinates of components A and B in the MagIC $`i`$ band image were calculated by the single Gaussian fit. The results agree with those obtained by the imexamine task in IRAF<sup>18</sup><sup>18</sup>18IRAF is the Image Reduction and Analysis Facility, a general purpose software system for the reduction and analysis of astronomical data. IRAF is written and supported by the IRAF programming group at the National Optical Astronomy Observatories (NOAO) in Tucson, Arizona. NOAO is operated by the Association of Universities for Research in Astronomy (AURA), Inc. under cooperative agreement with the National Science Foundation.. The $`i`$ band PSF-subtracted image is shown in the right panel of Figure 4. Although it is not obvious in Figure 4, we find a faint, extended object ($``$ 21.8 mag in 0$`\stackrel{}{\mathrm{.}}`$5 aperture radius) between the two stellar components after the PSF subtraction. For further discussion of this object, see the Keck results in §2.3 and/or the HST results in §2.4. The angular separation (in the $`i`$ band image) between components A and B is 1$`\stackrel{}{\mathrm{.}}`$042$`\pm `$0$`\stackrel{}{\mathrm{.}}`$003. The $`ugri`$ band PSF magnitudes of components A and B and the $`i`$ band aperture magnitude of the extended objects between components A and B (component G, see below) are summarized in Table 1. Here we used star S (see Figure 2) as a photometric star<sup>19</sup><sup>19</sup>19 Note that there are small (a few percent) differences between the filter response functions of the WB6.5m and those of the SDSS. .
### 2.3 Follow-up Observation using the Keck Telescope
We conducted a near-infrared imaging observation of SDSS J0246$``$0825 using the Keck I telescope, on 2002 August 21. Although the night was not photometric, we obtained $`K^{}`$ band data with the Near InfraRed Camera (NIRC; Matthews & Soifer, 1994). The seeing (FWHM) was variable, averaging about 0$`\stackrel{}{\mathrm{.}}`$4, and the total exposure time was 900 sec. The deconvolved $`K^{}`$ band image is shown in Figure 5; the deconvolution was done using the method which is described in Magain, Courbin, & Sohy (1998). The pixel size is half that of the raw image, 0$`\stackrel{}{\mathrm{.}}`$075. In Figure 5, we can clearly see an extended object (named component G) at the same position of the extended object found in the PSF-subtracted MagIC $`i`$ band image (we can see component G more clearly in the PSF-subtracted images, the inset of Figure 5). This object is probably the lensing galaxy of this lensed quasar system; this hypothesis is further supported by the spectroscopic observation of the two stellar components using the Keck telescope and the high-resolution imaging observation using the HST (see §2.4). The flux ratio between components B and A is 0.29 in the $`K^{}`$ band image; this is consistent with the mean flux ratio of the optical bands.
Spectra of the two components of SDSS J0246$``$0825 were acquired using the Keck II telescope on 2002 December 4. We conducted a spectroscopic follow-up observation using the Echellette Spectrograph and Imager (ESI; Sutin, 1997; Sheinis et al., 2002) with the MIT-LL 2048$`\times `$4096 CCD camera. We used the Echellette mode. The spectral range covers 3900 Å to 11,000 Å, and the spectral resolution is R $``$ 27000. We set the slit width 1$`\stackrel{}{\mathrm{.}}`$0 and the slit orientation so that the two stellar components (components A and B; see Figure 3) were on the slit (thus, component G should be also on the slit, but we were not able to extract the spectrum of component G because it is much fainter than the quasar components). The exposure time was 900 sec. Since the seeing (FWHM) was about 0$`\stackrel{}{\mathrm{.}}`$7 and the angular separation between the two components is only 1$`\stackrel{}{\mathrm{.}}`$04, we can clearly see there are two objects on the slit but these are not well resolved. Therefore we extracted the spectra using a method of summing the fluxes in a window around the position of the brighter component and in another window shifted by 1$`\stackrel{}{\mathrm{.}}`$04. Sky was subtracted using neighboring windows on either side of the trace. There are bad columns on the ESI chip between 4460 Å and 4560 Å, and therefore, we excluded the data of this bad region from the spectra. The binned spectra of components A and B taken with ESI on Keck II are shown in Figure 6. In both spectra, C IV, C III\], and Mg II emission lines are clearly seen at the same wavelength positions. The redshifts of components A and B are 1.6820$`\pm `$0.0001 and 1.6816$`\pm `$0.0004, which were calculated by fitting the Voigt functions with the C IV emission lines. The velocity difference between the two quasar components is 20 km sec$`{}_{}{}^{1}\pm 40`$ km sec<sup>-1</sup>. The redshifts and the widths of the emission lines are summarized in Table 2. In addition to the identical redshift, the spectral energy distributions themselves are also quite similar. The bottom line in Figure 6 shows the spectral flux ratio between components B and A. We find that this is almost constant for a wide range of wavelengths and this is almost consistent with the mean flux ratio (0.33) of the $`ugri`$ images. A large difference is seen around the C IV emission line; this might be caused by microlensing of the broad emission line region (Richards et al., 2004).
In addition to the emission lines, we found a Mg II/Mg I absorption system redshifted to $`z=0.724`$, as shown in Figure 7. This absorption system probably arises from a lensing galaxy, both because galaxies can produce such Mg II absorption systems (Bergeron, 1988; Bergeron & Boisse, 1991) and because the Mg II/Mg I absorption system ($`z=0.724`$) is not so far from the half of the angular diameter distance to the quasar components ($`z=1.688`$) that has the maximum lensing efficiency (Ofek et al., 2003). This prediction is further supported by the comparison of the expected magnitude by the mass models with the observed magnitude of the lensing galaxy (see §3). Furthermore, we found two Mg II absorption systems, whose redshifts are 1.537 and 1.531 (both at $``$ 7100 Å), with their Fe II absorption lines (at $``$ 6600 Å, $``$ 6000 Å, and so on) and weaker absorption lines of other elements. Moreover, at $``$ 7100 Å, we found two possible Mg II absorption systems whose redshifts are 1.540 and 1.533, although their Fe II absorption lines and other element absorption lines were not detected in the spectra. However these 4 absorption systems are unlikely to be associated with the lensing system, because the difference between the redshift of the source quasar and that of the absorption systems is small.
### 2.4 Follow-up Observation using the HST
Our final observation of SDSS J0246$``$0825 was conducted using the Hubble Space Telescope, under the “HST Imaging of Gravitational Lenses” program (principal investigator: C. S. Kochanek, proposal ID: Cycle12–9744). We used the Advanced Camera for Surveys (ACS; Clampin, 2000) and the Near Infrared Camera and Multi-Object Spectrometer (NICMOS; Thompson, 1992), installed on the HST. The ACS and NICMOS observations were conducted on 2003 October 13 and 2003 December 18, respectively. We used F555W filter ($`V`$-band) and the F814W filter ($`I`$-band) in the ACS observation and the F160W filter ($`H`$-band) in the NICMOS observation.
The ACS imaging observation consists of two dithered exposures taken in ACCUM mode. Total exposure time of each dithered exposure for the F555W filter was 1094 sec and that for the F814W filter was 1144 sec. The reduced (drizzled and calibrated) images were extracted using the CALACS calibration pipeline which includes the PyDrizzle algorithm. The combined image of the F555W filter and that of the F814W filter are shown in the left panels of Figure 8. As well as the NIRC $`K^{}`$ band image (Figure 5), we find the extended object (component G) between the two quasar components in the ACS images, particularly in the F814W image. The extended nature and the red color of component G suggest that this object is a galaxy. In addition to the identical nature of the spectral energy distributions of components A and B (described in §2.3), the existence of the galaxy strongly supports the lensing hypothesis; we conclude that SDSS 0246$``$0825 is certainly a lensed quasar system. Although the spectrum of component G has not been obtained yet, the color of component G ($`VI2.8`$, estimated from the PSF-subtracted images, the right panels of Figure 8) is consistent with it being a galaxy at $`z=0.724`$ (Fukugita, Shimasaku, & Ichikawa, 1995), which is the redshift of Mg II/Mg I absorption system seen in the spectra of components A and B (see Figure 7).
In addition to the central extended object (component G), we find a highly distorted object (named component C) in the upper left of component A. Similar “ring” features of lensed (host) galaxies have been found in the optical and/or near-infrared band images (e.g., Impey et al., 1998; King et al., 1998; Warren et al., 1996), and therefore this object is probably the lensed host galaxy of the source quasar of SDSS 0246$``$0825. To improve the visibility of this distorted feature, we subtract components A and B using the quasar PSFs produced by the Tiny Tim software (version 6.1a; Krist & Hook, 2003). The PSFs of quasars were constructed with the Keck spectra of components A and B. The result is shown in the right panels of Figure 8. Component C is bright even in the F555W filter image as well as the F814W filter image (see Figure 8); the blue color of component C indicates that the host galaxy might be a starburst galaxy. The (magnification) center of component C is 0$`\stackrel{}{\mathrm{.}}`$50 away from that of component A in the image plane, but the similar phenomenon have been observed in another lensed quasar system (Inada et al., 2005). See §3 for the comparison of this “ring” and mass models.
The NICMOS imaging observation consists of four dithered exposures taken in MULTIACCUM mode, through the F160W filter ($`H`$-band). The exposure time was 640 sec for three dithered exposures and 704 sec for one dithered exposure. The reduced (calibrated) images were extracted using the CALNICA pipeline. The central bad columns of each dithered image were corrected by linear interpolation. The combined image is shown in the left panel of Figure 9. First, we confirm the existence of the central extended object (component G) between components A and B, and the highly distorted object (component C) near component A. In addition, we find a further highly distorted object, which is similar to component C, near component B in the F160W image. To improve the visibility of them, we subtract components A and B using the quasar PSFs produced by the Tiny Tim software. The PSF of a quasar was constructed with an $`\alpha _\nu =0.5`$ power law spectrum. The peak flux and the center coordinates of components A and B in the NICMOS F160W image were calculated using the imexamine task in IRAF. The result is shown in the right panel of Figure 9; one can clearly see the distorted objects form a nearly perfect ring. This prefect ring feature can be also seen in the PSF-subtracted NIRC $`K^{}`$ band image (Figure 5). The PSF magnitudes of components A and B and the aperture magnitude of component G in F160W filter ($`H`$-band) are summarized in Table 1. In estimating the aperture magnitude of component G, we used the aperture radius of 0$`\stackrel{}{\mathrm{.}}`$375 in the PSF subtracted image.
## 3 Mass Modeling
We model the lensed quasar system using two mass models. One is the Singular Isothermal Sphere with external shear (SIS+shear) model, with a projected potential of
$$\psi (r,\theta )=\alpha _er+\frac{\gamma }{2}r^2\mathrm{cos}2(\theta \theta _\gamma ),$$
(1)
where $`\alpha _e`$ is the Einstein radius in arcseconds, $`r`$ and $`\theta `$ are the radial and the angular parts of the angular position on the sky, respectively, and $`\theta _\gamma `$ is the position angle of the shear, measured East of North. The other is the Singular Isothermal Ellipsoid (SIE) model, with a projected potential of
$$\psi (r,\theta )=\alpha _er[1+((1q^2)/(1+q^2))\mathrm{cos}2(\theta \theta _e)]^{1/2},$$
(2)
where $`q`$ is the lens axis ratio and the $`\theta _e`$ is the position angle of the ellipse, measured East of North. Since component A is saturated in both the ACS F555W image and the F814W image, we adopted the positions and flux of components A and B in the MagIC image and the position of component G in the ACS (F814W) image. In order to obtain the position of G relative to A, we used the positions of components B and G in the ACS F814W image and the position angle of the ACS image, 120.37 degree. These observables are summarized in Table 1. We use the lensmodel software (Keeton, 2001b) to model the lensed quasar system. Both mass models have eight parameters; $`\alpha _e`$, $`\gamma `$, $`\theta _\gamma `$, the lensing galaxy position, the source quasar position, and the source quasar flux for the SIS+shear model; $`\alpha _e`$, ellipticity $`e`$ ($`e=1q`$), $`\theta _e`$, the lensing galaxy position, the source quasar position, and the source quasar flux for the SIE model. Since we have only eight constraints (positions of two lensed components and the lensing galaxy, and the fluxes of lensed components), the number of degrees of freedom is zero. The results are summarized in Table 3 and Table 4. The critical curve in the image plane and the caustics in the source plane of the SIS+shear model are shown in the left panel of Figure 10, and those of the SIE model are shown in the right panel of Figure 10, respectively. The best-fit models predict the time delay between the two lensed images to be $`\mathrm{\Delta }t=7.1h_0^1`$ day for the SIS+shear model and $`\mathrm{\Delta }t=8.9h_0^1`$ day for the SIE model ($`H_0=100h_0`$ km sec<sup>-1</sup> Mpc<sup>-1</sup>), assuming that the redshift of the lensing galaxy is 0.724. Both the mass models predict the Einstein radius $`\alpha _e`$ of about 0$`\stackrel{}{\mathrm{.}}`$55, which corresponds to the velocity dispersion (in one direction) of the lensing galaxy of $`\sigma 265`$ km sec<sup>-1</sup> at $`z=0.724`$. Using the Faber-Jackson law (Faber & Jackson, 1976), $`M_{i}^{}{}_{}{}^{}=21.3`$ and $`\sigma ^{}=225`$ km sec<sup>-1</sup> (Blanton et al., 2001; Kochanek, 1996), and K- and evolution-correction for elliptical galaxies (Inada 2004; 1.0 and $``$0.1, respectively, derived from the model of Poggianti 1997 and the SDSS filter response function), the $`i`$-band magnitude of the lensing galaxy is predicted to be 22.1. This is in good agreement with the $`i`$-band magnitude of component G, 21.8. Thus these mass models further support that the redshift of the lensing galaxy is 0.724.
Although the number of degrees of freedom is zero in the above mass models, we were not able to reproduce the observables perfectly. The observed flux ratio between components B and A are $`0.33`$, while the predicted flux ratios are 0.58 for the SIS+shear model and 0.54 for the SIE model (see Table 3). The differences between the observed flux ratio and the predicted flux ratios are very large and therefore cannot be explained by the observed data errors. To explore the discrepancy, we speculate that an unknown object near component G (we named this hypothetical second lensing object “G”) causes the discrepancy. In order to minimize the difference between the numbers of constraints and parameters, we simply applied the SIS model for both G and G (i.e., we did not consider the shear and the ellipticity). This model, however, has nine parameters and is under-constrained (the number of the difference between the constraints and the parameters is 1), thus we required that the Einstein radius of G should not be so large, smaller than $``$ 0$`\stackrel{}{\mathrm{.}}`$20; this is natural since G should not be the dominant contribution. We varied $`\alpha _e`$ of G between 0$`\stackrel{}{\mathrm{.}}`$01 and 0$`\stackrel{}{\mathrm{.}}`$20, and search for the best-fit models for each value of $`\alpha _e`$ of G. We found that well-fitted models ($`\chi ^20`$) exist for each $`\alpha _e`$ of G between $`0\stackrel{}{\mathrm{.}}05`$ and $`0\stackrel{}{\mathrm{.}}20`$. A result for the case that the Einstein radius of G is 0$`\stackrel{}{\mathrm{.}}`$07 is summarized in Table 3 and Table 4; the predicted flux ratio between components B and A is 0.33. If we simply apply the Faber-Jackson law, the Einstein radius of the SIS model is proportional to $`L^{1/2}`$; thus that the Einstein radius of G is 0$`\stackrel{}{\mathrm{.}}`$07 corresponds to G $`5`$ mag fainter than G if we assume that G is at the same redshift of G. This is consistent with the fact that we cannot see any objects at the predicted position of G in the ACS images even after the PFS subtraction. The critical curve in the image plane and the caustics in the source plane of this model (2 SIS lens model for $`\alpha _e`$(G)$`=`$0$`\stackrel{}{\mathrm{.}}`$07) are shown in Figure 11, and the predicted time delay of this model is $`\mathrm{\Delta }t=8.7h_0^1`$ day. We note that the total magnification factors of all mass models (especially 2 SIS lens model) are much higher than those of the previously known two-image lensed quasars, e.g., HE1104$``$1805 (Wisotzki et al., 1993).
In the above lens models, it is expected that the shape of the lensed host galaxy, if it exists, is similar to that of the critical curve, because the quadruple moments are negligible and the source is close to the center of the lens system (Kochanek et al., 2001). To compare the ring feature observed in the NICMOS image and the critical curve of the best-fit mass model, we overplot the critical curve predicted by the 2 SIS lens model on the NICMOS ring, after subtracting the central lensing galaxy and excluding the pixels around components A and B. The result is shown in Figure 12. The red dot-dashed line in Figure 12 connects the pixels that have the maximum count rate in each column; thus it almost corresponds to the structure of the observed ring. In addition to the 2 SIS lens model, we also find that the critical curves of the SIS+shear model and the SIE model are in good agreement with the observed ring. A possible objection of this interpretation is that the magnification center of the observed ring is largely different from the center of component A in both the ACS and the NICMOS images; the ring might not be associated to the source quasar but be a background galaxy lensed by the same lensing galaxy.
## 4 Summary and Conclusion
We have reported the discovery of the doubly-imaged gravitationally lensed quasar at $`z=1.68`$ with the separation angle 1$`\stackrel{}{\mathrm{.}}`$04, SDSS J0246$``$0825, which was selected from the SDSS data. The redshift of the lensing galaxy is likely to be $`z=0.724`$, judged from the absorption lines in the spectra of the quasar components and the apparent magnitude of the galaxy combined with the expected absolute magnitude from the Faber-Jackson relation.
We found that the simple mass models (SIS+shear and SIE) with reasonable parameters well reproduce the lensing geometry of SDSS J0246$``$0825 but do not reproduce the flux ratio between the lensed quasar components. Thus, we speculated that there is another lensing object near the lensed quasar system; we found possible lens models which can reproduce both the lensing geometry and the flux ratio, by introducing a faint second lensing object. All of the best fit models (including the SIS + shear model and the SIE model) predicts the time delay between the two lensed images to be $`10h^1`$ days. In addition to the two lensed images, we have found a highly distorted object in the HST ACS and NICMOS images. In particular, in the near infrared (NICMOS) image this object forms a nearly perfect ring. Since the critical curves of the best fit models show a good agreement with the ring, this is likely to be the lensed host galaxy of the source quasar.
A portion of this work was supported by NASA HST-GO-09744.20. N. I. and M. O. are supported by JSPS through JSPS Research Fellowship for Young Scientists. A portion of this work was also performed under the auspices of the U.S. Department of Energy, National Nuclear Security Administration by the University of California, Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48. Some of the data presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute f Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation. Funding for the creation and distribution of the SDSS Archive has been provided by the Alfred P. Sloan Foundation, the Participating Institutions, the National Aeronautics and Space Administration, the National Science Foundation, the U.S. Department of Energy, the Japanese Monbukagakusho, and the Max Planck Society. The SDSS Web site is http://www.sdss.org/. The SDSS is managed by the Astrophysical Research Consortium (ARC) for the Participating Institutions. The Participating Institutions are The University of Chicago, Fermilab, the Institute for Advanced Study, the Japan Participation Group, The Johns Hopkins University, the Korean Scientist Group, Los Alamos National Laboratory, the Max-Planck-Institute for Astronomy (MPIA), the Max-Planck-Institute for Astrophysics (MPA), New Mexico State University, University of Pittsburgh, University of Portsmouth, Princeton University, the United States Naval Observatory, and the University of Washington.
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# Wave-vector dependent intensity variations of the Kondo peak in photoemission from CePd3
## Acknowledgement
This work was supported by the Deutsche Forschungsgemeinschaft, SFB 463, projects B2, B4, B11, and B16. Expert assistance by R. Follath and other staff members of BESSY is acknowledged.
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# Lorentz-Covariant Quantization of Massive Non-Abelian Gauge Fields in The Hamiltonian Path-Integral Formalism
## I Introduction
It was shown in that the massive non-Abelian gauge fields can well be quantized in the Hamiltonian path-integral formalism following the procedure initiated in for the massless gauge fields. In this quantization, the massive Yang-Mills Lagrangian density written below was chosen to be the starting point.
$$=\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a$$
(1)
where $`A_\mu ^a`$ are the vector potentials for a massive gauge field,
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c$$
(2)
are the field strengths and m is the mass of gauge bosons. Working with the Lagrangian in equation (1.1), the massive non-Abelian gauge fields appear to be second class systems. In the quantization, it is necessary to introduce a primary constraint of second class defined by
$$\mathrm{\Pi }_0^a(x)=\frac{}{\dot{A}_0^a}=0$$
(3)
which was incorporated into the Lagrangian (1.1) by the Lagrange multiplier method. The result of the quantization amounts to the one given in so-called unitary gauge. An alternative quantization within the Hamiltonian path-integral formalism was subsequently performed in by the approach proposed initially in . In this approach, the second class non-Abelian system is converted to a first class system by enlarging the phase space with introducing a series of extra fields. The results of this quantization can be compared with those obtained early in the Lagrangian path-integral formalism by using the gauge-invariant Stückelberg Lagrangian \[7-11\]. Due to the introduction of the extra field variables which have a correspondence with the Stückelberg scalar functions, the quantized result looks much complicated.
In this paper, a Lorentz-covariant quantization of the massive non-Abelian gauge fields will be carried out within the framework of Hamiltonian path-integral along the line described in our previous paper for the quantization of massless non-Abelian gauge fields . The essential point of this quantization which is different from the previous is that the Lorentz condition
$$\phi ^a^\mu A_\mu ^a=0$$
(4)
, as a necessary constraint, is introduced from the beginning and imposed on the massive Yang-Mills Lagrangian. This treatment is consistent with the fact that a massive gauge field has only three polarization states and can completely be described by the Lorentz-covariant (four-dimensionally) transverse part of the vector potential, $`A_T^{a\mu }(x)`$. Whereas, the Lorentz-covariant longitudinal part of the vector potential, $`A_L^{a\mu }`$, appears to be a redundant unphysical variable which must be constrained by introducing the Lorentz condition which implies $`A_L^{a\mu }=0`$. Conventionally, the Lorentz condition is viewed as a consequence of the field equations of motion
$$^\mu F_{\mu \nu }^a+m^2A_\nu ^a=j_\nu ^a$$
(5)
where $`j_\mu ^a`$ is the current generated by the gauge field itself. The argument of this viewpoint is as follows. When we take divergence of the both sides of equation (1.5) and notice the current conservation, it is found that
$$m^2^\mu A_\mu ^a=0$$
(6)
Since $`m0`$, the above equation leads to the Lorentz condition. This seems to imply that the Lorentz condition has already been included in the massive Yang-Mills Lagrangian. If so, when the Lagrangian is written in the first order form, we should see a term in the Lagrangian which is given by incorporating the Lorentz condition by the Lagrange multiplier method. Nevertheless, as will be seen in the next section, there is no such a term to appear in the Lagrangian. Therefore, the viewpoint stated above is not reasonable. The correct procedure is to treat the Lorentz condition as a primary constraint imposed on the massive Yang-Mills Lagrangian. In this case, due to the Lorentz condition introduced, equation (1.6), as a trivial identity, naturally holds, exhibiting the self-consistency of the theory. This procedure coincides with the aforementioned procedure of the quantization performed in where the condition in equation (1.2), as a primary constraint, is necessarily introduced and is also consistent with the conventional canonical quantization. In the latter quantization, to derive the free propagator of the gauge boson, one only uses the transverse free field operator
$$𝐀_\mu ^c(x)=\frac{d^3k}{(2\pi )^3\sqrt{2\omega (k)}}\underset{\lambda =1}{\overset{3}{}}ϵ_\mu ^\lambda (k)[𝐚_\lambda ^c(k)e^{ikx}+𝐚_\lambda ^{c+}(k)e^{ikx}]$$
(7)
where $`ϵ_\mu ^\lambda (k)`$ are the polarization vectors satisfying the transversality condition
$$k^\mu ϵ_\mu ^\lambda (k)=0$$
(8)
which follows directly from the Lorentz condition.
We have a great interest to note that in the physical subspace defined by the Lorentz condition, i.e., spanned by the transverse vector potential $`A_T^{a\mu },`$ the dynamics of massive gauge fields is gauge-invariant in contrast to that the massive Yang-Mills Lagrangian itself is not gauge-invariant. In fact, when we make an infinitesimal gauge transformation\[13 \]:
$$\delta A_\mu ^a=D_\mu ^{ab}\theta ^b$$
(9)
where
$$D_\mu ^{ab}=\delta ^{ab}_\mu gf^{abc}A_\mu ^c$$
(10)
to the action given by the Lagrangian in equation (1.1), noticing the identity $`f^{abc}A^{a\mu }A_\mu ^b=0`$ and applying the Lorentz condition, it can be found that
$$\delta S=m^2d^4x\theta ^a^\mu A_\mu ^a=0$$
(11)
Alternatively, the gauge-invariance may also be formulated by means of the Lagrangian represented in terms of the transverse fields
$$=\frac{1}{4}F_T^{a\mu \nu }F_{T\mu \nu }^a+\frac{1}{2}m^2A_T^{a\mu }A_{T\mu }^a$$
(12)
which is obtained from equation (1.1) by applying the solution of equation (1.4): $`A_L^{a\mu }=0`$. Under the gauge transformation taking place in the physical subspace
$$\delta A_{T\mu }^a=D_{T\mu }^{ab}\theta ^b$$
(13)
where $`D_{T\mu }^{ab}`$ is defined as that in equation (1.10) with $`A_\mu ^c`$ being replaced by $`A_{T\mu }^c,`$ it is easy to prove that the action given by the Lagrangian (1.12) is invariant
$$\delta S=m^2d^4x\theta ^a^\mu A_{T\mu }^a=0$$
(14)
where the transversality condition: $`^\mu A_{T\mu }^a=0,`$ as an identity, has been noticed. The gauge-invariance of the action in the physical subspace suggests that the quantum non-Abelian gauge field theory may also be set up on the basis of gauge-invariance principle as will be specified in section 4.
The remainder of this paper is arranged as follows. In section 2, we will start from the Lagrangian written in the first order formulation. This Lagrangian is given by incorporating the Lorentz condition into the massive Yang-Mills Lagrangian by the Lagrange undetermined multiplier method and will be used to derive canonical equations of motion for the fields. It will be demonstrated that the equations of motion derived are self-consistent and complete for describing the field dynamics. In section3, the quantization of the massive gauge field will be Lorentz-covariantly performed in the Hamiltonian path-integral formalism. The result of this quantization is different from those given in . In section 4, for comparison, the quantization will be alternatively carried out in the Lagrangian path-integral formalism by applying the Lagrangian multiplier method and the gauge-invariance principle. It will be shown that the Lagrangian multiplier method is equivalent to the Faddeev-Popov approach . The last section serves to make conclusions and comments on the problems of renormalizability and unitarity of the theory.
## II First order formulation and equations of motion
According to the general procedure, the Lorentz condition (1.4) may be incorporated into the Lagrangian (1.1) by the Lagrange undetermined multiplier method to give a generalized Lagrangian. In the first order formalism, this Lagrangian is written as
$$=\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a\frac{1}{2}F^{a\mu \nu }(_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c)+\frac{1}{2}m^2A^{a\mu }A_\mu ^a+\lambda ^a^\mu A_\mu ^a$$
(15)
where $`A_\mu ^a`$ and $`F_{\mu \nu }^a`$ are now treated as the mutually independent variables and $`\lambda ^a`$ are chosen to represent the Lagrange multipliers. Using the canonically conjugate variables defined by
$$\mathrm{\Pi }_\mu ^a(x)=\frac{}{\dot{A}^{a\mu }}=F_{\mu 0}^a+\lambda ^a\delta _{\mu 0}=\{\begin{array}{c}F_{k0}^a=E_k^a,\text{ if }\mu =k=1,2,3;\hfill \\ \lambda ^a=E_0^a,\text{ if }\mu =0.\hfill \end{array}$$
(16)
the Lagrangian (2.1) may be rewritten in the canonical form
$$=E^{a\mu }\dot{A}_\mu ^a+A_0^aC^aE_0^a\phi ^a$$
(17)
where
$$C^a=^\mu E_\mu ^a+gf^{abc}A_k^bE^{ck}+m^2A_0^a$$
(18)
$$=\frac{1}{2}(E_k^a)^2+\frac{1}{4}(F_{ij}^a)^2+\frac{1}{2}m^2[(A_0^a)^2+(A_k^a)^2]$$
(19)
here $`E_\mu ^a=(E_0^a,E_k^a)`$ is a Lorentz vector, $``$ is the Hamiltonian density in which $`F_{ij}^a`$ are defined by
$$F_{ij}^a=_iA_j^a_jA_i^a+gf^{abc}A_i^bA_j^c$$
(20)
In the above, the four-dimensional and the spatial indices are respectively denoted by the Greek and Latin letters. From the stationary condition of the action constructed by the Lagrangian (2.3), one may derive the equations of motion as follows
$$\dot{A}_k^a=_kA_0^a+gf^{abc}A_k^bA_0^cE_k^a$$
(21)
$$\dot{E}_k^a=^iF_{ik}^a+gf^{abc}(E_k^bA_0^c+F_{ki}^bA^{ci})+m^2A_k^a+_kE_0^a$$
(22)
$$C^a(x)^\mu E_\mu ^a+gf^{abc}A_k^bE^{ck}+m^2A_0^a=0$$
(23)
and equation (1.4) where $`k=1,2,3`$. Equations (2.7) and (2.8) act as the equations of motion satisfied by the independent canonical variables $`A_k^a`$ and $`E_k^a(k=1,2.3)`$ which precisely describe the three degrees of freedom of polarization for a massive gauge field with a given group index, while, equations (1.4) and (2.9) can only be regarded as the constraint equations obeyed by the constrained variables $`A_0^a`$ and $`E_0^a`$ because in these equations, there are no time-derivatives of the dynamical variables $`A_k^a`$ and $`E_k^a`$. Equation (2.3) clearly shows that these constraints have already been incorporated in the Lagrangian by the Lagrange undetermined multiplier method. Especially, the Lagrange multipliers are just the constrained variables themselves in this case. It is clear to see that in equations (2.7)-(2.9) and (1.4) there are altogether eight equations for a given group index. They are sufficient to determine the eight variables including three pairs of the dynamical canonical variables $`A_k^a`$ and $`E_k^a(k=1,2.3)`$ and one pair constrained variables $`A_0^a`$ and $`E_0^a`$, showing the completeness of the equations.
Along the general line by Dirac , we shall examine the evolution of the constraints $`\phi ^a`$ and $`C^a`$ with time. Taking the derivative of the both equations (1.4) and (2.9) with respect to time and making use of the equations of motion :
$$\dot{A}_\mu ^a(x)=\frac{\delta H}{\delta E^{a\mu }(x)}d^4y[A_0^b(y)\frac{\delta C^b(y)}{\delta E^{a\mu }(x)}E_0^b(y)\frac{\delta \phi ^b(y)}{\delta E^{a\mu }(x)}]$$
(24)
$$\dot{E}_\mu ^a(x)=\frac{\delta H}{\delta A^{a\mu }(x)}+d^4y[A_0^b(y)\frac{\delta C^b(y)}{\delta A^{a\mu }(x)}E_0^b(y)\frac{\delta \phi ^b(y)}{\delta A^{a\mu }(x)}]$$
(25)
which are obtained from the stationary condition of the action given by the Lagrangian (2.3), one may derive the following consistency equations
$$\{H,\phi ^a(x)\}+d^4y[\{\phi ^a(x),C^b(y)\}A_0^b(y)\{\phi ^a(x),\phi ^b(y)\}E_0^b(y)]=0$$
(26)
$$\{H,C^a(x)\}+d^4y[\{C^a(x),C^b(y)\}A_0^b(y)\{C^a(x),\phi ^b(y)\}E_0^b(y)]=0$$
(27)
where equations (1.4) and (2.9) have been used. In the above, $`H`$ is the Hamiltonian defined by an integral of the Hamiltonian density shown in equation (2.5) over the coordinate x and $`\{F,G\}`$ represents the Poisson bracket which is Lorentz-covariantly defined as
$$\{F,G\}=d^4x\{\frac{\delta F}{\delta A_\mu ^a(x)}\frac{\delta G}{\delta E^{a\mu }(x)}\frac{\delta F}{\delta E_\mu ^a(x)}\frac{\delta G}{\delta A^{a\mu }(x)}\}$$
(28)
The Poisson brackets in equations (2.12) and (2.13) are easily calculated. The results are
$$\{C^a(x),\phi ^b(y)\}=D_\mu ^{ab}(x)_x^\mu \delta ^4(xy)$$
(29)
$$\{\phi ^a(x),\phi ^b(y)\}=0$$
(30)
$$\{C^a(x),C^b(y)\}=m^2[gf^{abc}A_0^c(x)2\delta ^{ab}_0^x]\delta ^4(xy)$$
(31)
$$\{H,\phi ^a(x)\}=_x^kE_k^a(x)$$
(32)
$$\{H,C^a(x)\}=m^2[_0^xA_0^a(x)+_k^xA_k^a(x)]$$
(33)
where $`D_\mu ^{ab}(x)`$ was defined in equation (1.10). It is pointed out that by the requirement of Lorentz-covariance, in the computation of the above brackets, the second term in equation (2.9) has been written in a Lorentz-covariant form $`gf^{abc}A^{b\mu }E_\mu ^c`$. We are allowed to do it because the added term $`gf^{abc}A_0^bE_0^c`$ only gives a vanishing contribution to the term $`A_0^aC^a`$ in equation (2.3) due to the identity $`f^{abc}A_0^aA_0^b=0`$. Particularly, the determinant of the matrix which is constructed by the Poisson bracket denoted in equation (2.15) is not singular. This indicates that equations (2.12) and (2.13) are solvable to determine the Lagrange multipliers $`A_0^a(x)`$ and $`E_0^a(x)`$. There is no necessity of taking other subsidiary constraint conditions into account further. This reveals the consistency of the constraints (1.4) and (2.9). On substituting equations (2.15)-(2.19) into equations (2.12) and (2.13), we find
$$\mathrm{}_xA_0^a(x)gf^{abc}_x^\mu [A_0^b(x)A_\mu ^c(x)]_x^kE_k^a(x)=0$$
(34)
$$[\delta ^{ab}\mathrm{}_xgf^{abc}A_\mu ^c(x)_x^\mu ]E_0^b(x)=0$$
(35)
These equations are compatible with the equations (1.4) and (2.7)-(2.9). In fact, as can easily be verified, equations.(2.20) and (2.21) may directly be derived from equations (1.4) and (2.7-(2.9). In addition, we mention that if the Hamiltonian density is defined by $`\overline{}=A_0^aC^a+E_0^a\phi ^a`$, the equations (2.20) and (2.21) can also be obtained from the equations $`\{\overline{H},\phi ^a(x)\}=0`$ and $`\{\overline{H},C^a(x)\}=0`$ respectively where $`\overline{H}`$ is the Hamiltonian defined by an integral of the $`\overline{\text{ }}`$ over the coordinate x.
## III Quantization in the Hamiltonian path-integral formalism
This section serves to formulate the quantization performed in the Hamiltonian path-integral formalism for the massive non-Abelian gauge fields. In accordance with the general procedure of the quantization , we firstly write the generating functional of Green’s functions via the independent canonical variables which are now chosen to be the transverse parts of the vectors $`A_\mu ^a`$ and $`E_\mu ^a`$
$$Z[J]=\frac{1}{N}D(A_T^{a\mu },E_T^{a\mu })exp\{id^4x[E_T^{a\mu }\dot{A}_{T\mu }^a^{}(A_T^{a\mu },E_T^{a\mu })+J_T^{a\mu }A_{T\mu }^a]\}$$
(36)
where $`^{}(A_T^{a\mu },E_T^{a\mu })`$ is the Hamiltonian which is obtained from the Hamiltonian (2.5) by replacing the constrained variables $`A_L^{a\mu }`$ and $`E_L^{a\mu }`$ with the solutions of equations (1.4) and (2.9). As mentioned before, equation (1.4) leads to $`A_L^{a\mu }=0`$. Noticing this solution and the decomposition $`E^{a\mu }(x)=E_T^{a\mu }+E_L^{a\mu }(x),`$ when setting $`E_L^{a\mu }(x)=_x^\mu Q^a(x)`$ where $`Q^a(x)`$ is a scalar function, one may get from equation (2.9) an equation obeyed by the scalar function $`Q^a(x)`$ for a given group index
$$K^{ab}(x)Q^b(x)=W^a(x)$$
(37)
where
$$K^{ab}(x)=\delta ^{ab}\mathrm{}_xgf^{abc}A_T^{c\mu }(x)_\mu ^x$$
(38)
and
$$W^a(x)=gf^{abc}E_T^{b\mu }(x)A_{T\mu }^c(x)m^2A_T^{a0}(x)$$
(39)
With the aid of the Green’s function $`G^{ab}(xy)`$ (the ghost particle propagator) which satisfies the following equation
$$K^{ac}(x)G^{cb}(xy)=\delta ^{ab}\delta ^4(xy)$$
(40)
one may find the solution to the equation (3.2) as follows
$$Q^a(x)=d^4yG^{ab}(xy)W^b(y)$$
(41)
With the expressions given in equations (3.4) and (3.6), we see that the $`E_L^{a\mu }(x)`$ is a complicated functional of the variables $`A_T^{a\mu }`$ and $`E_T^{a\mu }`$ so that the Hamiltonian $`^{}(A_T^{a\mu },E_T^{a\mu })`$ is of much more complicated functional structure which is not convenient for constructing the diagram technique in perturbation theory. Therefore, it is better to express the generating functional in equation (3.1) in terms of the variables $`A_\mu ^a`$ and $`E_\mu ^a`$. For this purpose, it is necessary to insert the following delta-functional into equation (3.1)
$$\delta [A_L^{a\mu }]\delta [E_L^{a\mu }E_L^{a\mu }(A_T^{a\mu },E_T^{a\mu })]=detM\delta [C^a]\delta [\phi ^a]$$
(42)
where $`M`$ is the matrix whose elements are
$$M^{ab}(x,y)=\{C^a(x),\phi ^b(y)\}$$
(43)
which was given in equation (2.15). The relation in equation (3.7) is easily derived from equations (1.4) and (2.9) by applying the property of delta-functional . Upon inserting equation (3.7) into equation (3.1) and utilizing the Fourier representation of the delta-functional
$$\delta [C^a]=D(\eta ^a/2\pi )e^{i{\scriptscriptstyle d^4x\eta ^aC^a}}$$
(44)
we have
$`Z[J]`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle }D(A_\mu ^a,E_\mu ^a,\eta ^a)detM\delta [^\mu A_\mu ^a]exp\{i{\displaystyle }d^4x[E^{a\mu }\dot{A}_\mu ^a`$ ()
$`+\eta ^aC^a(A_\mu ^a,E_\mu ^a)+J^{a\mu }A_\mu ^a]\}`$
In the above exponential, there is a $`E_0^a`$-related term $`E_0^a(_0A_0^a_0\eta ^a)`$ which permits us to perform the integration over $`E_0^a`$, giving a delta-functional $`\delta [_0A_0^a_0\eta ^a]=det|_0|^1\delta [A_0^a\eta ^a]`$. The determinant $`det|_0|^1`$, as a constant, may be put in the normalization constant $`N`$ and the delta-functional $`\delta [A_0^a\eta ^a]`$ will disappear when the integration over $`\eta ^a`$ is carried out. The integral over $`E_k^a`$ is of Gaussian-type and hence easily calculated. After these manipulations, we arrive at
$`Z[J]`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle }D(A_\mu ^a)detM\delta [^\mu A_\mu ^a]exp\{i{\displaystyle }d^4x[{\displaystyle \frac{1}{4}}F^{a\mu \nu }F_{\mu \nu }^a`$ ()
$`+{\displaystyle \frac{1}{2}}m^2A^{a\mu }A_\mu ^a+J^{a\mu }A_\mu ^a]\}`$
When employing the familiar expression
$$detM=D(\overline{C}^a,C^a)e^{i{\scriptscriptstyle d^4xd^4y\overline{C}^a(x)M^{ab}(x,y)C^b(y)}}$$
(47)
where $`\overline{C}^a(x)`$ and $`C^a(x)`$ are the mutually conjugate ghost field variables and the following limit for the Fresnel functional
$$\delta [^\mu A_\mu ^a]=\underset{\alpha 0}{lim}C[\alpha ]e^{\frac{i}{2\alpha }{\scriptscriptstyle d^4x(^\mu A_\mu ^a)^2}}$$
(48)
where $`C[\alpha ]_x(\frac{i}{2\pi \alpha })^{1/2}`$ and supplementing the external source terms for the ghost fields, the generating functional in equation (3.11) is finally given in the form
$$Z[J,\overline{\xi },\xi ]=\frac{1}{N}D(A_\mu ^a,\overline{C}^a,C^a)exp\{id^4x[_{eff}+J^{a\mu }A_\mu ^a+\overline{\xi }^aC^a+\overline{C}^a\xi ^a]\}$$
(49)
where
$$_{eff}=\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a\frac{1}{2\alpha }(^\mu A_\mu ^a)^2^\mu \overline{C}^aD_\mu ^{ab}C^b$$
(50)
which is the effective Lagrangian for the quantized massive gauge field in which the first two terms are the Yang-Mills Lagrangian , the third and fourth terms are the so-called gauge-fixing term and the ghost term respectively. In equation (3.14), the limit $`\alpha 0`$ is implied. Certainly, the theory may be given in general gauges $`(\alpha 0)`$. In this case, as will be seen in the next section, the ghost particle will acquire a spurious mass $`\mu =\sqrt{\alpha }m`$.
## IV Quantization in the Lagrangian path-integral formalism
To help understanding of the result of the quantization given in the preceding section, in this section, we attempt to quantize the massive non-Abelian gauge fields in the Lagrangian path-integral formalism. For later convenience, the massive Yang-Mills Lagrangian in equation (1.1) and the Lorentz constraint condition in equation (1.4) are respectively generalized to the following forms
$$_\lambda =\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a\frac{1}{2}\alpha (\lambda ^a)^2$$
(51)
and
$$^\mu A_\mu ^a+\alpha \lambda ^a=0$$
(52)
where $`\lambda ^a(x)`$ are the extra functions which will be identified with the Lagrange multipliers and $`\alpha `$ is an arbitrary constant playing the role of gauge parameter. Now, according to the general procedure for constrained systems, equation (4.2) may be incorporated into equation (4.1) by the Lagrange multiplier method, giving a generalized Lagrangian
$$_\lambda =\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a+\lambda ^a^\mu A_\mu ^a+\frac{1}{2}\alpha (\lambda ^a)^2$$
(53)
This Lagrangian is obviously not gauge-invariant. However, for building up a correct gauge field theory, it is necessary to require the dynamics of the gauge field to be gauge-invariant. In other words, the action given by the Lagrangian (4.3) is required to be invariant under the gauge transformations shown in equation (1.9). By this requirement, noticing the identity $`f^{abc}A^{a\mu }A_\mu ^b=0`$ and applying the constraint condition (4.2), we have
$$\delta S_\lambda =\frac{1}{\alpha }d^4x^\nu A_\nu ^a(x)^\mu (𝒟_\mu ^{ab}(x)\theta ^b(x))=0$$
(54)
where
$$𝒟_\mu ^{ab}(x)=\delta ^{ab}\frac{\mu ^2}{\mathrm{}_x}_\mu ^x+D_\mu ^{ab}(x)$$
(55)
in which $`\mu ^2=\alpha m^2`$ and $`D_\mu ^{ab}(x)`$ was defined in equation (1.10). From equation (4.2) we see $`\frac{1}{\alpha }^\nu A_\nu ^a=\lambda ^a0`$. Therefore, to ensure the action to be gauge-invariant, the following constraint condition on the gauge group is necessary to be required
$$_x^\mu (𝒟_\mu ^{ab}(x)\theta ^b(x))=0$$
(56)
These are the coupled equations satisfied by the parametric functions $`\theta ^a(x)`$ of the gauge group. Since the Jacobian is not singular
$$detM0$$
(57)
where
$`M^{ab}(x,y)`$ $`=`$ $`{\displaystyle \frac{\delta (_x^\mu 𝒟_\mu ^{ac}(x)\theta ^c(x))}{\delta \theta ^b(y)}}`$ (58)
$`=`$ $`\delta ^{ab}(\mathrm{}_x+\mu ^2)\delta ^4(xy)gf^{abc}_x^\mu (A_\mu ^c(x)\delta ^4(xy))`$ ()
the above equations are solvable and would give a set of solutions which express the functions $`\theta ^a(x)`$ as functionals of the vector potentials $`A_\mu ^a(x)`$. The constraint conditions in equation (4.6) may also be inserted into the Lagrangian (4.3) by the Lagrange undetermined multiplier method. In doing this, it is convenient, as usually done, to introduce ghost field variables $`C^a(x)`$ in such a fashion
$$\theta ^a(x)=\zeta C^a(x)$$
(59)
where $`\zeta `$ is an infinitesimal Grassmann’s number. In accordance with equation (4.9), the constraint condition (4.6) can be rewritten as
$$^\mu (𝒟_\mu ^{ab}C^b)=0$$
(60)
where the number $`\zeta `$ has been dropped. This constraint condition usually is called ghost equation. When the condition (4.10) is incorporated into the Lagrangian (4.3) by the Lagrange multiplier method, we obtain a more generalized Lagrangian as follows
$$_\lambda =\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a+\lambda ^a^\mu A_\mu ^a+\frac{1}{2}\alpha (\lambda ^a)^2+\overline{C}^a^\mu (𝒟_\mu ^{ab}C^b)$$
(61)
where $`\overline{C}^a(x)`$, acting as Lagrange undetermined multipliers, are the new scalar variables conjugate to the ghost variables $`C^a(x).`$
At present, we are ready to formulate the quantization of the massive gauge field . As we learn from the Lagrange undetermined multiplier method, the dynamical and constrained variables as well as the Lagrange multipliers in the Lagrangian (4.11) can all be treated as free ones, varying arbitrarily. Therefore, we are allowed to use this kind of Lagrangian to construct the generating functional of Green’s functions
$`Z[J^{a\mu },\overline{\xi }^a,\xi ^a]`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle }D(A_\mu ^a,\overline{C}^a,C^a,\lambda ^a)exp\{i{\displaystyle }d^4x[_\lambda (x)+J^{a\mu }(x)A_\mu ^a(x)`$ ()
$`+\overline{\xi }^a(x)C^a(x)+\overline{C}^a(x)\xi ^a(x)]\}`$
where $`D(A_\mu ^a,\mathrm{},\lambda ^a)`$ denotes the functional integration measure, $`J_\mu ^a,\overline{K}^a`$ and $`K^a`$ are the external sources coupled to the gauge and ghost fields and $`N`$ is a normalization constant. Looking at the expression of the Lagrangian (4.11), we see, the integral over $`\lambda ^a(x)`$ is of Gaussian-type. Upon completing the calculation of this integral, we finally arrive at
$`Z[J^{a\mu },\overline{\xi }^a,\xi ^a]`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle }D(A_\mu ^a,\overline{C}^a,C^a,)exp\{i{\displaystyle }d^4x[_{eff}(x)`$ ()
$`+J^{a\mu }(x)A_\mu ^a(x)+\overline{\xi }^a(x)C^a(x)+\overline{C}^a(x)\xi ^a(x)]\}`$
where
$$_{eff}=\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a\frac{1}{2\alpha }(^\mu A_\mu ^a)^2^\mu \overline{C}^a𝒟_\mu ^{ab}C^b$$
(64)
is the effective Lagrangian given in the general gauges. In the Landau gauge ($`\alpha 0`$), The Lagrangian (4.14) goes over to the one given in equation (3.15). When the mass m tends to zero, equation (4.14) is immediately converted to the Lagrangian encountered in the massless gauge field theory.
Now let us turn to show that the quantization described above is equivalent to the quantization by the Faddeev-Popov approach. According to the procedure of the latter approach, we need to insert the following identity
$$\mathrm{\Delta }[A]D(\theta ^a)\delta [^\mu A_\mu ^\theta +\alpha \lambda ^\theta ]=1$$
(65)
into the vacuum-to-vacuum transition amplitude, obtaining
$$Z[0]<0^+0^{}>=\frac{1}{N}D(A_\mu ^a,\lambda ^a,\theta ^a)\mathrm{\Delta }[A]\delta [^\mu A_\mu ^\theta +\alpha \lambda ^\theta ]exp(iS_\lambda )$$
(66)
where $`S_\lambda `$ is the action given by the Lagrangian (4.1)
$$S_\lambda =d^4x[\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2A^{a\mu }A_\mu ^a\frac{1}{2}\alpha (\lambda ^a)^2]$$
(67)
The delta-functional in equation (4.16) implies
$$^\mu (A^\theta )_\mu ^a+\alpha (\lambda ^\theta )^a=^\mu A_\mu ^a+\alpha \lambda ^a=0$$
(68)
where $`(A^\theta )_\mu ^a`$ and $`(\lambda ^\theta )^a`$ represent the gauge-transformed vector potentials and Lagrange multipliers, $`(A^\theta )_\mu ^a=A_\mu ^a+\delta A_\mu ^a`$ here $`\delta A_\mu ^a`$ was denoted in equation (1.9) and $`(\lambda ^\theta )^a=\lambda ^a+\delta \lambda ^a`$. Equation (4.18) holds naturally because the constraint condition (4.2) is required to be satisfied for all the field variables including the ones before and after gauge transformation. The $`\delta \lambda ^a`$ may be determined by the requirement that the action $`S_\lambda `$ is to be gauge-invariant with respect to the gauge transformations of $`A_\mu ^a`$ and $`\lambda ^a`$. By this requirement, it is easy to find
$$\delta S_\lambda =d^4x^\mu A_\mu ^a(\delta \lambda ^am^2\theta ^a)=0$$
(69)
where the condition in equation (4.2) has been considered. From the above equation, noticing $`^\mu A_\mu ^a0`$ in the general gauges, we see, it must be
$$\delta \lambda ^a=m^2\theta ^a.$$
(70)
Here we see that in the case of massless gauge fields, $`\delta \lambda ^a=0`$ so that equation (4.19) becomes a trivial identity. When the gauge transformations given in equations (1.9) and (4.20) are inserted into equation (4.18), we may obtain a constraint condition which is identical to that denoted in equation (4.6). Hence, from equation (4.15) we get $`\mathrm{\Delta }[A]=det`$M\[A\] in which the matrix M\[A\] is completely the same as given in equation (4.8). It is easy to verify that the determinant det M\[A\] , the integration measure and the action in equation (4.17) are all invariant with respect to the gauge transformations of the functions $`A_\mu ^a`$ and $`\lambda ^a.`$ Therefore, when we make the gauge transformations: $`A_\mu ^a(A^\theta )_\mu ^a`$ and $`\lambda ^a(\lambda ^\theta )^a`$ to the functional integral in equation (4.16), the integral over $`\theta ^a(x)`$ , as a constant, may be factored out from the functional integral over $`A_\mu ^a`$ and $`\lambda ^a`$ and put in the normalization constant N. Thus, we have
$$Z[0]=\frac{1}{N}D(A_\mu ^a,\lambda ^a)detM[A]\delta [^\mu A_\mu +\alpha \lambda ]exp(iS_\lambda )$$
(71)
On completing the integration over $`\lambda ^a`$, we obtain
$$Z[0]=\frac{1}{N}D(A_\mu ^a)𝑑etM[A]exp\{iS\frac{i}{2\alpha }d^4x(^\mu A_\mu ^a)^2\}$$
(72)
where S is the massive Yang-Mills action. By making use of the representation of det M\[A\] as shown in equation (3.12) and introducing the external sources, we exactly recover the generating functional written in equations (4.13) and (4.14).
In the end, we would like to emphasize that under the Lorentz condition which has been introduced into the Lagrangian in equation (4.3) and the transition amplitude in equation (4.16), the gauge-invariance of the action shown in equation (4.4) and equation (4.19) was formulated only for the infinitesimal gauge transformations denoted in equation (1.9). The quantized result shown in equations (4.13) and (4.14) was derived just by the use of such gauge transformations. In the Landau gauge ($`\alpha =0)`$, this result is exactly the same as that obtained in section 3 by the quantization performed in the Hamiltonian path-integral formalism for which we only need to calculate the classical Poisson brackets without concerning any gauge transformation. This fact indicates that to get the correct quantized result by the methods formulated in this section, the infinitesimal gauge transformations are only necessary to be taken into account. This implies that in the physical subspace restricted by the Lorentz condition, only the infinitesimal gauge transformations are possible to exist.
## V Conclusions and comments
The new features of this paper consist in: (1) The Lorentz-covariant quantization is successfully achieved in the Hamiltonian path-integral formalism; (2) The result of this quantization is confirmed by the quantization performed by the Lagrange multiplier method in the Lagrangian path-integral formalism. The latter method is proposed first in this paper and shown to be equivalent to the Faddeev-Popov approach; (3) The established quantum massive non-Abelian gauge field theory can be straightforwardly converted to the massless theory in the zero-mass limit in contrast to the previous quantization in the Hamiltonian formulation for which the zero-mass limit does not exist. The essential points to achieve these results are clarified in this paper. They include: (1) The Lorentz condition, as a necessary constraint, must be introduced initially and imposed on the massive Yang-Mills Lagrangian so as to give a complete formulation of the dynamics of massive gauge fields. That is to say, in the whole space of the full vector potential, the massive non-Abelian gauge fields must be considered as constrained systems and the massive Yang-Mills Lagrangian itself is not complete for describing the dynamics; (2) The Lorentz condition may be incorporated into the massive Yang-Mills Lagrangian by the Lagrange multiplier method so that each component of a vector potential acquires its canonically conjugate counterpart. This makes the Lorentz-covariant formulation of the Hamiltonian path-integral quantization become possible; (3) In the physical subspace restricted by the Lorentz condition, the action other than the Lagrangian for the massive gauge field is gauge-invariant. Since the action is of more fundamental dynamical meaning than the Lagrangian, the gauge-invariance of the action implies that the dynamics of the massive gauge field is gauge-invariant. Therefore , the quantum massive non-Abelian gauge field theory may also be set up on the basis of gauge-invariance; (4) In the physical subspace, only infinitesimal gauge transformations are possibly allowed. This point explicated clearly in the end of the former section was already pointed out in reference . In the reference, after introducing the identity denoted in equation (4.15) into the vacuum-to vacuum transition amplitude, the authors said that ” We must know $`\mathrm{\Delta }[A]`$ is only for the transverse fields and in this case all contributions to the last integral are given in the neighborhood of the unite element of the group”. By this point, it can be understood why in the ordinary quantum gauge field theories such as the standard model, the BRST-transformations are all taken to be infinitesimal.
Now we are in a position to comment on the previous works for the massive gauge field theories. Originally, the theory was set up from the massive Yang-Mills Lagrangian alone which is viewed as complete for describing the field dynamics \[15-17\]. To overcome the gauge-non-invariance of the mass term in the Lagrangian, the Stückelberg Lagrangian in which the mass term is written in a gauge-invariant form by introducing the extra Stückelberg field functions was proposed and widely chosen to be the starting point to establish the quantum theory \[7-11\]. All these theories were considered to be nonrenormalizable due to the appearance of the unphysical gauge degrees of freedom in the theories which arises from the finite gauge transformations. We would like to point out that according to the basic ideas stated before, the Lorentz condition should be , as a primiry constraint, introduced initially into the theories and thereby only the infinitesimal gauge transformations are necessary to be taken into account. Based on this point, let us analysize the previous works mentioned above. In reference , the authors proved an equivalence theorem by which they gave a Hamiltonian derived from the massive Yang-Mills Lagrangian by introducing an auxiliary Stückelberg field functions. When making an unitary transformation to the Schrödinger equation they derived, the mass term in the Hamiltonian becomes dependent on the auxiliary field and contains an infinite number of terms in its expansion of power series which lead to bad unrenormalizability. In reference , the equivalence theorem was given in the form of S-matrix. The author also introduced the Stückelberg field and used it to make a finite gauge transformation to the fields involved in the theory. As a result, the mass term in the S-matrix contains an exponential function of the auxiliary field which gives rise to an infinite variety of distinct primitively divergent graphs that can not be eliminated by the introduced renormalizability conditions imposed on the gauge transformation (see equations (39) and (40) in the reference). It is noted here that when infinitesimal gauge transformations are, as they should be, considered only in the theory, it is easy to find that the gauge degrees of freedom are reduced and, particularly, the introduced renormalizability conditions are perfectly satisfied. In reference , the author found a relation by which any full vector potential may be represented as a gauge transformation of the physical transverse vector potential and tried to separate the gauge degrees of freedom from the transverse ones as it is able to be done for the renormalizable theory of the massive neutral gauge boson coupled to a conserved current. He eventually failed to do it because in the non-Abelian case the coupling between the both degrees of freedom does not vanish upon integration and thereof he concluded that the theory is nonrenormalizable. However, this conclusion is invalid for the infinitesimal gauge transformations required by the Lorentz condition because in this case the coupling mentioned above vanishes and, therefore, the gauge degree of freedom can also be separated out. As a result, the generating functional, as in the case of massive neutral gauge boson theory, may also be expressed through the transverse fields like equation (3.1). In reference , the quantization of the massive non-Abelian gauge field was implemented by the Faddeev-Popov operation with incorporation of the Lorentz condition in the generating functional. Nevertheless, the incorporation of the Lorentz condition is only for the purpose of improving the behavior of the massive gauge boson propagator, being not thought of a necessary procedure. In addition, different from the procedure stated in section 4, the quantization does not meet the requirement of gauge-invariance so that the ghost particle loses a mass term in the general gauges ($`\alpha 0`$). In particular, opposite to the procedure that the ghost term in the effective Lagrangian is introduced by applying the infinitesimal gauge transformations, the finite gauge transformations are made to the effective Lagrangian. As a result of these transformations, the mass term depends on the parametric functions of the gauge group and contains various unrenormalizable infinities. As pointed out before, when the Lorentz condition is incorporated in the theory, a consistent procedure is that the infinitesimal gauge transformations are required only. In this case, the unrenormalizable infinities in the mass term could not appear. Let us turn to the gauge-invariant Stückelberg Lagrangian \[8-10\]
$$L=\frac{1}{4}F^{a\mu \nu }F_{\mu \nu }^a+\frac{1}{2}m^2(A_\mu ^a\omega _\mu ^a)(A^{a\mu }\omega ^{a\mu })$$
(73)
where $`\omega _\mu ^a`$ are the group-valued Stückelberg functions defined as
$$\omega _\mu ^a=\frac{i}{g}(_\mu U^1U)^a=(\frac{e^{i\omega }1}{i\omega })^{ab}_\mu \varphi ^b$$
(74)
in which $`U=e^{ig\varphi ^aT^a}`$ and $`\omega ^{ab}=igf^{abc}\varphi ^c.`$ The function $`\omega _\mu ^a`$ for each group index a contains an infinite number of terms when we expand the exponential function $`e^{i\omega }`$ as a series. The quantum theory built by the Lagrangian (5.1) was proved to be nonrenormalizable due to the nonpolynomial nature of the function $`\omega _\mu ^a`$. It is noted that in the Lagrangian (5.1), the vector field $`\omega _\mu ^a`$ is also given a mass. Therefore, if it is considered to be physical, having three polarized states, we should also impose on it the Lorentz condition
$$^\mu \omega _\mu ^a=0$$
(75)
Substituting equation (5.2) in equation (5.3), it will be found that a necessary and sufficient condition to meet equation (5.3) is $`_\mu \varphi ^a=0`$ which according to equation (5.2) leads to $`\omega _\mu ^a=0.`$ This implies that the Stückelberg functions are unnecessary to be introduced. On the other hand, if the functions $`\omega _\mu ^a`$ are treated as free variables, not receiving any constraint, there should be an integral over them in the generating functional. The integral is of Gaussian-type and hence easily calculated. As a result of the integration, the mass term of the gauge fields will completely be cancelled out from the Lagrangian. Thus, we are left with only a massless gauge field theory. In a word, if the quantization of the massive gauge fields respets the basic requirements stated in the beginning of this section as was done in the preceding sections, the gauge degrees of freedom will be reduced and the problem of the nonrenormalizability will disappear.
From the perturbative expansion of the generating functional given in equations (4,13) and (4.14), it is easy to derive the bare vertices and free propagators. The vertices are the same as those given in the massless theory. The gauge boson propagator and the ghost particle one are respectively, in the momentum space, of the forms
$$iD_{\mu \nu }^{ab}(k)=i\delta ^{ab}\{\frac{g_{\mu \nu }k_\mu k_\nu /k^2}{k^2m^2+i\epsilon }+\frac{\alpha k_\mu k_\nu /k^2}{k^2\mu ^2+i\epsilon }\}$$
(76)
and
$$i\mathrm{\Delta }^{ab}(q)=\frac{i\delta ^{ab}}{q^2\mu ^2+i\epsilon }$$
(77)
These propagators show good renormalizable behavior in the large momentum limit, indicating that the power counting argument is applicable in this case to support the renormalizability of th theory. The renormalizability of the theory originates from the fact that the unphysical degrees of freedom contained in the massive Yang-Mills Lagrangian, i e., the longitudinal component of the gauge field and the residual gauge degrees of freedom existing in the physical subspace are respectively constrained by the introduced Lorentz condition and the ghost equation. The both constraints respectively give rise to the gauge-fixing term and the ghost term in the effective Lagrangian (4.14) which just play the role of counteracting the unphysical degrees of freedom in the massive Yang-Mills Lagrangian. This cancellation also ensures the unitarity of the theory. It is easy to check that for the theory of the gauge field coupled to a vector current such as QCD, all the S-matrix elements given in the tree diagram approximation are unitary because the unphysical longitudinal part of the gauge boson propagator gives a vanishing contribution to the S-matrix elements. For the higher order perturbative S-matrix elements, it can also be proved that the unphysical intermediate states are cancelled out with each other so that the unitarity of the theory is still preserved ( The proofs will be reported later). The problem arises for the theory of the charged gauge boson coupled to a charged chiral current as we met in the electroweak theory . For this kind of theory , it seems that the tree-unitarity is violated if without recourse to the Higgs mechanism . We do not concern this kind of theory in this paper. But, we would like to note that even for the charged gauge boson theory without introducing the Higgs mechanism, the unitarity of the theory may also be preserved by means of the limiting procedure proposed originally in reference . The limiting procedure means that the gauge boson propagator in equation (5.4) which is written in the $`\alpha `$-gauge will, in the limit: $`\alpha \mathrm{}`$, be converted to the one given in the unitary gauge
$$iD_{\mu \nu }^{ab}(k)=i\delta ^{ab}\frac{g_{\mu \nu }k_\mu k_\nu /m^2}{k^2m^2+i\epsilon }$$
(78)
which was originally derived in the canonical quantization by making use of the Fourier representation of the transverse vector potential denoted in equation (1.7). Since the $`\alpha `$-gauge propagator has good renormalizable behavior, one may employ this kind of propagator to calculate the S-matrix and then by the limiting procedure to obtain the unitary gauge result so as to guarantee the unitarity of the theory. In doing this, the theory given in the $`\alpha `$-gauge can be viewed as a regularization of the theory given in the unitary gauge.
## VI ACKNOWLEDGMENT
The author would like to thank professor Shi-Shu Wu for useful discussions. This subject was supported in part by National Natural Science Foundation of China and the Research Fund for the Doctoral Program of Higher Education.
## VII References
P. Senjanovic, Ann. Phys. (New York) 100, 227 (1976).
C.Grosse-Knetter, Phys. Rev. D48, 2854 (1993).
P. A. M. Dirac, Lectures on Quantum Mechanics, Yeshiva University Press, New York (1964).
L. D. Faddeev, Theor. Math. Phys. 1, 1; L. D. Faddeev and A. A. Slavnov,
Gauge Fields: Introduction to Quantum Theory, The Benjamin Cummings Publishing Company Inc.(1980).
N. Banerjee, R. Banejee and Subir Ghosh, Ann. Phys. 241, 237 (1995).
I. A. Batalin and I. V. Tyutin, Nucl. Phys. B279, 514 (1987); Inter. J. Mod. Phys. A6, 3255 (1991).
E. C. G. Stückelberg, Helv. Phys. Acta. 30, 225 (1957).
T. Kunimasa and T. Goto, Prog. Theor. Phys. 37, 452 (1967).
K.Shizuya, Nucl. Phys. B87, 255 (1975); B94, 260 (1975); B121, 125 (1977).
T. Fukuda ,M. Monoa , M. Takeda and K. Yokoyama, Prog. Theor. Phys. 66, 1827 (1981); 67, 1206 (1982);
70, 284 (1983).
A. Burnel, Phys. Rev. D33, 2981 (1986); D33, 2985 (1986).
J. C. Su, J. Phys. G: Nucl. Part. Phys. 27, 1493 (2001).
C. Itzykson and F-B. Zuber, Quantum Field Theory, McGraw-Hill, New York (1980).
L. D. Faddeev and V. N. Popov, Phys. Lett. B25, 29 (1967).
H. Umezawa and S. Kamefuchi, Nucl. Phys. 23, 399 (1961).
A. Salam, Phys. Rev. 127, 331 (1962).
D .G. Boulware, Ann. Phys. 56, 140 (1970).
C. H. Llewellyn-Smith, Phys. Lett. B46, 233 (1973).
J. M. Cornwall, D. N. Levin and G. Tiktopoulos, Phys. Rev. Lett. 30, 1268 (1973).
T. D. Lee and C. N. Yang, Phys. Rev. 128, 885 (1962).
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# Moderate deviations and laws of the iterated logarithm for the renormalized self-intersection local times of planar random walks
## 1 Introduction
Let $`\{S_n\}`$ be a symmetric random walk on $`^2`$ with covariance matrix $`\mathrm{\Gamma }`$. Let
(1.1)
$$B_n=\underset{1j<kn}{}\delta (S_j,S_k)$$
where
(1.2)
$$\delta (x,y)=\{\begin{array}{cc}1\hfill & \text{ if }x=y\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}$$
is the usual Kroenecker delta. We refer to $`B_n`$ as the self-intersection local time up to time $`n`$. We call
$$\gamma _n=:B_n𝔼B_n$$
the renormalized self-intersection local time of the random walk up to time $`n`$.
In it was shown that $`\gamma _n`$, appropriately scaled, converges to the renormalized self-intersection local time of planar Brownian motion. Renormalized self-intersection local time for Brownian motion was originally studied by Varadhan for its role in quantum field theory. Renormalized self-intersection local time turns out to be the right tool for the solution of certain “classical” problems such as the asymptotic expansion of the area of the Wiener sausage in the plane and the range of random walks, , , .
One of the applications of self-intersection local time is to polymer growth. If $`S_n`$ is a planar random walk and $``$ is its law, one can construct self-repelling and self-attracting random walks by defining
$$d_n/d=c_ne^{\zeta B_n/n},$$
where $`\zeta `$ is a parameter and $`c_n`$ is chosen to make $`_n`$ a probability measure. When $`\zeta <0`$, more weight is given to those paths with a small number of self-intersections, hence $`_n`$ is a model for a self-repelling random walk. When $`\zeta >0`$, more weight is given to paths with a large number of self-intersections, leading to a self-attracting random walk. Since $`𝔼B_n`$ is deterministic, by modifying $`c_n`$, we can write
$$d_n/d=c_ne^{\zeta (B_n𝔼B_n)/n}.$$
It is known that for small positive $`\zeta `$ the self-attracting random walk grows with $`n`$ while for large $`\zeta `$ it “collapses,” and its diameter remains bounded in mean square. It has been an open problem to determine the critical value of $`\zeta `$ at which the phase transition takes place. The work suggested that the critical value $`\zeta _c`$ could be expressed in terms of the best constant of a certain Gagliardo-Nirenberg inequality, but that work was for planar Brownian motion, not for random walks. In the current paper we obtain moderate deviations estimates for $`\gamma _n`$ and these are in terms of the best constant of the Gagliard-Nirenberg inequality; see Theorem 1.1. However the critical constant $`\zeta _c`$ is different (see Remark 4.3) and it is still an open problem to determine it. See and for details and further information on these models.
In the present paper we study moderate deviations of $`\gamma _n`$. Before stating our main theorem we recall one of the Gagliardo-Nirenberg inequalities:
$$f_4Cf_2^{1/2}f_2^{1/2},$$
which is valid for $`fC^1`$ with compact support, and can then be extended to more general $`f`$’s. We define $`\kappa (2,2)`$ to be the infimum of those values of $`C`$ for which the above inequality holds. In particular, $`0<\kappa (2,2)<\mathrm{}`$. For further details, see .
In this paper we will always assume that the smallest group which supports $`\{S_n\}`$ is $`^2`$. For simplicity we assume further that our random walk is strongly aperiodic.
###### Theorem 1.1
Let $`\{b_n\}`$ be a positive sequence satisfying
(1.3)
$$\underset{n\mathrm{}}{lim}b_n=\mathrm{}\text{and}\text{ }b_n=o(n).$$
For any $`\lambda >0`$,
(1.4)
$$\underset{n\mathrm{}}{lim}\frac{1}{b_n}\mathrm{log}\left\{B_n𝔼B_n\lambda nb_n\right\}=\lambda \sqrt{det\mathrm{\Gamma }}\kappa (2,2)^4.$$
We call Theorem 1.1 a moderate deviations theorem rather than a large deviations result because of the second restriction in (1.3). Our techniques do not apply when this restriction is not present, and and in fact it is not hard to show that the value on the right hand side of (1.4) should be different when $`b_nn`$; see Remark 4.3.
Moderate deviations for $`\gamma _n`$ are more subtle. In the next theorem we obtain the correct rate, but not the precise constant.
###### Theorem 1.2
Suppose $`𝔼|S_1|^{2+\delta }<\mathrm{}`$ for some $`\delta >0`$. There exist $`C_1,C_2>0`$ such that for any $`\theta >0`$ and sequence $`b_n\mathrm{}\text{with}\text{ }b_n=o(n^{1/\theta })`$
(1.5) $`C_1`$ $``$ $`\underset{n\mathrm{}}{lim\; inf}b_n^\theta \mathrm{log}\left\{𝔼B_nB_n\theta (2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$``$ $`\underset{n\mathrm{}}{lim\; sup}b_n^\theta \mathrm{log}\left\{𝔼B_nB_n\theta (2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$``$ $`C_2.`$
Here are the corresponding laws of the iterated logarithm for $`\gamma _n`$.
###### Theorem 1.3
(1.6)
$$\underset{n\mathrm{}}{lim\; sup}\frac{B_n𝔼B_n}{n\mathrm{log}\mathrm{log}n}=det(\mathrm{\Gamma })^{1/2}\kappa (2,2)^4a.s.$$
and if $`𝔼|S_1|^{2+\delta }<\mathrm{}`$ for some $`\delta >0`$,
(1.7)
$$\underset{n\mathrm{}}{lim\; inf}\frac{B_n𝔼B_n}{n\mathrm{log}\mathrm{log}\mathrm{log}n}=(2\pi )^1det(\mathrm{\Gamma })^{1/2}a.s.$$
In this paper we deal exclusively with the case where the dimension $`d`$ is 2. We note that in dimension $`1`$ no renormalization is needed, which makes the results much simpler. See . When $`d3`$, the renormalized intersection local time is in the domain of attraction of a centered normal random variable. Consequently the tails of the weak limit are expected to be of Gaussian type, and in particular, the tails are symmetric; see .
Theorems 1.1-1.3 are the analogues of the theorems proven in for the renormalized self-intersection local time of planar Brownian motion. Although the proofs for the random walk case have some elements in common with those for Brownian motion, the random walk case is considerably more difficult. The major difficulty is the fact that we do not have Gaussian random variables. Consequently, the argument for the lower bound of Theorem 1.1 needs to be very different from the one given in \[2, Lemma 3.4\]. This requires several new tools, such as Theorem 4.1, which we expect will have applications beyond the specific needs of this paper.
## 2 Integrability
Let $`\{S_n^{}\}`$ be an independent copy of the random walk $`\{S_n\}`$. Let
(2.1)
$$I_{m,n}=\underset{j=1}{\overset{m}{}}\underset{k=1}{\overset{n}{}}\delta (S_j,S_k^{})$$
and set $`I_n=I_{n,n}`$. Thus
(2.2)
$$I_n=\mathrm{\#}\{(j,k)[1,n]^2;S_j=S_k^{}\}.$$
###### Lemma 2.1
(2.3)
$$𝔼I_{m,n}c\left((m+n)\mathrm{log}(m+n)m\mathrm{log}mn\mathrm{log}n\right).$$
In particular
(2.4)
$$𝔼(I_n)cn.$$
We also have
(2.5)
$$𝔼I_{m,n}c\sqrt{mn}.$$
ProofUsing symmetry and independence
$`𝔼I_{m,n}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}𝔼\delta (S_j,S_k^{})`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}𝔼\delta (S_jS_k^{},0)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}𝔼\delta (S_{j+k},0)={\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}p_{j+k}(0).`$
By \[17, p. 75\],
(2.7)
$$p_m(0)=\frac{1}{2\pi \sqrt{det\mathrm{\Gamma }}}\frac{1}{m}+o\left(\frac{1}{m}\right)$$
so that
(2.8)
$$𝔼I_{m,n}c\underset{j=1}{\overset{m}{}}\underset{k=1}{\overset{n}{}}\frac{1}{j+k}c_{r=0}^m_{s=0}^n\frac{1}{r+s}𝑑r𝑑s$$
and (2.3) follows. (2.2) is then immediate. (2.5) follows from (2.8) and the bound $`(r+s)^1(\sqrt{rs})^1`$.
It follows from the proof of \[8, Lemma 5.2\] that for any integer $`k1`$
(2.9)
$$𝔼(I_n^k)(k!)^2(1+𝔼(I_n))^k$$
Furthermore, by \[13, (5.k)\] we have that $`I_n/n`$ converges in distribution to a random variable with finite moments. Hence for any integer $`k1`$
(2.10)
$$\underset{n\mathrm{}}{lim}\frac{𝔼(I_n^k)}{n^k}=c_k<\mathrm{}.$$
###### Lemma 2.2
There is a constant $`c>0`$ such that
(2.11)
$$\underset{n}{sup}𝔼\mathrm{exp}\left\{\frac{c}{n}I_n\right\}<\mathrm{}.$$
Proof. For any $`m1`$ write $`l(m,n)=[n/m]+1`$. Using \[8, Theorem 5.1\] with $`p=2`$ and $`a=m`$, and then (2.4), (2.9) and (2.10), we obtain
$`\left(𝔼I_n^m\right)^{1/2}`$ $``$ $`{\displaystyle \underset{\stackrel{k_1+\mathrm{}+k_m=m}{k_1,\mathrm{},k_m0}}{}}{\displaystyle \frac{m!}{k_1!\mathrm{}k_m!}}\left(𝔼I_{l(m,n)}^{k_1}\right)^{1/2}\mathrm{}\left(𝔼I_{l(m,n)}^{k_m}\right)^{1/2}`$
$``$ $`{\displaystyle \underset{\stackrel{k_1+\mathrm{}+k_m=m}{k_1,\mathrm{},k_m0}}{}}{\displaystyle \frac{C^mm!}{k_1!\mathrm{}k_m!}}k_1!\mathrm{}k_m!\left(𝔼I_{l(m,n)}\right)^{k_1/2}\mathrm{}\left(𝔼I_{l(m,n)}\right)^{k_m/2}`$
$``$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{2m1}{m}}\right)m!C^m\left({\displaystyle \frac{n}{m}}\right)^{m/2}\left({\displaystyle \genfrac{}{}{0pt}{}{2m}{m}}\right)m!C^m\left({\displaystyle \frac{n}{m}}\right)^{m/2}`$
where $`C>0`$ can be chosen independently of $`m`$ and $`n`$. Hence
(2.13)
$$𝔼I_n^m\left(\genfrac{}{}{0pt}{}{2m}{m}\right)^2C^m(m!)^2\left(\frac{n}{m}\right)^m\left(\genfrac{}{}{0pt}{}{2m}{m}\right)^2C^mm!n^m.$$
Notice that
(2.14)
$$\left(\genfrac{}{}{0pt}{}{2m}{m}\right)4^m.$$
The conclusion then follows using the power series for $`e^x`$.
For any random variable $`X`$ we define
$$\overline{X}=:X𝔼X.$$
We write
(2.15)
$$(m,n]_<^2=\{(j,k)(m,n]^2;j<k\}$$
For any $`A\left\{(j,k)(^+)^2;j<k\right\}`$, write
(2.16)
$$B(A)=\underset{(j,k)A}{}\delta (S_j,S_k)$$
In our proofs we will use several decompositions of $`B_n`$. If $`J_1,\mathrm{},J_{\mathrm{}}`$ are consecutive disjoint blocks of integers whose union is $`\{1,\mathrm{},n\}`$, we have
$$B_n=\underset{i}{}B((J_i\times J_i)(0,n]_<)+\underset{i<j}{}B(J_i\times J_j)$$
and also
$$B_n=\underset{i}{}B((J_i\times J_i)(0,n]_<)+\underset{i}{}B(_{j=1}^{i1}J_j)\times J_i).$$
###### Lemma 2.3
There is a constant $`c>0`$ such that
(2.17)
$$\underset{n}{sup}𝔼\mathrm{exp}\left\{\frac{c}{n}|\overline{B}_n|\right\}<\mathrm{}.$$
Proof. We first prove that there is $`c>0`$ such that
(2.18)
$$M\underset{n}{sup}𝔼\mathrm{exp}\left\{\frac{c}{2^n}|\overline{B}_{2^n}|\right\}<\mathrm{}.$$
We have
(2.19) $`B_{2^n}`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}B\left(((2k2)2^{nj},(2k1)2^{nj}]\times ((2k1)2^{nj},(2k)2^{nj}]\right)`$
Write
(2.20) $`\alpha _{j,k}=B\left(((2k2)2^{nj},(2k1)2^{nj}]\times ((2k1)2^{nj},(2k)2^{nj}]\right)`$
$`𝔼B\left(((2k2)2^{nj},(2k1)2^{nj}]\times ((2k1)2^{nj},(2k)2^{nj}]\right)`$
For each $`1jn`$, the random variables $`\alpha _{j,k}`$, $`k=1,\mathrm{},2^{j1}`$ are i.i.d. with common distribution $`I_{2^{nj}}𝔼I_{2^{nj}}`$. By the previous lemma there exists $`\delta >0`$ such that
(2.21)
$$\underset{n}{sup}\underset{jn}{sup}𝔼\mathrm{exp}\left\{\delta \frac{1}{2^{nj}}\left|\alpha _{j,1}\right|\right\}<\mathrm{}.$$
By \[3, Lemma 1\], there exists $`\theta >0`$ such that
$`C(\theta )`$ $``$ $`\underset{n}{sup}\underset{jn}{sup}𝔼\mathrm{exp}\left\{\theta 2^{j/2}{\displaystyle \frac{1}{2^n}}\left|{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}\alpha _{j,k}\right|\right\}`$
$`=`$ $`\underset{n}{sup}\underset{jn}{sup}𝔼\mathrm{exp}\left\{\theta 2^{j/2}{\displaystyle \frac{1}{2^{nj}}}\left|{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}\alpha _{j,k}\right|\right\}<\mathrm{}.`$
Write
(2.23)
$$\lambda _N=\underset{j=1}{\overset{N}{}}\left(12^{j/2}\right)\text{and}\text{ }\lambda _{\mathrm{}}=\underset{j=1}{\overset{\mathrm{}}{}}\left(12^{j/2}\right).$$
Using Hölder’s inequality with $`1/p=12^{n/2},\mathrm{\hspace{0.17em}1}/q=2^{n/2}`$ we have
(2.24) $`𝔼\mathrm{exp}\left\{\lambda _n{\displaystyle \frac{\theta }{2^n}}\left|{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}\alpha _{j,k}\right|\right\}`$
$`\left(𝔼\mathrm{exp}\left\{\lambda _{n1}{\displaystyle \frac{\theta }{2^n}}\left|{\displaystyle \underset{j=1}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}\alpha _{j,k}\right|\right\}\right)^{12^{n/2}}`$
$`\times \left(𝔼\mathrm{exp}\left\{2^{n/2}\lambda _n{\displaystyle \frac{\theta }{2^n}}\left|{\displaystyle \underset{k=1}{\overset{2^{n1}}{}}}\alpha _{n,k}\right|\right\}\right)^{2^{n/2}}`$
$`𝔼\mathrm{exp}\left\{\lambda _{n1}{\displaystyle \frac{\theta }{2^n}}\left|{\displaystyle \underset{j=1}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}\alpha _{j,k}\right|\right\}C(\theta )^{2^{n/2}}`$
Repeating this procedure,
(2.25) $`𝔼\mathrm{exp}\left\{\lambda _n{\displaystyle \frac{\theta }{2^n}}\left|{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{k=1}{\overset{2^{j1}}{}}}\alpha _{j,k}\right|\right\}`$
$`C(\theta )^{2^{1/2}+\mathrm{}+2^{n/2}}C(\theta )^{2^{1/2}(12^{1/2})^1}`$
So we have
(2.26)
$$\underset{n}{sup}𝔼\mathrm{exp}\left\{\lambda _{\mathrm{}}\frac{\theta }{2^n}|\overline{B}_{2^n}|\right\}<\mathrm{}$$
We now prove our lemma for general $`n`$. Given an integer $`n2`$, we have the following unique representation:
(2.27)
$$n=2^{m_1}+2^{m_2}+\mathrm{}+2^{m_l}$$
where $`m_1>m_2>\mathrm{}m_l0`$ are integers. Write
(2.28)
$$n_0=0\text{and}\text{ }n_i=2^{m_1}+\mathrm{}+2^{m_i},i=1,\mathrm{},l.$$
Then
(2.29) $`{\displaystyle \underset{1j<kn}{}}\delta (S_j,S_k)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \underset{n_{i1}<j<kn_i}{}}\delta (S_j,S_k)+{\displaystyle \underset{i=1}{\overset{l1}{}}}B\left((n_{i1},n_i]\times (n_i,n]\right)`$
$`=`$ $`:{\displaystyle \underset{i=1}{\overset{l}{}}}B_{2^{m_i}}^{(i)}+{\displaystyle \underset{i=1}{\overset{l1}{}}}A_i.`$
By Hölder’s inequality, with $`M`$ as in (2.18)
(2.30) $`𝔼\mathrm{exp}\left\{{\displaystyle \frac{c}{n}}\left|{\displaystyle \underset{i=1}{\overset{l}{}}}(B_{2^{m_i}}^{(i)}𝔼B_{2^{m_i}}^{(i)})\right|\right\}`$
$`{\displaystyle \underset{i=1}{\overset{l}{}}}\left(𝔼\mathrm{exp}\left\{{\displaystyle \frac{c}{2^{m_i}}}|B_{2^{m_i}}^{(i)}𝔼B_{2^{m_i}}^{(i)}|\right\}\right)^{\frac{2^{m_i}}{n}}{\displaystyle \underset{i=1}{\overset{l}{}}}M^{2^{m_i}/n}=M.`$
Using Hölder’s inequality,
(2.31)
$$𝔼\mathrm{exp}\left\{\frac{c}{n}\underset{i=1}{\overset{l1}{}}A_i\right\}\underset{i=1}{\overset{l1}{}}\left(𝔼\mathrm{exp}\left\{\frac{c}{2^{m_i}}A_i\right\}\right)^{\frac{2^{m_i}}{n}}.$$
Notice that for each $`1il1`$,
(2.32)
$$A_i\stackrel{d}{=}\underset{j=1}{\overset{2^{m_i}}{}}\underset{k=1}{\overset{nn_i}{}}\delta (S_j,S_k^{})\underset{j=1}{\overset{2^{m_i}}{}}\underset{k=1}{\overset{2^{m_i}}{}}\delta (S_j,S_k^{}),$$
where the inequality follows from
(2.33)
$$nn_i=2^{m_{i+1}}+\mathrm{}+2^{m_l}2^{m_i}.$$
Using (2.32) and Lemma 2.1, we can take $`c>0`$ so that
(2.34)
$$𝔼\mathrm{exp}\left\{\frac{c}{2^{m_i}}A_i\right\}\underset{n}{sup}𝔼\mathrm{exp}\left\{\frac{c}{n}I_n\right\}N<\mathrm{}.$$
Consequently,
(2.35)
$$𝔼\mathrm{exp}\left\{\frac{c}{n}\underset{i=1}{\overset{l1}{}}A_i\right\}\underset{i=1}{\overset{l1}{}}N^{2^{m_i}/n}N.$$
In particular, this shows that
(2.36)
$$𝔼\left\{\frac{c}{n}\underset{i=1}{\overset{l1}{}}A_i\right\}N.$$
Combining (2.35) and (2.36) with (2.30) we have
(2.37)
$$\underset{n}{sup}𝔼\mathrm{exp}\left\{\frac{c}{2n}|\overline{B}_n|\right\}<\mathrm{}.$$
###### Lemma 2.4
(2.38)
$$𝔼B_n=\frac{1}{2\pi \sqrt{det\mathrm{\Gamma }}}n\mathrm{log}n+o(n\mathrm{log}n),$$
and if $`𝔼|S_1|^{2+2\delta }<\mathrm{}`$ for some $`\delta >0`$ then
(2.39)
$$𝔼B_n=\frac{1}{2\pi \sqrt{det\mathrm{\Gamma }}}n\mathrm{log}n+O(n).$$
Proof.
(2.40)
$$𝔼B_n=𝔼\underset{1j<kn}{}\delta (S_j,S_k)=\underset{1j<kn}{}p_{kj}(0)$$
where $`p_m(x)=𝔼(S_m=x)`$. If $`𝔼|S_1|^{2+2\delta }<\mathrm{}`$, then by \[12, Proposition 6.7\],
(2.41)
$$p_m(0)=\frac{1}{2\pi \sqrt{det\mathrm{\Gamma }}}\frac{1}{m}+o\left(\frac{1}{m^{1+\delta }}\right).$$
Since the last term is summable, it will contribute $`O(n)`$ to (2.40). Also,
(2.42)
$$\underset{1j<kn}{}\frac{1}{kj}=\underset{m=1}{\overset{n}{}}\underset{i=1}{\overset{nm}{}}\frac{1}{m}=\underset{m=1}{\overset{n}{}}\frac{nm}{m}=n\underset{m=1}{\overset{n}{}}\frac{1}{m}n$$
and our Lemma follows from the well known fact that
(2.43)
$$\underset{m=1}{\overset{n}{}}\frac{1}{m}=\mathrm{log}n+\gamma +O\left(\frac{1}{n}\right)$$
where $`\gamma `$ is Euler’s constant.
If we only assume finite second moments, instead of (2.41) we use (2.7) and proceed as above.
###### Lemma 2.5
For any $`\theta >0`$
(2.44)
$$\underset{n}{sup}𝔼\mathrm{exp}\left\{\frac{\theta }{n}(𝔼B_nB_n)\right\}<\mathrm{}$$
and for any $`\lambda >0`$
(2.45)
$$\underset{n\mathrm{}}{lim}\frac{1}{b_n}\mathrm{log}\left\{𝔼B_nB_n\lambda nb_n\right\}=\mathrm{}.$$
Proof. By Lemma 2.3 this is true for some $`\theta _o>0`$. For any $`\theta >\theta _o`$, take an integer $`m1`$ such that $`\theta m^1<\theta _o`$. We can write any $`n`$ as $`n=rm+i`$ with $`1i<m`$. Then
(2.46) $`𝔼B_nB_n`$
$`{\displaystyle \underset{j=1}{\overset{m}{}}}\left[𝔼{\displaystyle \underset{(j1)r<k,ljr}{}}\delta (S_k,S_l){\displaystyle \underset{(j1)r<k,ljr}{}}\delta (S_k,S_l)\right]+𝔼B_nm𝔼B_r.`$
We claim that
(2.47)
$$𝔼B_nm𝔼B_r=O(n).$$
To see this, write
(2.48)
$$𝔼B_nm𝔼B_r=𝔼B_n\underset{l=1}{\overset{m}{}}𝔼B(((l1)r,lr]_<^2)$$
Notice that
(2.49) $`B_n{\displaystyle \underset{l=1}{\overset{m}{}}}B(((l1)r,lr]_<^2)`$
$`={\displaystyle \underset{l=1}{\overset{m}{}}}B(((l1)r,lr]\times (lr,mr])+B((mr,n]_<^2)`$
$`+B((0,mr]\times (mr,n])`$
Since
(2.50)
$$B(((l1)r,lr]\times (lr,mr])\stackrel{d}{=}I_{r,(ml)r}$$
by (2.3) we have
(2.51) $`𝔼B(((l1)r,lr]\times (lr,mr])`$
$`C\{(m(l1))r)\mathrm{log}(m(l1))r)`$
$`((ml)r)\mathrm{log}((ml)r)r\mathrm{log}r\}`$
Therefore
(2.52) $`{\displaystyle \underset{l=1}{\overset{m}{}}}𝔼B(((l1)r,lr]\times (lr,mr])`$
$`C{\displaystyle \underset{l=1}{\overset{m}{}}}\{(m(l1))r)\mathrm{log}(m(l1))r)`$
$`((ml)r)\mathrm{log}((ml)r)r\mathrm{log}r\}`$
$`=C\left\{mr\mathrm{log}mrmr\mathrm{log}r\right\}=Cmr\mathrm{log}m.`$
Using (2.5) for $`𝔼B((0,mr]\times (mr,n])=𝔼I_{mr,i}`$ and (2.38) for $`𝔼B((mr,n]_<^2)`$ then completes the proof of (2.47).
Note that the summands in (2.46) are independent. Therefore, for some constant $`C>0`$ depending only on $`\theta `$ and $`m`$,
(2.53)
$$𝔼\mathrm{exp}\left\{\frac{\theta }{n}(𝔼B_nB_n)\right\}C\left(𝔼\mathrm{exp}\left\{\frac{\theta }{n}(𝔼B_rB_r)\right\}\right)^m$$
which proves (2.44), since $`\theta /n\theta /mr<\theta _o/r`$ and $`r\mathrm{}`$ as $`n\mathrm{}`$.
Then, by Chebychev’s inequality, for any fixed $`h>0`$
(2.54)
$$\left\{𝔼B_nB_n\lambda nb_n\right\}e^{h\lambda b_n}𝔼\mathrm{exp}\left\{\frac{h}{n}(𝔼B_nB_n)\right\}$$
so that by (2.44)
(2.55)
$$\underset{n\mathrm{}}{lim}\frac{1}{b_n}\mathrm{log}\left\{𝔼B_nB_n\lambda nb_n\right\}h\lambda .$$
Since $`h>0`$ is arbitrary, this proves (2.45).
## 3 Proof of Theorem 1.1
By the Gärtner-Ellis theorem ( \[11, Theorem 2.3.6\]), we need only prove
(3.1)
$$\underset{n\mathrm{}}{lim}\frac{1}{b_n}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{\frac{b_n}{n}}|B_n𝔼B_n|^{1/2}\right\}=\frac{1}{4}\kappa (2,2)^4\theta ^2det(\mathrm{\Gamma })^{1/2}.$$
Indeed, by the Gärtner-Ellis theorem the above implies that
(3.2)
$$\underset{n\mathrm{}}{lim}\frac{1}{b_n}\mathrm{log}\left\{|B_n𝔼B_n|\lambda nb_n\right\}=\lambda \sqrt{det(\mathrm{\Gamma })}\kappa (2,2)^4.$$
Using (2.45) we will then have Theorem 1.1. It thus remains to prove (3.1).
Let $`f`$ be a symmetric probability density function in the Schwartz space $`𝒮(^2)`$ of $`C^{\mathrm{}}`$ rapidly decreasing functions. Let $`ϵ>0`$ be a small number and write
(3.3)
$$f_ϵ(x)=ϵ^2f(ϵ^1x),x^2$$
and
(3.4)
$$l(n,x)=\underset{k=1}{\overset{n}{}}\delta (S_k,x),l(n,x,ϵ)=\underset{k=1}{\overset{n}{}}f_{ϵ(b_n^1n)^{1/2}}(S_kx).$$
By \[8, Theorem 3.1\],
(3.5) $`\underset{n\mathrm{}}{lim}{\displaystyle \frac{1}{b_n}}\mathrm{log}𝔼\mathrm{exp}\left\{{\displaystyle \frac{\theta }{\sqrt{2}}}\sqrt{{\displaystyle \frac{b_n}{n}}}\left({\displaystyle \underset{x^2}{}}l^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`=\underset{g_2}{sup}\left\{{\displaystyle \frac{\theta }{\sqrt{2}}}\left({\displaystyle _^2}|g^2f_ϵ(x)|^2𝑑x\right)^{1/2}{\displaystyle \frac{1}{2}}{\displaystyle _^2}g,\mathrm{\Gamma }g𝑑x\right\}`$
where
(3.6)
$$_2=\left\{gW^{1,2}(^2)\right|g_2=1\}.$$
As in the proof of \[10, Theorem 1\], (3.1) will follow from (3.5) and the next Theorem.
###### Theorem 3.1
For any $`\theta >0`$,
$$\underset{ϵ0}{lim}\underset{n\mathrm{}}{lim}\frac{1}{b_n}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{\frac{b_n}{n}}|B_n𝔼B_n\frac{1}{2}\underset{x^2}{}l^2(n,x,ϵ)|^{1/2}\right\}=0.$$
Proof. Let $`l>1`$ be a large but fixed integer. Divide $`[1,n]`$ into $`l`$ disjoint subintervals $`D_1,\mathrm{},D_l`$, each of length $`[n/l]`$ or $`[n/l]+1`$. Write
(3.7)
$$D_i^{}=\{(j,k)D_i^2;j<k\}i=1,\mathrm{},l$$
With the notation of (2.16) we have
(3.8)
$$B_n=\underset{i=1}{\overset{l}{}}B(D_i^{})+\underset{1j<kl}{}B(D_j\times D_k)$$
Notice that
(3.9) $`B(D_j\times D_k)`$ $`=`$ $`{\displaystyle \underset{n_1D_j,n_2D_k}{}}\delta (S_{n_1},S_{n_2})`$
$`=`$ $`{\displaystyle \underset{n_1D_j,n_2D_k}{}}\delta ((S_{n_1}S_{b_j})+S_{b_j},S_{a_k}+(S_{n_2}S_{a_k}))`$
$`=`$ $`{\displaystyle \underset{n_1D_j,n_2D_k}{}}\delta ((S_{n_1}S_{b_j}),Z+(S_{n_2}S_{a_k}))`$
with $`Z\stackrel{d}{=}S_{a_k}S_{b_j}`$, so that $`Z,S_{n_1}S_{b_j},S_{n_2}S_{a_k}`$ are independent. Then as in (2)
(3.10)
$$𝔼B(D_j\times D_k)=𝔼\underset{n_1D_j,n_2D_k}{}p_{n_1+n_2}(Z)\underset{n_1D_j,n_2D_k}{}p_{n_1+n_2}(0)$$
since $`sup_xp_j(x)=p_j(0)`$ for a symmetric random walk. Then as in the proof of (2.4) we have that
(3.11)
$$𝔼B(D_j\times D_k)cn/l.$$
Hence,
(3.12) $`B_n𝔼B_n`$
$`={\displaystyle \underset{i=1}{\overset{l}{}}}\left[B(D_i^{})𝔼B(D_i^{})\right]+{\displaystyle \underset{1j<kl}{}}B(D_j\times D_k)𝔼{\displaystyle \underset{1j<kl}{}}B(D_j\times D_k)`$
$`={\displaystyle \underset{i=1}{\overset{l}{}}}\left[B(D_i^{})𝔼B(D_i^{})\right]+{\displaystyle \underset{1j<kl}{}}B(D_j\times D_k)+O(n)`$
where the last line follows from (3.11).
Write
(3.13)
$$\xi _i(n,x,ϵ)=\underset{kD_i}{}f_{ϵ(b_n^1n)^{1/2}}(S_kx).$$
Then
(3.14)
$$\underset{x^2}{}l^2(n,x,ϵ)=\underset{i=1}{\overset{l}{}}\underset{x^2}{}\xi _i^2(n,x,ϵ)+2\underset{1jkl}{}\underset{x^2}{}\xi _j(n,x,ϵ)\xi _k(n,x,ϵ).$$
Therefore, by (3.12)
(3.15) $`\left|(B_n𝔼B_n){\displaystyle \frac{1}{2}}{\displaystyle \underset{x^2}{}}l^2(n,x,ϵ)\right|`$
$`{\displaystyle \underset{i=1}{\overset{l}{}}}\left|B(D_i^{})𝔼B(D_i^{})\right|+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)`$
$`+{\displaystyle \underset{1j<kl}{}}\left|B(D_j\times D_k){\displaystyle \underset{x^2}{}}\xi _j(n,x,ϵ)\xi _k(n,x,ϵ)\right|+O(n).`$
The proof of Theorem 3.1 is completed in the next two lemmas.
###### Lemma 3.2
For any $`\theta >0`$,
(3.16) $`\underset{n\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{b_n}}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}\left({\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`l^1{\displaystyle \frac{1}{2}}\kappa (2,2)^4\theta ^2det(\mathrm{\Gamma })^{1/2}`$
and
(3.17)
$$\underset{n\mathrm{}}{lim\; sup}\frac{1}{b_n}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{\frac{b_n}{n}}\left(\underset{i=1}{\overset{l}{}}\left|B(D_i^{})𝔼B(D_i^{})\right|\right)^{1/2}\right\}l^1H\theta ^2,$$
where
(3.18)
$$H=\left(sup\left\{\lambda >0;\underset{n}{sup}𝔼\mathrm{exp}\left\{\lambda \frac{1}{n}|B_n𝔼B_n|\right\}<\mathrm{}\right\}\right)^2.$$
Proof. Replacing $`\theta `$ by $`\theta /\sqrt{l}`$, $`n`$ by $`n/l`$, and $`b_n`$ by $`b_n^{}=b_{ln}`$ (notice that $`b_{n/l}^{}=b_n`$)
(3.19) $`\underset{n\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{b_n}}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}\left({\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`=\underset{n\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{b_{n/l}^{}}}\mathrm{log}𝔼\mathrm{exp}\left\{{\displaystyle \frac{\theta }{\sqrt{l}}}\sqrt{{\displaystyle \frac{b_{n/l}^{}}{n/l}}}\left({\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\}`$
Applying Jensen’s inequality on the right hand side of (3.5),
$`{\displaystyle _^2}|g^2f_\epsilon (x)|^2`$ $`=`$ $`{\displaystyle _^2}\left[{\displaystyle _^2}g^2(xy)f_\epsilon (y)𝑑y\right]^2𝑑x`$
$``$ $`{\displaystyle g^4(xy)f_\epsilon (y)𝑑y𝑑x}={\displaystyle f_\epsilon (y)\left[g^4(xy)𝑑x\right]𝑑y}`$
$`=`$ $`\left[{\displaystyle g^4(x)𝑑x}\right]{\displaystyle f_\epsilon (y)𝑑y}={\displaystyle _^2}g^4(y)𝑑y.`$
Combining the last two displays with (3.5) we have that
(3.20) $`\underset{n\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{b_n}}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}\left({\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`\underset{g_2}{sup}\left\{{\displaystyle \frac{\theta }{\sqrt{l}}}\left({\displaystyle _^2}|g(x)|^4𝑑x\right)^{1/2}{\displaystyle \frac{1}{2}}{\displaystyle _^2}g(x),\mathrm{\Gamma }g(x)𝑑x\right\}`$
$`=l^1\theta ^2\underset{h_2}{sup}\left\{\left({\displaystyle _^2}|h(x)|^4𝑑x\right)^{1/2}{\displaystyle \frac{1}{2}}{\displaystyle _^2}|h(x)|^2𝑑x\right\}`$
$`={\displaystyle \frac{1}{2}}l^1det(\mathrm{\Gamma })^{1/2}\kappa (2,2)^4\theta ^2,`$
where the third line follows from the substitution $`g(x)=\sqrt{|det(A)|}f(Ax)`$ with a $`2\times 2`$ matrix $`A`$ satisfying
(3.21)
$$A^\tau \mathrm{\Gamma }A=\frac{1}{2}\theta ^2\sqrt{det(\mathrm{\Gamma })}𝐈_2$$
and the last line in \[8, Lemma A.2 \]; here $`𝐈_2`$ is the $`2\times 2`$ identity matrix.
Given $`\delta >0`$, there exist $`\overline{a}_1=(a_{1,1},\mathrm{},a_{1,l}),\mathrm{},\overline{a}_m=(a_{m,1},\mathrm{},a_{m,l})`$ in $`^l`$ such that $`|\overline{a}_1|=\mathrm{}=|\overline{a}_m|=1`$ and
(3.22)
$$|z|(1+\delta )\mathrm{max}\{\overline{a}_1z,\mathrm{},\overline{a}_mz\},z^l.$$
In particular, with
(3.23)
$$z=(\left(\underset{x^2}{}\xi _1^2(n,x,ϵ)\right)^{1/2},\mathrm{},\left(\underset{x^2}{}\xi _l^2(n,x,ϵ)\right)^{1/2})$$
we have
(3.24)
$$\left(\underset{i=1}{\overset{l}{}}\underset{x^2}{}\xi _i^2(n,x,ϵ)\right)^{1/2}(1+\delta )\underset{1jm}{\mathrm{max}}\underset{i=1}{\overset{l}{}}a_{j,i}\left(\underset{x^2}{}\xi _i^2(n,x,ϵ)\right)^{1/2}.$$
Hence
(3.25) $`𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}\left({\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`{\displaystyle \underset{j=1}{\overset{m}{}}}𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}(1+\delta ){\displaystyle \underset{i=1}{\overset{l}{}}}a_{j,i}\left({\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`={\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \underset{i=1}{\overset{l}{}}}𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}(1+\delta )a_{j,i}\left({\displaystyle \underset{x^2}{}}\xi _i^2(n,x,ϵ)\right)^{1/2}\right\},`$
where the last line follows from independence of $`\xi _i(n,x,ϵ)_{L^2(^2)}`$, $`i=1,\mathrm{},l`$. Therefore
(3.26) $`\underset{n\mathrm{}}{lim\; sup}{\displaystyle \frac{1}{b_n}}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}\left({\displaystyle \underset{k=1}{\overset{l}{}}}{\displaystyle \underset{x^2}{}}\xi _k^2(n,x,ϵ)\right)^{1/2}\right\}`$
$`\underset{1jm}{\mathrm{max}}{\displaystyle \frac{1}{2}}l^1\kappa (2,2)^4(1+\delta )^2\theta ^2\left({\displaystyle \underset{i=1}{\overset{l}{}}}a_{j,i}^2\right)`$
$`={\displaystyle \frac{1}{2}}l^1det(\mathrm{\Gamma })^{1/2}\kappa (2,2)^4(1+\delta )^2\theta ^2.`$
Letting $`\delta 0^+`$ proves (3.16).
By the inequality $`aba^2+b^2`$ we have that
(3.27) $`𝔼\mathrm{exp}\left\{\theta \sqrt{{\displaystyle \frac{b_n}{n}}}|B_n𝔼B_n|^{1/2}\right\}`$
$`\mathrm{exp}\left\{c^2\theta ^2b_n\right\}𝔼\mathrm{exp}\left\{c^2{\displaystyle \frac{1}{n}}|B_n𝔼B_n|\right\},`$
and taking $`c^2H^2`$ we see that for any $`\theta >0`$,
(3.28)
$$\underset{n\mathrm{}}{lim\; sup}\frac{1}{b_n}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{\frac{b_n}{n}}|B_n𝔼B_n|^{1/2}\right\}H^2\theta ^2.$$
Notice that for any $`1il`$,
(3.29)
$$B(D_i^{})𝔼B(D_i^{})\stackrel{d}{=}B_{\mathrm{\#}(D_i)}𝔼B_{\mathrm{\#}(D_i)}.$$
We have
$$𝔼\mathrm{exp}\{\theta \sqrt{\frac{b_n}{n}}|B(D_i^{})𝔼B(D_i^{})|^{1/2}\}=𝔼\mathrm{exp}\{\frac{\theta }{\sqrt{l}}\sqrt{\frac{b_n}{n/l}}|B(D_i^{}𝔼B(D_i^{})|\}.$$
Replacing $`\theta `$ by $`\theta /\sqrt{l}`$, $`n`$ by $`n/l`$, and $`b_n`$ by $`b_n^{}=b_{ln}`$ (notice that $`b_{n/l}^{}=b_n`$) gives
(3.30)
$$\underset{n\mathrm{}}{lim\; sup}\frac{1}{b_n}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{\frac{b_n}{n}}|B(D_i^{})𝔼B(D_i^{})|^{1/2}\right\}l^1H^2\theta ^2.$$
Thus (3.17) follows by the same argument we used to prove (3.16).
###### Lemma 3.3
For any $`\theta >0`$ and any $`1j<kl`$,
(3.31)
$$\underset{ϵ0^+}{lim\; sup}\underset{n\mathrm{}}{lim\; sup}\frac{1}{b_n}\mathrm{log}𝔼\mathrm{exp}\left\{\theta \sqrt{\frac{b_n}{n}}\left|B(D_j\times D_k)\underset{x^2}{}\xi _j(n,x,ϵ)\xi _k(n,x,ϵ)\right|^{1/2}\right\}=0.$$
Proof. Define $`a_j,b_j`$ so that $`D_j=(a_j,b_j]`$ $`(1jl)`$. We now fix $`1j<kl`$ and estimate
(3.32)
$$B(D_j\times D_k)\underset{x^2}{}\xi _j(n,x,ϵ)\xi _k(n,x,ϵ).$$
Without loss of generality we may assume that $`v=:[n/l]=\mathrm{\#}(D_j)=\mathrm{\#}(D_k)`$. For $`yZ^2`$ set
(3.33)
$$I_n(y)=\underset{n_1,n_2=1}{\overset{n}{}}\delta (S_{n_1},S_{n_2}^{}+y).$$
Note that $`I_n=I_n(0)`$. By (3.9) we have that
(3.34)
$$B(D_j\times D_k)\stackrel{d}{=}I_v(Z)$$
with $`Z`$ independent of $`S,S^{}`$.
Similarly, we have
$`{\displaystyle \underset{x^2}{}}\xi _j(n,x,ϵ)\xi _k(n,x,ϵ)`$
$`={\displaystyle \underset{x^2}{}}{\displaystyle \underset{n_1D_k,n_2D_k}{}}f_{ϵ(b_n^1n)^{1/2}}(S_{n_1}x)f_{ϵ(b_n^1n)^{1/2}}(S_{n_2}x)`$
$`={\displaystyle \underset{x^2}{}}{\displaystyle \underset{n_1D_k,n_2D_k}{}}f_{ϵ(b_n^1n)^{1/2}}(x)f_{ϵ(b_n^1n)^{1/2}}(S_{n_2}S_{n_1}x)`$
$`={\displaystyle \underset{n_1D_k,n_2D_k}{}}f_{ϵ(b_n^1n)^{1/2}}f_{ϵ(b_n^1n)^{1/2}}(S_{n_2}S_{n_1})`$
(3.35) $`={\displaystyle \underset{n_1D_j,n_2D_k}{}}f_{ϵ(b_n^1n)^{1/2}}f_{ϵ(b_n^1n)^{1/2}}((S_{n_2}S_{a_k})(S_{n_1}S_{b_j})+Z)`$
where
(3.36)
$$ff(y)=\underset{x^2}{}f(x)f(yx)$$
denotes convolution in $`L^1(^2)`$. It is clear that if $`f𝒮(^2)`$ so is $`ff`$. For $`yZ^2`$, define the link
(3.37)
$$L_{n,ϵ}(y)=\underset{n_1,n_2=1}{\overset{n}{}}f_ϵf_ϵ(S_{n_2}^{}S_{n_1}+y).$$
By (3.35) we have that
(3.38)
$$\underset{x^2}{}\xi _j(n,x,ϵ)\xi _k(n,x,ϵ)\stackrel{d}{=}L_{v,(b_n^1n)^{1/2}ϵ}(Z)$$
with $`Z`$ independent of $`S,S^{}`$.
###### Lemma 3.4
Let $`f𝒮(^2)`$ with Fourier transform $`\widehat{f}`$ supported on $`(\pi ,\pi )^2`$. Then for any $`r1`$
(3.39)
$$e^{i\lambda y}(f_rf_r)(y)𝑑y=\widehat{f}^2(r\lambda ),\lambda ^2.$$
Proof. We have
$`{\displaystyle e^{i\lambda y}(ff)(y)𝑑y}`$ $`=`$ $`{\displaystyle \underset{x^2}{}}f(x){\displaystyle e^{i\lambda y}f(yx)𝑑y}`$
$`=`$ $`\widehat{f}(\lambda ){\displaystyle \underset{x^2}{}}f(x)e^{i\lambda x}`$
$`=`$ $`\widehat{f}(\lambda ){\displaystyle \underset{x^2}{}}\left({\displaystyle e^{ipx}\widehat{f}(p)𝑑p}\right)e^{i\lambda x}.`$
For $`x^2`$
(3.41) $`{\displaystyle e^{ipx}\widehat{f}(p)𝑑p}={\displaystyle \underset{u^2}{}}{\displaystyle _{[\pi ,\pi ]^2}}e^{ipx}\widehat{f}(p+2\pi u)𝑑p`$
and using Fourier inversion
$`{\displaystyle \underset{x^2}{}}\left({\displaystyle e^{ipx}\widehat{f}(p)𝑑p}\right)e^{i\lambda x}`$
(3.42) $`={\displaystyle \underset{u^2}{}}{\displaystyle \underset{x^2}{}}\left({\displaystyle _{[\pi ,\pi ]^2}}e^{ipx}\widehat{f}(p+2\pi u)𝑑p\right)e^{i\lambda x}`$
$`={\displaystyle \underset{u^2}{}}\widehat{f}(\lambda +2\pi u).`$
Thus from (3) we find that
(3.43)
$$e^{i\lambda y}ff(y)𝑑y=\widehat{f}(\lambda )\underset{u^2}{}\widehat{f}(\lambda +2\pi u).$$
Since $`\widehat{f}_r(\lambda )=\widehat{f}(r\lambda )`$ we see that for any $`r>0`$
(3.44)
$$e^{i\lambda y}(f_rf_r)(y)𝑑y=\widehat{f}(r\lambda )\underset{u^2}{}\widehat{f}(r\lambda +2\pi ru).$$
Then if $`r1`$, using the fact that $`\widehat{f}(\lambda )`$ is supported in $`(\pi ,\pi )^2`$, we obtain (3.39).
Taking $`f𝒮(R^2)`$ with $`\widehat{f}(\lambda )`$ supported in $`(\pi ,\pi )^2`$, Lemma 3.3 will follow from Theorem 4.1 of the next section.
## 4 Intersections of Random Walks
Let $`S_1(n),S_2(n)`$ be independent copies of the symmetric random walk $`S(n)`$ in $`Z^2`$ with a finite second moment.
Let $`f`$ be a positive symmetric function in the Schwartz space $`𝒮(R^2)`$ with $`f𝑑x=1`$ and $`\widehat{f}`$ supported in $`(\pi ,\pi )^2`$. Given $`ϵ>0`$, and with the notation of the last section, let us define the link
(4.1)
$$I_{n,ϵ}(y)=\underset{n_1,n_2=1}{\overset{n}{}}f_{(b_n^1n)^{1/2}ϵ}f_{(b_n^1n)^{1/2}ϵ}(S_2(n_2)S_1(n_1)+y))$$
with $`I_{n,ϵ}=I_{n,ϵ}(0)`$.
###### Theorem 4.1
For any $`\lambda >0`$
(4.2) $`\underset{ϵ0}{lim\; sup}\underset{n\mathrm{}}{lim\; sup}\underset{y}{sup}`$
$`{\displaystyle \frac{1}{b_n}}\mathrm{log}E\left(\mathrm{exp}\left\{\lambda \left|{\displaystyle \frac{I_n(y)I_{n,ϵ}(y)}{b_n^1n}}\right|^{1/2}\right\}\right)=0.`$
Proof of Theorem 4.1. We have
(4.3) $`{\displaystyle \frac{1}{b_n^1n}}I_n(y)`$
$`={\displaystyle \frac{1}{b_n^1n}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}\delta (S_1(n_1),S_2(n_2)+y)`$
$`={\displaystyle \frac{1}{b_n^1n(2\pi )^2}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}\left[{\displaystyle _{[\pi ,\pi ]^2}}e^{ip(S_2(n_2)+yS_1(n_1))}𝑑p\right]`$
where from now on we work modulo $`\pm \pi `$. Then by scaling we have
(4.4) $`{\displaystyle \frac{1}{b_n^1n}}I_n(y)`$
$`={\displaystyle \frac{1}{(b_n^1n)^2(2\pi )^2}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}\left[{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^2}}e^{ip(S_2(n_2)+yS_1(n_1))/(b_n^1n)^{1/2}}𝑑p\right]`$
As in (4.3)-(4.4), using Lemma 3.4, the fact that $`ϵ(b_n^1n)^{1/2}1`$ for $`ϵ>0`$ fixed and large enough $`n`$, and abbreviating $`\widehat{h}=\widehat{f}^2`$
(4.5) $`{\displaystyle \frac{1}{b_n^1n}}I_{n,ϵ}(y)`$
$`={\displaystyle \frac{1}{b_n^1n(2\pi )^2}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}\left[{\displaystyle _{R^2}}e^{ip(S_2(n_2)+yS_1(n_1))}\widehat{h}(ϵ(b_n^1n)^{1/2}p)𝑑p\right]`$
$`={\displaystyle \frac{1}{(b_n^1n)^2(2\pi )^2}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}\left[{\displaystyle _{R^2}}e^{ip(S_2(n_2)+yS_1(n_1))/(b_n^1n)^{1/2}}\widehat{h}(ϵp)𝑑p\right].`$
Using our assumption that $`\widehat{h}`$ supported in $`[\pi ,\pi ]^2`$, and that $`ϵ^1(b_n^1n)^{1/2}`$ for $`ϵ>0`$ fixed and large enough $`n`$, we have that
(4.6) $`{\displaystyle \frac{1}{b_n^1n}}I_{n,ϵ}(y)`$
$`={\displaystyle \frac{1}{(b_n^1n)^2(2\pi )^2}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}`$
$`\left[{\displaystyle _{ϵ^1[\pi ,\pi ]^2}}e^{ip(S_2(n_2)+yS_1(n_1))/(b_n^1n)^{1/2}}\widehat{h}(ϵp)𝑑p\right]`$
$`={\displaystyle \frac{1}{(b_n^1n)^2(2\pi )^2}}{\displaystyle \underset{n_1,n_2=1}{\overset{n}{}}}`$
$`\left[{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^2}}e^{ip(S_2(n_2)+yS_1(n_1))/(b_n^1n)^{1/2}}\widehat{h}(ϵp)𝑑p\right].`$
To prove (4.2) it suffices to show that for each $`\lambda >0`$ we have
(4.7) $`\underset{y}{sup}E\left(\mathrm{exp}\left\{\lambda \left|{\displaystyle \frac{I_n(y)I_{n,ϵ}(y)}{b_n^1n}}\right|^{1/2}\right\}\right)`$
$`Cb_n(1C\lambda ϵ^{m/4})^1(1+C\lambda ϵ^{1/4}b_n^{1/2})e^{C\lambda ^2ϵ^{1/2}b_n}.`$
for some $`C<\mathrm{}`$ and all $`ϵ>0`$ sufficiently small.
We begin by expanding
(4.8) $`E\left(\mathrm{exp}\left\{\lambda \left|{\displaystyle \frac{I_n(y)I_{n,ϵ}(y)}{b_n^1n}}\right|^{1/2}\right\}\right)`$
$`={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda ^m}{m!}}E\left(\left|{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right|^{m/2}\right)`$
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\lambda ^m}{m!}}\left(E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^{2m}\right)\right)^{1/4}`$
By (4.4), (4.6) and the symmetry of $`S_1`$ we have
(4.9) $`E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^m\right)`$
$`={\displaystyle \frac{1}{(b_n^1n)^{2m}(2\pi )^{2m}}}{\displaystyle \underset{\stackrel{n_{1,j},n_{2,j}=1}{j=1,\mathrm{},m}}{\overset{n}{}}}{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`E\left(e^{i_{j=1}^mp_j(S_2(n_{2,j})+y+S_1(n_{1,j}))/(b_n^1n)^{1/2}}\right){\displaystyle \underset{j=1}{\overset{m}{}}}(1\widehat{h}(ϵp_j))dp_j.`$
Then
(4.10) $`\left|E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^m\right)\right|`$
$`{\displaystyle \frac{1}{(b_n^1n)^{2m}(2\pi )^{2m}}}{\displaystyle \underset{\stackrel{n_{1,j}=1}{j=1,\mathrm{},m}}{\overset{n}{}}}{\displaystyle \underset{\stackrel{n_{2,j}=1}{j=1,\mathrm{},m}}{\overset{n}{}}}{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`\left|E\left(e^{i_{j=1}^mp_jS_1(n_{1,j})/(b_n^1n)^{1/2}}\right)\right|`$
$`\left|E\left(e^{i_{j=1}^mp_jS_2(n_{2,j})/(b_n^1n)^{1/2}}\right)\right|{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j.`$
By the Cauchy-Schwarz inequality
(4.11) $`{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}\left|E\left(e^{i_{j=1}^mp_jS_1(n_{1,j})/(b_n^1n)^{1/2}}\right)\right|`$
$`\left|E\left(e^{i_{j=1}^mp_jS_2(n_{2,j})/(b_n^1n)^{1/2}}\right)\right|{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j`$
$`{\displaystyle \underset{i=1}{\overset{2}{}}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
(4.12) $`|E\left(e^{i_{j=1}^mp_jS(n_{i,j})/(b_n^1n)^{1/2}}\right)|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}.`$
Thus
(4.13) $`\left|E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^m\right)\right|^{1/2}`$
$`{\displaystyle \underset{\stackrel{n_j=1}{j=1,\mathrm{},m}}{\overset{n}{}}}{\displaystyle \frac{1}{(b_n^1n)^m(2\pi )^m}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`|E\left(e^{i_{j=1}^mp_jS(n_j)/(b_n^1n)^{1/2}}\right)|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}`$
For any permutation $`\pi `$ of $`\{1.\mathrm{},m\}`$ let
(4.14)
$$D_m(\pi )=\{(n_1,\mathrm{},n_m)|\mathrm{\hspace{0.17em}1}n_{\pi (1)}\mathrm{}n_{\pi (m)}n\}.$$
Using the (non-disjoint) decomposition
$$\{1,\mathrm{},n\}^m=\underset{\pi }{}D_m(\pi )$$
we have from (4.13) that
(4.15) $`\left|E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^m\right)\right|^{1/2}`$
$`{\displaystyle \underset{\pi }{}}{\displaystyle \underset{D_m(\pi )}{}}{\displaystyle \frac{1}{(b_n^1n)^m(2\pi )^m}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`|E\left(e^{i_{j=1}^mp_jS(n_j)/(b_n^1n)^{1/2}}\right)|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}.`$
where the first sum is over all permutations $`\pi `$ of $`\{1.\mathrm{},m\}`$.
Set
(4.16)
$$\varphi (u)=E\left(e^{iuS(1)}\right).$$
It follows from our assumptions that $`\varphi (u)C^2`$, $`\frac{}{u_i}\varphi (0)=0`$ and $`\frac{^2}{u_iu_j}\varphi (0)=E\left(S_{(i)}(1)S_{(j)}(1)\right)`$ where $`S(1)=(S_{(1)}(1),S_{(2)}(1))`$ so that for some $`\delta >0`$
(4.17)
$$\varphi (u)=1E\left((uS(1))^2\right)/2+o(|u|^2),|u|\delta .$$
Then for some $`c_1>0`$
(4.18)
$$\varphi (u)e^{c_1|u|^2},|u|\delta .$$
Strong aperiodicity implies that $`|\varphi (u)|<1`$ for $`u0`$ and $`u[\pi ,\pi ]^2`$. In particular, we can find $`b<1`$ such that $`|\varphi (u)|b`$ for $`\delta |u|`$ and $`u[\pi ,\pi ]^2`$. But clearly we can choose $`c_2>0`$ so that $`be^{c_2|u|^2}`$ for $`u[\pi ,\pi ]^2`$. Setting $`c=\mathrm{min}(c_1,c_2)>0`$ we then have
(4.19)
$$\varphi (u)e^{c|u|^2},u[\pi ,\pi ]^2.$$
On $`D_m(\pi )`$ we can write
(4.20)
$$\underset{j=1}{\overset{m}{}}p_jS(n_j)=\underset{j=1}{\overset{m}{}}(\underset{i=j}{\overset{m}{}}p_{\pi (i)})(S(n_{\pi (j)})S(n_{\pi (j1)})).$$
Hence on $`D_m(\pi )`$
(4.21)
$$E\left(e^{i_{j=1}^mp_jS(n_j)/(b_n^1n)^{1/2}}\right)=\underset{j=1}{\overset{m}{}}\varphi ((\underset{i=j}{\overset{m}{}}p_{\pi (i)})/(b_n^1n)^{1/2})^{(n_{\pi (j)}n_{\pi (j1)})}.$$
Now it is clear that
(4.22) $`{\displaystyle \underset{D_m(\pi )}{}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`|{\displaystyle \underset{j=1}{\overset{m}{}}}\varphi (({\displaystyle \underset{i=j}{\overset{m}{}}}p_{\pi (i)})/(b_n^1n)^{1/2})^{(n_{\pi (j)}n_{\pi (j1)})}|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}`$
$`={\displaystyle \underset{1n_{\pi (1)}\mathrm{}n_{\pi (m)}n}{}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`|{\displaystyle \underset{j=1}{\overset{m}{}}}\varphi (({\displaystyle \underset{i=j}{\overset{m}{}}}p_{\pi (i)})/(b_n^1n)^{1/2})^{(n_{\pi (j)}n_{\pi (j1)})}|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}`$
is independent of the permutation $`\pi `$. Hence writing
(4.23)
$$u_j=\underset{i=j}{\overset{m}{}}p_u$$
we have from (4.15) that
(4.24) $`\left|E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^m\right)\right|^{1/2}`$
$`m!{\displaystyle \underset{1n_1\mathrm{}n_mn}{}}{\displaystyle \frac{1}{(b_n^1n)^m(2\pi )^m}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`|{\displaystyle \underset{j=1}{\overset{m}{}}}\varphi (u_j/(b_n^1n)^{1/2})^{(n_jn_{j1})}|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}.`$
For each $`A\{2,3,\mathrm{},m\}`$ we use $`D_m(A)`$ to denote the subset of $`\{1n_1\mathrm{}n_mn\}`$ for which $`n_j=n_{j1}`$ if and only if $`jA`$. Then we have
(4.25) $`\left|E\left(\left\{{\displaystyle \frac{1}{b_n^1n}}(I_n(y)I_{n,ϵ}(y))\right\}^m\right)\right|^{1/2}`$
$`m!{\displaystyle \underset{A\{2,3,\mathrm{},m\}}{}}{\displaystyle \underset{D_m(A)}{}}{\displaystyle \frac{1}{(b_n^1n)^m(2\pi )^m}}\{{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}`$
$`|{\displaystyle \underset{j=1}{\overset{m}{}}}\varphi (u_j/(b_n^1n)^{1/2})^{(n_jn_{j1})}|^2{\displaystyle \underset{j=1}{\overset{m}{}}}|1\widehat{h}(ϵp_j)|dp_j\}^{1/2}.`$
For any $`uR^d`$ let $`\stackrel{~}{u}`$ denote the representative of $`u\text{ mod }(b_n^1n)^{1/2}2\pi Z^2`$ of smallest absolute value. We note that
(4.26)
$$|\stackrel{~}{u}|=|\stackrel{~}{u}|,\text{ and }|\stackrel{~}{u+v}|=|\stackrel{~}{u}+\stackrel{~}{v}||\stackrel{~}{u}|+|\stackrel{~}{v}|.$$
Using the periodicity of $`\varphi `$ we see that (4.19) implies that for all $`u`$
(4.27)
$$|\varphi (u/(b_n^1n)^{1/2})|e^{c|\stackrel{~}{u}|^2/(b_n^1n)}.$$
Then we have that on $`\{1n_1\mathrm{}n_mn\}`$
(4.28)
$$\left|\underset{j=1}{\overset{m}{}}\varphi (u_j/(b_n^1n)^{1/2})^{(n_jn_{j1})}\right|^2\underset{j=1}{\overset{m}{}}e^{c|\stackrel{~}{u}_j|^2(n_jn_{j1})/(b_n^1n)}$$
Using $`|1\widehat{h}(ϵp_j)|cϵ^{1/2}p_j^{1/2}`$ we bound the integral in (4.25) by
(4.29)
$$c^mϵ^{m/2}_{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}\underset{j=1}{\overset{m}{}}e^{c|\stackrel{~}{u}_j|^2(n_jn_{j1})/(b_n^1n)}|p_j|^{1/2}dp_j.$$
Using (4.23) and (4.26) we have that
(4.30)
$$\underset{j=1}{\overset{m}{}}|p_j|^{1/2}\underset{j=1}{\overset{m}{}}(|\stackrel{~}{u}_j|^{1/2}+|\stackrel{~}{u}_{j+1}|^{1/2})$$
and when we expand the right hand side as a sum of monomials we can be sure that no factor $`|\stackrel{~}{u}_k|^{1/2}`$ appears more than twice. Thus we see that we can bound (4.29) by
(4.31)
$$C^mϵ^{m/2}\underset{h(j)}{\mathrm{max}}_{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}\underset{j=1}{\overset{m}{}}e^{c|\stackrel{~}{u}_j|^2(n_jn_{j1})/(b_n^1n)}|\stackrel{~}{u}_j|^{h(j)/2}dp_j$$
where the $`\mathrm{max}`$ runs over the the set of functions $`h(j)`$ taking values $`0,1\text{ or }2`$ and such that $`_jh(j)=m`$. Changing variables, we thus need to bound
(4.32)
$$_{\mathrm{\Lambda }_n}\underset{j=1}{\overset{m}{}}e^{c|\stackrel{~}{u}_j|^2(n_jn_{j1})/(b_n^1n)}|\stackrel{~}{u}_j|^{h(j)/2}du_j$$
where, see (4.23),
(4.33)
$$\mathrm{\Lambda }_n=\{(u_1,\mathrm{},u_m)|u_ju_{j+1}(b_n^1n)^{1/2}[\pi ,\pi ]^2,j\}.$$
Let $`C_n`$ denote the rectangle $`(b_n^1n)^{1/2}[\pi ,\pi ]^2`$ and let us call any rectangle of the form $`2\pi k+C_n`$, where $`kZ^2`$, an elementary rectangle. Note that any rectangle of the form $`v+C_n`$, where $`vR^2`$, can be covered by $`4`$ elementary rectangles. Hence for any $`vR^2`$ and $`1sn`$
(4.34) $`{\displaystyle _{v+C_n}}e^{c\frac{s}{b_n^1n}|\stackrel{~}{u}|^2}|\stackrel{~}{u}|^a𝑑u`$
$`4{\displaystyle _{R^2}}e^{c\frac{s}{(b_n^1n)}|u|^2}|u|^{h/2}𝑑u`$
$`C\left({\displaystyle \frac{s}{b_n^1n}}\right)^{(1+h/4)}.`$
Similarly
(4.35)
$$_{v+C_n}|\stackrel{~}{u}|^{h/2}𝑑uC(b_n^1n)^{(1+h/4)}.$$
We now bound (4.32) by bounding successively the integration with respect to $`u_1,\mathrm{},u_m`$. Consider first the $`du_1`$ integral, fixing $`u_2,\mathrm{},u_m`$. By (4.33) the $`du_1`$ integral is over the rectangle $`u_2+C_n`$, hence the factors involving $`u_1`$ can be bounded using (4.34). Proceeding inductively, using (4.33) when $`n_jn_{j1}>0`$ and (4.35) when $`n_j=n_{j1}`$, leads to the following bound of (4.32), and hence of (4.29) on $`D_m(A)`$:
(4.36) $`c^mϵ^{m/2}{\displaystyle _{(b_n^1n)^{1/2}[\pi ,\pi ]^{2m}}}{\displaystyle \underset{j=1}{\overset{m}{}}}e^{c|\stackrel{~}{u}_j|^2(n_jn_{j1})/(b_n^1n)}|p_j|^{1/2}dp_j`$
$`C^mϵ^{m/2}{\displaystyle \underset{jA}{}}(b_n^1n)^{(1+h(j)/4)}{\displaystyle \underset{jA^c}{}}\left({\displaystyle \frac{(n_jn_{j1})}{b_n^1n}}\right)^{(1+h(j)/4)}.`$
Here $`A^c`$ means the complement of $`A`$ in $`\{1,\mathrm{},m\}`$, so that $`A^c`$ always contains $`1`$.
Note that
(4.37)
$$(b_n^1n)^{(1+h(j)/4)/2}\frac{1}{b_n^1n}0\text{ as }n\mathrm{}.$$
If $`A=\{i_1,\mathrm{},i_k\}`$ where $`i_1<\mathrm{}<i_k`$ we then obtain for the sum in (4.25) over $`D_m(A)`$, the bound
(4.38) $`C^mϵ^{m/4}\underset{h(j)}{\mathrm{max}}{\displaystyle \underset{1n_{i_1}<\mathrm{}<n_{i_k}n}{}}{\displaystyle \underset{jA^c}{}}\left({\displaystyle \frac{(n_jn_{j1})}{b_n^1n}}\right)^{(1+h(j)/4)/2}{\displaystyle \frac{1}{b_n^1n}}`$
$`C^mϵ^{m/4}\underset{h(j)}{\mathrm{max}}{\displaystyle _{0r_{i_1}<\mathrm{}<r_{i_k}b_n}}{\displaystyle \underset{jA^c}{}}(r_jr_{j1})^{(1/2+h(j)/8)}dr_j`$
$`C^mϵ^{m/4}\underset{h(j)}{\mathrm{max}}{\displaystyle \frac{b_n^{_{jA^c}(1/2h(j)/8)}}{\mathrm{\Gamma }(_{jA^c}(1/2h(j)/8))}}`$
Using this together with (4.25), but with $`m`$ replaced by $`2m`$, and the fact that $`(2m!)^{1/2}/m!2^m`$, we see that (4.7) is bounded by
(4.39) $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}C^m\lambda ^mϵ^{m/4}\left({\displaystyle \underset{A\{2,3,\mathrm{},2m\}}{}}\underset{h(j)}{\mathrm{max}}{\displaystyle \frac{b_n^{_{jA^c}(1/2h(j)/8)}}{\mathrm{\Gamma }(_{jA^c}(1/2h(j)/8))}}\right)^{1/2}.`$
We have $`_{A\{1,2,3,\mathrm{},2m\}}1=2^{2m}`$, and the number of ways to choose the $`\{h(j)\}`$ is bounded by the number of ways of dividing $`2m`$ objects into $`3`$ groups, which is $`3^{2m}`$. Then noting that $`_{jA^c}(1/2h(j)/8)`$ is an integer multiple of $`1/8`$ which is always less than $`m`$, we can bound the last line by
(4.40) $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \underset{m=l}{\overset{\mathrm{}}{}}}C^m\lambda ^mϵ^{m/4}\right){\displaystyle \underset{j=0}{\overset{7}{}}}\left({\displaystyle \frac{b_n^{l+j/8}}{\mathrm{\Gamma }(l+j/8)}}\right)^{1/2}`$
$`Cb_n{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \underset{m=l}{\overset{\mathrm{}}{}}}C^m\lambda ^mϵ^{m/4}\right)\left({\displaystyle \frac{b_n^l}{\mathrm{\Gamma }(l)}}\right)^{1/2}`$
$`Cb_n(1C\lambda ϵ^{m/4})^1{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}C^l\lambda ^l|ϵ|^{l/4}b_n^{l/2}\left({\displaystyle \frac{1}{\mathrm{\Gamma }(l)}}\right)^{1/2}`$
for $`ϵ>0`$ sufficiently small.
(4.7) then follows from the fact that for any $`a>0`$
(4.41) $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}a^l\left({\displaystyle \frac{1}{\mathrm{\Gamma }(l)}}\right)^{1/2}`$
$`={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left(a^{2m}\left({\displaystyle \frac{1}{\mathrm{\Gamma }(2m)}}\right)^{1/2}+a^{2m+1}\left({\displaystyle \frac{1}{\mathrm{\Gamma }(2m+1)}}\right)^{1/2}\right)`$
$`C(1+a){\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}a^{2m}\left({\displaystyle \frac{1}{\mathrm{\Gamma }(2m)}}\right)^{1/2}`$
$`C(1+a)e^{Ca^2}.`$
###### Remark 4.2
It follows from the proof that in fact for $`\rho >0`$ sufficiently small, for any $`\lambda >0`$
(4.42) $`\underset{ϵ0}{lim\; sup}\underset{n\mathrm{}}{lim\; sup}\underset{y}{sup}`$
$`{\displaystyle \frac{1}{b_n}}\mathrm{log}E\left(\mathrm{exp}\left\{\lambda \left|{\displaystyle \frac{I_n(y)I_{n,ϵ}(y)}{ϵ^\rho b_n^1n}}\right|^{1/2}\right\}\right)=0.`$
###### Remark 4.3
Without the the restriction that $`b_n=o(n)`$, Theorem 1.1 is not true. To see this, let $`N`$ be an arbitrarily large integer, let $`\epsilon =2/N^2`$, and let $`X_i`$ be be an i.i.d. sequence of random vectors in $`^2`$ that take the values $`(N,0),(N,0),(0,N)`$, and $`(0,N)`$ with probability $`\epsilon /4`$ and $`(X_1=(0,0))=1\epsilon `$. The covariance matrix of the $`X_i`$ will be the identity. Let $`b_n=(1\epsilon )n`$. Then the event that $`S_i=S_0`$ for all $`in`$ will have probability at least $`(1\epsilon )^n`$, and on this event $`B_n=n(n1)/2`$. This shows that
$$\mathrm{log}(B_n𝔼B_n>nb_n/2)n\mathrm{log}(1\epsilon ),$$
which would contradict (1.4).
The same example shows that the critical constant in the polymer model is different than the one in . Then
$$𝔼\mathrm{exp}\left\{C\frac{B_n𝔼B_n}{n}\right\}\mathrm{exp}\left\{C\frac{𝔼B_n}{n}\right\}(1\epsilon )^n\mathrm{exp}\left\{C\frac{n1}{2}\right\}.$$
This shows that the critical constant is no more than $`2\mathrm{log}\frac{1}{1\epsilon }`$.
## 5 Theorem 1.2: Upper bound for $`𝔼B_nB_n`$
Proof of Theorem 1.2.
We first prove (1.5) for $`\theta =1`$:
(5.1) $`C_1`$ $``$ $`\underset{n\mathrm{}}{lim\; inf}b_n^1\mathrm{log}\left\{𝔼B_nB_n(2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$``$ $`\underset{n\mathrm{}}{lim\; sup}b_n^1\mathrm{log}\left\{𝔼B_nB_n(2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}C_2`$
for any $`\{b_n\}`$ satisfying (1.3). $`b_n^\theta `$ in (5.1).
In this section we prove the upper bound for (5.1). Let $`t>0`$ and write $`K=[t^1b_n]`$. Divide $`[1,n]`$ into $`K>1`$ disjoint subintervals $`(n_0,n_1],\mathrm{},(n_{K1},n_K]`$, each of length $`[n/K]`$ or $`[n/K]+1`$. Notice that
(5.2) $`𝔼B_nB_n{\displaystyle \underset{i=1}{\overset{K}{}}}\left[𝔼B\left((n_{i1},n_i]_<^2\right)B\left((n_{i1},n_i]_<^2\right)\right]`$
$`+𝔼B_n{\displaystyle \underset{i=1}{\overset{K}{}}}𝔼B\left((n_{i1},n_i]_<^2\right)`$
By (2.39),
(5.3) $`{\displaystyle \underset{i=1}{\overset{K}{}}}𝔼B\left((n_{i1},n_i]_<^2\right)={\displaystyle \underset{i=1}{\overset{K}{}}}𝔼B_{n_in_{i1}}`$
$`={\displaystyle \underset{i=1}{\overset{K}{}}}\left[{\displaystyle \frac{1}{(2\pi )\sqrt{det\mathrm{\Gamma }}}}(n/K)\mathrm{log}(n/K)+O(n/K)\right]`$
$`={\displaystyle \frac{1}{(2\pi )\sqrt{det\mathrm{\Gamma }}}}n\mathrm{log}(n/K)+O(n)`$
With $`K>1`$, the error term can be taken to be independent of $`t`$ and $`\{b_n\}`$. Thus, by (2.39), there is constant $`\mathrm{log}a>0`$ independent of $`t`$ and $`\{b_n\}`$ such that
(5.4) $`𝔼B_n{\displaystyle \underset{j=1}{\overset{K}{}}}𝔼B\left((n_{i1},n_i]_<^2\right)`$
$`{\displaystyle \frac{1}{(2\pi )\sqrt{det\mathrm{\Gamma }}}}n\left(\mathrm{log}(t^1b_n)+\mathrm{log}a\right).`$
It is here that we use the condition that $`𝔼|S_1|^{2+\delta }<\mathrm{}`$ for some $`\delta >0`$, needed for (2.39).
By first using Chebyshev’s inequality, then using (5.2), (5.4) and the independence of the $`B\left((n_{i1},n_i]_<^2\right)`$, for any $`\varphi >0`$,
(5.5) $`\left\{𝔼B_nB_n(2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$`\mathrm{exp}\left\{\varphi b_n\mathrm{log}b_n\right\}𝔼\mathrm{exp}\left\{2\pi \varphi \sqrt{det\mathrm{\Gamma }}{\displaystyle \frac{b_n}{n}}(B_n𝔼B_n)\right\}`$
$`\mathrm{exp}\left\{\varphi b_n(\mathrm{log}a\mathrm{log}t)\right\}\left(𝔼\mathrm{exp}\left\{2\pi \varphi \sqrt{det\mathrm{\Gamma }}{\displaystyle \frac{b_n}{n}}(B_{[n/K]}𝔼B_{[n/K]})\right\}\right)^K`$
By \[16, Theorem 1.2\],
(5.6)
$$\sqrt{det\mathrm{\Gamma }}\frac{b_n}{n}(B_{[n/K]}𝔼B_{[n/K]})\stackrel{d}{}\gamma _t,(n\mathrm{})$$
where $`\gamma _t`$ is the renormalized self-intersection local time of planar Brownian motion $`\{W_s\}`$ up to time $`t`$. By Lemma 2.5 and the dominated convergence theorem,
(5.7)
$$𝔼\mathrm{exp}\left\{2\pi \varphi \sqrt{det\mathrm{\Gamma }}\frac{b_n}{n}(B_{[n/K]}𝔼B_{[n/K]})\right\}𝔼\mathrm{exp}\left\{2\pi \varphi t\gamma _1\right\},(n\mathrm{})$$
where we used the scaling $`\gamma _t\stackrel{d}{=}t\gamma _1`$.
Thus,
(5.8) $`\underset{n\mathrm{}}{lim\; sup}b_n^1\mathrm{log}\left\{𝔼B_nB_n(2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$`\theta (\mathrm{log}a\mathrm{log}t)+{\displaystyle \frac{1}{t}}\mathrm{log}𝔼\mathrm{exp}\left\{2\pi \varphi t\gamma _1\right\}`$
$`=\theta \mathrm{log}(a\varphi )+{\displaystyle \frac{1}{t}}\mathrm{log}𝔼\mathrm{exp}\left\{(\varphi t)\mathrm{log}(\theta t)2\pi (\varphi t)\gamma _1\right\}`$
By \[2, p. 3233\], the limit
(5.9)
$$C\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{log}𝔼\mathrm{exp}\left\{t\mathrm{log}t2\pi t\gamma _1\right\}$$
exists. Hence
(5.10) $`\underset{n\mathrm{}}{lim\; sup}b_n^1\mathrm{log}\left\{𝔼B_nB_n(2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$`\varphi \mathrm{log}(a\varphi )+C\varphi .`$
Taking the minimizer $`\varphi =a^1e^{(1+C)}`$ we have
(5.11) $`\underset{n\mathrm{}}{lim\; sup}b_n^1\mathrm{log}\left\{𝔼B_nB_n(2\pi )^1det(\mathrm{\Gamma })^{1/2}n\mathrm{log}b_n\right\}`$
$`a^1e^{(1+C)}.`$
This proves the upper bound for (5.1).
## 6 Theorem 1.2: Lower bound for $`𝔼B_nB_n`$
In this section we complete the proof of Theorem 1.2 by proving the lower bound for (5.1).
Let $`B(x,r)`$ be the ball of radius $`r`$ centered at $`x`$. Let $`_k=\sigma \{X_i:ik\}`$. Let us assume for simplicity that the covariance matrix for the random walk is the identity; routine modifications are all that are needed for the general case. We write $`\mathrm{\Theta }`$ for $`(2\pi )^1\text{ det }(\mathrm{\Gamma })^{1/2}=(2\pi )^1`$. We write $`D(x,r)`$ for the disc of radius $`r`$ in $`^2`$ centered at $`x`$.
Let $`K=[b_n]`$ and $`L=n/K`$. Let us divide $`\{1,2,\mathrm{},n\}`$ into $`K`$ disjoint contiguous blocks, each of length strictly between $`L/2`$ and $`3L/2`$. Denote the blocks $`J_1,\mathrm{},J_K`$. Let $`v_i=\mathrm{\#}(J_i)`$, $`w_i=_{j=1}^iv_j`$. Let
(6.1)
$$B_{v_i}^{(i)}=\underset{j,kJ_i,j<k}{}\delta (S_j,S_k),A_i=\underset{jJ_{i1},kJ_i}{}\delta (S_j,S_k).$$
Define the following sets:
$`F_{i,1}`$ $`=`$ $`\{S_{w_i}D(i\sqrt{L},\sqrt{L}/16)\},`$
$`F_{i,2}`$ $`=`$ $`\{S(J_i)[(i1)\sqrt{L}\sqrt{L}/8,i\sqrt{L}+\sqrt{L}/8]\times [\sqrt{L}/8,\sqrt{L}/8]\},`$
$`F_{i,3}`$ $`=`$ $`\{B_{v_i}^{(i)}𝔼B_{v_i}^{(i)}\kappa _1L\},`$
$`F_{i,4}`$ $`=`$ $`\{{\displaystyle \underset{jJ_i}{}}1_{D(x,r\sqrt{L})}(S_j)\kappa _2rL\text{ for all }xD(i\sqrt{L},3\sqrt{L}),1/\sqrt{L}<r<2\},`$
$`F_{i,5}`$ $`=`$ $`\{A_i<\kappa _3L\},`$
where $`\kappa _1,\kappa _2,\kappa _3`$ are constants that will be chosen later and do not depend on $`K`$ or $`L`$. Let
(6.2)
$$C_i=F_{i,1}F_{i,2}F_{i,3}F_{i,4}F_{i,5}$$
and
(6.3)
$$E=_{i=1}^KC_i.$$
We want to show
(6.4)
$$(C_i_{w_{i1}})c_1>0$$
on the event $`C_1\mathrm{}C_{i1}`$. Once we have (6.4), then
(6.5)
$$(_{i=1}^mC_i)=𝔼((C_m_{w_{m1}});_{i=1}^{m1}C_i)c_1(_{i=1}^{m1}C_i),$$
and by induction
(6.6)
$$(E)=(_{i=1}^KC_i)c_1^K=e^{K\mathrm{log}c_1}=e^{c_2K}.$$
On the set $`E`$, we see that $`S(B_i)S(B_j)=\mathrm{}`$ if $`|ij|>1`$. So we can write
(6.7)
$$B_n=\underset{k=1}{\overset{K}{}}(B_{v_k}^{(k)}𝔼B_{v_k}^{(k)})+\underset{k=1}{\overset{K}{}}𝔼B_{v_k}^{(k)}+\underset{k=1}{\overset{K}{}}A_k.$$
On the event $`E`$, each $`B_{v_k}^{(k)}𝔼B_{v_k}^{(k)}`$ is bounded by $`\kappa _1L`$ and each $`A_k`$ is bounded by $`\kappa _3L`$. By (2.38), each $`𝔼B_{v_k}^{(k)}=\mathrm{\Theta }v_k\mathrm{log}v_k+O(L)=\mathrm{\Theta }v_k\mathrm{log}L+O(v_k)`$. Therefore
(6.8)
$$B_n\kappa _1KL+\mathrm{\Theta }KL\mathrm{log}L+O(n)+\kappa _3KL,$$
and using (2.38) again,
$`𝔼B_nB_n`$ $``$ $`\mathrm{\Theta }n\mathrm{log}nc_3n\mathrm{\Theta }n\mathrm{log}(n/b_n)`$
$`=`$ $`\mathrm{\Theta }n\mathrm{log}b_nc_3n`$
on the event $`E`$. We conclude that
(6.10)
$$(𝔼B_nB_n\mathrm{\Theta }n\mathrm{log}b_nc_3n)e^{c_2b_n}.$$
We apply (6.10) with $`b_n`$ replaced by $`b_n^{}=c_4b_n`$, where $`\mathrm{\Theta }\mathrm{log}c_4=c_3`$. Then
(6.11)
$$\mathrm{\Theta }n\mathrm{log}b_n^{}c_3n=\mathrm{\Theta }n\mathrm{log}b_n+\mathrm{\Theta }n\mathrm{log}c_4c_3n=\mathrm{\Theta }n\mathrm{log}b_n.$$
We then obtain
(6.12)
$$(𝔼B_nB_n\mathrm{\Theta }n\mathrm{log}b_n)=(𝔼B_nB_n\mathrm{\Theta }n\mathrm{log}b_n^{}c_3n)e^{c_2b_n^{}},$$
which would complete the proof of the lower bound for (5.1), hence of Theorem 1.2.
So we need to prove (6.4). By scaling and the support theorem for Brownian motion (see \[1, Theorem I.6.6\]), if $`W_t`$ is a planar Brownian motion and $`|x|\sqrt{L}/16`$, then
(6.13) $`^x(W_{v_i}D(\sqrt{L},\sqrt{L}/16)\text{ and}`$
$`\{W_s;0sv_i\}[\sqrt{L}/8,9\sqrt{L}/8]\times [\sqrt{L}/8,\sqrt{L}/8])>c_5,`$
where $`c_5`$ does not depend on $`L`$. Using Donsker’s invariance principle for random walks with finite second moments together with the Markov property,
(6.14)
$$(F_{i,1}F_{i,2}F_{w_{i1}})>c_6.$$
By Lemma 2.5, for $`L/2\mathrm{}3L/2`$
(6.15)
$$(B_{\mathrm{}}𝔼B_{\mathrm{}}>\kappa _1L)c_6/2$$
if we choose $`\kappa _1`$ large enough. Again using the Markov property,
(6.16)
$$(F_{i,1}F_{i,2}F_{i,3}F_{w_{i1}})>c_6/2.$$
Now let us look at $`F_{i,4}`$. By \[17, p. 75\], $`(S_j=y)c_7/j`$ with $`c_7`$ independent of $`y^2`$ so that
(6.17)
$$(S_jD(x,r\sqrt{L}))=\underset{yD(x,r\sqrt{L})}{}(S_j=y)\frac{c_8r^2L}{j}.$$
Therefore
$`𝔼{\displaystyle \underset{jJ_1}{}}1_{D(x,r\sqrt{L})}(S_j)`$ $``$ $`{\displaystyle \underset{j=1}{\overset{[2L]}{}}}(S_jD(x,r\sqrt{L}))`$
$``$ $`r^2L+{\displaystyle \underset{j=r^2L}{\overset{[2L]}{}}}{\displaystyle \frac{c_9r^2L}{j}}`$
$``$ $`r^2L+c_{10}Lr^2\mathrm{log}(1/r)c_{11}Lr^2\mathrm{log}(1/r)`$
if $`1/\sqrt{L}r2`$. Let $`C_m=_{j<m}1_{D(x,r\sqrt{L})}(S_j)`$ for $`m[2L]+1`$ and let $`C_m=C_{[2L]+1}`$ for $`m>L`$. By the Markov property and independence,
(6.19) $`𝔼[C_{\mathrm{}}C_m_m]1+𝔼[C_{\mathrm{}}C_{m+1}_m]`$
$`1+𝔼^{S_m}C_{\mathrm{}}c_{12}Lr^2\mathrm{log}(1/r).`$
By \[1, Theorem I.6.11\], we have
(6.20)
$$𝔼\mathrm{exp}\left(c_{13}\frac{C_{[2L]+1}}{c_{12}Lr^2\mathrm{log}(1/r)}\right)c_{14}$$
with $`c_{13},c_{14}`$ independent of $`L`$ or $`r`$. We conclude
(6.21)
$$\left(\underset{jJ_1}{}1_{D(x,r\sqrt{L})}(S_j)>c_{15}Lr^2\mathrm{log}(1/r)\right)c_{16}e^{c_{17}c_{15}}.$$
Suppose $`2^sr<2^{s+1}`$ for some $`s0`$. If $`xD(0,3\sqrt{L})`$, then each point in the disc $`D(x,r\sqrt{L})`$ will be contained in $`D(x_i,2^{s+3}\sqrt{L})`$ for some $`x_i`$, where each coordinate of $`x_i`$ is an integer multiple of $`2^{s2}\sqrt{L}`$. There are at most $`c_{18}2^{2s}`$ such balls, and $`Lr^2\mathrm{log}(1/r)c_{19}2^{s/2}Lr`$, so
(6.22)
$$\left(\underset{xD(0,3\sqrt{L}),2^sr<2^{s+1}}{sup}\underset{jJ_1}{}1_{D(x,r\sqrt{L})}(S_j)>c_{20}rL\right)c_{21}2^{2s}e^{c_{22}c_{20}2^{s/2}}.$$
If we now sum over positive integers $`s`$ and take $`\kappa _2=c_{20}`$ large enough, we see that
(6.23)
$$(F_{1,4})c_6/4.$$
By the Markov property, we then obtain
(6.24)
$$(F_{i,1}F_{i,2}F_{i,3}F_{i,4}F_{w_{i1}})>c_6/4.$$
Finally, we examine $`F_{i,5}`$. We will show
(6.25)
$$(F_{i,5}_{w_{i1}})c_6/8$$
on the set $`_{j=1}^{i1}C_j`$ if we take $`\kappa _3`$ large enough. By the Markov property, it suffices to show
(6.26)
$$\left(\underset{j=1}{\overset{[2L]}{}}1_{(S_jG)}\kappa _3L\right)c_6/8$$
whenever $`G^2`$ is a fixed nonrandom set consisting of $`[2L]`$ points satisfying the property that
(6.27)
$$\mathrm{\#}(GD(x,r\sqrt{L}))\kappa _2rL,xD(0,3\sqrt{L}),1/\sqrt{L}r2.$$
We compute the expectation of
(6.28)
$$\underset{j=1}{\overset{[2L]}{}}1_{(S_jG(D(0,2^k\sqrt{L})D(0,2^{k+1}\sqrt{L})))}.$$
When $`j2^{2k}L`$, then the fact that the random walk has finite second moments implies that the probability that $`|S_j|`$ exceeds $`2^{k+1}\sqrt{L}`$ is bounded by $`c_{23}j/(2^{2k+2}L)`$. When $`j>2^{2k}L`$, we use \[17, p. 75\], and obtain
(6.29)
$$(S_jG(D(0,2^k\sqrt{L})c_{24}\frac{\kappa _22^kL}{j}.$$
So
(6.30) $`𝔼{\displaystyle \underset{j=1}{\overset{[2L]}{}}}1_G(S_j)`$
$`{\displaystyle \underset{k}{}}{\displaystyle \underset{[2L]j>2^{2k}L}{}}c_{24}{\displaystyle \frac{\kappa _22^kL}{j}}+{\displaystyle \underset{k}{}}{\displaystyle \underset{j2^{2k}L}{}}c_{23}{\displaystyle \frac{j}{2^{2k+2}L}}`$
$`{\displaystyle \underset{k}{}}(c_{25}\kappa _2k2^kL+c_{26}2^{2k}L)c_{27}L.`$
So if take $`\kappa _3`$ large enough, we obtain (6.26).
This completes the proof of (6.4), hence of Theorem 1.2.
## 7 Laws of the iterated logarithm
### 7.1 Proof of the LIL for $`B_n𝔼B_n`$
First, let $`S_j,S_j^{}`$ be two independent copies of our random walk. Let
(7.1)
$$\mathrm{}(n,x)=\underset{i=1}{\overset{n}{}}\delta (S_i,x),\mathrm{}^{}(n,x)=\underset{i=1}{\overset{n}{}}\delta (S_i^{},x)$$
and note that
(7.2)
$$I_{k,n}=\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{n}{}}\delta (S_i,S_j^{})=\underset{x^2}{}\mathrm{}(k,x)\mathrm{}^{}(n,x).$$
###### Lemma 7.1
There exist constants $`c_1,c_2`$ such that
(7.3)
$$(I_{k,n}>\lambda \sqrt{kn})c_1e^{c_2\lambda }.$$
Proof. Clearly
(7.4)
$$(I_{k,n})^m=\underset{x_1^2}{}\mathrm{}\underset{x_m^2}{}\left(\underset{i=1}{\overset{m}{}}\mathrm{}(k,x_i)\right)\left(\underset{i=1}{\overset{m}{}}\mathrm{}^{}(n,x_i)\right)$$
Using the independence of $`S`$ and $`S^{}`$,
(7.5)
$$𝔼\left((I_{k,n})^m\right)=\underset{x_1^2}{}\mathrm{}\underset{x_m^2}{}𝔼\left(\underset{i=1}{\overset{m}{}}\mathrm{}(k,x_i)\right)𝔼\left(\underset{i=1}{\overset{m}{}}\mathrm{}^{}(n,x_i)\right).$$
By Cauchy-Schwarz, this is less than
(7.6) $`\left[{\displaystyle \underset{x_1^2}{}}\mathrm{}{\displaystyle \underset{x_m^2}{}}\left(𝔼\left({\displaystyle \underset{i=1}{\overset{m}{}}}\mathrm{}(k,x_i)\right)\right)^2\right]^{1/2}`$
$`\left[{\displaystyle \underset{x_1^2}{}}\mathrm{}{\displaystyle \underset{x_m^2}{}}\left(𝔼\left({\displaystyle \underset{i=1}{\overset{m}{}}}\mathrm{}^{}(n,x_i)\right)\right)^2\right]^{1/2}`$
$`=:J_1^{1/2}J_2^{1/2}.`$
We can rewrite
(7.7)
$$J_1=\underset{x_1^2}{}\mathrm{}\underset{x_m^2}{}𝔼\left(\underset{i=1}{\overset{m}{}}\mathrm{}(k,x_i)\right)𝔼\left(\underset{i=1}{\overset{m}{}}\mathrm{}^{}(k,x_i)\right)=𝔼\left((I_k)^m\right),$$
and similarly $`J_2=𝔼\left((I_n)^m\right)`$.
Therefore,
$`𝔼\mathrm{exp}(aI_{k,n}/\sqrt{kn})`$
$`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{a^m}{k^{m/2}n^{m/2}m!}}𝔼\left((I_{k,n})^m\right)`$
$``$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{a^m}{k^{m/2}n^{m/2}m!}}(𝔼\left((I_k)^m\right))^{1/2}(𝔼\left((I_n)^m\right))^{1/2}`$
$``$ $`\left({\displaystyle \frac{a^m}{m!}𝔼\left(\frac{I_k}{k}\right)^m}\right)^{1/2}\left({\displaystyle \frac{a^m}{m!}𝔼\left(\frac{I_n}{n}\right)^m}\right)^{1/2}`$
$``$ $`\left(𝔼e^{aI_k/k}\right)^{1/2}\left(𝔼e^{aI_n/n}\right)^{1/2}.`$
By Lemma 2.2 this can be bounded independently of $`k`$ and $`n`$ if $`a`$ is taken small, and our result follows.
We are now ready to prove the upper bound for the LIL for $`B_n𝔼B_n`$. Write $`\mathrm{\Xi }`$ for $`\sqrt{det\mathrm{\Gamma }}\kappa (2,2)^4`$. Recall that for any integrable random variable $`Z`$ we let $`\overline{Z}`$ denote $`Z𝔼Z`$. Let $`\epsilon >0`$ and let $`q>1`$ be chosen later. Our first goal is to get an upper bound on
$$(\underset{n/2kn}{\mathrm{max}}\overline{B}_k>(1+\epsilon )\mathrm{\Xi }^1n\mathrm{log}\mathrm{log}n).$$
Let $`m_0=2^N`$, where $`N`$ will be chosen later to depend only on $`\epsilon `$. Let $`𝒜_0`$ be the integers of the form $`nkm_0`$ that are contained in $`\{n/4,\mathrm{},n\}`$. For each $`i`$ let $`𝒜_i`$ be the set of integers of the form $`nkm_02^i`$ that are contained in $`\{n/4,\mathrm{},n\}`$. Given an integer $`k`$, let $`k_j`$ be the largest element of $`𝒜_j`$ that is less than or equal to $`k`$. For any $`k\{n/2,\mathrm{},n\}`$, we can write
(7.9)
$$\overline{B}_k=\overline{B}_{k_0}+(\overline{B}_{k_1}\overline{B}_{k_0})+\mathrm{}+(\overline{B}_{k_N}\overline{B}_{k_{N1}}).$$
If $`\overline{B}_k(1+\epsilon )\mathrm{\Xi }^1n\mathrm{log}\mathrm{log}n`$ for some $`n/2kn`$, then either (a) $`\overline{B}_{k_0}(1+\frac{\epsilon }{2})\mathrm{\Xi }^1n\mathrm{log}\mathrm{log}n`$ for some $`k_0𝒜_0`$; or else (b) for some $`i1`$ and some pair of consecutive elements $`k_i,k_i^{}𝒜_i`$, we have
(7.10)
$$\overline{B}_{k_i^{}}\overline{B}_{k_i}\frac{\epsilon }{40i^2}n\mathrm{log}\mathrm{log}n.$$
For each $`k_0`$, using Theorem 1.1 and the fact that $`k_0n/4`$, the probability in (a) is bounded by
(7.11)
$$\mathrm{exp}((1+\frac{\epsilon }{4})\mathrm{log}\mathrm{log}k_0)c_1(\mathrm{log}n)^{(1+\frac{\epsilon }{4})}.$$
There are at most $`n/m_0`$ elements of $`𝒜_0`$, so the probability in (a) is bounded by
(7.12)
$$\frac{n}{m_0}\frac{c_1}{(\mathrm{log}n)^{1+\frac{\epsilon }{4}}}.$$
Now let us examine the probability in (b). Fix $`i`$ for the moment. Any two consecutive elements of $`𝒜_i`$ are $`2^im_0`$ apart. Recalling the notation (2.16) we can write
(7.13)
$$\overline{B}_k\overline{B}_j=\overline{B}([j+1,k]_<^2)+\overline{B}([1,j]\times [j+1,k]),$$
So
(7.14) $`(\overline{B}_k\overline{B}_j\frac{\epsilon }{40i^2}n\mathrm{log}\mathrm{log}n)`$ $``$ $`(\overline{B}([j+1,k]_<^2)\frac{\epsilon }{80i^2}n\mathrm{log}\mathrm{log}n)`$
$`+\left(B([1,j]\times [j+1,k])\frac{\epsilon }{80i^2}n\mathrm{log}\mathrm{log}n\right).`$
We bound the first term on the right by Theorem 1.1, and get the bound
(7.15)
$$\mathrm{exp}\left(\frac{\epsilon \mathrm{\Xi }}{80i^2}\frac{n\mathrm{log}\mathrm{log}n}{2^im_0}\right)\mathrm{exp}\left(\frac{\epsilon \mathrm{\Xi }}{80i^2}2^i(n/m_0)\mathrm{log}\mathrm{log}n\right)$$
if $`j`$ and $`k`$ are consecutive elements of $`𝒜_i`$. Note that $`B([1,j]\times [j+1,k])`$ is equal in law to $`I_{j1,kj}`$. Using Lemma 7.1, we bound the second term on the right hand side of (7.14) by
$`c_1\mathrm{exp}\left(c_2{\displaystyle \frac{\epsilon }{80i^2}}{\displaystyle \frac{n\mathrm{log}\mathrm{log}n}{\sqrt{2^im_0}\sqrt{j}}}\right)`$
(7.16) $`c_1\mathrm{exp}\left(c_2{\displaystyle \frac{\epsilon }{80i^2}}2^{i/2}(n/m_0)^{1/2}\mathrm{log}\mathrm{log}n\right).`$
The number of pairs of consecutive elements of $`𝒜_i`$ is less than $`2^{i+1}(n/m_0)`$. So if we add (7.15) and (7.16) and multiply by the number of pairs, the probability of (b) occurring for a fixed $`i`$ is bounded by
(7.17)
$$c_3\frac{n}{m_0}2^i\mathrm{exp}\left(c_42^{i/2}(n/m_0)^{1/2}\mathrm{log}\mathrm{log}n/(80i^2)\right).$$
If we now sum over $`i1`$, we bound the probability in (b) by
(7.18)
$$c_5\frac{n}{m_0}\mathrm{exp}\left(c_6(n/m_0)^{1/2}\mathrm{log}\mathrm{log}n\right).$$
We now choose $`m_0`$ to be the largest power of $`2`$ so that $`c_6(n/m_0)^{1/2}>2`$; recall $`n`$ is big.
Let us use this value of $`m_0`$ and combine (7.12) and (7.18). Let $`n_{\mathrm{}}=q^{\mathrm{}}`$ and
(7.19)
$$C_{\mathrm{}}=\{\underset{n_\mathrm{}1kn_{\mathrm{}}}{\mathrm{max}}\overline{B}_k(1+\epsilon )\mathrm{\Xi }^1n_{\mathrm{}}\mathrm{log}\mathrm{log}n_{\mathrm{}}\}.$$
By our estimates, $`(C_{\mathrm{}})`$ is summable, so for $`\mathrm{}`$ large, by Borel-Cantelli we have
(7.20)
$$\underset{n_\mathrm{}1kn_{\mathrm{}}}{\mathrm{max}}\overline{B}_k(1+\epsilon )\mathrm{\Xi }^1n_{\mathrm{}}\mathrm{log}\mathrm{log}n_{\mathrm{}}.$$
By taking $`q`$ sufficiently close to 1, this implies that for $`k`$ large we have $`\overline{B}_k(1+2\epsilon )\mathrm{\Xi }^1k\mathrm{log}\mathrm{log}k`$. Since $`\epsilon `$ is arbitrary, we have our upper bound.
The lower bound for the first LIL is easier. Let $`\delta >0`$ be small and let $`n_{\mathrm{}}=[e^{\mathrm{}^{1+\delta }}]`$. Let
(7.21)
$$D_{\mathrm{}}=\{\overline{B}([n_\mathrm{}1+1,n_{\mathrm{}}]_<^2)(1\epsilon )\mathrm{\Xi }^1n_{\mathrm{}}\mathrm{log}\mathrm{log}n_{\mathrm{}}\}.$$
Using Theorem 1.1, and the fact that $`n_{\mathrm{}}/(n_{\mathrm{}}n_\mathrm{}1)ce^{c^{}l^\delta }`$ we see that $`_{\mathrm{}}(D_{\mathrm{}})=\mathrm{}`$. The $`D_{\mathrm{}}`$ are independent, so by Borel-Cantelli
(7.22)
$$\overline{B}([n_\mathrm{}1+1,n_{\mathrm{}}]_<^2)(1\epsilon )\mathrm{\Xi }^1n_{\mathrm{}}\mathrm{log}\mathrm{log}n_{\mathrm{}}$$
infinitely often with probability one. Note that as in (7.13) we can write
(7.23)
$$\overline{B}_n_{\mathrm{}}=\overline{B}([n_\mathrm{}1+1,n_{\mathrm{}}]_<^2)+\overline{B}_{n_\mathrm{}1}+\overline{B}([1,n_\mathrm{}1]\times [n_\mathrm{}1+1,n_{\mathrm{}}]).$$
By the upper bound,
$$\underset{\mathrm{}\mathrm{}}{lim\; sup}\frac{\overline{B}_{n_\mathrm{}1}}{n_\mathrm{}1\mathrm{log}\mathrm{log}n_\mathrm{}1}\mathrm{\Xi }^1$$
almost surely, which implies
(7.24)
$$\underset{\mathrm{}\mathrm{}}{lim\; sup}\frac{\overline{B}_{n_\mathrm{}1}}{n_{\mathrm{}}\mathrm{log}\mathrm{log}n_{\mathrm{}}}=0.$$
Since $`B([1,n_\mathrm{}1]\times [n_\mathrm{}1+1,n_{\mathrm{}}])0`$ and by (2.5)
(7.25)
$$𝔼B([1,n_\mathrm{}1]\times [n_\mathrm{}1+1,n_{\mathrm{}}])c_1\sqrt{n_\mathrm{}1}\sqrt{n_{\mathrm{}}n_\mathrm{}1}=o(n_{\mathrm{}}\mathrm{log}\mathrm{log}n_{\mathrm{}}),$$
using (7.22)-(7.25) yields the lower bound.
### 7.2 LIL for $`𝔼B_nB_n`$
Let $`\mathrm{\Delta }=2\pi \sqrt{det\mathrm{\Gamma }}`$. Let us write $`J_n=𝔼B_nB_n`$.
First we do the upper bound. Let $`m_0`$, $`𝒜_i`$, and $`k_j`$ be as in the previous subsection. We write, for $`n/2kn`$,
(7.26)
$$J_k=J_{k_0}+(J_{k_1}J_{k_0})+\mathrm{}+(J_{k_N}J_{k_{N1}}).$$
If $`\mathrm{max}_{n/2kn}J_k(1+\epsilon )\mathrm{\Delta }^1n\mathrm{log}\mathrm{log}\mathrm{log}n`$, then either (a) $`J_{k_0}(1+\frac{\epsilon }{2})\mathrm{\Delta }^1n\mathrm{log}\mathrm{log}\mathrm{log}n`$ for some $`k_0𝒜_0`$, or else (b) for some $`i1`$ and $`k_i,k_i^{}`$ consecutive elements of $`𝒮_i`$ we have
(7.27)
$$J_{k_i^{}}J_{k_i}\frac{\epsilon }{40i^2}n\mathrm{log}\mathrm{log}\mathrm{log}n.$$
There are at most $`n/m_0`$ elements of $`𝒜_0`$. Using Theorem 1.2, the probability of (a) is bounded by
(7.28)
$$c_1\frac{n}{m_0}e^{(1+\frac{\epsilon }{4})\mathrm{log}\mathrm{log}n}.$$
To estimate the probability in (b), suppose $`j`$ and $`k`$ are consecutive elements of $`𝒜_i`$. There are at most $`2^{i+1}(n/m_0)`$ such pairs. We have
$`J_kJ_j`$ $`=`$ $`\overline{B}([j+1,k]_<^2)\overline{B}([1,j]\times [j+1,k])`$
$`\overline{B}([j+1,k]_<^2)+𝔼B([1,j]\times [j+1,k])`$
$`\overline{B}([j+1,k]_<^2)+c_2\sqrt{j}\sqrt{kj},`$
as in the previous subsection. Provided $`n`$ is large enough, $`c_2\sqrt{j}\sqrt{kj}=c_2\sqrt{j}\sqrt{2^im_0}`$ will be less than $`\frac{\epsilon }{80i^2}n\mathrm{log}\mathrm{log}\mathrm{log}n`$ for all $`i`$. So in order for $`J_kJ_j`$ to be larger than $`\frac{\epsilon }{40i^2}n\mathrm{log}\mathrm{log}\mathrm{log}n`$, we must have $`\overline{B}([j+1,k]_<^2)`$ larger than $`\frac{\epsilon }{80i^2}n\mathrm{log}\mathrm{log}\mathrm{log}n`$. We use Theorem 1.2 to bound this. Then multiplying by the number of pairs and summing over $`i`$, the probability is (b) is bounded by
(7.30) $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}2^{i+1}{\displaystyle \frac{n}{m_0}}e^{\frac{\epsilon }{80i^2}\frac{n}{2^im_0}\mathrm{log}\mathrm{log}n}c_3{\displaystyle \frac{n}{m_0}}e^{c_4(n/m_0)\mathrm{log}\mathrm{log}n}.`$
We choose $`m_0`$ to be the largest possible power of 2 such that $`c_4(n/m_0)>2`$.
Combining (7.28) and (7.30), we see that if we set $`q>1`$ close to 1, $`n_{\mathrm{}}=[q^{\mathrm{}}]`$, and
(7.31)
$$E_{\mathrm{}}=\{\underset{n_{\mathrm{}}/2kn_{\mathrm{}}}{\mathrm{max}}J_k(1+\epsilon )\mathrm{\Delta }^1n_{\mathrm{}}\mathrm{log}\mathrm{log}\mathrm{log}n_{\mathrm{}}\},$$
then $`_{\mathrm{}}(E_{\mathrm{}})`$ is finite. So by Borel-Cantelli, the event $`E_{\mathrm{}}`$ happens for a last time, almost surely. Exactly as in the previous subsection, taking $`q`$ close enough to 1 and using the fact that $`\epsilon `$ is arbitrary leads to the upper bound.
The proof of the lower bound is fairly similar to the previous subsection. Let $`n_{\mathrm{}}=[e^{\mathrm{}^{1+\delta }}]`$. Theorem 1.2 and Borel-Cantelli tell us that $`F_{\mathrm{}}`$ will happen infinitely often, where
(7.32)
$$F_{\mathrm{}}=\{\overline{B}([n_\mathrm{}1+1,n_{\mathrm{}}]_<^2)(1\epsilon )\mathrm{\Delta }^1n_{\mathrm{}}\mathrm{log}\mathrm{log}\mathrm{log}n_{\mathrm{}}\}.$$
We have
(7.33)
$$J_n_{\mathrm{}}\overline{B}([n_\mathrm{}1+1,n_{\mathrm{}}]_<^2)+J_{n_\mathrm{}1}A(1,n_\mathrm{}1;n_\mathrm{}1,n_{\mathrm{}}).$$
By the upper bound,
(7.34)
$$J_{n_\mathrm{}1}=O(n_\mathrm{}1\mathrm{log}\mathrm{log}\mathrm{log}n_\mathrm{}1)=o(n_{\mathrm{}}\mathrm{log}\mathrm{log}\mathrm{log}n_{\mathrm{}}).$$
By Lemma 7.1,
(7.35)
$$(B([1,n_\mathrm{}1]\times [n_\mathrm{}1+1,n_{\mathrm{}}])\epsilon n_{\mathrm{}}\mathrm{log}\mathrm{log}\mathrm{log}n_{\mathrm{}})c_1\mathrm{exp}\left(c_2\frac{\epsilon n_{\mathrm{}}\mathrm{log}\mathrm{log}\mathrm{log}n_{\mathrm{}}}{\sqrt{n_\mathrm{}1}\sqrt{n_{\mathrm{}}n_\mathrm{}1}}\right).$$
This is summable in $`\mathrm{}`$, so
(7.36)
$$\underset{\mathrm{}\mathrm{}}{lim\; sup}\frac{B([1,n_\mathrm{}1]\times [n_\mathrm{}1+1,n_{\mathrm{}}])}{n_{\mathrm{}}\mathrm{log}\mathrm{log}\mathrm{log}n_{\mathrm{}}}\epsilon $$
almost surely. This is true for every $`\epsilon `$, so the limsup is 0. Combining this with (7.34) and substituting in (7.33) completes the proof.
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# Discovery of an M4 Spectroscopic Binary in Upper Scorpius: A Calibration Point for Young Low-Mass Evolutionary Models
## 1 Introduction
The study of very low mass stars and brown dwarfs has experienced rapid growth during the past 15 years, primarily due to the advent of the next generation of large telescopes. Many questions are beginning to be answered about the atmospheric and evolutionary state of these objects, yet many fundamental questions remain. Thus far we have not gathered much empirical data on mass, one of the most fundamental quantities of interest. The evolutionary models that are currently used to assign almost all masses, based on age and luminosity, are largely untested. Modelers warn that such models cannot be trusted too far (Baraffe et al.,, 2003) for very low masses and very young ages. The search for binary systems that would provide fundamental stellar parameters is only now beginning to bear fruit. Indeed, the frequency of occurrence and physical configurations of very low mass binaries are themselves not well understood. These systems appear to have a somewhat, but not drastically, lower binary frequency than solar-type stars, and the scale of the systems seems to be smaller, along with more equal mass ratios (Bouy et al.,, 2003).
Recently there have been indications that the evolutionary models substantially under-predict masses for young objects in the 0.03-0.5 M range. Hillenbrand & White, (2004) compiled a collection of 115 low-mass stars with masses determined from orbital dynamics. Comparing them to masses derived from different evolutionary models, they found the masses of low-mass pre-main sequence stars to be generally under-predicted. The mass at which model predictions are close to dynamically derived masses depends on the model, but agreement tends to be better above 0.5 M. The lowest tested pre-main sequence mass, however, was 0.3 M. At main-sequence ages, Hillenbrand & White, (2004) found masses to be under-predicted in the 0.1-0.6 M regime in most of the models they used.
The first study focusing on less massive objects at very young ages was done by Mohanty et al., 2004a ; Mohanty et al., 2004b , who derived masses for 11 low mass objects in the young Upper Scorpius association. These mass estimates are independent of evolutionary models, relying instead on comparison of high resolution spectra with synthetic model spectra to obtain effective temperatures and surface gravities, then combining these with measured luminosities and the known distance and age of the star-forming region. They conclude that models under-predict the mass of young low-mass objects with a given temperature and luminosity, similar to the results of Hillenbrand & White, (2004) at higher masses and ages. More recently, a study using the more familiar orbital dynamical technique also suggested that the models under-predict mass for the object AB Dor C (Close et al.,, 2005). In that case, the determined mass of 0.09 M is almost twice that which models predict for an object of that luminosity, temperature, and age. Neither the age nor temperature are determined very precisely by this study, however (see also Luhman & Stauffer, (2005)). Close et al., (2005) overextend their result to claim that low mass brown dwarfs are being mistaken for “free-floating planets” due to the mass calibration problem. Actually, Mohanty et al., 2004b find that models *over-predict* masses below about 0.03 M (a regime not otherwise tested); the slope of their empirical mass-luminosity relation is shallower than predicted by models throughout the brown dwarf range.
Astronomers are much more used to calibrating masses by dynamical studies of binary systems than by surface gravity analysis. Indeed, the dynamical method looks much more straightforward at first glance, although when used to calibrate models it suffers from the same need to accurately measure a star’s luminosity and temperature as does the surface gravity method. In this paper, we report the discovery of a double-lined spectroscopic binary among the same population of low-mass members of Upper Sco studied by Mohanty et al., 2004a ; Mohanty et al., 2004b . This system is substantially younger than AB Dor C (studied by Close et al.,, 2005). We are able to determine an orbit for this system, UScoCTIO 5, which allows us to place a lower limit on its mass. We show that its mass appears to be underpredicted by the models, based on its measured luminosity and known age of the Upper Scorpius association. We find the same result using temperature, obtained in the usual way (via spectral type), which is not as precise as we would desire.
Our current result tends to confirm the results suggested by the surface gravity analysis of Mohanty et al., 2004b , which are in the same sense as found by Hillenbrand & White, (2004) and Close et al., (2005). This strengthens the case that evolutionary models for young very low-mass objects under-predict their masses at the bottom of the main sequence and in the heavier half of brown dwarfs. This has further implications for the initial mass functions derived in young clusters and star-forming regions, and provides a direction for modelers to improve their initial conditions. That they might need to do so comes as no surprise; they have been saying so themselves (e.g. Baraffe et al.,, 2003).
## 2 Observations
UScoCTIO 5 was the warmest object we observed in an extension of our original program to study very young brown dwarfs (Jayawardhana et al.,, 2002, 2003). We made the first observation of it on 11 June 2003 with HIRES at Keck I. It was immediately clear from the raw spectrum on that night that the lines were doubled. We then undertook a program to determine the orbital parameters. The region around the KI lines in the first spectrum is plotted in Fig. 1. Both binary components are clearly visible with comparable line intensities. After this exposure, twenty further spectra were taken in 2004, and one in 2005. We (and other helpful observers) carried out high-resolution spectroscopy using the HIRES spectrograph at Keck I, the echelle spectrograph at the CTIO 4m telescope, and the MIKE-Red echelle on Magellan II at Las Campanas Observatory. The instruments used, dates and times of observations, exposure times and names of observers are compiled in Table 1. For every instrumental setup, one or more spectra of the slowly rotating M6-star Gl 406 (CN Leo) were also taken to allow production of cross-correlation profiles. The resolution in all our spectra is at least $`R=\mathrm{25\hspace{0.17em}000}`$.
It is essential to the utility of our analysis that UScoCTIO 5 be a member of the Upper Scorpius Association. It is this membership which gives us the distance and age of the object. There are a number of observations which strongly point to its membership. It was discovered because it sits on the pre-main sequence for that association based on its photometry (and of course is spatially within the association). If it is that young, it should display strong lithium lines, and indeed both components are easily visible in lithium. At 0.3 M, the lifetime of lithium is only around 30 Myr. This rules out UScoCTIO 5 being a foreground main sequence star. There is further evidence against that in the narrow Na line profiles (main sequence stars have higher gravity and broader lines). The system velocity is consistent with our other targets in Upper Sco; the central velocity matches that of most other Upper Sco targets to within 3 km s<sup>-1</sup>. Finally, the H$`\alpha `$ lines in these objects are quite similar to those in several of the other non-accretors in Upper Scorpius shown in Mohanty et al., 2004b , though that by itself would not be definitive.
## 3 Differential Radial Velocities
We infer the orbital elements of UScoCTIO 5 from the radial velocities of the binary components relative to each other (differential radial velocities). Cross-correlation functions are computed using spectra of Gl 406 as a template. For each spectrum, we obtained a template spectrum using identical instrument setups for UScoCTIO 5 and Gl 406. In case of HIRES spectra, we used one order covering approximately 100 Å between 7850 Å and 7950 Å. From the CTIO data we merged four echelle orders and calculated the correlation functions in about 600 Å between 8350 Å and 8950 Å. The region between 8600 Å and 8800 Å was used in the Magellan data. None of these wavelength regions is contaminated by significant telluric lines.
Using wavelength regions as large as 600 Å, the wavelength dependency of Doppler shifts becomes important for the calculation of precise radial velocities. At differential radial velocities of 50 km s<sup>-1</sup>, for example, the velocity difference per pixel is of the order of 3 km s<sup>-1</sup> from one end of the linear wavelength scale used in the CTIO data to the other end. For highest accuracy, we therefore calculate the correlation function using logarithmic wavelengths. Differential radial velocities are then derived from the separation of the two peaks belonging to the different components. Both components were easily distinguishable in all correlation functions. Differential radial velocities $`\mathrm{\Delta }v_{\mathrm{rad}}`$ for all observations are given in column five of Table 1. An estimation of our measurement errors will be given in the next section.
It turns out that all observations taken after our first one in 2003 reveal significantly lower differential radial velocities, i.e., the components are more seriously blended than in the example given in Fig. 1. While both components are always distinguishable, an assignment of each of the two correlation peaks to individual stellar components is ambiguous, even in the cross-correlation functions derived from large spectral regions. We therefore did not individually identify the components in every spectrum, i.e., the sign of $`\mathrm{\Delta }v_{\mathrm{rad}}`$ is not derived from the data. However, as shown in the next section, we were able to eliminate the ambiguity in the sign of $`\mathrm{\Delta }v_{\mathrm{rad}}`$ when determining the orbit of UScoCTIO 5.
## 4 Determining the Orbit
As mentioned above, we do not have information about the sign of $`\mathrm{\Delta }v_{\mathrm{rad}}`$. Our data were taken during five observing campaigns, indicated by horizontal lines in Table 1. Radial velocities from spectra obtained during the same campaigns show clear correlations, and it is very unlikely that the sign of $`\mathrm{\Delta }v_{\mathrm{rad}}`$ changes within one campaign. However, we looked at the periodograms of all possible permutations of the signs of $`\mathrm{\Delta }v_{\mathrm{rad}}`$ and found no hints of a short orbit with the sign of $`\mathrm{\Delta }v_{\mathrm{rad}}`$ alternating within consecutive observations of any campaign. The radial velocity curve of a spectroscopic binary can be uniquely specified by the following five parameters: Period, $`P`$; total projected mass $`(M_1+M_2)\mathrm{sin}i`$; eccentricity of the orbit, $`e`$; longitude of periastron, $`\mathrm{\Omega }`$; and Epoch, $`T`$. From differential radial velocities it is not possible to obtain the inclination $`i`$. Thus we only derive a minimum mass sum from our data. The mass fraction $`M_1/M_2`$ could be determined in principle from either the individual component line ratios or the amplitude of their excursions about the system velocity. In our case, we could not confidently find any difference between the two components. This places a minimum mass limit on the primary (it cannot contain much more than half of the total mass).
We searched for the best orbital solution in the five free parameters by minimizing the rms scatter of $`\mathrm{\Delta }v_{\mathrm{rad}}`$, $`\sigma ^2=(\mathrm{\Delta }v_{\mathrm{rad}}\mathrm{\Delta }v_{\mathrm{rad},\mathrm{orbit}})^2/N`$, with $`N=22`$ the number of measurements. The rms scatter of $`\mathrm{\Delta }v_{\mathrm{rad}}`$ around the best fit is $`\sigma _{\mathrm{min}}=480`$ m s<sup>-1</sup>. The uncertainties in $`\mathrm{\Delta }v_{\mathrm{rad}}`$ are not well determined due to the variety of instruments used, and they are primarily systematic rather than of statistical origin. Thus we cannot give a strict statistical confidence limit on our result. We conservatively estimate the uncertainty in the orbital parameters by searching for the intervals for which $`\sigma <2\sigma _{\mathrm{min}}`$, when varying each parameter independently. The parameters of our best fit are given in Table 2 with the derived errors. Our best solution is plotted over the whole range of observations and over phase in the left and right panels of Fig. 2, respectively.
We find a minimum projected total mass of $`(M_1+M_2)\mathrm{sin}i=0.64`$ M. The solution is unique in the sense that for the second best local minimum in $`\sigma `$ we get a significantly higher value of $`\sigma _2=7.5\sigma _{\mathrm{min}}`$. This solution would provide $`(M_1+M_2)\mathrm{sin}i=0.70`$ M, and all other parameters are comparable to our best solution as well. We ran an extensive survey in parameters considering different signs in individual values and groups of $`\mathrm{\Delta }v_{\mathrm{rad}}`$. No other parameter set yielding comparable fit quality could be found.
## 5 Comparison to Evolutionary Tracks
In order to compare the derived minimum mass to theoretical evolutionary tracks, we calculate bolometric luminosity $`L_{\mathrm{bol}}/L_{}`$ and effective temperature $`T_{\mathrm{eff}}`$ from photometry and spectral type, respectively. A spectral type of M4 has been determined for UScoCTIO 5 by Ardila et al., (2000). Photometry, effective temperatures and bolometric luminosity are given in Table 3. We calculate bolometric luminosity from $`J`$ and $`K`$ colors taken from 2MASS (Cutri et al.,, 2003), and determine $`L_{\mathrm{bol}}/L_{}`$ independently from both colors. From the $`J`$-color we determine $`V`$ according to dwarf calibrations given in Kenyon & Hartmann, (1995), $`BC_V`$ is also taken from that paper. A calibration of $`BC_K`$ in the UKIRT photometric system can be found in Leggett et al., (2001); the transformation of 2MASS colors to the UKIRT system is adopted from Carpenter, (2001).
We calculate the extinction by comparing $`JH`$ and $`JK`$ to the young-disk calibration given in Legget, (1992) using the extinction law from Schlegel et al., (1998); colors are converted to the $`CRI`$-system for that purpose (Carpenter,, 2001). This yields $`A_V(JH)=1.0`$ and $`A_V(JK)=0.5`$, while a value of $`A_V=0.5\pm 0.5`$ is reported in Ardila et al., (2000). Including the uncertainty in spectral type, our estimates of $`A_V`$ have uncertainties of 0.3 mag, hence the values of $`A_V`$ are consistent within the uncertainties. We choose $`A_V=0.75\pm 0.4`$ for the calculation of bolometric luminosity, and we derive a bolometric luminosity of $`\mathrm{log}L_{\mathrm{bol}}/L_{}=1.17\pm 0.08`$ from the $`J`$band and $`\mathrm{log}L_{\mathrm{bol}}/L_{}=1.16\pm 0.06`$ from the $`K`$band. The two values derived from different colors agree very well, and in the following we will use the luminosity derived from the $`K`$-band since it has smaller uncertainties. For the uncertainties, we took into account a spectral class uncertainty of half a subclass ($`\pm 0.5`$), age ($`5\pm 2`$ Myr) and distance ($`145\pm 20`$ pc) for the cluster adopted from Mohanty et al., 2004b . These errors have been included in the luminosity error. For the bolometric correction $`BC_V`$ an uncertainty of 0.1 mag has been assumed, and 0.05 mag has been employed for $`BC_K`$ (Leggett et al.,, 2001).
In Fig. 3, UScoCTIO 5 is plotted in a luminosity-age diagram with evolutionary tracks from Baraffe et al., (1998) and Chabrier et al., (2000). Assuming that both components of UScoCTIO 5 contribute equal light, the most probable mass of one component according to its position in the diagram is $`M=0.23`$ M. Our dynamically determined total system projected minimum mass is $`(M_1+M_2)\mathrm{sin}i=0.64\pm 0.02`$ M. The color and spectral type of a spectroscopic binary is primarily determined by the hotter component, i.e., the minimum mass of the more massive component in the binary system is at least
$$M_{\mathrm{min}}=0.5(\mathrm{M}_1+\mathrm{M}_2)=\frac{0.32}{\mathrm{sin}i}\mathrm{M}_{}.$$
(1)
The real mass is very likely to be higher because of the inclination effect. On the other hand, the inclination is probably fairly close to edge-on, because the lower limit on mass is already rather high for the spectral type (this luck allows us to find a meaningful lower limit). The hatched region in Fig. 3 is defined by the known age of the association and the lower limit to the dynamical mass from the orbit. The upper limit to the mass is set from the late spectral type. Even on the main sequence, an M2 dwarf has a mass of about $`M=0.4`$ M (Delfosse et al.,, 2000), the upper limit we adopt. Our object is both cooler and younger than an M2 dwarf, making our upper limit a safely conservative estimate.
Temperatures of young dwarfs have been investigated by Luhman, (1999) and Luhman et al., (2003). Investigating members of the young cluster IC 348 and the young quadruple system GG Tau, they provide a temperature scale calibrated by evolutionary models under the premises that coeval stars fall on the same isochrones. These temperatures are intermediate between older giant and dwarf temperature scales, leading to an estimate for the temperature(s) of UScoCTIO 5 of $`T_{\mathrm{eff}}3270\pm 100`$ K. The error here is that cited by Luhman, and is his estimate of possible systematic errors in his scale (which was constructed to fill a specific purpose, relevant in our context). Mohanty et al., 2004a have discussed the temperatures of the cool components of GG Tau in detail (their Appendix B). They have the advantage of a spectroscopically determined temperature in addition to newer spectral types. They find that the temperatures of these young objects (comparable to UScoCTIO 5) are consistent with newly determined dwarf scales for main sequence M-dwarfs given in Leggett et al., (2000) or Mohanty & Basri, (2003). This leads to an estimate for the temperature of UScoCTIO 5 (M4) of $`T_{\mathrm{eff}}3175\pm 100`$ K.
In Fig. 4, UScoCTIO 5 is plotted in an HR-diagram with evolutionary tracks from the same models. With the effective temperature derived from a scale that is designed to fit isochrones of these evolutionary models, it is not surprising that UScoCTIO 5 is located near the 5 Myr isochrone. If we adopt the temperature scale based on the work of Mohanty & Basri, (2003) and Mohanty et al., 2004a which is about 100K cooler, the effect is to lower the inferred mass to about $`M=0.19`$ M and the age to 4 Myr (still fully consistent with the age estimated for the association). An equally weak link in setting the temperature is the assumed spectral type. A change to M3.5 from M4, for example, would have a substantial effect on the implied temperature, and raise the inferred mass to better agreement with the models. We intend to determine the temperature of UScoCTIO 5 directly from the high resolution spectra in a follow-up paper. It is therefore the case that the mass discrepancy cannot be demonstrated as conclusively in the HR-diagram with the current uncertainties in temperature, and our conclusions rest primarily on the (similar) result in the luminosity-age diagram.
The hatched region in Fig. 4 displays the area consistent with our previously adopted mass and age limits. The position of UScoCTIO 5 inferred for its dynamically determined mass from evolutionary models is too hot to be consistent with its measured luminosity. With the error bars on temperature (assuming a spectral type of M4) adopted from Luhman, (1999), UScoCTIO 5’s position in the HRD is barely consistent with even the minimum mass at $`i=90^{}`$ and equal mass components. The conclusion from Fig. 4 is that the mass of UScoCTIO 5 is probably significantly underestimated by the evolutionary model parameters that would fit the observational temperature and luminosity. The mass estimated by models in the luminosity-temperature diagram is consistent with that estimated from the luminosity-age diagram, because the temperature scale employed positions objects near the 5 Myr isochrone by construction. This probably means that the temperature is not too far off.
Besides a problem with the evolutionary models themselves, the underestimation of mass could only be due to (1) uncertainty in the dynamical mass, (2) underestimation of the temperature, and/or (3) mis-estimation of the spectral type by a half-subclass cooler or more. The dynamical mass does not have a large error, except that the unknown inclination can easily make the mass discrepancy worse. The estimation of luminosity is unlikely to have a large error, and this only plays a minor role since evolutionary tracks are largely independent of luminosity at a given mass; they are nearly vertical in the hatched region in Fig. 4. An error in the spectral type, or in the conversion of spectral type to temperature, could move the observed point leftward in Fig. 4. Note, however, that moving it to the value of 3450 K which is suggested by the dynamical mass, would also make the age uncomfortably large (as the luminosity would be too low). We leave the detailed analysis of these issues to another paper, wherein we will obtain an independent spectroscopic determination of the temperature. Our current result from the luminosity-age analysis is not affected by errors in spectral type or its conversion to temperature in any case.
## 6 Conclusions
We have derived the orbit of the M4 binary system UScoCTIO 5, a member of the Up Sco OB association with an age of $`5`$ Myr. From our radial velocity measurements we infer a projected minimum total mass of the system of $`(M_1+M_2)\mathrm{sin}i=(0.64\pm 0.02)`$ M, i.e. a minimum mass of the primary of $`M=0.32`$ M. The luminosity and effective temperature of UScoCTIO 5 have been derived from photometry and spectral type. As our main result, we compare the position of UScoCTIO 5 to theoretical evolutionary tracks of Baraffe et al., (1998) and Chabrier et al., (2000) in a luminosity-age diagram. These evolutionary tracks predict a mass of $`M0.23`$ M at the luminosity and age observed for UScoCTIO 5. Taking into account our ignorance of the inclination $`i`$ of the orbit, and the mass ratio $`M_2/M_1`$ (both of which can only lead to higher dynamical mass), we conclude that there is a significant discrepancy between the “true” mass and the mass derived from current commonly-used evolutionary models in this part of parameter space.
This conclusion is the same in our analysis using the HR-diagram, although there it is subject to further uncertainties in the true spectral type, and the conversion from spectral type to temperature at this age. Uncertainties in the derivation of $`T_{\mathrm{eff}}`$ lead to lower masses and an even larger discrepancy with the evolutionary models if the temperature is cooler than we estimate here; there is some suggestion for that from Mohanty et al., 2004a . One could obtain higher mass (and better agreement with the models) if the spectral type were earlier than the current value by half a subclass or more (using the current temperature scale). The same issues are present in the analysis by Close et al., (2005), who also found that models underestimated the mass for an object that is lower in mass but substantially older than this one. These uncertainties, however, do not affect the conclusions of the first paragraph.
These problems with temperatures are not applicable to the similar conclusions about model discrepancies reached earlier in the analysis of a number of objects in Upper Sco having even lower masses by Mohanty et al., 2004b , because they determined the temperatures of those objects more directly (from high resolution spectra). However, their masses are found by a surface gravity analysis rather than dynamically. The effect of the current analysis, therefore, is to support the methodologies of the earlier paper with an independent check. It seems fair to say at this point that the evidence is accumulating from several directions that commonly-used evolutionary models require adjustments to higher mass at a given luminosity and/or temperature for stars or brown dwarfs in the 0.03-0.3 solar mass range which are younger than a hundred million years or so. This has implications for all the conclusions drawn about the initial mass function in star-forming regions at low masses, and for other papers which rely on evolutionary mass estimates to draw conclusions about what sort of low-mass objects are being studied.
We would like to acknowledge the kind assistance of several observers who willingly aided this project during time they had for other purposes. They include Wallace Sargent and Masahide Takada-Hidai from CIT, Angelle Tanner from SAO, and Diane Paulsen from U Mich. This work is based on observations obtained from the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration, on data taken at Las Campanas Observatory, and at CTIO which is operated by AURA for the NSF. We would like to acknowledge the great cultural significance of Mauna Kea for native Hawaiians and express our gratitude for permission to observe from atop this mountain. GB thanks the NSF for grant support through AST00-98468. AR has received research funding from the European Commission’s Sixth Framework Programme as an Outgoing International Fellow (MOIF-CT-2004-002544).
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# 1 Emission rate normalized to the vacuum rate versus frequency for an 𝑥-oriented dipole in the central hole of a PC membrane (details in text). Black lines correspond to dipoles in the slab (𝑧0), as listed in the legend.
## Abstract
We show theoretically that photonic crystal membranes cause large variations in the spontaneous emission rate of dipole emitters, not only inside but also in the near-field above the membranes. Our three-dimensional finite difference time-domain calculations reveal an inhibition of more than five times and an enhancement of more than ten times for the spontaneous emission rate of emitters with select dipole orientations and frequencies. Furthermore we demonstrate theoretically, the potential of a nanoscopic emitter attached to the end of a glass fiber tip as a local probe for mapping the large spatial variations of the photonic crystal local radiative density of states. This arrangement is promising for on-command modification of the coupling between an emitter and the photonic crystal in quantum optical experiments.
Spontaneous emission in the near-field of 2D photonic crystals
A. Femius Koenderink
Laboratory of Physical Chemistry, Swiss Federal Institute of Technology (ETH), 8093 Zurich, Switzerland
Maria Kafesaki
IESL, Foundation for Research and Technology Hellas (FORTH), P.O. Box 1527, 71110 Heraklion, Crete, Greece
Costas M. Soukoulis
Ames Laboratory, Iowa State University, Ames, Iowa 50011, and IESL, FORTH, P.O. Box 1527, 71110 Heraklion, Greece, and Dept. of Materials Science and Technology, University of Crete, Greece
Vahid Sandoghdar
Laboratory of Physical Chemistry, Swiss Federal Institute of Technology (ETH), 8093 Zurich, Switzerland
OCIS codes: 130.130, 160.0160, 180.5810, 270.5580
It is well known that the rate of spontaneous emission can be controlled by the geometry of the medium surrounding the fluorescent species. In particular, many recent research efforts have been devoted to studying spontaneous emission in photonic crystals. The quantitative interpretation of these experiments, however, remains frustrated by lack of detailed information about many parameters that strongly affect the emission dynamics. These include the exact position of the emitters on the subwavelength scale and the orientation of the emission dipole moments, as well as systematic effects such as surface-induced quenching or other chemical or electronic surface phenomena. An ideal arrangement would require accurate placement of a single emitter at an arbitrary location in a photonic crystal (PC). Very recently Badolato *et al*. have achieved this by precise fabrication of a PC structure around a given semiconductor emitter. In this Letter we discuss the *in-situ* control of the position and thereby modification of the spontaneous emission rate of a single emitter close to or in a two-dimensional PC slab.
Two-dimensional (2D) photonic crystals fabricated in thin semiconductor membranes promise to achieve many of the long-standing goals of photonic band gap materials. Indeed, recently it has been demonstrated that it is possible to achieve very high-Q and low mode volume cavities in these structures. Due to their planar nature, PC membranes can be easily accessed by subwavelength probes such as optical fiber or atomic force microscope tips. Motivated by this opportunity, we investigate the prospects of coupling between a PC and nanoscopic optical emitters located at the end of sharp probes. Although 2D crystals do not yield a zero density of states, we show that both inside and in the near-field above a PC membrane the emission rate of properly oriented dipoles can be strongly modified. We show that the nanometer accuracy in scanning probe positioning allows direct mapping of the dependence of the emission rate on the spatial coordinates of the subwavelength emitter.
We have used the three-dimensional Finite-Difference Time-Domain (FDTD) method to calculate the local radiative density of states (LRDOS), accounting for the position dependence of the photon states available for fluorescent decay of a quantum emitter. This calculation relies on the fact that the LRDOS appearing in the formulation of Fermi’s Golden Rule for the spontaneous emission rate, also describes the total emitted power of a classical point-dipole antenna run at a fixed current. We consider semiconductor membranes with dielectric constant $`ϵ=11.76`$ and thickness $`d=250`$ nm, surrounded by up to $`1\mu `$m of air above and below. We take the membrane to contain a hexagonal array of holes with radius $`r=0.3a`$ at a lattice spacing of $`a=420`$ nm. Such a structure possesses a band gap in the range $`a/\lambda =0.25`$ to 0.33 for the transverse electric (TE) mode where the electric field is parallel to the plane of the membrane. The ratio $`a/\lambda `$ is used as normalized frequency units throughout our work. We used discretization with 14 or 20 points per lattice constant and employed volume-averaging of the dielectric constant to reduce staircasing errors. We considered finite hexagonally shaped PC structures up to $`25`$ holes across, terminated by the unperforated slab extending into Liao’s absorbing boundary conditions. By broadband temporal excitation of the dipole, we simulated the emission power spectrum over a wide frequency range. After dividing the resulting spectrum by that of an identically excited dipole in vacuum, we obtained the LRDOS normalized to the vacuum LRDOS .
For an emitter half way deep in a PC membrane, the solid black spectrum in Fig. 1 shows a strong inhibition of fluorescence by over a factor of 7 in the band gap as compared to its vacuum rate. In this case the dipole was laterally centered in the structure, and its orientation was chosen in the $`x`$-direction (cf. Fig. 1). The slab was taken to be $`13a`$ across, and we verified that no significant further reduction of the emission rate was obtained if we increased the size of the structure. However, the magnitude of the enhancement at the blue edge of the gap, as well as the Fabry-Pérot oscillations at frequencies below the gap depend on the finite size of the PC structure. For all tested structures wider than 7 holes across, we find emission enhancements larger than a factor of 15, representing a jump of two orders of magnitude as compared to the LRDOS for frequencies in the gap.
Given the inherent strong modulation of the dielectric constant in a PC structure, it is particularly interesting to examine the lateral dependence of the spontaneous emission rate. Figure 2(a) shows a contour plot of the LRDOS modification for an $`x`$-oriented dipole midway in the slab depth versus emission frequency and for lateral locations along a trajectory that traces the irreducible part of the unit cell (Fig. 2(d)). The emission is inhibited in the band gap at all positions whereas outside the gap we observe Fabry-Pérot modulations together with enhancement at the low and high frequency edges. The enhancement of the emission occurs especially on the high frequency edge of the gap (the ‘air band’) for dipoles in air holes, and predominantly at the low frequency edge (the ‘dielectric band’) for dipoles in the dielectric.
Next, we ask whether it is possible to capture these effects by scanning an emitter just above the PC slab. Different spectra in Fig. 1 show the LRDOS modification of a dipole laterally centered in the structure but at various heights $`z`$ above the membrane. In addition Fig. 2(b) and (c) display the modification of the LRDOS for the dipole right at the crystal-air interface and at $`110`$ nm above this plane. These data reveal that as $`z`$ increases, the inhibition and enhancement reduce in size. To examine this distance dependence more closely, in Fig. 2 we plot the normalized emission rate as a function of the distance between the dipole and the membrane surface for three key frequencies $`a/\lambda =0.23,0.28`$ and $`0.34`$ just below, in, and just above the band gap, respectively. Evidently, the inhibition diminishes for emitters located above the slab. In contrast, enhancements persist at the blue edge of the gap even if the dipole is lifted into air above the membrane. Figures 2 and 3 let us conclude that it is possible to enhance the spontaneous emission rate by a factor of $`5`$ to $`10`$ if the emitter position is controlled to within $`50`$ nm above the PC membrane. Note that the emission of a dipole near a simple homogeneous dielectric slab is also enhanced, due to coupling to the guided modes. However, at the gap edges the PC causes a further strong enhancement of the LRDOS.
The required resolution and control for mapping the modification of the emission rate can be achieved by scanning a subwavelength emitter attached to the end of a sharp tip. A crucial question arises as to the effect of the tip on the LRDOS. To estimate this effect, we have calculated the LRDOS for sources embedded inside cylindrical tips of diameter 125 nm pushed 130 nm into the central hole of the PC structure. We find that the influence of the PC structure on the spectrum of LRDOS enhancement and inhibition is unchanged for the system of emitters embedded in glass ($`ϵ2.25`$). In contrast, tips of very high-index material such as silicon fundamentally change the LRDOS spectrum, leading to the creation of a low-Q localized defect mode from the band edge due to the addition of dielectric material. The presence of a silicon tip causes an overall reduction of the gap depth and an increase and red-shift of the rate enhancement at the blue edge of the gap. Suitable probes of the LRDOS in photonic crystal membranes are therefore emitters inside low index tips.
Although optical detection of single emitters has become possible for some systems, many applications such as the realization of a nanolaser would benefit from coupling an ensemble of emitters to a PC. Furthermore, nanoscopic ensembles are more readily available than single emitter systems. Thus, we have also considered the modification of the spontaneous emission rate for a subwavelength ensemble of randomly oriented dipoles. Note that because the LRDOS is essentially unchanged for the TM polarization, dipolar components normal to the slab reduce the visibility of lifetime effects. We have considered over 25 in-equivalent dipole orientations (corresponding to over 300 orientations in $`2\pi `$ solid angle) in the central air hole and have calculated the corresponding LRDOS and luminescence extraction efficiency. We find that in general, the time-resolved flux of fluorescence photons extracted from the slab follows a significantly non-single exponential decay behavior. Nonetheless, the mean decay constant reveals inhibition by a factor three, and enhancement by a factor five compared to vacuum. The observation of inhibition is facilitated by the increase of the emission extraction efficiency for in-plane dipoles from $`20\%`$ for frequencies below the gap to $`>80\%`$ in the gap.
In conclusion, we have shown that strong inhibition and enhancement of emission can be achieved for emitters well inside photonic crystal membranes while a significant level of enhancement persists even in the near field above the structures. Since these results also hold for emitters embedded in nanoscopic dielectric probes, scanning probe technologies can be promising for on-command spontaneous emission control. An important advantage of emitters inside such probes is that they are shielded from unwanted interactions, and can be calibrated by simply retracting the probe from the structure.
This work was funded by the Deutsche Forschungsgemeinschaft (DFG) through focus program SP1113 and by ETH Zürich. Vahid Sandoghdar can be reached by e-mail at vahid.sandoghdar@ethz.ch.
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# Electrostatic self-energy and Bekenstein entropy bound in the massive Schwinger model
## 1 Introduction
Using the generalized second law for black holes, a universal upper bound on the entropy of a macroscopic charged object can be obtained by a gedanken experiment known as Geroch process. In this process the object is lowered adiabatically toward the horizon and then is assimilated from a small proper distance into the hole ,. On the basis of the no-hair theorem it is claimed that this bound cannot be improved by other quantum numbers which are carried by the object, e.g. the baryon number . A direct proof for the special case of a particle with a scalar charge, is given in . This proof is based on the fact that the self-energy receives no contribution from the scalar charge and the entropy of the assimilated object consists only of its gravitational energy.
A quantum effect relevant to this subject is the radiative correction involving the loop graphs, the so called vacuum polarization. Because of computational difficulties in studying the quantum effects in four dimensional curved space-times, one can consider lower dimensional models as a framework to explore these effects and to obtain the physical clues to the real four dimensional cases.
In it is shown that in the massless Schwinger model , the Bekenstein entropy upper bound is not affected by vacuum polarization. This lies on the fact that the effect of vacuum polarization is only appeared in the mass gained by the gauge field. This can not change the entropy upper bound since the gauge field behaves like a conformal massless field near the horizon.
In this paper, using the generalized second law for black-holes, we compute the entropy upper bound of a charged macroscopic object in the massive Schwinger model -.
The paper is organized as follows: The section two, is devoted to a brief introduction of massive Schwinger model in curved space-time. In the section three we discuss the confining and screening aspects of this model in static black-hole background. We show that our results asymptotically reduce to ones obtained in where the WKB approximation has been used. Also we compute the electrostatic self-energy of a charged particle very close to the horizon in terms of Bessel functions. These results are used to check the more complicated expression of self-energy derived in the section four where a term proportional to scalar curvature is included in the metric expansion. Our method is based on Taylor expansion of the metric in Schwarzschild gauge near the horizon, leading to Rindler (section three) or AdS space-time (section four) depending on the order of our approximation. Finally in the section five, we find out when an external charge is assimilated into the hole, the massive dynamical fermions affects the hole’s entropy via the self-energy of the charge. This effect disappears when dynamical fermions are massless.
In this paper the charges of the hole and the object are assumed to be very small in the scale of the hole mass. The units $`c=G=\mathrm{}=1`$ are used throughout the paper.
## 2 preliminaries
Since all two dimensional spaces are conformally flat, the most general stationary metric reads
$$ds^2=\sqrt{g(x)}\left(dt^2dx^2\right),$$
(1)
where $`g(x)=|\mathrm{det}g_{\mu \nu }(x)|`$, and $`g_{\mu \nu }`$ depends only on the spatial coordinate $`x`$. On this space-time, the quantum electrodynamics of fermions of mass $`m`$ and charge $`e`$ is described by the action ,
$$S=\left[\frac{1}{4}g^{\mu \nu }g^{\lambda \beta }F_{\mu \lambda }F_{\nu \beta }+i\overline{\psi }\gamma ^\mu \left(_\mu ieA_\mu \right)\psi m\overline{\psi }\psi \right]\sqrt{g(x)}𝑑x𝑑t.$$
(2)
$`\gamma ^\mu =e^{\mu a}\gamma _a`$ is the curved space-time counterparts of Dirac gamma matrices $`\gamma _a`$ and the zweibeins $`e^{\mu a}`$’s satisfy $`e^{\mu a}e_{\nu a}=\delta _\nu ^\mu `$. The covariant derivative, acting on fermions, is given by $`_\mu =_\mu +\frac{1}{2}\omega _\mu ^{ab}\sigma _{ab}`$ in which $`\sigma _{ab}=\frac{1}{4}[\gamma _a,\gamma _b]`$ and $`\omega _\mu ^{ab}`$’s are spin-connections. $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ are the field strength tensor components. The vacuum angle is assumed to be zero.
The corresponding bosonic action, in the presence of covariant conserved current
$$J^\mu =\frac{e^{}}{\sqrt{g(x)}}\left[\delta (xa)\delta (xb)\right]\delta _0^\mu ,$$
(3)
describing two opposite external charges $`e^{}`$ and $`e^{}`$ located at $`x=a`$ and $`x=b`$, respectively, is
$`S`$ $`=`$ $`{\displaystyle }[{\displaystyle \frac{1}{2}}\sqrt{g}g^{\mu \nu }_\mu \varphi _\nu \varphi +{\displaystyle \frac{e}{\sqrt{\pi }}}F\varphi +{\displaystyle \frac{m}{\pi }}g^{\frac{1}{4}}(x)\mathrm{exp}[2\pi G(𝐱,𝐱)]`$ (4)
$`N_\mu \mathrm{cos}(2\sqrt{\pi }\varphi )+{\displaystyle \frac{1}{2\sqrt{g(x)}}}F^2+{\displaystyle \frac{e}{\sqrt{\pi }}}\eta F]dxdt.`$
The dual field strength $`F`$ is defined through $`F=\widehat{ϵ}^{\mu \nu }_\mu A_\nu `$, where $`\widehat{ϵ}^{01}=\widehat{ϵ}^{10}=1`$. $`N_\mu `$ is normal ordering with respect to scale
$$\mu =\frac{e}{\sqrt{\pi }},$$
(5)
and $`G(𝐱,𝐱):=G_\mu (𝐱,𝐱)D(𝐱,𝐱)`$, $`𝐱=(t,x)`$. $`G_\mu (𝐱,𝐱)`$ is the Green function of a massive scalar field of mass $`\mu `$ and
$$D(𝐱_i,𝐱_j)=\frac{1}{4\pi }\mathrm{ln}|𝐱_i𝐱_j|^2.$$
In terms of step functions, $`\eta `$ is defined as
$$\eta =\frac{\sqrt{\pi }e^{}}{e}\left[\theta (xb)\theta (xa)\right].$$
In the Schwarzschild gauge, the metric
$$ds^2=f(r)dt^2\frac{1}{f(r)}dr^2,r>h$$
(6)
describes a black hole whose the horizon is located at $`r=h`$. $`h`$ is defined through $`f(h)=0`$. We assume that $`f(r)`$ is a positive $`𝒞^{\mathrm{}}`$ function for $`r>h`$. In tortoise coordinates $`(t,x)`$, defined by $`dx^2=dr^2/f^2(r)`$, the metric (6) reduces to (1) for $`f(r)=\sqrt{g(x)}`$. By integrating over the gauge fields of action (4), we obtain
$$S_{\mathrm{eff}.}=\left[\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi +m\mathrm{\Sigma }N_\mu \mathrm{cos}(2\sqrt{\pi }\varphi )\frac{\mu ^2}{2}(\varphi +\eta )^2\right]𝑑r𝑑t,$$
(7)
where $`g_{tt}=1/g_{rr}=f(r)`$ and $`\mathrm{\Sigma }:=\mathrm{exp}[2\pi G(𝐱,𝐱)]/(\pi f^{\frac{1}{2}}(r))`$ is the chiral condensate $`<\overline{\psi }\psi >`$ .
The electrostatic energy of external charges can be calculated by either performing a typical Wilson loop calculation, or by computing the ground state expectation value of the Hamiltonian in the presence of external source
$$E_{\mathrm{elec}.}:=<\mathrm{\Omega }_Q|H_Q|\mathrm{\Omega }_Q><\mathrm{\Omega }_0|H_0|\mathrm{\Omega }_0>.$$
(8)
$`H_0(H_Q)`$ and $`|\mathrm{\Omega }_0>(|\mathrm{\Omega }_Q>)`$ are the Hamiltonian and the ground state in the absence (presence) of the test charges, respectively . In the static case, the energy of the system measured by an observer whose two-velocity is parallel to the global time-like Killing vector of the space-time, i.e. $`u=(f^{\frac{1}{2}},0)`$, is
$$E=T_0^0𝑑r=L𝑑r.$$
(9)
$`T_\mu ^\nu `$ is the energy-momentum tensor and $`L`$ is the Lagrangian density. For small $`\varphi `$, where $`\mathrm{cos}(2\sqrt{\pi }\varphi )12\pi \varphi ^2`$, the action (4) becomes quadratic with respect to the gauge fields (this can be seen by integrating out the fermionic degrees of freedom in (2) or the bosonic ones in (4)). Hence in the static case, the electrostatic energy of charges in terms of the gauge field Green function $`𝒢(r,r^{})`$ becomes
$$E_{\mathrm{elec}.}=\frac{e^2}{2}\left[𝒢(r_a,r_a)+𝒢(r_b,r_b)2𝒢(r_a,r_b)\right].$$
(10)
$`e^2𝒢(r_a,r_b)`$ is the interaction energy, $`E_{\mathrm{int}.}`$, of two opposite charges $`e^{}`$ and $`e^{}`$ located at $`r_a:=r(a)`$ and $`r_b:=r(b)`$, and $`\left(e^2/2\right)𝒢(r,r)`$ is the self-energy $`E_{\mathrm{self}.}`$. In this paper, by self-energy of a charge we mean the change of energy of the system when the charge is added to it. Note that in contrast to higher dimensional cases, $`𝒢(x,x)`$, where $`x`$ is the spatial coordinate of the two dimensional space-time, is not infinite. For example in the massless Schwinger model on flat space-time, we have $`𝒢(x,y)=1/\left(2\mu \right)\mathrm{exp}[\mu |xy|]`$, which becomes the constant $`𝒢(x,x)=1/\left(2\mu \right)`$ in the coincident limit .
Up to the first order of $`m`$ and when $`m\mathrm{\Sigma }\mu ^2`$, the electrostatic energy of widely separated charges, i.e. $`\eta =(e^{}/e)\sqrt{\pi }`$, is ,
$$E_{\mathrm{elec}.}=E(\eta )E(0)=m\left[1\mathrm{cos}\left(2\pi \frac{e^{}}{e}\right)\right]_{r_a}^{r_b}\mathrm{\Sigma }𝑑r.$$
(11)
Note that the action is not quadratic here and the equation (10) is not valid. For constant values of $`\eta `$, using (4), one can show that the gauge fields with zero winding number ( $`F𝑑r𝑑t=0`$) do not contribute to the electrostatic energy of external charges.
For finitely separated charges, $`\eta `$ is no more a constant and the solution of the classical equation of motion for the field $`|\varphi |1`$ is
$$\varphi =\frac{\mu ^2}{_rg^{rr}_r+\mu ^2+4\pi m\mathrm{\Sigma }}\eta .$$
(12)
In this regime, the action becomes Gaussian and the classical solutions of the action coincide with the quantum ones. Putting (12) back into Lagrangian (7), eq.(9) results
$$E_{\mathrm{elec}.}=\frac{\mu ^2}{2}\left[_{r_a}^{r_b}\left(\eta ^2\mu ^2\eta \frac{1}{_rg^{rr}_r+\mu ^2+4\pi m\mathrm{\Sigma }}\eta \right)𝑑r\right].$$
(13)
This equation can be rewritten as
$$E_{\mathrm{elec}.}=\left(\frac{e^{}}{\mu }\right)^2\left[\frac{\mu ^2}{2}(r_br_a)+\mathrm{\Pi }\right],$$
(14)
where
$`\mathrm{\Pi }`$ $`:=`$ $`{\displaystyle \frac{\mu ^4}{2}}{\displaystyle _{r_a}^{r_b}}dr[{\displaystyle _{r^{}=r_a}^{r^{}=r}}G(r_>=r,r_<=r^{})dr^{}`$ (15)
$`+`$ $`{\displaystyle _{r^{}=r}^{r^{}=r_b}}G(r_>=r^{},r_<=r)dr^{}],`$
and
$$\left(_rg^{rr}_r+\mu ^2+4\pi m\mathrm{\Sigma }\right)G(r,r^{})=\delta (r,r^{}).$$
(16)
We denote the smaller (bigger) argument of the Green function by $`>`$ ($`<`$).
In , for a slowly varying metric $`df^{\frac{1}{2}}(r)/dr\overline{\mu }`$, where
$$\overline{\mu }^2=\mu ^2+4\pi m\mathrm{\Sigma },$$
(17)
$`E_{\mathrm{elec}.}`$ has been derived as
$`E_{\mathrm{elec}.}`$ $`=`$ $`{\displaystyle \frac{e^2}{2}}(1{\displaystyle \frac{\mu ^2}{\overline{\mu }^2}})(r_br_a)+{\displaystyle \frac{e^2\mu ^2}{4\overline{\mu }^3}}[f^{\frac{1}{2}}\left(r_a\right)`$ (18)
$`+`$ $`f^{\frac{1}{2}}\left(r_b\right)2f^{\frac{1}{4}}\left(r_a\right)f^{\frac{1}{4}}\left(r_b\right)\mathrm{exp}[{\displaystyle _{r_a}^{r_b}}\overline{\mu }f^{\frac{1}{2}}(u)du]].`$
The above equation has been obtained in the zeroth order of WKB approximation . In this small curvature limit, (18) is the leading term of the energy. If the metric were a constant, the self-force (the derivative of the self-energy) would vanish and the energy should become the same as the flat case, shifted by the metric factor $`f^{1/2}`$. This result is consistent with the measurement of a constantly accelerated observer. The condition of validity of (18) is $`e^{}e`$ and the first expression in (18) reproduces (11) in this limit.
Near the horizon, WKB approximation fails . This can be related to zero frequency modes of massive scalar fields appeared in the Schwinger model . So, to obtain the energy of external charges near the horizon, we must use another approximation. Note that the solutions of (16) in regions near and far from the horizon must be matched asymptotically.
## 3 $`q\overline{q}`$ potential near horizon in the massive Schwinger model
Near the horizon of a non-extremal black-hole, we can expand the metric as
$$f(r)=\kappa (rh)+O(rh)^2,$$
(19)
where $`\kappa :=df/dr`$ is twice of the surface gravity. This metric describes Rindler space-time. In this region, $`G(r,r^{})`$ satisfies
$$\left[_r\kappa (rh)_r+\mu ^2+4\pi m\mathrm{\Sigma }\right]G(r,r^{})=\delta (r,r^{}),$$
(20)
where $`\mathrm{\Sigma }=\mu \mathrm{exp}(\gamma )/\left(2\pi \right)`$ and $`\gamma `$ is the Euler constant. Two independent solutions of the corresponding homogeneous equation are $`I_0\left(2\overline{\mu }\sqrt{(rh)/\kappa }\right)`$ and $`K_0\left(2\overline{\mu }\sqrt{(rh)/\kappa }\right)`$, where $`I_0`$ and $`K_0`$ are modified Bessel functions. The well-defined Green function is then
$$G(r,r^{})=\frac{2}{\kappa }K_0\left(2\overline{\mu }\sqrt{\frac{r_>h}{\kappa }}\right)I_0\left(2\overline{\mu }\sqrt{\frac{r_<h}{\kappa }}\right).$$
(21)
Defining $`𝒦`$ and $``$ through
$`{\displaystyle \frac{d𝒦}{dr}}`$ $`=`$ $`K_0\left(2\overline{\mu }\sqrt{{\displaystyle \frac{rh}{\kappa }}}\right),`$
$`{\displaystyle \frac{d}{dr}}`$ $`=`$ $`I_0\left(2\overline{\mu }\sqrt{{\displaystyle \frac{rh}{\kappa }}}\right),`$ (22)
then $`\mathrm{\Pi }`$ in (15) reduces to
$`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \frac{\mu ^4}{\kappa }}\{{\displaystyle _{r_a}^{r_b}}W[(r),𝒦(r)]dr`$ (23)
$`+`$ $`[(r_b)𝒦(r_b)+(r_a)𝒦(r_a)2(r_a)𝒦(r_b)]\},`$
where $`W`$ is the Wronskian. Considering the recurrence formulas
$$\frac{d}{dx}\left[x^nI_n(x)\right]=x^nI_{n1}(x),\frac{d}{dx}\left[x^nK_n(x)\right]=x^nK_{n1}(x),$$
one can show
$`𝒦(r)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\kappa }}{\overline{\mu }}}\left(rh\right)^{\frac{1}{2}}K_1\left(2\overline{\mu }\sqrt{{\displaystyle \frac{rh}{\kappa }}}\right),`$
$`(r)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\kappa }}{\overline{\mu }}}\left(rh\right)^{\frac{1}{2}}I_1\left(2\overline{\mu }\sqrt{{\displaystyle \frac{rh}{\kappa }}}\right).`$ (24)
Using $`W[K_1(x),I_1(x)]=1/x`$, we arrive at $`W[(r),𝒦(r)]=\kappa /(2\overline{\mu }^2)`$. Therefore
$`E_{\mathrm{elec}.}`$ $`=`$ $`{\displaystyle \frac{e^2}{2}}\left(1\left({\displaystyle \frac{\mu }{\overline{\mu }}}\right)^2\right)\left(r_br_a\right)`$ (25)
$``$ $`{\displaystyle \frac{e^2\mu ^2}{\kappa }}\left[(r_b)𝒦(r_b)+(r_a)𝒦(r_a)2(r_a)𝒦(r_b)\right].`$
The first and the last terms describe the interaction of opposite charges. While the first term corresponds to confinement aspects of the model, the last term illustrates the screening effect. Note that both the screening and confining phenomenon appear in the same problem, similar to what happened in the flat case . For $`m=0`$ ($`\mu =\overline{\mu }`$), the confining term disappears. The second and the third terms are the self-energies of the charges. Writing eq.(25) in the form (10), we get
$$𝒢(r,r^{})=\frac{1}{2}\left[1\left(\frac{\mu }{\overline{\mu }}\right)^2\right]\left(r_>r_<\right)+\frac{2\mu ^2}{\kappa }(r_<)𝒦(r_>).$$
(26)
$`\left(2\mu ^2/\kappa \right)(r_<)𝒦(r_>)`$ is the only term which contributes to the self-energy in (25). As this part of the Green function satisfies Dirichlet boundary condition at the horizon, we find that $`E_{\mathrm{self}.}(rh)=0`$. Far from the horizon, WKB approximation is applicable and from (18) we obtain $`E_{\mathrm{self}.}(rh)=e^2\mu ^2/\left(4\overline{\mu }^3\right)`$.
We can use the global method of Smith and Will to determine the self-force. In a free falling system, the work done by the force $`F`$ to displace slowly (such that the location of the event horizon remains unchanged) the test charge by an infinitesimal distance $`\delta \overline{r}`$ toward the horizon is
$$\delta \overline{W}=F\delta \overline{r}.$$
(27)
The corresponding energy detected by an observer at asymptotic infinity will be red-shifted
$$\delta E=\sqrt{g_{tt}(r)}\delta \overline{W}.$$
(28)
This change will be manifested by a change in the asymptotic mass M of the system, given by the total mass variation law of Carter
$$\delta \text{M}=\delta _h^{\mathrm{}}T_t^t𝑑r,$$
(29)
where $`T_t^t`$ is the component of energy momentum tensor and is the same as the effective Lagrangian in (7). Therefore $`\delta \text{M}=\delta E_{\mathrm{self}.}`$. Hence by transforming locally the flat coordinates denoted by $`\overline{r}`$ to the Schwarzschild ones, we obtain
$$F=\frac{\delta E_{\mathrm{self}.}}{\delta r}.$$
(30)
To derive (25) we have assumed $`|\varphi (r)|1`$. Let us check this assumption more carefully. Note that $`\varphi `$ is the solution of the equation (12), obtained using eqs.(21) and (21) as:
$$\varphi (r)=\{\begin{array}{cc}\frac{2\mu e^{}}{\kappa }I_0\left(2\overline{\mu }\sqrt{\frac{rh}{\kappa }}\right)\left[𝒦(r_b)𝒦(r_a)\right]\hfill & \text{ }r<r_a<r_b,\hfill \\ \frac{2\mu e^{}}{\kappa }K_0\left(2\overline{\mu }\sqrt{\frac{rh}{\kappa }}\right)\left[(r_b)(r_a)\right]\hfill & r>r_b>r_a,\hfill \\ \frac{2\mu e^{}}{\kappa }[K_0\left(2\overline{\mu }\sqrt{\frac{rh}{\kappa }}\right)(r)I_0\left(2\overline{\mu }\sqrt{\frac{rh}{\kappa }}\right)𝒦(r)+\hfill & \\ +I_0\left(2\overline{\mu }\sqrt{\frac{rh}{\kappa }}\right)𝒦(r_b)K_0\left(2\overline{\mu }\sqrt{\frac{rh}{\kappa }}\right)(r_a)]\hfill & r_a<r<r_b.\hfill \end{array}$$
The expressions appeared in $`\varphi (r)`$ are in the forms
$$_1=\frac{2\mu e^{}}{\kappa }K_0\left(2\overline{\mu }\sqrt{\frac{r_>h}{\kappa }}\right)(r_<),$$
and
$$_2=\frac{2\mu e^{}}{\kappa }I_0\left(2\overline{\mu }\sqrt{\frac{r_<h}{\kappa }}\right)𝒦(r_>).$$
Using the following expansions near the horizon ,
$`K_0(x)`$ $`\left[\gamma +\mathrm{ln}\left({\displaystyle \frac{x}{2}}\right)\right]\left(1+{\displaystyle \frac{x^2}{4}}\right)+{\displaystyle \frac{x^2}{4}}+O(x^4),`$
$`K_1(x)`$ $`{\displaystyle \frac{1}{x}}+{\displaystyle \frac{x}{2}}\left[\mathrm{ln}\left({\displaystyle \frac{x}{2}}\right){\displaystyle \frac{\psi (1)+\psi (2)}{2}}\right]+O(x^3),`$
$`I_0(x)`$ $`1+{\displaystyle \frac{x^2}{4}}+O(x^4),`$
$`I_1(x)`$ $`{\displaystyle \frac{x}{2}}+O(x^3),`$ (31)
where $`\psi `$ is digamma function, one finds
$`_1`$ $``$ $`{\displaystyle \frac{2\mu e^{}(r_<h)}{\kappa }}\left[\gamma +\mathrm{ln}\left(\overline{\mu }\sqrt{{\displaystyle \frac{r_>h}{\kappa }}}\right)\right],`$
$`_2`$ $`\mu e^{}\{{\displaystyle \frac{1}{\overline{\mu }^2}}+{\displaystyle \frac{r_<h}{\kappa }}`$ (32)
$`+`$ $`{\displaystyle \frac{2(r_>h)}{\kappa }}[\mathrm{ln}\left(\overline{\mu }\sqrt{{\displaystyle \frac{r_>h}{\kappa }}}\right){\displaystyle \frac{\psi (1)+\psi (2)}{2}}]\}.`$
This shows that the assumption $`|\varphi |1`$ is applicable near the horizon, in which $`rh`$.
From the asymptotic expansions
$`I_\nu (z)`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}z^\nu (z^2)^{\frac{1}{4}(2\nu +1)}[e^{i\left[\sqrt{z^2}\frac{1}{4}(2\nu +1)\pi \right]}(1+O(z^1))`$
$`+`$ $`\mathrm{c}.\mathrm{c}.]`$
$`K_\nu (z)`$ $`\sqrt{{\displaystyle \frac{\pi }{2z}}}e^z\left[1+O(z^1)\right],`$ (33)
one can show that the eq.(25) asymptotically reduces to (18), as expected. This is the requirement explained at the end of section 2.
## 4 Electrostatic self-energy in the second order approximation
To study the effect of the curvature in higher order approximation of self-energy, we must improve and refine the approximation (19) by including terms proportional to the scalar curvature in the metric expansion near the horizon $`rh`$. We put
$$f(r)=\kappa (rh)+\frac{R}{2}(rh)^2+O((rh)^3),$$
(34)
where $`\kappa =f^{}(h)>0`$ and $`R=f^{\prime \prime }(h)`$. This describes a space-time with constant curvature $`R`$ (locally AdS or dS, depending on whether the sign of $`R`$ is negative or positive, respectively). In terms of the coordinate $`u=x+b`$, where
$$b=\frac{\kappa }{R},x=rh,$$
the equation (16) becomes
$$\left[_u\frac{R}{2}\left(u^2b^2\right)_u+\overline{\mu }^2\right]G(u,u^{})=\delta (u,u^{}).$$
(35)
In space-times with a constant positive curvature, $`\mathrm{\Sigma }`$ is a constant :
$$\mathrm{\Sigma }=\frac{\mu e^\gamma }{2\pi }\mathrm{exp}\left[\frac{1}{2}\left\{\mathrm{ln}\left(\frac{R}{2\mu ^2}\right)+\psi (\frac{1}{2}+\alpha )+\psi (\frac{1}{2}\alpha )\right\}\right],$$
(36)
where $`\alpha ^2=1/4+2\mu ^2/R`$. We restrict ourselves to small negative $`R`$, which we encounter in the next section where a dilatonic black hole with a large mass is considered. For $`|R|\mu ^2`$,the chiral condensate becomes
$$\mathrm{\Sigma }=\frac{e^\gamma \mu }{2\pi }\mathrm{exp}\left[\frac{R}{12\mu ^2}\right].$$
(37)
Defining
$$z:=\frac{u}{b},0<z<1;\stackrel{~}{\mu }^2:=\frac{2\overline{\mu }^2}{R}$$
the homogeneous counterpart of the equation (35) becomes
$$\left[\left(1z^2\right)_z^22z_z\stackrel{~}{\mu }^2\right]G_h(z)=0.$$
(38)
The real solutions of (38) can be expressed in terms of conical functions defined by
$`p_{\stackrel{~}{\nu }}^n(z)`$ $`=`$ $`P_{\stackrel{~}{\nu }}^n(z),`$
$`q_{\stackrel{~}{\nu }}^n(z)`$ $`=`$ $`{\displaystyle \frac{(1)^n}{2}}\left[Q_{\stackrel{~}{\nu }}^n(z)+Q_{\stackrel{~}{\nu }1}^n(z)\right]`$ (39)
$`=`$ $`{\displaystyle \frac{\pi }{2\mathrm{sin}(\pi \stackrel{~}{\nu })}}P_{\stackrel{~}{\nu }}^n(z).`$
$`P_{\stackrel{~}{\nu }}^n(z)`$ and $`Q_{\stackrel{~}{\nu }}^n(z)`$ are associated Legendre functions and $`\stackrel{~}{\nu }`$ is defined through $`\stackrel{~}{\nu }(\stackrel{~}{\nu }+1)=\stackrel{~}{\mu }^2`$. Note that in the small curvature limit $`2\overline{\mu }^2/R>\frac{1}{4}`$ or $`\stackrel{~}{\mu }^2>\frac{1}{4}`$, $`\stackrel{~}{\nu }`$ is a complex number. Near the horizon, $`z=1`$, we have
$`\underset{z1}{lim}q_{\stackrel{~}{\nu }}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}\left({\displaystyle \frac{1z}{2}}\right)\left[1+O(1z)\right]`$
$`{\displaystyle \frac{1}{2}}\left[\psi (1+\stackrel{~}{\nu })+\psi (\stackrel{~}{\nu })+2\gamma \right]\left[1+O(1z)\right],`$
$`\underset{z1}{lim}p_{\stackrel{~}{\nu }}(z)`$ $`=`$ $`1{\displaystyle \frac{\stackrel{~}{\nu }(\stackrel{~}{\nu }+1)}{2}}(1z)+O(1z)^2.`$ (40)
Hence the real well behaved Green function is
$$G(r,r^{})=\frac{2}{\kappa }p_{\stackrel{~}{\nu }}(z_>)q_{\stackrel{~}{\nu }}(z_<).$$
(41)
To arrive at this relation, we have used $`W[p_{\stackrel{~}{\nu }}^n(z),q_{\stackrel{~}{\nu }}^n(z)]=1/\left(1z^2\right)`$ .
Let us compare this result with one obtained in the previous section. In terms of $`\theta `$ defined by $`z=\mathrm{cos}\theta `$, (41) becomes
$$G(z,z^{})=\frac{2}{\kappa }p_{\stackrel{~}{\nu }}(\mathrm{cos}\theta _<)q_{\stackrel{~}{\nu }}(\mathrm{cos}\theta _>).$$
(42)
One can make use of the relations
$`\underset{\theta 0}{lim}\stackrel{~}{\nu }^np_{\stackrel{~}{\nu }}^n(\mathrm{cos}\theta )`$ $`=`$ $`i^nI_n(n\lambda ),`$
$`\underset{\theta 0}{lim}\stackrel{~}{\nu }^nq_{\stackrel{~}{\nu }}^n(\mathrm{cos}\theta )`$ $`=`$ $`i^nK_n(n\lambda ),`$ (43)
where $`n\lambda /\left(\mathrm{sin}\theta \right)=\stackrel{~}{\mu }`$, to obtain
$$G(z,z^{})=\frac{2}{\kappa }I_0(\stackrel{~}{\mu }\theta _<)K_0(\stackrel{~}{\mu }\theta _>),$$
(44)
in the vicinity of the horizon, i.e. $`\theta 0`$. In this region we have also $`x\theta ^2b/2`$ which yields $`\theta \stackrel{~}{\mu }=2\sqrt{x\overline{\mu }^2/\kappa }`$. Therefore in the limit $`x0`$, (41) tends to (21).
Putting back (41) into (15) we get
$$\mathrm{\Pi }=\frac{\mu ^4}{\kappa }W[𝒫,𝒬](r_br_a)\frac{\mu ^4}{\kappa }\left[𝒫(r_a)𝒬(r_a)+𝒫(r_b)𝒬(r_b)2𝒫(r_a)𝒬(r_b)\right],$$
(45)
where
$`{\displaystyle \frac{d𝒬(r)}{dr}}`$ $`=`$ $`q_{\stackrel{~}{\nu }}\left(1+{\displaystyle \frac{rh}{b}}\right),`$
$`{\displaystyle \frac{d𝒫(r)}{dr}}`$ $`=`$ $`p_{\stackrel{~}{\nu }}\left(1+{\displaystyle \frac{rh}{b}}\right).`$ (46)
With the help of the recurrence formulas ,
$`\left(1z^2\right){\displaystyle \frac{dP_{\stackrel{~}{\nu }}^1(z)}{dz}}`$ $`=`$ $`\stackrel{~}{\nu }zP_{\stackrel{~}{\nu }}^1(z)+\left(\stackrel{~}{\nu }1\right)P_{\stackrel{~}{\nu }1}^1(z),`$
$`\left(1z^2\right)^{\frac{1}{2}}P_{\stackrel{~}{\nu }}(z)`$ $`=`$ $`z\left(\stackrel{~}{\nu }+1\right)P_{\stackrel{~}{\nu }}^1(z)\left(\stackrel{~}{\nu }1\right)P_{\stackrel{~}{\nu }1}^1(z),`$ (47)
and
$`\left(1z^2\right){\displaystyle \frac{dq_{\stackrel{~}{\nu }}^1(z)}{dz}}`$ $`=`$ $`\stackrel{~}{\nu }zq_{\stackrel{~}{\nu }}^1(z)+\left(\stackrel{~}{\nu }+1\right)q_{\stackrel{~}{\nu }1}^1(z),`$
$`\stackrel{~}{\nu }\left(1z^2\right)^{\frac{1}{2}}q_{\stackrel{~}{\nu }}(z)`$ $`=`$ $`q_{\stackrel{~}{\nu }1}^1(z)zq_{\stackrel{~}{\nu }}^1(z),`$ (48)
one obtains
$`𝒬(r)`$ $`=`$ $`{\displaystyle \frac{b}{\stackrel{~}{\nu }\left(\stackrel{~}{\nu }+1\right)}}\left(1z^2\right)^{\frac{1}{2}}q_{\stackrel{~}{\nu }}^1(z),`$
$`𝒫(r)`$ $`=`$ $`b\left(1z^2\right)^{\frac{1}{2}}p_{\stackrel{~}{\nu }}^1(z).`$ (49)
The first equation of (4) is derived directly from the first identity of (4). To verify the second equation, we write it in terms of associated Legendre function $`P`$ as
$$P_{\stackrel{~}{\nu }1}^1(z)+zP_{\stackrel{~}{\nu }}^1(z)=\stackrel{~}{\nu }(1z^2)^{\frac{1}{2}}P_{\stackrel{~}{\nu }}(z).$$
(50)
Using
$$zP_{\stackrel{~}{\nu }}^\mu (z)P_{\stackrel{~}{\nu }+1}^\mu (z)=(\mu +\stackrel{~}{\nu })(1z^2)^{\frac{1}{2}}P_{\stackrel{~}{\nu }}^{\mu 1}(z),$$
(51)
(50) reduces to
$$P_{\stackrel{~}{\nu }1}^1(z)+zP_{\stackrel{~}{\nu }}^1(z)=\frac{\stackrel{~}{\nu }}{\stackrel{~}{\nu }+1}\left[zP_{\stackrel{~}{\nu }}^1(z)P_{\stackrel{~}{\nu }+1}^1(z)\right].$$
(52)
But (52) is nothing but the known equation:
$$P_\nu ^\mu (z)=\frac{2\nu +3}{\mu +\nu +1}zP_{\nu +1}^\mu (z)+\frac{\mu \nu 2}{\nu +\mu +1}P_{\nu +2}^\mu (z),$$
(53)
when one takes $`\nu =\stackrel{~}{\nu }1`$ and $`\mu =1`$.
Using $`W[𝒫,𝒬]=b/\left[\stackrel{~}{\nu }\left(\stackrel{~}{\nu }+1\right)\right]`$, $`E_{\mathrm{elec}.}`$ from (14) and (45) becomes
$`E_{\mathrm{elec}.}`$ $`=`$ $`{\displaystyle \frac{e^2}{2}}\left(1\left({\displaystyle \frac{\mu }{\overline{\mu }}}\right)^2\right)\left(r_br_a\right)`$ (54)
$``$ $`{\displaystyle \frac{e^2\mu ^2}{\kappa }}\left[𝒫(r_b)𝒬(r_b)+𝒫(r_a)𝒬(r_a)2𝒫(r_a)𝒬(r_b)\right].`$
From $`\mathrm{\Gamma }(\stackrel{~}{\nu })\mathrm{\Gamma }(1\stackrel{~}{\nu })=\pi /\left[\mathrm{sin}\left(\pi \stackrel{~}{\nu }\right)\right]`$, it follows that
$`\underset{z1}{lim}p_{\stackrel{~}{\nu }}^1(z)`$ $`=`$ $`(1+z)^{\frac{1}{2}}\left[(1z)^{\frac{1}{2}}{\displaystyle \frac{\stackrel{~}{\nu }(\stackrel{~}{\nu }+1)}{2\mathrm{\Gamma }(3)}}(1z)^{\frac{3}{2}}\right]+O(1z)^{\frac{5}{2}},`$
$`\underset{z1}{lim}q_{\stackrel{~}{\nu }}^1(z)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\nu }(\stackrel{~}{\nu }+1)}{2}}(1+z)^{\frac{1}{2}}(1z)^{\frac{1}{2}}[\mathrm{ln}\left({\displaystyle \frac{1z}{2}}\right)\psi (2)+\gamma +`$
$`\psi (\stackrel{~}{\nu })`$ $`+`$ $`\psi (\stackrel{~}{\nu }+1)\left]\right[1+O(1z)]+(1+z)^{\frac{1}{2}}(1z)^{\frac{1}{2}}.`$ (55)
We have also
$`\psi (\stackrel{~}{\nu })+\psi (\stackrel{~}{\nu }+1)`$ $`=`$ $`\psi ({\displaystyle \frac{1}{2}}+{\displaystyle \frac{i\lambda }{2}})+\psi ({\displaystyle \frac{1}{2}}{\displaystyle \frac{i\lambda }{2}})`$ (56)
$`=`$ $`2\mathrm{}\left(\psi ({\displaystyle \frac{1}{2}}+{\displaystyle \frac{i\lambda }{2}})\right),`$
where $`\lambda =\sqrt{4\stackrel{~}{\mu }^21}`$. For slowly varying metrics $`\lambda 1`$, one can use the expansion $`\psi (w)\mathrm{ln}w\left(1/2w\right)1/\left(12w^2\right)+O\left(w^4\right)`$, , to arrive at
$$\psi (\stackrel{~}{\nu })+\psi (\stackrel{~}{\nu }+1)=\mathrm{ln}\frac{2\overline{\mu }^2}{R}+\frac{8\overline{\mu }^2RR^2}{48\overline{\mu }^4}+O(\frac{R^3}{\overline{\mu }^6}).$$
(57)
Combining (54)-(57) one finds
$`E_{\mathrm{self}.}(x)`$ $`=`$ $`{\displaystyle \frac{e^2}{2}}{\displaystyle \frac{\mu ^2}{\overline{\mu }^2}}\{x+{\displaystyle \frac{\overline{\mu }^2}{\kappa }}x^2[\mathrm{ln}{\displaystyle \frac{\overline{\mu }^2x}{\kappa }}+\gamma \psi (2)+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{R}{6\overline{\mu }^2}}`$ (58)
$``$ $`{\displaystyle \frac{R^2}{48\overline{\mu }^4}}+O(R^3)]\}+O(x^3).`$
Note that at $`m=0`$, where $`\overline{\mu }=\mu `$, the above expression reduces to one obtained in .
## 5 Bekenstein entropy bound
Our aim is now to use the relation (58) to obtain the upper entropy bound of a charged object. To do so, we allow the black hole to carry a charge $`q`$ and assume $`q`$, $`e^{}`$ and $`\mu `$ to be very small with respect to the black hole mass. We consider two dimensional dilatonic charged black hole with mass $`M`$ and charge $`q`$:
$$ds^2=(12Me^r+q^2e^{2r})dt^2(12Me^r+q^2e^{2r})^1dr^2,$$
(59)
This metric emerges in the heterotic string theory as a solution of the action
$$S[g,\phi ,A]=\frac{1}{2}_{}d^2x\sqrt{g}e^\phi (R+(\phi )^2\frac{1}{2}F_{\mu \nu }F^{\mu \nu })_{}𝑑x\sqrt{I}e^\phi K,$$
(60)
describing $`2d`$ gravity coupled to dilatonic field $`\phi `$. $`K`$ is the extrinsic curvature and $`I`$ is the induced metric on $``$, where $``$ is the surface under study. The boundary term is added to make the variation procedure self-consistent. The metric (59) can support an electrostatic test charge . Thermodynamical quantities for this system can be obtained using the Massieu function, expressed in term of grand canonical partition function corresponding to the action (60). In $`2d`$, the horizon surface of the black hole is a point and we cannot consider the area. Nevertheless it might be useful to think about the value of dilatonic field at the horizon, $`\phi _h`$, as a quantity playing the rôle of the logarithm of an effective area. The event horizon of the black hole is located at
$$h=\mathrm{ln}[M+(M^2q^2)^{\frac{1}{2}}].$$
(61)
For $`q>M`$ we have a naked singularity. The entropy of the system is obtained as
$`S=2\pi e^{\phi _h}=2\pi [M+(M^2q^2)^{\frac{1}{2}}].`$ (62)
Comparing (59) with (6), gives $`f(r)=12Me^r+q^2e^{2r}`$. Expanding $`f(r)`$ near $`r=h`$, results $`\kappa >0`$ and $`R<0`$ for $`qM`$. $`\kappa `$ and $`R`$ are defined through eq.(34).
We now consider a charged object of rest mass $`m`$ and charge $`e`$ whose gravitational field is negligible on this background. This object is slowly (adiabatically) descended toward the black hole. This process causes no change in the horizon location and the entropy of the black hole remains unchanged . To find the change in black hole entropy caused by assimilation of the object, one should evaluate the energy at the point of capture,which is at a proper distance $`l`$ outside the horizon
$$l=_0^a\frac{dy}{[\kappa y+\frac{Ry^2}{2}]^{\frac{1}{2}}}=2\sqrt{\frac{a}{\kappa }}\frac{1}{6}R\left(\frac{a}{\kappa }\right)^{\frac{3}{2}}+O\left(\frac{a}{\kappa }\right)^{\frac{5}{2}}.$$
(63)
$`y`$ is the coordinate distance from the horizon and $`a`$ is the position of the center of mass of the object. In fact $`l`$ is the proper radius of the object at the point of capture. For small $`l`$ we get $`\sqrt{a/\kappa }=l/2+\left(R/96\right)l^3+O(l^5)`$. When the object is assimilated, its charge modifies the hole’s charge to $`q+e^{}`$, and its total energy, which we denote by $`E_T`$, augments the hole’s mass from $`M`$ to $`M^{}=M+E_T`$.
The energy of the object, $`E_T`$, is constituted of the energy of the body’s mass, $`m^{}`$, shifted by the gravitational field, $`E_m^{}m^{}\kappa l/2+O(l^3)`$, and the electrostatic self-energy $`E_{\mathrm{self}.}`$ (58) which in terms of $`l`$ is:
$`E_{\mathrm{self}.}`$ $`=`$ $`{\displaystyle \frac{e^2}{2}}{\displaystyle \frac{\mu ^2}{\overline{\mu }^2}}\{{\displaystyle \frac{\kappa l^2}{4}}+{\displaystyle \frac{\kappa Rl^4}{96}}+{\displaystyle \frac{\overline{\mu }^2\kappa l^4}{16}}[2\mathrm{ln}\left(\overline{\mu }({\displaystyle \frac{l}{2}}+{\displaystyle \frac{Rl^3}{96}})\right)+\gamma \psi (2)`$ (64)
$`+`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{R}{6\overline{\mu }^2}}{\displaystyle \frac{R^2}{48\overline{\mu }^4}}+O(R^3)]\}+O(l^6).`$
The final entropy of the black hole is
$$S_f=2\pi \left(2M+2E_T\frac{q^2+e^2+2qe^{}}{2M}+O(M^2)\right).$$
(65)
In the massive Schwinger model, the fermionic parameters $`\mu `$ and $`m`$ are appeared in the final black hole entropy through the self-energy $`E_T`$. In contrast, in the massless Schwinger model, where $`\overline{\mu }=\mu `$, any information about fermionic field is lost (up to the order $`l^4`$) . This result agrees with .
Assuming the validity of the generalized second law of thermodynamics
$$S_f(S_i+S),$$
(66)
where $`S_{f(i)}`$ is the black hole entropy in final (initial) state and $`S`$ is the object entropy, we obtain
$$S4\pi E_T\frac{\pi e^2}{M}\frac{2\pi e^{}q}{M}.$$
(67)
When $`M`$ is large with respect to $`e`$, $`q`$, $`m^{}`$, and $`e^{}`$, we can expand the surface gravity
$$\kappa =2\sqrt{M^2q^2}/\left(M+\sqrt{M^2q^2}\right),$$
(68)
as $`\kappa 1q^2/\left(4M^2\right)+O(M^4)`$. In this limit the leading terms in the inequality (67) are independent of black hole’s parameters
$$S2\pi m^{}l+\frac{\pi e^2}{2}\frac{\mu ^2}{\overline{\mu }^2}l^2+O(\frac{1}{M}).$$
(69)
This shows that the electrostatic self-energy modifies the Bekenstein upper bound.
In the massless Schwinger model, $`\mu =\overline{\mu }`$, the eq.(69) becomes
$$S2\pi m^{}l+\frac{\pi e^2}{2}l^2+O(\frac{1}{M}),$$
(70)
which indicates the absence of the fermionic information in the upper entropy bound of the object. This agrees to the fact that near the horizon, massive gauge fields act like conformal massless fields. Besides the rôle of vacuum polarization, a main difference with respect to the result obtained in $`QED_4`$, , is the sign and the order of terms in right hand side of (69), in other words the bound is not tightened here. For $`e^{}=0`$ we obtain the well known result
$$S2\pi m^{}l.$$
(71)
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# Fully three dimensional breather solitons can be created using Feshbach resonance
## Abstract
We investigate the stability properties of breather solitons in a three-dimensional Bose-Einstein Condensate with Feshbach Resonance Management of the scattering length and confined only by a one dimensional optical lattice. We compare regions of stability in parameter space obtained from a fully 3D analysis with those from a quasi two-dimensional treatment. For moderate confinement we discover a new island of stability in the 3D case, not present in the quasi 2D treatment. Stable solutions from this region have nontrivial dynamics in the lattice direction, hence they describe fully 3D breather solitons. We demonstrate these solutions in direct numerical simulations and outline a possible way of creating robust 3D solitons in experiments in a Bose Einstein Condensate in a one-dimensional lattice. We point other possible applications.
The creation of Bose–Einstein condensates (BEC) in vapors of alkali metals has opened an excellent opportunity to investigate nonlinear interactions of atomic matter waves. An important challenge of practical interest is to develop methods to create and control matter-wave solitons. One dimensional dark 1fsol , bright 2bright , and gap-mode ober solitons have already been observed. A promising approach to obtain multidimensional solitons consists in varying the scattering length a of interatomic collisions. This can be achieved by means of sweeping an external magnetic field through the $`a=0`$ point (the point, where the scattering length vanishes). This point occurs close to the Feshbach resonance 3Inouye . The application of an ac magnetic field may induce a periodic modulation of a, opening a way to “Feshbach-resonance management” (FRM) 4Greece . A noteworthy FRM-induced effect is the possibility of creating self-trapped oscillating BEC solitons (breathers) without an external trap in the 2D case. The underlying mechanism is fast modulations creating an effective potential on a slower timescale. This potential can stabilize the soliton. The BEC model based on the Gross-Pitaevskii equation (GPE) with harmonic modulation of a was investigated in Refs. 6Saito ; 8Abdul ; 9VPG . The conclusion was that FRM renders it possible to stabilize 2D breather solitons even without the use of an external trap. According to these references, 3D breathers require at least a tight, one dimensional harmonic trap Gerlitz , practically reducing the problem to 2D. Later on we will call this approach a quasi two dimensional (Q2D) treatment. It will be defined more precisely below.
Recently 10EL we demonstrated that quasi 2D solitons can be stabilized by a combination of FRM and a strong 1D optical lattice (1D OL), instead of a 1D harmonic trap 6Saito ; 8Abdul ; 9VPG . By a “strong lattice” we mean one in which the atoms in neighboring cells cannot interact. This issue has practical relevance, as a 1D OL can easily be created, illuminating the BEC by a pair of counterpropagating laser beams that form a periodic interference pattern 11OL . The lattice will be weak or strong depending on the intensity of the laser light. In fact, it is easier to realize a tight confinement configuration in an optical lattice than in an harmonic trap. Hence this environment may be more friendly for creating quasi 2D solitons.
In this Letter, we demonstrate that the combined OL-FRM stabilization of 3D solitons is possible even in a weak lattice, when atoms confined in different cells interact. By analyzing the stability charts in configuration space we discover two distinct regions where stable solutions exist. The first of these regions has its counterpart in the Q2D treatment. The other region appears when the frequency of modulation exceeds a critical value dependent on the confining potential. It is not present in the Q2D treatment; hence it corresponds to fully 3D solitons. In the limit of tight confinement the latter region moves to extremely high frequencies and the Q2D stability chart is recovered.
We describe our system by the GPE in physical units, including a time-dependent (FRM-controlled) scattering length $`a(\tau )`$ and an external potential $`\stackrel{~}{U}(𝐫,\tau )`$
$$i\mathrm{}\frac{\mathrm{\Psi }}{\tau }=\left[\frac{\mathrm{}^2}{2m}^2+\stackrel{~}{U}(𝐫,\tau )+\frac{4\pi a(\tau )\mathrm{}^2}{m}|\mathrm{\Psi }|^2\right]\mathrm{\Psi }.$$
(1)
Initially the BEC is in the ground state of a radial (2D) parabolic trap with frequency $`\stackrel{~}{\omega }_{}`$, supplemented, in the longitudinal direction, by “end caps” induced by transverse light sheets. The configuration is much like the one used to create soliton trains in a Li<sup>7</sup> condensate 2bright . A 1D lattice potential in the axial direction is adiabatically turned on from $`\stackrel{~}{\epsilon }=0`$ to $`\stackrel{~}{\epsilon }=\stackrel{~}{\epsilon }_\mathrm{f}`$, see Fig. 1. Thus, the full potential is
$`\stackrel{~}{U}(𝐫,\tau )`$ $`=`$ $`\stackrel{~}{\epsilon }(\tau )\left[1\mathrm{cos}(2\pi z/\lambda )\right]+`$ (2)
$`+`$ $`f(\tau )\left[(m/2)\stackrel{~}{\omega }_{}^2\varrho ^2+\stackrel{~}{U}_0(z)\right],`$
where $`\lambda `$ is the lattice spacing, $`\varrho `$ is the radial variable in the plane transverse to $`z`$, and the axial “end-cap” potential, $`\stackrel{~}{U}_0(z)`$ is approximated by a sufficiently deep one dimensional rectangular potential well. The width of the well determines the number of peaks in the finally established structure. The $`f(\tau )`$ is a switching-off function (see Fig. 1). We introduce dimensionless variables $`𝐱=(\pi /\lambda )𝐫`$, $`t=\tau \omega `$, $`\psi =\mathrm{\Psi }\sqrt{\lambda ^3/(N\pi ^3)}`$, where $`\omega =\pi ^2\mathrm{}/(m\lambda ^2)`$ and $`N`$ is the number of atoms. We also define $`g=4\pi ^2Na/\lambda `$, $`\epsilon =\stackrel{~}{\epsilon }/(\omega \mathrm{})`$, $`\omega _{}=\stackrel{~}{\omega }_{}/\omega `$, and $`U_0=\stackrel{~}{U}_0/(\omega \mathrm{})`$ to obtain
$$i\frac{\psi }{t}=\left[\frac{^2}{2}+U(𝐱,t)+g(t)|\psi |^2\right]\psi ,$$
(3)
and the potential in the form
$$U(𝐱,t)=\epsilon (t)\left[1\mathrm{cos}(2z)\right]+f(t)\left[\frac{1}{2}\omega _{}^2\varrho ^2+U_0(z)\right].$$
(4)
The nonlinear interaction coupling is described by
$$g(t)=g_0(t)+g_1(t)\mathrm{sin}(\mathrm{\Omega }t),$$
(5)
and the dimensionless modulation frequency by $`\mathrm{\Omega }=\stackrel{~}{\mathrm{\Omega }}/\omega `$. Initially $`g_1(0)=0`$ and $`g(0)=g_0(0)>0`$. At some moment $`t_1`$, we begin to linearly decrease $`g_0(t)`$. It vanishes at time $`t_2`$, and remains zero up to $`t_3`$, when we start to gradually switch on the rapid FRM modulation of a. In the interval $`[t_3,t_4]`$, $`g_0(t)`$ decreases linearly from zero to a negative $`g_{0\mathrm{f}}`$ and the amplitude of the modulation $`g_1(t)`$ increases from zero to $`g_{1\mathrm{f}}`$, see Fig. 1. Simultaneously, both the radial confinement and end-caps are gradually switched off. At times $`t>t_4`$, $`g(t)`$ oscillates with a constant amplitude $`g_{1\mathrm{f}}`$ around a negative average value $`g_{0\mathrm{f}}`$. Consequently, a soliton so created, if any, is supported by the combination of the 1D lattice and FRM.
Numerical experiments following the path outlined in Fig. 1 indicate that it is possible to create stable solitons 10EL (see inset to Fig.2). Before showing the results, we first resort to the variational approximation (VA) in order to predict conditions on the modulation frequency and the size of the negative average nonlinear coefficient $`g_{0\mathrm{f}}`$, necessary to support 3D solitons.
The VA can be applied to the description of BEC dynamics under diverse circumstances 6Saito ; 8Abdul ; 9VPG ; 10EL ; 12BBB ; Progress . Equation (3) is derived from the Lagrangian density
$`=i(\psi _t^{}\psi \psi _t^{}\psi )|\psi _\varrho |^2|\psi _z|^2g(t)|\psi |^42U|\psi |^2`$ (6)
We use VA for $`t>t_3`$ and choose a complex Gaussian ansatz for the solution for one lattice cell. The amplitude is $`A(t)`$, radial and axial widths are $`W(t)`$ and $`V(t)`$ respectively, and $`b(t)`$ and $`\beta (t)`$ are the corresponding chirps
$$\psi (𝐫,t)=Ae^{\left[\varrho ^2(1/2W^2+ib)z^2(1/2V^2+i\beta )+i\varphi \right]}.$$
(7)
The reduced Lagrangian can be found upon substituting (7) into (6) and integrating over space. By varying this reduced Lagrangian with respect to $`\varphi `$ we obtain the constant $`E=A^2W^2V=\pi ^{3/2}_{cell}|\psi |^2𝑑𝐫=\pi ^{3/2}/n`$, where the integral extends over one cell of the lattice and $`n`$ is the number of occupied lattice cells. Notice that the total number of atoms is included in the definition of the nonlinear coupling $`g(t)`$ and the total wavefunction is normalized to unity. When the other four variational equations are derived, we can deduce two dynamical equations for the widths:
$`\ddot{W}`$ $`=`$ $`{\displaystyle \frac{1}{W^3}}f(t)\omega _{}^2W+{\displaystyle \frac{Eg(t)}{\sqrt{8}W^3V}},`$ (8)
$`\ddot{V}`$ $`=`$ $`{\displaystyle \frac{1}{V^3}}4\epsilon _\mathrm{f}V\mathrm{exp}\left(V^2\right)+{\displaystyle \frac{Eg(t)}{\sqrt{8}W^2V^2}}.`$ (9)
These equations describe the dynamics of a single peak, and with one term slightly altered can be applied to the problem of a BEC confined in a 1D harmonic trap 6Saito ; 8Abdul ; 9VPG .
In the corresponding Q2D treatment we drop the $`z`$ dimension in Eq. (3). It is assumed that in this direction the profile of the wavefunction is fixed and reproduces the ground state $`\psi _0`$ of the single lattice cell (Wannier function) or harmonic potential, as in Refs. 6Saito ; 9VPG . The reduced potential in 2D GP will take the form $`U(\varrho ,t)=f(t)(1/2)\omega _{}^2\varrho ^2`$. In VA we take $`VV_0`$ from $`4\epsilon _\mathrm{f}V_0^4\mathrm{exp}\left(V_0^2\right)=1`$, and only solve equation (8). In numerical simulations we rescale the nonlinear coupling coefficient $`g_{2\mathrm{D}}=g\times (|\psi _0|^2\psi _0𝑑z)/(\psi _0𝑑z)`$.
We simulated both the full GPE, Eq. (3), using an axisymmetric code (for 3D), a Cartesian code (for 2D), and the variational equations for comparison. Numerical simulations followed the path outlined in Fig. 1. The parameters used in the simulations would correspond, for <sup>85</sup>Rb atoms and an OL period of $`\lambda =3\mu `$m, to an initial radial-confinement frequency of $`\omega _{}=2\pi \times 39`$ Hz, a FRM frequency of $`\mathrm{\Omega }=2\pi \times 2.87`$ kHz. A good candidate is the Feshbach resonance at 155 G for <sup>85</sup>Rb. Since in this case the $`a=0`$ point is far from the resonance, both two and three body losses are negligible 155G . A lattice depth of $`\epsilon _\mathrm{f}=10.25`$ recoil energies, and an effective nonlinear coefficient of $`(N/n)a=\pm \mathrm{\hspace{0.17em}10}^5`$m are required. Here $`N/n`$ is the number of atoms per lattice cell, with a total number of atoms in the range of $`10^410^6`$. The respective values of the normalized parameters are given in the figure captions.
Examples of the numerical results, which are generic, are displayed in Fig. 2, which shows the evolution of the central-peak’s amplitude versus time in dimensionless units, defined by Eq. (3). After an initial transient, a stable structure is established, featuring breathing without any systematic decay. The figure reveals the influence of neighboring solitons on the amplitude of the central peak. The thin curve corresponds to the evolution of the full multipeak structure, shown in the inset for a fixed moment of time. To obtain the thick curve we repeated the above calculations up to the time $`t=7500`$ and then removed all but the central peak. In this case the amplitude gradually decreased. The interaction between neighboring solitons can be explained by the difference between oscillation periods from cell to cell due to small deviations in the numbers of atoms. It is essentially the same mechanism as in a Josephson junction. When the lattice potential is weak, interaction between neighboring solitons is possible.
In Figs. 3 and 4 we have collected results of a systematic scan of parameter space based on GPE simulations and compared them with the predictions of the VA (a similar analysis can be performed if we replace 1D OL with a 1D harmonic trap - the conclusion does not depend on the form of confinement in the $`z`$ direction). By stability we mean shape preservation during one run. When this was not the case, after a short period of time we observed clear collapse or spreadout in transverse directions.
The agreement between VA and direct simulations is very good. As seen from Fig. 3, in our fully 3D treatment we found two islands of stability. Note the similarity of the lower region to that of Fig. 4, which portrays the results of the Q2D treatment. This region corresponds to Q2D solitons. On the other hand, the upper region contains fully 3D solitons. It appears when the frequency of modulation exceeds the lowest excitation frequency of the confining potential, $`\mathrm{\Omega }_0=26.76`$. If the strength of the lattice $`\epsilon _\mathrm{f}`$ is increased, this region moves towards higher frequencies, the Q2D region expands upwards, and becomes more and more like in Fig. 4. This will be demonstrated in a fuller version of this work future .
As we saw, new region of stability appears in the 3D treatment as compared to the Q2D treatment. We are more used to the effect of adding a new dimension simply shrinking or abolishing the basin of stability of solitons or waves. For example, water waves are unstable with respect to perturbations along their direction of propagation only when the depth exceeds a critical value TB . When, however, two dimensional perturbations are allowed, there will always be an unstable angle regardless of the depth WD . Another example is that of 1D solitons of the Nonlinear Schrödinger equation (NLS) with constant coefficients. They are stable in 1D, but unstable in 2D or 3D. This is also true for some NLS waves EI .
In the problem treated in this Letter this is not the case. For much of the quasi two dimensional stability chart adding a degree of freedom stabilizes the soliton solution. The key to this dichotomy seems to be the presence of a periodic modulation, absent in the above mentioned classical examples. This can be illustrated by a simple case involving oscillators. Take as the one dimensional version a forced oscillator problem:
$$\ddot{x}+\omega _0^2x=y\mathrm{cos}(\omega _0t).$$
(10)
If $`y`$ is fixed, the solution has a secular component $`x=yt/(2\omega _0)\mathrm{sin}(\omega _0t)+F(t)`$, where $`F(t)`$ is a periodic function, and so the amplitude will grow as $`t`$. If however we allow a second degree of freedom, such that $`y`$ also oscillates (for instance $`\ddot{y}+ϵ^2y=0`$) the solution stabilizes, unless $`ϵ=\pm 2\omega _0`$. In general this can be the case when there are periodic modulations. This fact, obvious in oscillator theory, is perhaps less well known in the soliton context.
The main result of this letter is the possibility of creating fully 3D breather solitons in a BEC confined by a 1D optical lattice potential, corresponding to the upper region in Fig. 3. The stable patterns may feature a multi-cell structure, which in the case studied here forms a set of weakly interacting fundamental solitons. The scheme proposed in this Letter is based on a combination of FRM and a 1D optical lattice, and could be implemented in an experiment, as outlined here. This would open the way to the creation of robust 3D solitons (breathers) in BECs. A similar idea can apply in the field of nonlinear optics. The Feshbach resonance could be replaced by a nonlinear periodic structure, for example ferromagnetic domains as used in contemporary quasi phase matching Arie , or by optically induced photorefractive lattices Wiesio .
M.M. acknowledges support from the KBN grant 2P03 B4325, M.T. was supported by the Polish Ministry of Scientific Research and Information Technology under grant PBZ MIN-008/P03/2003 and E.I. from that of grant 2P03B09722. The work of B.A.M. was partially supported by the Israel Science Foundation through grant No. 8006/03.
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# Multiwavelength studies of WR 21a and its surroundings
## 1 Introduction
Massive O and Wolf-Rayet (WR) stars lose large amounts of mass through dense, energetic winds. These winds can interact with the surrounding interstellar medium (ISM) to create cavities in the H i distribution (e.g. Benaglia & Cappa 1999) or produce strong shock fronts that, in principle, can accelerate charged particles up to relativistic energies (e.g. Völk & Forman 1982). In fact, non-thermal radio emission has been detected in a number of early-type stars, indicating the presence of populations of locally accelerated electrons (e.g. Chapman et al. 1999, Benaglia et al. 2001b, Benaglia & Koribalski 2004).
Particle acceleration mediated by strong shocks could be taking place in a number of different regions of a massive stellar system such as the outer boundary of the interaction between the wind and the ISM (e.g. Cassé & Paul 1980), the base of the wind, where line-driven instabilities are thought to drive strong shocks (e.g. White 1985), or, in massive binary systems, the colliding-wind region (CWR, e.g. Eichler & Usov 1993). In addition to the expected non-thermal radio emission from relativistic electrons, high-energy radiation could be produced through, for example, inverse-Compton scattering of UV stellar photons or hadronic interactions from shock-accelerated ions and ambient material (e.g. Pollock 1987, White & Chen 1992, Romero et al. 1999, Benaglia et al. 2001a, Mücke & Pohl 2002, Benaglia & Romero 2003).
Only a small number of non-thermal radio WR binaries have been identified so far (see for instance Dougherty et al. 2003) and, even then, the connection between the radio and X-ray properties of these systems is far from clear. The archetype system WR 140 is bright in both regimes with a well-defined radio light curve that shares the 7.94-year period of the optical radial-velocity orbit. The X-ray light curve is not so well determined but shows strong variability consistent with the same period. WR 140, however, is the exception. The radio brightest star, WR 146, is a relatively faint X-ray source whereas WR 147 seems to be intermediate in both regimes. Sharing some characteristics with these systems, the star WR 21a \[($`\alpha ,\delta )_{\mathrm{J2000}}=10^\mathrm{h}25^\mathrm{m}56^\mathrm{s}.49,57^{}48^{}44.4\mathrm{"}`$, $`(l,b)=284^{}.52,0^{}.24`$\], is particularly interesting for various reasons. In the past, it coincided within the fairly large positional uncertainties with the COS B source 2CG 284-00 (Goldwurm et al. 1987). More recently, the unidentified EGRET gamma-ray source 3EG J1027–5817 was detected southwest of WR 21a (Hartman et al. 1999), being the star 30 arcmins apart from the center of the gamma-ray source. The 95% probability radius of the EGRET source is about 0.3.
Caraveo et al. (1989) established that WR 21a is associated with the X-ray source 1E 1024.0–5732, and reported X-ray pulses. Its X-ray emission could be explained if the star is part of a binary system with either a compact companion forming a high-mass X-ray binary (Caraveo et al. 1989) or another massive, early-type star, with a CWR (Reig 1999). According to its stellar spectrum, it was early classified as an O5 star (Caraveo et al. 1989), and as a likely binary (WN5-6+O3f), at a maximum distance of 3 kpc (Reig 1999).
With the current investigation we aim at determining the radio properties of the star and its surroundings. Specifically, using interferometric observations at 3 and 6 cm, we looked for non-thermal radiation that might be interpreted as evidence for sites where electrons are being accelerated in colliding winds or terminal shocks. By means of low-resolution 21cm-line observations, we have also studied the distribution of neutral hydrogen around the star. The detection of non-thermal emission coincident with the stellar position would help in the identification of the system components. Analysis of all available X-ray data of WR 21a reveals a variability history that sheds fresh light on the nature of the binary. Finally, we discuss the possibility of a physical link between the stellar system and the nearby EGRET source, which can help to unveil the nature of the latter.
The contents of the paper are as follows: Section 2 reviews the main sources detected in this region of the sky that are relevant for the present study. Section 3 describes the observations carried out and the data reduction; Section 4 outlines the new observational results, whereas Section 5 presents the corresponding analysis. In Section 6 we discuss the X-ray data. Section 7 contains a comment on the possibility of a physical association between the star and the EGRET source 3EG J1027–5817. Section 8 closes with the summary.
## 2 WR 21a and its surroundings
The region of the Galactic plane towards $`(l,b)=(284^{},0^{})`$ has been widely observed, from radio to gamma rays. The optical images of Rodgers et al. (1960) revealed various H ii complexes, of which RCW 49 is the most prominent. The COS B satellite discovered the unidentified $`\gamma `$-ray source 2CG 284-00 (Bignami & Hermsen 1983, Caraveo 1983). The zone was then investigated by means of Einstein X-ray measurements in which were detected emission from RCW 49 and a point source named 1E 1024.0–5732 (Goldwurm et al. 1987). The point source was tentatively linked to WR 21a (Hertz & Grindlay 1984), an emission-line star with $`m_V=12.8`$ (Wackerling 1969), also called Th35$``$42 and Wack 2134.
Caraveo et al. (1989) gathered enough evidence definitely to associate the X-ray source with the star. They took optical stellar spectra with the 3.6-m telescope at La Silla, classified the star as O5, set an upper limit for the stellar distance of 3 kpc, and studied the Einstein X-ray emission, reporting an X-ray periodicity of $`60`$ ms. They explained the X-ray behaviour as a binary with a compact companion, likely an accreting neutron star (NS). Dieters et al. (1990) looked for optical pulsations with the Tasmania 1-m telescope. They set an upper limit around 19.7 mag for any optical pulsation.
Mereghetti et al. (1994) observed the region with the ROSAT PSPC instrument, finding no pulsations, and obtained a further spectrum with the CTIO 1.6m telescope, revising the spectral classification to WN6. Based on the small equivalent widths of the emission lines, they suggested that the object could be a binary with an O star, rather than a compact companion, in which case the X-rays could come from colliding winds, in common with other Wolf-Rayet binary systems.
In 1999, Reig presented RXTE data and a new spectrum taken with the 1.9m telescope at SAAO. He stated that the lack of pulsations and the relatively soft and low X-ray emission seem to exclude the presence of a NS as responsible for the observed X-rays. Claiming that the spectrum shows features of both WR and O stars, he suggested that the system is formed by a WN6 and a possibly supergiant O3 companion, favoring the hypothesis of the colliding-wind binary (CWB). Roberts et al. (2001) observed the region towards WR 21a with ASCA, and interpreted the hard X-ray emission detected as produced by shocks from colliding winds in the stellar system. Very recently, Niemela et al.’s (2005) optical radial measurements have finally confirmed that the system is formed by two massive stars.
The distance $`d`$ to WR 21a is not well established. We shall adopt $`d=3`$ kpc throughout this paper, a value that seems to be consistent, as we will see, with all current observations.
Table 1 lists the adopted parameters that make up the Wack 2134/WR 21a binary system. Very few of the properties have so far been measured, so we have to assume parameters from similar stars or theoretical predictions. The stellar luminosity, effective temperature, and stellar mass of the WR component were estimated as an average of the same variables given by Hamann et al. (1995) for WN6 stars. The WR mean molecular weight of the ions $`\mu `$ was assumed the same as in Leitherer et al. (1997) for a WN6 star. The WR wind terminal velocity is taken as a mean value between data listed by Hamman et al. (1995) and by Prinja et al. (1990) for WN6 stars. The WR predicted mass loss rate ($`\dot{M}`$) is taken as the lowest value tabulated for WN6 stars by Nugis & Lamers (2000), in their compilation of observable $`\dot{M}`$, and it can be considered a lower limit. The O3 (I) mass was taken as the spectroscopic mass from the tables of Vacca et al. (1996), and its terminal velocity from the averaged values listed by Prinja et al. (1990). We assumed $`\mu =1.5`$ for the O star because of its evolved stage. A predicted mass loss rate was estimated using the recipe derived by Vink et al. (2000) (astro.ic.ac.uk/$``$jvink/). The expected (WR+Of) mass-loss rate would imply a 6-cm flux density of 0.24 mJy at 3 kpc, if thermal emission from ionized winds in both stars is assumed.
Close to WR 21a there are two interesting sources: the H ii region RCW 49, and the gamma-ray source 3EG J1027–5817. They are discussed below.
### 2.1 The H ii region RCW 49
RCW 49 is a southern H ii region, located at $`(l,b)=(284.3^{},0.3^{})`$, and extended over an area of $`90^{}\times 70^{}`$. Values for its distance range between 2.3 and 7.9 kpc (Manchester et al. 1970, Moffat et al. 1991, Brand & Blitz 1993, etc). H i spectra towards RCW 49 taken by Goss et al. (1972) show prominent H i absorption from about –20 to +5 km s<sup>-1</sup> (see also Figs. 2 and 3). Whiteoak & Uchida (1997) have imaged RCW 49 at radio continuum with MOST at 0.843 GHz, and the central region using ATCA at 1.38 and 2.38 GHz, attaining an angular resolution of 7 arcsec. They found two shells, and ascribe the formation of the northern one to the Westerlund 2 cluster containing the binary star WR 20a, and the southern one to the star WR 20b (see their Fig. 2d). Recent infrared images of RCW 49 obtained with the Spitzer Space Telescope (Churchwell et al. 2004) show the intricate filamentary structure of the nebula in the inner 5 arcmin shaped by stellar winds and radiation.
### 2.2 The gamma-ray source 3EG J1027–5817
After the analysis of the EGRET data, Hartman et al. (1999) were the first to point at the proximity between the gamma-ray source 3EG J1027–5817 and the X-ray source 1E 1024.0–5732, associated with WR 21a (see Figure 3). The averaged measured flux at $`E>100`$ MeV is 65.9$`\pm 0.70\times 10^8`$ photons cm<sup>-2</sup> s<sup>-1</sup>. The source is constant within errors on timescales of months (variability index $`I`$ = 1.6, see Torres et al. 2001) and with a photon spectral index $`\mathrm{\Gamma }=1.94\pm 0.09`$ ($`dN/dEE^\mathrm{\Gamma }`$). The more recent variability analysis by Nolan et al. (2003), who calculated a likelihood function for the flux of each source in each observation, also suggests that this source is not variable on short, monthly timescales.
## 3 Observations and data reduction
In order to search for non-thermal radio emission in the direction of WR 21a we have carried out interferometric radio continuum observations at high angular resolution ($`1^{\prime \prime }2^{\prime \prime }`$). These were complemented with single-dish H i-line observations to map the distribution of neutral material in the neighbourhood and study its kinematic behaviour.
### 3.1 Radio continuum observations
Radio continuum data were obtained in September 2001 with the Australia Telescope Compact Array in the 6D array, observing simultaneously at 3 and 6 cm, or 8.64 and 4.8 GHz, respectively. The total bandwidth used was 128 MHz. The primary calibrator was PKS 1934-638, with flux densities of 5.83 and 2.84 Jy, at 4.8 and 8.64 GHz. WR 21a was tracked 12 h –full synthesis– to gain maximum $`uv`$ coverage, interleaving with observations of the phase calibrator 1039–47. The theoretical r.m.s. noise after $``$ 9 h on source is 0.03 mJy at both frequencies, taking all baselines into account.
The data were reduced and analyzed with miriad routines. The images built using “robust” weighting showed the best signal to noise ratio, and minimized sidelobes. The diffuse emission from extended sources was removed by taking out the visibilities corresponding to the shortest baselines. The resulting beams were $`0.83\mathrm{"}\times 0.70\mathrm{"}`$ at 3 cm, and $`1.73\mathrm{"}\times 1.49\mathrm{"}`$ at 6 cm. The r.m.s. noise of the final maps is 0.1 mJy beam<sup>-1</sup> at 3 cm and 0.06 mJy beam<sup>-1</sup> at 6 cm. The difference between the theoretical and the observed r.m.s. noise is mainly due to flagging of bad data – specially at 3 cm – as well as short baselines contributing with confusing emission from extended sources both in and outside the main beam.
The observations were set to optimize the detection of point-like features. Maps at two frequencies would allow the determination of spectral indices.
### 3.2 H i-line observations
The 21cm-line data were obtained with a 30m-single dish radiotelescope at the Instituto Argentino de Radioastronomía (IAR, Villa Elisa, Argentina) in June 2000. The observations were done in total power mode, covering a total field of $`5^{}\times 5^{}`$, in Galactic coordinates with a cell size of $`15^{}\times 15^{}`$. The HPBW at 1420 MHz is 30’. The receiver’s system temperature was $``$ 35 K. The velocity coverage was (–450, +450) km s<sup>-1</sup>; the 1008 channel autocorrelator allows a maximum velocity resolution of 1.05 km s<sup>-1</sup>. The r.m.s. noise of the brightness-temperature ($`T_\mathrm{B}`$) of a single spectral point is $``$ 0.1 K. The $`T_\mathrm{B}`$ scale was calibrated with the standard region S9 (Morras & Cappa 1995). A series of ($`l`$, $`b`$)-$`T_\mathrm{B}`$ maps were built every 1.05 km s<sup>-1</sup> to proceed with the analytical stage.
## 4 Results
### 4.1 ATCA radio continuum data
The ATCA images at 3 and 6 cm are shown in Figure 1. A point source positionally coincident with WR 21a is visible at 6 cm, at $`S/N`$ $`>`$ 4. A flux density $`S_{6\mathrm{c}\mathrm{m}}=0.25`$ mJy was derived using imfit after a gaussian fit (see Table 2). The r.m.s. noise in the 3 cm image is 0.1 mJy beam<sup>-1</sup>. No radio source is detected at the stellar position above 3 r.m.s.. This non-detection imposes an upper limit for the spectral index of $`\alpha <0.3`$ ($`S\nu ^\alpha `$). The deviation from purely thermal emission, characterized by $`\alpha =0.60.8`$, indicates the presence of a non-thermal contribution.
In a similar way as in Benaglia & Koribalski (2004), we derive a mass loss rate of WR 21a, from the flux density at 6 cm. In a first approximation, the mean number of electrons per ion ($`\gamma `$), and the rms ionic charge ($`Z`$) were taken equal to unity. The gas temperature is computed from the stellar effective temperature (see Table 1) as $`0.4T_{\mathrm{eff}}`$. The Gaunt factor results in 5.6. At the adopted distance of 3 kpc, this would imply a radio-derived mass loss rate for the WR star of $`\dot{M}=f\times \mathrm{\hspace{0.17em}4.8}\times \mathrm{\hspace{0.17em}10}^5`$ M yr<sup>-1</sup>, where $`f`$ is the fraction of thermal to total radio emission. From the relation between the flux density values measured with ATCA at 3 and 6 cm we know that at 6 cm there is a non-thermal contribution to the emission. Thus, in deriving a mass loss rate from the flux density at 6 cm, the result with $`f=1`$ is an upper limit.
A second source, called S2, is detected 10” away from WR 21a, at both frequencies. The measured fluxes are 0.55 mJy at 3 cm, and 0.36 mJy at 6 cm. The corresponding spectral index of $`\alpha _{\mathrm{S2}}=+0.72`$ points to an ultracompact H ii region or the thermal wind of another, as yet unidentified, early-type star on the field of WR 21a. Table 2 lists the position and flux density of the two sources detected.
### 4.2 Large-scale H i-line data
Studies of neutral hydrogen were performed to look for evidences of gas features such as shells, filaments, bubbles, etc., that could be physically related to the object under the present study. The H i gas kinematics in the line-of-sight to a Galactic source can be used to constrain its distance (e.g. Koribalski et al. 1995) using the Galactic rotation curve (e.g. Fich et al. 1989) and the Galactic velocity field (Brand & Blitz 1993). Here we searched for the signatures of an interstellar H i bubble, created by the action of the stellar winds of WR 21a. Because of the large angular size of the IAR telescope beam (HPBW = 30 arcmin; see Figs. 2 and 3) and the proximity of WR 21a to the extended H ii region RCW 49, which has a total radio continuum flux of $``$210 Jy at 843 MHz (Whiteoak & Uchida 1997), H i spectra in this direction are completely dominated by H i absorption against RCW 49 (Goss et al. 1972; McClure-Griffiths et al. 2001). We also investigated the high-resolution (130 arcsec) H i data cubes from the Southern Galactic Plane Survey (SGPS; McClure-Griffiths et al. 2001) in the region of RCW 49 and WR 21a. Unfortunately, the region around RCW 49, including the H i line emission at the position of WR 21a, suffers from artifacts caused by the strong radio continuum emission from RCW49 and potential H i structures associated with WR 21a cannot be distinguished.
Figure 2 displays the H i brightness temperature maps each 4 km s<sup>-1</sup> built from the IAR data.
If we were able to see an H i bubble around WR 21a, its size would give us some information about the energetics of the stellar wind and its velocity would give us an estimate of its kinematic distance using the Galactic rotation curve and velocity field. These issues are explored in the following section.
## 5 Analysis of the line features
The derivation of gas kinematic distances is difficult for this particular region of the Galactic plane. The line of sight goes tangential to the Carina arm at this Galactic longitude, and velocity crowding becomes very important. Brand & Blitz (1993) showed that the measured velocities deviate strongly from Galactic rotation: gas in the velocity range from –20 to –10 km s<sup>-1</sup> can be located at distances between 2 and 6 kpc for $`l285^{}`$, where both WR 21a and RCW 49 are likely to be located.
By means of CO observations, Grabelsky et al. (1987) studied molecular gas associated with the Carina arm. They interpreted the velocity-longitude behaviour of the gas in terms of material at different heliocentric distances (see their Fig. 5), separating local clouds from Carina arm gas. According to their results it is possible to determine that toward $`l284^{}`$ gas showing velocities around $``$15 km s<sup>-1</sup> belongs to the Carina arm and is located at about 3 kpc, i.e., the distance of WR 21a. In Fig. 2 it can be appreciated that the gas distribution changes at about $`v14`$ to –12 km s<sup>-1</sup>. We are going to focus on gas with $`\mathrm{\Delta }v`$ = –20 to –12 km s<sup>-1</sup> because, according to the CO results, its kinematical distance is compatible with that of the target star. The H i column density at the mentioned velocity interval is presented in Fig. 3. From the plot of Grabelsky et al. (1987), CO gas with velocities between –20 and –10 km s<sup>-1</sup> is located between 2 and $`3.5`$ kpc. This fact helps to constrain an approximate error in the H i distance of $``$ 1 kpc. At larger velocities, Grabelsky et al. claimed that gas related to RCW 49 shows velocities of –5 km s<sup>-1</sup>, and placed it at 4 kpc. A distance of $``$5 kpc can be derived, for Carina gas with velocities near 0 km s<sup>-1</sup>. If the H i follows the motions proposed by Grabelsky et al., it is reasonable to suggest that gas with a velocity of –14 km s<sup>-1</sup> lies at 3 kpc.
Figure 3 shows the presence of gas at a distance compatible with that of the stellar system. Due to the strong continuum source RCW 49, part of the HI gas is not emitting but absorbing. At a distance of 3 kpc, the visible neutral gas coincident with the position of the EGRET source would sum up about 1500 $`M_{}`$, which can be considered as a lower limit for the masses of clouds at the distance we are interested on.
By means of Antenna I at IAR (HPBW = 12.4) we also measured the H125$`\alpha `$ (3326.9880 MHz) radio recombination line (RRL) towards $`(l,b)=(284.32^{},0.34^{})`$. Since RRLs are typically produced by H ii regions, the detected emission (see Fig. 4) is likely to be mostly from RCW 49. We measure a center velocity of $`v=+0.6\pm 1.4`$ km s<sup>-1</sup>, similar to those found at other RRLs in RCW 49 (e.g. Caswell & Haynes 1987).
The kinetic energy needed to form a typical bubble around WR 21a can be computed as $`E_\mathrm{k}=0.5M_{\mathrm{sh}}v_{\mathrm{exp}}^2=10^{49}`$ erg if we assume a shell mass of $`M_{\mathrm{sh}}`$ = 10000 M and an expansion velocity of $`v_{\mathrm{exp}}`$ = 10 km s<sup>-1</sup>, which are typical parameters of neutral shells detected around massive Of stars (Cappa & Benaglia 1998, Benaglia & Cappa 1999).
The stellar wind luminosity can be expressed as $`L_\mathrm{w}=0.5\dot{M}v_{\mathrm{}}^2`$. Using the values of $`\dot{M}`$ and $`v_{\mathrm{}}`$ given in Table 1, we find $`L_{\mathrm{w},\mathrm{WR}}=3.8\times 10^{37}`$ erg s<sup>-1</sup>, and $`L_{\mathrm{w},\mathrm{O}}=4.4\times 10^{37}`$ erg s<sup>-1</sup>.
Finally, a wind mechanical energy $`E_\mathrm{w}=L_\mathrm{w}t=2.6\times 10^{51}`$ erg per Myr is obtained if both stars are considered, and $`E_\mathrm{w}=1.2\times 10^{51}`$ erg per Myr for only WN6. It can be seen that the energy deposited by the wind is much larger than the energy needed to create a typical bubble.
## 6 The X-ray history of WR 21a
Among the Wolf-Rayet stars, WR 21a is especially prominent in X-rays: despite its modest optical magnitude, it is one of the five or six brightest sources in both apparent and absolute terms. Since the discovery of 1E 1024.0–5732 in 1979 with the Einstein Observatory by Goldwurm et al. (1987), through its identification with a Wolf-Rayet star by Mereghetti et al. (1994), WR 21a has been observed on ten separate occasions in X-rays as shown in Table 3. All of these data are available through the HEASARC. As we mentioned before, Reig (1999) compared his 1997-RXTE data with earlier Einstein-IPC and ROSAT-PSPC measurements to show that the X-ray luminosity had apparently been steadily rising by about a factor of five in the eighteen years since its discovery and argued that this showed WR 21a is a long-period CWB of the type exemplified by WR 140 (Williams et al. 1990).
Though in general terms the latter is probably correct, the details are more complicated. First, the X-ray model which Reig used had too high an absorbing column density ($`N_\mathrm{H}=2.9\times 10^{22}\mathrm{cm}^2`$) to be consistent with the the soft X-ray data; and second, the faintness of the pair of ROSAT-HRI measurements made in 1994 (Belloni & Mereghetti 1994) contradicts the apparent inexorable rise in luminosity since 1979. Since the RXTE observation, two archived sets of ASCA data have also become available that are useful for bridging the soft X-ray images and the hard X-ray collimator data. The first ASCA observation took place two months after Reig’s RXTE pointing though WR 21a was not the main objective and thus appeared near the edge of the GIS field-of-view, one of the two ASCA instruments that provided imaging X-ray spectroscopy. As a consequence there are no data from the SIS, the other instrument. On the other hand, data are available from the full set of ASCA instruments from an observation performed nearly a year later at the end of 1998.
We have obtained and analysed with XSPEC v11.2 all the archived spectra and associated response matrices available from the HEASARC. Despite the obvious changes in overall luminosity, we could find no evidence of any changes in the shape of the spectrum which, within the limited energy resolution available, seems to be consistent with the relatively hot few keV thermal plasma observed from the Wolf-Rayet binaries exemplified by WR 140 (Pollock et al. 2005), in contrast with the cooler temperatures more typical of the intrinsic emission of single stars (see for example the studies on the presumably single WN stars WR 1 \[Ignace et al. 2003\], WR 6 \[Skinner et al. 2002b\], and WR 110 \[Skinner et al. 2002a\]). The spectrum was modeled as an absorbed Bremsstrahlung continuum with additional emission lines of Si, Mg and Ne. The best-fit values of column density and temperature were $`N_\mathrm{H}=7.0\pm 0.6\times 10^{21}\mathrm{cm}^2`$, about a factor of 4 lower than Reig’s (1999) value, and $`kT=3.3\pm 0.2`$ keV. This empirical approach has the natural advantage of reproducing the range of ionization species in the spectrum, notably the simultaneous presence of the lines of SiXIII and SiXIV. The alternative fits given by XSPEC’s plasma models were slightly worse but gave completely consistent best-fit values of column density and temperature.
The luminosities reported in Table 3 are for a joint fit to all the available spectra with only the luminosities free to vary between observations. The resulting lightcurve is shown in Figure 5. The X-ray variability is apparently irregular though the measurements are spaced at such large intervals with respect to the newly-discovered period of weeks (Niemela et al. 2005) that it will only be possible to tell if it is related to the binary orbit once a precise orbit is available. Some care is also required with the brightest point that came from RXTE, the only instrument here with no imaging capabilities. Though Reig made a correction for the emission from other nearby sources by adding an extra spectral component with fixed parameters, this is quite uncertain, because of the extensive diffuse emission and the strength of the point sources enumerated by Belloni & Mereghetti (1994), of which 1E1022.2-5730 is of similar spectral shape to WR 21a.
## 7 Connection with the EGRET source?
The nature of the EGRET unidentified gamma-ray sources (Hartman et al. 1999) has become one of the most intriguing questions in astrophysics (e.g. Romero 2001). In the particular case of 3EG J1027-5817, the nearby X-ray source associated with WR 21a is indicated in the Third EGRET catalog as a potential counterpart.
The presence of a non-thermal contribution to the radio spectrum of WR 21a implies the existence of a population of relativistic particles in the source, which are probably accelerated at the CWR by diffusive shock acceleration (e.g. Bell 1978a,b) and cool by synchrotron emission in the local magnetic field. Electrons can also cool in such an environment through inverse-Compton interactions with stellar UV photons, producing non-thermal X-rays and gamma-rays (e.g. Pollock 1987, Benaglia & Romero 2003). However, the facts that the X-ray emission from the system can be correctly modeled as thermal Bremsstrahlung and that the source has not been directly detected at gamma-rays suggest that the magnetic energy density in the CWR should largely exceed the photon energy density, implying significantly shorter cooling timescales for the synchrotron mechanism.
The same mechanism that accelerates the electrons should also operate on the ions. Synchrotron losses are not relevant for protons in the environment of the CWR. The maximum energy they can achieve will be determined by the photo-pion losses in the UV stellar field and by the size constraint imposed by the limited space available for the acceleration process. In the present case, where the CWR is not resolved and the geometry of the system remains unknown, we cannot calculate the high-energy cutoff for the non-thermal proton distribution. A value between 10 and 100 GeV seems not unreasonable (see Benaglia & Romero 2003).
It could be the case that some of these protons diffuse up to a nearby cloud where they might be trapped in the magnetic field, which is expected to be higher than the average value in the ISM (Crutcher 1999). Then they will interact there with the local material producing gamma-rays through $`p+pp+p+\pi ^0`$ interactions and the subsequent $`\pi ^0\gamma +\gamma `$ decays. The situation of a passive cloud irradiated by cosmic rays from some nearby accelerator has been discussed in detail by Black & Fazio (1973) and by Aharonian & Atoyan (1996). In the present case our ignorance on several basic parameters prevents accurate calculations, but there remains the possibility that a part of the flux detected from 3EG J1027-5817 could be originated in relativistic particles accelerated in the colliding wind region of WR 21a. Whether this is or not the case could be established through future observations of the gamma-ray source by instruments like AGILE and GLAST which could report the source position with higher accuracy.
## 8 Summary and Conclusions
We have detected radio emission from WR 21a at 4.8 GHz. The intensity of the source is $``$0.25 mJy. The non-detection at 8.64 GHz implies a spectral index of $`\alpha <0.3`$ ($`S_\nu \nu ^\alpha `$), which significantly departs from a typical Bremsstrahlung spectrum. Combined thermal/non-thermal spectra are usually found in colliding-wind binaries. We suggest that this is also the case here, since the latest spectral determinations show that WR 21a is a system formed by WN6 and an early O companion (Reig 1999, Niemela et al. 2005). An upper limit for the system mass loss rate of $`\dot{M}=4.8\times 10^5`$ M yr<sup>-1</sup> is derived from the 4.8-GHz radio flux density.
We have reanalyzed all X-ray observations of WR 21a in order to determine its time history in this waveband. Our results indicate, contrary to previous thought, that the X-ray flux has not been monotonically increasing since 1979 though the coverage is far too sparse to constrain the variability timescale. The X-ray spectrum is consistent with a few keV thermal plasma, with no obvious non-thermal contribution.
Locally accelerated relativistic electrons in the CWR probably mainly cool by synchrotron emission at radio frequencies, with small inverse-Compton losses, of which there was no X-ray evidence. If protons are also accelerated at the colliding-wind shocks, then they might diffuse through the ISM up to nearby clouds, where they might interact with an enhanced H i density to produce gamma-rays from $`\pi ^0`$-decays.
Future observations with both X-ray and gamma-ray instruments like CHANDRA, AGILE and GLAST can shed additional light on the nature of the high-energy emission in this interesting region. Detailed knowledge, on the other hand, of the orbital parameters of the system WR 21a will allow more sophisticated models to be built of the radiative processes taking place in the colliding-wind region.
###### Acknowledgements.
We thank Virpi Niemela for discussions on this source. This research has been supported by the Argentine agency ANPCyT through grant PICT 03-13291.
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# Reflections in abstract Coxeter groups
## 1. Introduction
The dihedral group of order 12 can be considered as Coxeter group of type $`I_2(6)`$ or as Coxeter group of type $`A_1\times I_2(3)`$. This example shows that, in general, the set of reflections in a Coxeter system is not determined by the abstract group $`W`$ alone, but does indeed depend on the choice of the Coxeter generating set $`R`$. However there are a lot of examples of Coxeter systems $`(W,R)`$ where the abstract group does determine the set of reflections or even the set $`R`$ up to $`W`$-conjugacy. The main motivation for the present paper is to show that the latter holds for infinite Coxeter groups having a finite, irreducible and 2-spherical Coxeter generating set, which is our Theorem 1 below.
In view of the main result of it suffices to show that these Coxeter groups determine the set of reflections. In order to achieve this goal we provide a handy criterion for an involution in an abstract Coxeter group $`W`$ to be a reflection with respect to any Coxeter generating set of $`W`$. Our principal observation is the following. Let $`(W,R)`$ be a Coxeter system and let $`wW`$ be an involution. If $`wR^W`$, then the centralizer of $`w`$ in $`W`$ contains a finite normal subgroup properly containing $`w`$. This is an immediate consequence of Richardson’s result in . Thus, if $`wW`$ is an involution having the property that $`w`$ is a maximal finite normal subgroup of its centralizer in $`W`$, then $`w`$ is a reflection with respect to any Coxeter generating set of $`W`$.
It turns out that it is more convenient to work with the finite continuation of an involution rather than to consider finite normal subgroups of its centralizer. The finite continuation of a finite order element $`w`$ in a Coxeter group is defined to be the intersection of all maximal finite subgroups containing it; we write $`\mathrm{FC}(w)`$ for the finite continuation of $`w`$. In this paper we restrict our attention to finitely generated Coxeter groups. For these it is a consequence of a result of Tits that every element of finite order is contained in some maximal finite subgroup; so $`\mathrm{FC}(w)`$ is a finite subgroup of $`W`$ (see Corollary 14 below). The main result of the present paper is a complete description of the finite continuation of a simple reflection in a Coxeter system of finite rank. Its proof constitutes the bulk of this paper.
Main Result: Let $`(W,R)`$ be a Coxeter system of finite rank. Then the following holds.
* For each $`rR`$ the finite continuation of $`r`$ can be described.
* Given an involution $`wW`$ such that $`\mathrm{FC}(w)=w`$, then $`wR^W`$.
Part a) of our main result is Theorem 7. Its precise statement requires some preparation. Part b) is Corollary 24.
The main result of this paper is in fact the first of two steps to reduce the isomorphism problem for Coxeter groups to its ‘reflection-preserving’ version. The second step is given in . We refer to for further information about the applications to the general isomorphism problem.
A special instance of the isomorphism problem for Coxeter groups is the question about their rigidity (see for further information). In combination with the main result of a consequence of our main result is the following rigidity result.
###### Theorem 1.
Let $`(W,R)`$ be an irreducible, non-spherical Coxeter system such that $`R`$ is finite and such that $`rr^{}`$ has finite order for all $`r,r^{}R`$. Then the following assertions hold.
* For each $`rR`$ we have $`\mathrm{FC}(r)=r`$.
* If $`SW`$ is such that $`(W,S)`$ is a Coxeter system, then there exists $`wW`$ such that $`S^w=R`$.
* All automorphisms of $`W`$ are inner-by-graph.
In the language of , Part b) of the previous theorem means that an infinite, irreducible, 2-spherical Coxeter system is strongly rigid. Part c), which is an immediate consequence of Part b), improves the result of .
To conclude this introduction we remark that characterizations results for reflections in even Coxeter groups have been obtained in . Some of the results there can be deduced as corollaries of our main result as well.
### Acknowledgement
The authors thank Frédéric Haglund for a helpful discussion on the subject.
## 2. Precise Statement of the Main Result
Recall that a Coxeter group is a group with a presentation of the form
(2.1)
$$W=\mathrm{gp}\{r_aa\mathrm{\Pi }\}(r_ar_b)^{m_{ab}}=1\text{ for all }a,b\mathrm{\Pi }$$
where $`\mathrm{\Pi }`$ is some indexing set, whose cardinality is called the rank of $`W`$ (relative to this presentation), and the $`m_{ab}`$ satisfy the following conditions: $`m_{ab}=m_{ba}`$, each $`m_{ab}`$ lies in the set $`\{mm1\}\{\mathrm{}\}`$, and $`m_{ab}=1`$ if and only if $`a=b`$. When $`m_{ab}=\mathrm{}`$ the relation $`(r_ar_b)^{m_{ab}}=1`$ is interpreted as vacuous. We shall restrict attention to Coxeter groups of finite rank.
A reduced expression for an element $`wW`$ is a minimal length word expressing $`w`$ as a product of elements of the distinguished generating set $`\{r_aa\mathrm{\Pi }\}`$. We define $`\mathrm{}(w)`$ to be the length of a reduced expression for $`w`$.
As is well known (and as we shall describe in Section 3 below), every Coxeter group W can be realized geometrically as a group generated by reflections. In this realization of $`W`$ the reflections in $`W`$ are the conjugates of the generators $`r_a`$.
The Coxeter graph associated with the presentation above is the graph with vertex set $`\mathrm{\Pi }`$ and edge set consisting of those pairs of vertices $`\{a,b\}`$ for which $`m_{ab}3`$. The edge $`\{a,b\}`$ is given the label $`m_{ab}`$. The components of $`\mathrm{\Pi }`$ are the connected components of the graph, and we say that $`W`$ is irreducible if the graph is connected.
For each $`I\mathrm{\Pi }`$ we define $`W_I`$ to be the subgroup of $`W`$ generated by the set $`\{r_aaI\}`$; we call these subgroups the visible subgroups of $`W`$. A parabolic subgroup of $`W`$ is any conjugate of a visible subgroup. We say that $`I\mathrm{\Pi }`$ is spherical if $`W_I`$ is finite, and we say that $`\mathrm{\Pi }`$ (or $`W`$) is $`k`$-spherical if all $`k`$-element subsets of $`\mathrm{\Pi }`$ are spherical.
The defitions given so far are fairly standard. In order to facilitate the precise statement of the main result, we introduce some nonstandard notation and terminology (in Definitions 2, 3, 4, 5 and 6 below).
###### Definition 2.
If $`wW`$ has finite order, define the finite continuation of $`w`$, written $`\mathrm{FC}(w)`$, to be the intersection of all maximal finite subgroups of $`W`$ containing $`w`$.
###### Definition 3.
The odd graph of $`W`$ is the graph $`\mathrm{\Omega }(\mathrm{\Pi })`$ obtained from the Coxeter graph by deleting the edges whose labels are infinite or even. For each $`a\mathrm{\Pi }`$ we define $`\mathrm{Odd}(a)`$ to be the connected component of $`\mathrm{\Omega }(\mathrm{\Pi })`$ containing $`a`$. For each connected component $`M`$ of $`\mathrm{\Omega }(\mathrm{\Pi })`$ we define $`\mathrm{E}(M)`$ to be the union of $`M`$ with the set of all $`b\mathrm{\Pi }`$ such that $`m_{cb}`$ is even for some $`cM`$. We also abbreviate $`\mathrm{E}(\mathrm{Odd}(a))`$ to $`\mathrm{EOdd}(a)`$.
In the discussions below, when we refer to the components of $`\mathrm{E}(M)`$ we regard $`\mathrm{E}(M)`$ as the full subgraph of the Coxeter graph spanned by the vertices in $`\mathrm{E}(M)`$. In other words, the edges with even and infinite labels, deleted when forming the odd graph, are restored in $`\mathrm{E}(M)`$.
Note that if $`aL\mathrm{\Pi }`$ and $`W_L`$ is finite then $`m_{ab}<\mathrm{}`$ for all $`bL`$. Whether $`m_{ab}`$ is odd or even it follows that $`b\mathrm{EOdd}(a)`$. Thus $`L\mathrm{EOdd}(a)`$.
###### Definition 4.
Let $`M\mathrm{\Pi }`$ be a connected component of $`\mathrm{\Omega }(\mathrm{\Pi })`$. We call $`b\mathrm{\Pi }M`$ a $`C_3`$-neighbour of $`M`$ if $`m_{bc}\{2,4\}`$ for all $`cE(M)`$, the case $`m_{bc}=4`$ occurring for at least one $`c`$, and for each $`cE(M)`$ with $`m_{bc}=4`$ there is an $`aM`$ such that the following conditions are satisfied:
1. $`m_{ba}=2`$ and $`m_{ca}=3`$, and $`m_{cd}=\mathrm{}`$ for all $`dM\{a,c\}`$;
2. for all $`e\mathrm{\Pi }(M\{b\})`$, either $`m_{ce}=\mathrm{}`$ or $`m_{ae}=m_{ce}=m_{be}=2`$.
###### Definition 5.
Let $`M\mathrm{\Pi }`$ be a connected component of $`\mathrm{\Omega }(\mathrm{\Pi })`$, and let $`aM`$ and $`b\mathrm{\Pi }M`$. We call the pair $`(a,b)`$ a focus of $`M`$ in $`\mathrm{\Pi }`$ if the following conditions all hold.
1. All the edge labels of $`M`$ are 3, and $`M`$ is a tree.
2. For each $`cM`$, the set $`C[b..c]\mathrm{\Pi }`$ consisting of $`b`$ and those elements of $`M`$ that form the path from $`a`$ to $`c`$ in $`M`$ constitutes a system of type $`C_k`$ (for some $`k`$ dependent on $`c`$).
3. If $`c,dM\{b\}`$ with $`cC[b..d]`$ and $`dC[b..c]`$ then $`m_{cd}=\mathrm{}`$.
4. If $`m_{ce}\mathrm{}`$ for some $`cM`$ and $`e\mathrm{\Pi }(M\{b\})`$, then $`m_{ce}=2=m_{de}`$ for all $`dC[b..c]`$.
5. The vertices of $`M\{b\}`$ do not form a spherical component of $`\mathrm{E}(M)`$.
###### Definition 6.
Let $`M\mathrm{\Pi }`$ be a connected component of $`\mathrm{\Omega }(\mathrm{\Pi })`$, and let $`a,bM`$. We call the two-element set $`\{a,b\}`$ a half focus of $`M`$ in $`\mathrm{\Pi }`$ if $`m_{ab}=2`$ and the following conditions all hold.
1. We have $`m_{ac}=m_{bc}\{2,3\}`$ for all $`cM\{a,b\}`$, and $`m_{ac}=m_{bc}\{2,\mathrm{}\}`$ for all $`c\mathrm{\Pi }M`$.
2. All the edge labels of $`M\{b\}`$ are 3, and $`M\{b\}`$ is a tree.
3. For each $`cM\{a,b\}`$, the set $`D[a,b..c]\mathrm{\Pi }`$ consisting of $`b`$ and those and those elements of $`M\{b\}`$ that form the path from $`a`$ to $`c`$ constitutes a system of type $`D_k`$ (for some $`k`$ dependent on $`c`$).
4. If $`c,dM\{a,b\}`$ with $`cD[a,b..d]`$ and $`dD[a,b..c]`$ then $`m_{cd}=\mathrm{}`$.
5. If $`m_{ce}\mathrm{}`$ for some $`cM\{a,b\}`$ and $`e\mathrm{\Pi }M`$, then $`m_{ce}=2=m_{de}`$ for all $`dD[a,b..c]`$.
6. The vertices of $`M`$ do not form a spherical component of $`\mathrm{E}(M)`$.
We are now able to give a precise statement of Part a) of our main result.
###### Theorem 7.
For each connected component $`M`$ of $`\mathrm{\Omega }(\mathrm{\Pi })`$ there is at least one $`aM`$ such that $`\mathrm{FC}(r_a)`$ is a visible subgroup of $`W`$. We have the following possibilities.
Case A: Suppose that the component of $`\mathrm{E}(M)`$ containing $`M`$ is spherical, and let $`aM`$ be arbitrary. Then $`\mathrm{FC}(r_a)=W_J`$, where $`J`$ is the union of the spherical components of $`\mathrm{E}(M)`$.
Case B: Suppose that the component of $`\mathrm{E}(M)`$ containing $`M`$ is not spherical, and $`M`$ does not have any focus or half-focus in $`\mathrm{\Pi }`$, and let $`J^{}`$ be the union of the spherical components of $`\mathrm{E}(M)`$ and the set of $`C_3`$-neighbours of $`M`$. If $`aM`$ is not adjacent in $`\mathrm{\Pi }`$ to any $`C_3`$-neighbour of $`M`$ then $`\mathrm{FC}(r_a)=W_{J^{}\{a\}}`$, and if $`aM`$ is adjacent in $`\mathrm{\Pi }`$ to a $`C_3`$-neighbour of $`M`$ then $`\mathrm{FC}(r_a)`$ is not visible.
Case C: Suppose that $`(a,b)`$ is a focus of $`M`$. Then $`\mathrm{FC}(r_a)=W_J`$ where $`J`$ is the union of $`\{a,b\}`$ and the spherical components of $`\mathrm{E}(M)`$, and $`\mathrm{FC}(r_c)`$ is not visible for any $`cM\{a\}`$.
Case D: Suppose that $`\{a,b\}`$ is a half-focus of $`M`$. Then $`\mathrm{FC}(r_a)=\mathrm{FC}(r_b)=W_J`$, where $`J`$ is the union of $`\{a,b\}`$ and the spherical components of $`\mathrm{E}(M)`$, and $`\mathrm{FC}(r_c)`$ is not visible for any $`cM\{a,b\}`$.
## 3. Reflections and root systems
Let $``$ be the real field, and $`V`$ the vector space over $``$ with basis $`\mathrm{\Pi }`$. Let $`B`$ the bilinear form on $`V`$ such that for all $`a,b\mathrm{\Pi }`$,
$$B(a,b)=\mathrm{cos}(\pi /m_{ab}).$$
To make our notation more compact we define $`uv=B(u,v)`$ for all $`u,vV`$. Note that $`aa=1`$ for all $`a\mathrm{\Pi }`$, since $`m_{aa}=1`$.
For each $`aV`$ such that $`aa=1`$, the reflection along $`a`$ is the transformation of $`V`$ given by $`vv2(av)a`$. It is well known (see, for example, Corollary 5.4 of ) that $`W`$ has a faithful representation on $`V`$ such that, for all $`a\mathrm{\Pi }`$, the element $`r_a`$ acts as the reflection along $`a`$. We shall identify elements of $`W`$ with their images in this representation. We also use the notation $`r_a`$ for the reflection along $`a`$ whenever $`aV`$ satisfies $`aa=1`$. It is straightforward to show that each reflection $`r_a`$ preserves the form $`B`$; hence all elements of $`W`$ preserve $`B`$. Furthermore, the equation $`gr_ag^1=r_{ga}`$ holds for all $`aV`$ such that $`aa=1`$ and all transformations $`g`$ that preserve $`B`$.
We write $`\mathrm{Ref}(W)`$ for the set of all reflections in $`W`$. It is immediate from the above comments that if $`\mathrm{\Phi }=\{wawW,a\mathrm{\Pi }\}`$ then $`\{r_bb\mathrm{\Phi }\}\mathrm{Ref}(W)`$.
The set $`\mathrm{\Phi }`$ is called the root system of $`W`$, and elements of $`\mathrm{\Phi }`$ are called roots. Elements of the basis $`\mathrm{\Pi }`$ are called simple roots, and the reflections $`r_a`$ for $`a\mathrm{\Pi }`$ are called simple reflections. A root is said to be positive if it has the form $`_{a\mathrm{\Pi }}\lambda _aa`$ with $`\lambda _a0`$ for all $`a\mathrm{\Pi }`$, and negative otherwise. We write $`\mathrm{\Phi }^+`$ for the set of all positive roots and $`\mathrm{\Phi }^{}`$ for the set of all negative roots.
###### Lemma 8.
With the notation as above, the following statements hold.
1. Every negative root has the form $`_{a\mathrm{\Pi }}\lambda _aa`$ with $`\lambda _a0`$ for all $`a\mathrm{\Pi }`$. Furthermore, $`\mathrm{\Phi }^{}=\{bb\mathrm{\Phi }^+\}`$.
2. If $`wW`$ and $`a\mathrm{\Pi }`$ then
$$\mathrm{}(wr_a)=\{\begin{array}{cc}\mathrm{}(w)+1\hfill & \text{if }wa\mathrm{\Phi }^+\text{,}\hfill \\ \mathrm{}(w)1\hfill & \text{if }wa\mathrm{\Phi }^{}\text{.}\hfill \end{array}$$
3. If $`t\mathrm{Ref}(W)`$ then $`t=r_b`$ for some $`b\mathrm{\Phi }`$.
4. The group $`W`$ is finite if and only if the bilinear form $`B`$ is positive definite.
5. The root system $`\mathrm{\Phi }`$ is finite if and only if the group $`W`$ is finite.
###### Proof.
Proofs of (1) and (2) can be found in \[14, Section 5.4\], Theorem 4.1 in includes both (4) and (5), and (3) is \[13, Lemma 2.2\]. ∎
The following result is well known.
###### Lemma 9.
Let $`a\mathrm{\Pi }`$. Then $`\mathrm{Odd}(a)=\mathrm{\Pi }Wa`$.
For each $`wW`$ we define $`N(w)=\{b\mathrm{\Phi }^+wb\mathrm{\Phi }^{}\}`$. By Part (2) of Lemma 8, if $`w1`$ then $`N(w)\mathrm{\Pi }\mathrm{}`$. An easy induction shows that $`N(w)`$ has exactly $`\mathrm{}(w)`$ elements. In particular, $`N(w)`$ is a finite set. It is also easily shown that if $`\mathrm{\Phi }`$ is finite then there is a unique $`wW`$ such that $`N(w)=\mathrm{\Phi }^+`$. This element, which we denote by $`w_\mathrm{\Pi }`$, is also the unique element of maximal length in $`W`$ (which is a finite group). Furthermore, $`w_\mathrm{\Pi }\mathrm{\Pi }=\mathrm{\Pi }`$.
For each $`\mathrm{\Gamma }\mathrm{\Phi }`$ the subgroup $`W_\mathrm{\Gamma }`$ generated by the set $`\{r_aa\mathrm{\Gamma }\}`$ is called a reflection subgroup of $`W`$. The set $`\mathrm{\Phi }_\mathrm{\Gamma }=\{a\mathrm{\Phi }r_aW_\mathrm{\Gamma }\}`$ is called the root subsystem generated by $`\mathrm{\Gamma }`$. Let $`\mathrm{\Phi }_\mathrm{\Gamma }^+=\mathrm{\Phi }_\mathrm{\Gamma }\mathrm{\Phi }^+`$ and $`\mathrm{\Phi }_\mathrm{\Gamma }^{}=\mathrm{\Phi }_\mathrm{\Gamma }\mathrm{\Phi }^{}`$, and define
$$\mathrm{\Pi }_\mathrm{\Gamma }=\{a\mathrm{\Phi }_\mathrm{\Gamma }^+N(r_a)\mathrm{\Phi }_\mathrm{\Gamma }=\{a\}\}.$$
The main theorem of Deodhar and Theorem (3.3) of Dyer yield the following result.
###### Theorem 10.
For each $`\mathrm{\Gamma }\mathrm{\Phi }`$ the group $`W_\mathrm{\Gamma }`$ is a Coxeter group on the generating set $`\{r_aa\mathrm{\Pi }_\mathrm{\Gamma }\}`$. The set $`\{aba,b\mathrm{\Pi }_\mathrm{\Gamma }\text{ and }ab\}`$ is a subset of $`𝒞=\{\mathrm{cos}(\pi /m)2m\}(\mathrm{},1]`$. Moreover, if $`\mathrm{\Delta }`$ is any subset of $`\mathrm{\Phi }^+`$ such that $`\{aba,b\mathrm{\Delta }\text{ and }ab\}𝒞`$ then $`W_\mathrm{\Delta }`$ is a Coxeter group on the generating set $`\{r_aa\mathrm{\Delta }\}`$.
Note that the notation $`W_\mathrm{\Gamma }`$ introduced above is an extension of the notation for visible subgroups introduced in Section 2. However, if $`\mathrm{\Gamma }\mathrm{\Pi }`$ then $`W_\mathrm{\Gamma }`$ need not be visible.
It is clear that if $`I\mathrm{\Pi }`$ then $`W_I`$ preserves the subspace $`V_I`$ of $`V`$ spanned by $`I`$, and acts on this subspace as a Coxeter group with $`I`$ as its set of simple roots. In this case $`\mathrm{\Phi }_I=\mathrm{\Phi }V_I`$ and $`\mathrm{\Pi }_I=I`$.
The following simple facts are well known.
###### Lemma 11.
In the above situation, $`\mathrm{\Phi }_I=\mathrm{\Phi }V_I`$. Furthermore, $`wW`$ normalizes $`W_I`$ if and only if $`w\mathrm{\Phi }_I=\mathrm{\Phi }_I`$. In particular, for all $`a\mathrm{\Phi }`$, the reflection $`r_a`$ normalizes $`W_I`$ if and only if $`a\mathrm{\Phi }_I`$ or $`ab=0`$ for all $`bI`$.
Suppose that $`I\mathrm{\Pi }`$ and $`a\mathrm{\Pi }I`$, and let $`L`$ be the component of (the Coxeter graph of) $`I\{a\}`$ to which $`a`$ belongs. If $`W_L`$ is finite we define $`v[a,I]=w_Lw_{L\{a\}}`$. It is easily seen that $`v[a,I]II\{a\}`$, and that $`v[a,I]b=b`$ for all $`bIL`$. In particular, $`v[a,I]I=\{J\mathrm{\Pi }J=wI\text{ for some }wW\text{ }\}`$. It was proved in (for finite Coxeter groups) and in (in the general case) that every element $`wW`$ satisfying $`wI\mathrm{\Pi }`$ can be expressed as a product of elements of the form $`v[a,I^{}]`$, with $`I^{}`$ and $`a\mathrm{\Pi }I^{}`$. That is,
(3.1)
$$w=v[a_1,I_1]v[a_2,I_2]\mathrm{}v[a_n,I_n]$$
for some $`I_i,a_i`$ such that (for each $`i`$) the component of $`I_i\{a_i\}`$ containing $`a_i`$ corresponds to a finite visible subgroup, $`v[a_i,I_i]I_i=I_{i1}`$ for $`1<in`$, and $`I_n=I`$. Furthermore, the following result holds.
###### Proposition 12.
Let $`I,J\mathrm{\Pi }`$. Then $`\{wWwW_Iw^1=W_J\}=N(J,I)W_I`$, where $`N(J,I)=\{wWwI=J\}`$. Furthermore, for each $`wN(J,I)`$ and each $`a\mathrm{\Pi }N(w)`$ there is an expression for $`w`$ of the form (3.1) above, with $`(a_n,I_n)=(a,I)`$ and $`\mathrm{}(w)=_{i=1}^n\mathrm{}(v[a_i,I_i])`$.
The following lemma, which appears in \[2, Exercise 2d, p. 130\], is fundamental to all of our arguments.
###### Lemma 13 (Tits).
If $`W`$ is a Coxeter group and $`HW`$ is finite, then $`H`$ is contained in a finite parabolic subgroup of $`W`$.
One immediate consequence of Lemma 13 is that every maximal finite subgroup of a Coxeter group is parabolic. Another consequence of the previous lemma is that each finite subgroup of $`W`$ is contained in a maximal finite parabolic subgroup. (Remember that we always assume that $`W`$ is finitely generated.) Thus the set of maximal finite subgroups of $`W`$ containing a given finite order element of $`W`$ is not empty, and hence we have the following fact.
###### Corollary 14.
If $`wW`$ has finite order, then $`\mathrm{FC}(w)`$ is a well-defined finite subgroup of $`W`$.
###### Lemma 15 (Kilmoyer).
Let $`I,J\mathrm{\Pi }`$. Then every $`(W_I,W_J)`$ double coset in $`W`$ contains a unique element of minimal length; moreover, if $`d`$ is the minimal length element of $`W_IdW_J`$ then $`W_IdW_Jd^1=W_K`$, where $`K=IdJ`$.
###### Proof.
See \[6, Theorem 2.7.4\]. ∎
###### Corollary 16.
The intersection of a finite number of parabolic subgroups is a parabolic subgroup.
The following consequence of Lemmas 13 and 15 is proved in \[10, Lemma 11\].
###### Lemma 17.
If $`J`$ is a maximal spherical subset of $`\mathrm{\Pi }`$ then $`W_J`$ is a maximal finite subgroup of $`W`$. Furthermore, $`W_J`$ is not conjugate to any other visible subgroup of $`W`$.
Another important tool in our analysis of automorphisms is the classification of involutions in Coxeter groups, due to Richardson .
###### Proposition 18.
Suppose that $`wW`$ is an involution. Then there is a $`tW`$ and a spherical $`I\mathrm{\Pi }`$ such that $`w=t^1w_It`$ with $`\mathrm{}(w)=\mathrm{}(w_I)+2\mathrm{}(t)`$, and $`w_I`$ is central in $`W_I`$.
###### Proof.
See \[10, Proposition 5\]. ∎
###### Definition 19.
We say that $`I\mathrm{\Pi }`$ is of ($``$1)-type if $`W_I`$ is finite and $`w_I`$ is central in $`W_I`$.
The reason for the terminology is that $`I`$ is of $`(1)`$-type if and only if there is an element of $`W_I`$ that acts on $`V_I`$ as multiplication by $`1`$.
We need the following lemma.
###### Lemma 20.
Suppose that $`I,J\mathrm{\Pi }`$ with $`I`$ of ($``$1)-type, and suppose that $`tW`$ has the property that $`tw_It^1W_J`$. Then $`tW_It^1W_J`$.
###### Proof.
Let $`aI`$. Then $`w_I(a)=a`$, and so $`(tw_It^1)(ta)=ta`$, whence it follows that either $`ta`$ or $`ta`$ is in the set $`N(tw_It^1)`$. But $`N(tw_It^1)\mathrm{\Phi }_J`$; so $`ta\mathrm{\Phi }_J`$, and therefore $`tr_at^1=r_{ta}W_J`$. Since $`W_I`$ is generated by $`\{r_aaI\}`$, the result follows. ∎
In particular, it follows from Lemma 20 that if $`I,J`$ are both of $`(1)`$-type and $`tw_It^1=w_J`$ then $`tW_It^1=W_J`$. Conversely, suppose that $`tW_It^1=W_J`$, so that in fact $`dW_Id^1=W_J`$ for all $`d`$ in $`W_JtW_I`$ (which equals $`tW_I`$). Taking $`d`$ to be the shortest element in $`tW_I`$, Lemma 15 yields that $`dI=J`$, and hence $`xdxd^1`$ is a length-preserving isomorphism $`W_IW_J`$; consequently $`dw_Id^1=w_J`$. If $`w_I,w_J`$ are central in $`W_I,W_J`$ we deduce that $`tw_It^1=w_J`$. So we have proved the following result.
###### Lemma 21.
Suppose that $`I,J`$ are subsets of $`\mathrm{\Pi }`$ that are both of ($``$1)-type. Then $`\{tWtw_It^1=w_J\}=\{tWtW_It^1=W_J\}`$.
###### Proposition 22.
Let $`I\mathrm{\Pi }`$ be of ($``$1)-type. Then $`W_I\mathrm{FC}(w_I)`$.
###### Proof.
Let $`F`$ be a maximal finite subgroup of $`W`$ such that $`w_IH`$. By Lemma 13 there exist $`tW`$ and $`J\mathrm{\Pi }`$ such that $`tFt^1=W_J`$. By Lemma 20 and the fact that $`w_IH`$ it follows that $`tW_It^1W_J`$. Hence $`W_It^1W_Jt=F`$. ∎
###### Proposition 23.
Let $`W,W^{}`$ be Coxeter groups of finite rank and $`\alpha :WW^{}`$ an isomorphism. Let $`\mathrm{\Pi }`$ be the set of simple roots corresponding to the distinguished generating set of $`W`$, and let $`a\mathrm{\Pi }`$. If $`r_a^\alpha `$ is not a reflection in $`W^{}`$ then the intersection of all maximal finite subgroups of $`W`$ containing $`r_a`$ is a parabolic subgroup of order greater than 2.
###### Proof.
Write $`\mathrm{\Pi }^{}`$ for the set of simple roots of $`W^{}`$. Observe that Lemma 13 and Corollary 16 trivially imply that $`\mathrm{FC}(r_a)`$ is a parabolic subgroup of $`W`$.
Since $`r_a^\alpha `$ is not a reflection it follows from Proposition 18 that $`r_a^\alpha =tw_It^1`$ for some $`tW^{}`$ and $`I\mathrm{\Pi }^{}`$ of $`(1)`$-type and of rank at least 2. Clearly $`\mathrm{FC}(r_a)^\alpha =t\mathrm{FC}(w_I)t^1`$, and by Proposition 22 we know that $`W_I\mathrm{FC}(w_I)`$. Hence $`(tW_It^1)^{\alpha ^1}\mathrm{FC}(r_a)`$, so that $`\mathrm{FC}(r_a)`$ has order greater than 2, as required. ∎
###### Corollary 24.
Let $`wW`$ be an involution such that $`\mathrm{FC}(w)=w`$ and let $`SW`$ be such that $`(W,S)`$ is a Coxeter system. Then $`wS^W`$.
## 4. The finite continuation of a reflection
Let $`r`$ be a reflection in $`W`$. Replacing $`r`$ by $`wrw^1`$ replaces $`\mathrm{FC}(r)`$ by $`w\mathrm{FC}(r)w^1`$, and so choosing $`w`$ suitably enables us to assume that $`\mathrm{FC}(r)=W_J`$, a visible parabolic subgroup. Furthermore, replacing $`r`$ by $`trt^1`$ for suitable $`tW_J`$ enables us to assume that $`r=r_a`$ for some $`aJ`$. (Note that these observations yield the first assertion of Theorem 7.)
###### Proposition 25.
Let $`aJ\mathrm{\Pi }`$, and suppose that $`W_J`$ is the intersection of all maximal finite subgroups of $`W`$ containing $`r_a`$. Then $`\{wWwr_aw^1W_J\}`$ is a subset of the normalizer of $`W_J`$ in $`W`$. Thus each $`W`$-conjugate of $`r_a`$ in $`W_J`$ is $`N_W(W_J)`$-conjugate to $`r_a`$, and $`C_W(r_a)N_W(W_J)`$. Moreover, if $`b\mathrm{\Pi }J`$ is such that $`W_{J\{b\}}`$ is infinite then $`m_{bc}=\mathrm{}`$ for all $`cJ`$ such that $`r_c`$ is conjugate to $`r_a`$ in $`W`$.
###### Proof.
Let $`𝒮`$ be the set of all maximal finite subgroups of $`W`$ containing $`r_a`$, so that $`W_J=\mathrm{FC}(r_a)=_{F𝒮}F`$. Suppose that $`wW`$ satisfies $`wr_aw^1W_J`$, and let $`F𝒮`$. Then $`wr_aw^1W_JF`$, and so $`r_aw^1Fw`$. Thus $`w^1Fw`$ is a maximal finite subgroup of $`W`$ containing $`r_a`$, whence $`w^1Fw𝒮`$. So
$$\underset{F𝒮}{}F\underset{F𝒮}{}w^1Fw$$
and so $`W_Jw^1W_Jw`$. Since $`W_J`$ is finite it follows that $`wN_W(W_J)`$.
Suppose that $`cJ`$ with $`r_c=wr_aw^1`$ for some $`wW`$. Clearly $`FwFw^1`$ is a bijection from the set of maximal finite subgroups of $`W`$ containing $`r_a`$ to the set of maximal finite subgroups of $`W`$ containing $`r_c`$, and so $`\mathrm{FC}(r_c)=w\mathrm{FC}(r_a)w^1`$. But $`w\mathrm{FC}(r_a)w^1=wW_Jw^1=W_J`$ by the first part of the proof, and so $`\mathrm{FC}(r_c)=W_J`$. Now suppose that $`b\mathrm{\Pi }J`$ with $`m_{cb}<\mathrm{}`$. Then $`W_{\{c,b\}}`$ is finite, and so contained in a maximal finite subgroup $`F`$. Since $`r_cF`$ we must have $`\mathrm{FC}(r_c)F`$. It follows that the finite group $`F`$ contains both $`W_J`$ and $`r_b`$, and therefore $`W_{J\{b\}}`$ is finite. ∎
Assume, as in Proposition 25, that $`aJ\mathrm{\Pi }`$ and $`W_J=\mathrm{FC}(r_a)`$, and suppose now that $`J\{a\}`$. Suppose that $`L\mathrm{\Pi }`$ is such that $`JL`$ and $`W_L`$ is finite. Then $`W_L`$ is a finite Coxeter group possessing a visible parabolic subgroup $`W_J`$ of rank greater than 1 that is normalized by the centralizer of some simple reflection $`r_aW_J`$. Indeed, $`W_J`$ is normalized by all $`wW_L`$ such that $`wr_aw^1W_J`$. Equivalently, by Lemma 8 (3), $`\{wW_Lwa\mathrm{\Phi }_J\}N_W(W_J)`$. This is a very restrictive condition, which we now proceed to examine with a case-by-case investigation of the different types of finite Coxeter groups. For the course of this investigation, we can (and shall) assume that $`L=\mathrm{\Pi }`$.
So we assume for now that $`W`$ is a finite Coxeter group of rank $`n`$, and our aim is to find all examples of the following phenomenon: there exist $`\{a\}J\mathrm{\Pi }`$ such that the set $`Q=\{wWwa\mathrm{\Phi }_J\}`$ is a subset of $`N_W(W_J)`$. We assume that $`J\mathrm{\Pi }`$, since the condition is trivially satisfied otherwise.
If $`K\mathrm{\Pi }`$ is a component of the Coxeter graph such that $`JK=\mathrm{}`$ then $`W_K`$ is a direct factor of $`N_W(W_J)`$; moreover, $`Q=(QW_{\mathrm{\Pi }K})W_K`$. So removing $`K`$ from the graph will have no bearing on whether or not the condition $`QN_W(W_J)`$ holds. So we assume that there are no such components of $`\mathrm{\Pi }`$. Exactly the same comments apply for a component $`K`$ of $`\mathrm{\Pi }`$ such that $`KJ`$. So we also assume that there are none of these.
Assume that $`\{a\}J\mathrm{\Pi }`$ and $`QN_W(W_J)`$. Suppose that $`K\mathrm{\Pi }`$ is a component of the Coxeter graph such that $`aK`$. Then $`r_ba=a`$ for all $`bK`$; so $`r_bQN_W(W_J)`$, and it follows that $`r_bc\mathrm{\Phi }_J`$ whenever $`cJ`$. If $`bc0`$ then $`b`$ is in the support of $`r_bc`$, and so $`r_bc\mathrm{\Phi }_J`$ implies $`bJ`$. Since $`K`$ is connected it follows that if $`K`$ contains any element of $`J`$ then $`KJ`$. So either $`KJ=\mathrm{}`$ or $`KJ`$. But we have assumed that there are no such components. So the component of $`\mathrm{\Pi }`$ that contains $`a`$ is the only component; that is, $`\mathrm{\Pi }`$ is irreducible.
Observe that the group $`\mathrm{Stab}(a)=\{wWwa=a\}`$ is a subset of $`Q`$ and hence of $`N_W(W_J)`$. Note also that $`N_W(W_J)=\{wWw\mathrm{\Phi }_J=\mathrm{\Phi }_J\}`$, which is also the stabilizer of the subspace $`V_J`$ (since $`V_J`$ is the subspace spanned by $`\mathrm{\Phi }_J`$ and $`\mathrm{\Phi }_J=V_J\mathrm{\Phi }`$). Now $`\mathrm{Stab}(a)`$ is a parabolic subgroup of $`W`$ whose root system is $`\mathrm{\Phi }a^{}`$, and the following table gives the type of this root system in all cases.
| $`W`$ | | $`\mathrm{Stab}(a)`$ |
| --- | --- | --- |
| $`A_n`$ | | $`A_{n2}`$ |
| $`C_n`$ | | $`C_{n2}+A_1`$ |
| $`C_n`$ | | $`C_{n1}`$ |
| $`D_n`$ | | $`D_{n2}+A_1`$ |
| $`F_4`$ | | $`C_3`$ |
| $`E_6`$ | | $`A_5`$ |
| $`W`$ | | $`\mathrm{Stab}(a)`$ |
| --- | --- | --- |
| $`E_7`$ | | $`D_6`$ |
| $`E_8`$ | | $`E_7`$ |
| $`H_3`$ | | $`A_1+A_1`$ |
| $`H_4`$ | | $`H_3`$ |
| $`I_2(2k)`$ | | $`A_1`$ |
| $`I_2(2k+1)`$ | | $`\mathrm{}`$ |
(For $`C_n`$ there are two $`W`$ orbits of roots, giving two possibilities for $`\mathrm{Stab}(a)`$. For $`F_4`$ and $`I_2(2k)`$ there are also two $`W`$-orbits of roots, but $`\mathrm{Stab}(a)`$ has the same type of root system whichever orbit $`a`$ belongs to.) Since each irreducible constituent of its root system spans an irreducible $`\mathrm{Stab}(a)`$-submodule of $`V`$, the table shows that as a $`\mathrm{Stab}(a)`$-module, $`V`$ has composition length two or three or (in one case only) four: $`a`$ itself spans a trivial $`\mathrm{Stab}(a)`$ submodule of dimension 1, and $`a^{}`$ is either irreducible of dimension $`n1`$ (for types $`F_4`$, $`E_6`$, $`E_7`$, $`E_8`$, $`H_4`$, $`I_2(2k)`$ and one of the $`C_n`$ possibilities), or the direct sum of irreducibles of dimensions 1 and $`n2`$ (for types $`A_n`$, $`C_n`$, $`D_n`$ when $`n>4`$, $`H_3`$ and $`I_2(2k+1)`$), or the direct sum of three irreducibles of dimension 1 (for type $`D_4`$). Furthermore, the summands of $`a^{}`$ are pairwise nonisomorphic as $`\mathrm{Stab}(a)`$ modules, since even if they are of the same type their centralizers in $`\mathrm{Stab}(a)`$ are different.
Since $`\{a\}J\mathrm{\Pi }`$ and $`V_J`$ is $`\mathrm{Stab}(a)`$-invariant, we see that $`a^{}=(V_Ja^{})V_J^{}`$, with both summands nonzero $`\mathrm{Stab}(a)`$-modules. So $`\mathrm{\Pi }`$ is of type $`A_n`$, $`C_n`$, $`D_n`$ or $`H_3`$. Furthermore, except in type $`D_4`$, the two direct summands of $`a^{}`$ are irreducible and not isomorphic, and are therefore the only proper $`\mathrm{Stab}(a)`$-submodules of $`a^{}`$. We conclude that $`V_J`$ is spanned by $`a`$ and one of the summands of $`a^{}`$, while $`V_J^{}`$ is the other summand. In type $`D_4`$ we similarly deduce that $`V_J`$ is spanned by $`a`$ and one or two of the three 1-dimensional summands of $`a^{}`$, and, correspondingly, $`V_J^{}`$ is of either of type $`A_1+A_1`$ or of type $`A_1`$.
If $`\mathrm{\Pi }`$ is of type $`A_n`$ then one of the summands of $`a^{}`$ is of type $`A_{n2}`$ while the other is a trivial 1-dimensional $`\mathrm{Stab}(a)`$-module. If $`V_J^{}`$ is of type $`A_{n2}`$ then $`V_J`$ must be of type $`A_1`$, since the orthogonal complement of a subsystem of type $`A_{n2}`$ in $`A_n`$ contains only a rank 1 root system. This contradicts the assumption that $`\{a\}J`$. So $`J`$ is of type $`A_1+A_{n2}`$. Since $`W_J`$ is visible, we deduce that $`a`$ is an end node of the $`A_n`$ diagram, and the node adjacent to $`a`$ is the unique simple root not in $`J`$. However, if $`n>3`$ then the maximal length element of $`W`$ is in $`Q`$ but not in the normalizer of $`W_J`$. So $`n=3`$ and $`J=\{a,c\}`$, where $`c`$ is the other end node. It is readily checked that $`Q`$ has 8 elements and coincides with $`N_W(W_J)`$ (which is generated by $`W_J`$ and an element that interchanges $`a`$ and $`c`$).
If $`\mathrm{\Pi }`$ is of type $`C_n`$ then one summand of $`a^{}`$ is of type $`C_{n2}`$ and the other of type $`A_1`$. The roots in the $`A_1`$ summand are in the same $`W`$-orbit as $`a`$. If $`V_J^{}`$ is the $`A_1`$ component of $`a^{}`$ then $`V_J=(V_J^{})^{}`$ is of type $`C_{n2}+A_1`$. This determines $`J`$ uniquely, since $`W_J`$ is visible. If $`n4`$ and $`w`$ is the longest element of the visible parabolic subgroup of type $`A_{n1}`$, then $`wa\mathrm{\Pi }\{b\}=J`$, but $`wN_W(W_J)`$. This contradicts the fact that $`QN_W(W_J)`$. So $`n=3`$, and the elements of $`J`$ are the end nodes $`a,c`$ of the $`C_3`$ diagram, the middle node $`b`$ being in the same $`W`$-orbit as $`a`$. Since $`\mathrm{\Phi }_J=\{\pm a,\pm c\}`$ and $`c`$ is not in the same $`W`$-orbit as $`a`$ and $`b`$ we deduce that $`Q=\{wWwa=\pm a\}`$. Furthermore, of the 6 roots in the $`W`$-orbit of $`c`$, only $`c`$ and $`c`$ are orthogonal to $`a`$. So if $`wa=\pm a`$ then $`wc=\pm c`$. Thus if $`wQ`$ then $`w\mathrm{\Phi }_J=\mathrm{\Phi }_J`$, as required.
Continuing the discussion of $`C_n`$, suppose now that $`V_J^{}`$ is the $`C_{n2}`$ component of $`a^{}`$. Then $`V_J=(V_J^{})^{}`$ is of type $`C_2`$. Writing $`J=\{a,b\}`$, the fact that $`\mathrm{Stab}(a)`$ is of type $`A_1+C_{n2}`$ means that it is $`b`$ rather than $`a`$ that is the end node of the $`C_n`$ diagram. If we put $`c=r_ba`$ then $`\{\pm c\}`$ is the component of $`\mathrm{\Phi }a^{}`$ of type $`A_1`$. It follows that $`\{\pm a\}=\{\pm r_bc\}`$ is the $`A_1`$-component of $`\mathrm{\Phi }(r_ba)^{}=\mathrm{\Phi }c^{}`$. We see that $`\mathrm{Stab}(a)=r_c\times W^{}`$ and $`\mathrm{Stab}(c)=r_a\times W^{}`$, where $`W^{}`$ is a parabolic (not visible) subgroup of $`W`$ of type $`C_{n2}`$. Indeed, the root system of $`W^{}`$ is $`\mathrm{\Phi }V_J^{}`$. The roots in $`\mathrm{\Phi }_J`$ that are in the same $`W`$-orbit as $`a`$ are $`\pm a`$ and $`\pm c`$, and so
$$Q=\{1,r_a,r_b,r_br_a\}\mathrm{Stab}(a)=\{1,r_a,r_b,r_br_a\}\{1,r_c\}W^{}=W_JW^{}.$$
Hence our requirement that $`Q`$ stabilizes $`\mathrm{\Phi }_J=\{\pm a,\pm b\}`$ is indeed satisfied.
If $`\mathrm{\Pi }`$ is of type $`D_n`$ with $`n>4`$ then one summand of $`a^{}`$ is of type $`D_{n2}`$ and the other of type $`A_1`$. The roots orthogonal to a $`D_{n2}`$ subsystem form a system of type $`A_1+A_1`$. There are in fact two $`W`$-orbits of parabolic $`A_1+A_1`$ subsystems, and the orthogonal complement of a $`D_{n2}`$ is perhaps better thought of as type $`D_2`$, since the visible parabolic in this orbit corresponds to the two nodes of the diagram that form the fork. So if $`V_J^{}`$ is the $`D_{n2}`$ summand of $`a^{}`$ then $`J=\{a,b\}`$ consists to the two nodes of valency 1 that are adjacent to $`c`$, the node of valency 3. A similar statement applies for $`D_4`$ in the case that $`V_J^{}`$ is of type $`A_1+A_1`$. In both cases the element $`w=r_cr_ar_br_cW`$ satisfies $`wa=b`$ and $`wb=a`$, and since $`\mathrm{\Phi }_J=\{\pm a,\pm b\}`$ we see that $`Q=\{1,r_a,w,wr_a\}\mathrm{Stab}(a)`$. But $`\mathrm{Stab}(a)=r_b\times W^{}`$ and $`\mathrm{Stab}(b)=r_a\times W^{}`$, where $`W^{}`$ is the parabolic subgroup corresponding to the subspace $`V_J^{}`$, and it follows readily that $`Q`$ stabilizes $`\mathrm{\Phi }_J=\{\pm a,\pm b\}`$, as required.
Continuing the discussion of $`D_n`$, where $`n4`$, suppose now that $`V_J^{}`$ is an $`A_1`$ component of $`a^{}`$. Then $`V_J=(V_J^{})^{}`$ is of type $`A_1+D_{n2}`$. But the maximal length element of a visible $`A_{n1}`$ subsystem containing $`a`$ takes $`a`$ to an element of $`\mathrm{\Phi }_J`$ but does not normalize $`W_J`$. So our requirement that $`QN_W(W_J)`$ is not met.
Finally, suppose that $`\mathrm{\Pi }`$ is of type $`H_3`$, so that $`\mathrm{Stab}(a)`$ is of type $`A_1+A_1`$. Then $`V_J^{}`$ is of type $`A_1`$, and hence $`J`$ is of type $`A_1+A_1`$. Let $`J=\{a,c\}`$, and note that $`c=wa`$ for some $`wW`$. Since $`N_W(W_J)`$ is generated by $`W_J`$ and the central involution of $`W`$, we see that $`c`$ is not in the $`N_W(W_J)`$-orbit of $`a`$. Hence the element $`w`$ above is in $`Q`$ but not in $`N_W(W_J)`$, and so our requirements are not met.
We have thus established the following result.
###### Proposition 26.
Let $`\mathrm{\Pi }`$ be the set of simple roots for the finite irreducible Coxeter group $`W`$, and suppose that $`aJ\mathrm{\Pi }`$. Then $`\{wWwa\mathrm{\Phi }_J\}`$ is a subset of $`N_W(W_J)`$ if and only if one of the following situations occurs:
1. $`J=\{a\}`$;
2. $`J=\mathrm{\Pi }`$;
3. $`\mathrm{\Pi }=\{a,b,c\}`$ is of type $`C_3`$, with $`m_{ab}=3`$ and $`m_{bc}=4`$, and $`J=\{a,c\}`$ of type $`A_1+A_1`$;
4. $`\mathrm{\Pi }`$ is of type $`D_n`$ or $`A_3`$, and $`J=\{a,b\}`$, where $`a`$ and $`b`$ are end nodes that are both adjacent to some $`c\mathrm{\Pi }`$;
5. $`\mathrm{\Pi }`$ is of type $`C_n`$ and $`J=\{a,b\}`$ is of type $`C_2`$, with $`b`$ an end node of $`\mathrm{\Pi }`$.
We return now to investigation of an arbitrary finite rank Coxeter group $`W`$. The next proposition is an immediate consequence of Proposition 26 and the discussion preceding it.
###### Proposition 27.
Let $`aJL\mathrm{\Pi }`$, and suppose that the group $`W_L`$ is finite and that $`W_J=\mathrm{FC}(r_a)`$. Let $`J_0`$ be the component of $`J`$ containing $`a`$ and $`L_0`$ the component of $`L`$ containing $`J_0`$. Then every component of $`J`$ that is not contained in $`L_0`$ is a component of $`L`$. Furthermore, if $`\{a\}JL_0L_0`$ then $`JL_0=\{a,b\}`$ for some $`b`$, and one of the following alternatives occurs:
1. $`L_0=\{a,c,b\}`$ is of type $`C_3`$, with $`m_{ac}=3`$ and $`m_{cb}=4`$;
2. $`L_0`$ is of type $`C_n`$ for some $`n3`$, with $`b`$ an end node and $`J_0=\{a,b\}`$ of type $`C_2`$;
3. $`L_0`$ is of type $`A_3`$ or type $`D_n`$ for some $`n4`$, the nodes $`a`$ and $`b`$ having valency 1 and sharing a common neighbour.
One of the ingredients of alternative (2) of Proposition 27 is that the component of $`\mathrm{FC}(r_a)`$ containing $`a`$ is of type $`C_2`$. We shall see that when this situation arises, $`\mathrm{Odd}(a)`$ has a focus in $`\mathrm{\Pi }`$.
###### Proposition 28.
Suppose that $`aJ\mathrm{\Pi }`$ with $`W_J=\mathrm{FC}(r_a)`$, and let $`J_0`$ be the component of $`J`$ containing $`a`$. Suppose that $`J_0=\{a,b\}`$ is of type $`C_2`$. Then either $`\mathrm{Odd}(a)\{b\}`$ is a spherical component of $`\mathrm{EOdd}(a)`$, or else $`(a,b)`$ is a focus of $`\mathrm{Odd}(a)`$ in $`\mathrm{\Pi }`$.
###### Proof.
We use induction on $`k`$ to prove that for all $`k2`$, if $`b=c_1,a=c_2,c_3,\mathrm{},c_k`$ are simple roots satisfying
1. $`2<m_{c_ic_{i+1}}<\mathrm{}`$ for all $`i\{1,2,\mathrm{},k1\}`$, and
2. $`c_1,c_2,\mathrm{},c_k`$ are distinct from each other,
then $`\{c_1,c_2,\mathrm{},c_k\}`$ forms a system of type $`C_k`$. The case $`k=2`$ is immediately true.
Suppose that $`k>2`$. The inductive hypothesis tells us that $`\{c_1,c_2,\mathrm{},c_{k1}\}`$ is of type $`C_{k1}`$. The element $`w=v[c_{k1},\{c_{k2}\}]\mathrm{}v[c_4,\{c_3\}]v[c_3,\{c_2\}]`$ has the property that $`wa=wc_2=c_{k1}`$, and so if we write $`d=wb`$ then
$$r_d=wr_bw^1w\mathrm{FC}(r_a)w^1=\mathrm{FC}(c_{k1}),$$
since it is given that $`b\mathrm{FC}(r_a)`$. But $`W_{\{c_{k1},c_k\}}`$ is finite, and so it follows that $`\{r_d,r_{c_{k1}},r_{c_k}\}`$ generates a finite group. Now $`dc_{k1}=ba=\mathrm{cos}(\pi /4)`$ and $`c_{k1}c_k=\mathrm{cos}(\pi /m)`$ for some $`m>2`$. If $`m4`$ then
$$c_kd=c_k\left(b+\sqrt{2}\underset{i=2}{\overset{k1}{}}c_i\right)\sqrt{2}(c_kc_{k1})1,$$
whence the reflection subgroup $`W_{\{r_d,r_{c_k}\}}`$ is infinite (by Theorem 10), a contradiction. So $`m=3`$. If $`m_{c_ic_k}>2`$ for any $`i\{1,2,\mathrm{},k2\}`$ then
$$c_kd\sqrt{2}(c_kc_{k1})+c_kc_i<1,$$
again giving a contradiction. So $`m_{c_kc_i}=2`$ for all $`i\{1,2,\mathrm{},k2\}`$ and $`m_{c_kc_{k1}}=3`$, and since $`\{c_1,c_2,\mathrm{},c_{k1}\}`$ is a system of type $`C_{k1}`$ it follows that $`\{c_1,c_2,\mathrm{},c_k\}`$ is a system of type $`C_k`$, as claimed.
If there were $`c,d\mathrm{Odd}(a)`$ with $`3<m_{cd}<\mathrm{}`$ then $`b`$ together with a minimal length odd-labelled path from $`a`$ to $`\{c,d\}`$ would yield $`c_1,c_2,\mathrm{},c_k\mathrm{\Pi }`$ satisfying (1) and (2) above and not forming a system of type $`C_k`$, contradicting the result proved above. The same argument yields a contradiction if $`c\mathrm{Odd}(a)`$ and $`d\mathrm{\Pi }\mathrm{Odd}(a)`$ with $`3<m_{cd}<\mathrm{}`$, unless $`\{c,d\}=\{a,b\}`$. So all edge labels in $`\mathrm{Odd}(a)`$ are $`3`$, if $`c,d\mathrm{Odd}(a)`$ are not adjacent in $`\mathrm{Odd}(a)`$ then $`m_{cd}\{2,\mathrm{}\}`$, and if $`c\mathrm{Odd}(a)`$ and $`d\mathrm{\Pi }\mathrm{Odd}(a)`$ then $`m_{cd}\{2,\mathrm{}\}`$ unless $`\{c,d\}=\{a,b\}`$. Furthermore, any circuit in $`\mathrm{Odd}(a)`$ would similarly yield a contradiction (by combining the circuit with a minimal finite-labelled path connecting it to $`b`$). So $`\mathrm{Odd}(a)`$ is tree.
For each $`c\mathrm{Odd}(a)`$ let $`C[b..c]\mathrm{\Pi }`$ consist of $`b`$ and the unique path from $`a`$ to $`c`$ in $`\mathrm{Odd}(a)`$. The discussion above shows that $`C[b..c]`$ is always of type $`C`$. Now suppose that $`c\mathrm{Odd}(a)`$ and $`e\mathrm{\Pi }C[b..c]`$ with $`m_{ce}=2`$. Write $`C[b..c]=\{c_1,c_2,\mathrm{},c_k\}`$, with $`c_1=b`$ and $`c_k=c`$, and let $`d=b+\sqrt{2}_{i=2}^kc_i`$. An argument similar to one used above shows that $`r_d\mathrm{FC}(c)`$, and hence $`W_{\{d,c,e\}}`$ is finite. So $`de>1`$. If $`c_ie0`$ then $`c_ie1/2`$; so it follows that there is at most one $`i`$ with $`c_ie0`$. Suppose, for a contradiction, that there is exactly one such $`i`$. If $`i>1`$ then $`de=\sqrt{2}(c_ie)`$, and so $`c_ie>1/\sqrt{2}`$. Hence $`m_{c_ie}=3`$, and $`de=1/\sqrt{2}`$. But this means that the edges $`\{c,d\}`$ and $`\{d,e\}`$ of the Coxeter graph of $`\{d,c,e\}`$ are both labelled $`4`$, contradicting the fact that $`W_{\{d,c,e\}}`$ is finite. So we must have $`i=1`$, and finiteness of $`W_{\{c,d,e\}}`$ forces $`be=de=1/2`$. But now if we put $`L=\{e\}J`$ then, in the notation of Proposition 27, we have that $`L_0=\{e,b,a\}`$ is of type $`C_3`$ with $`JL_0=\{b,a\}`$ of type $`C_2`$, and Proposition 27 shows that this is not possible. We conclude that if $`e\mathrm{\Pi }`$ has the property that $`m_{ce}=2`$ for some $`c\mathrm{Odd}(a)`$ then $`m_{de}=2`$ for all $`dC[b..c]`$. In particular, if $`e\mathrm{\Pi }(\mathrm{Odd}(a)\{b\})`$ and $`m_{ce}\mathrm{}`$ for some $`c\mathrm{Odd}(a)`$ then $`m_{ce}=2`$, as shown above, and so $`m_{de}=2`$ for all $`dC[b..c]`$.
All that remains to prove now is that if $`c,d\mathrm{Odd}(a)`$ with $`cC[b..d]`$ and $`dC[b..c]`$, then $`m_{cd}=\mathrm{}`$. Since $`c`$ and $`d`$ are not adjacent in $`\mathrm{Odd}(a)`$ the only alternative is that $`m_{cd}=2`$; so suppose, for a contradiction, that this holds. Choose the vertex $`e\mathrm{Odd}(a)`$ on the (unique) path from $`c`$ to $`d`$ such that the distance from $`e`$ to $`a`$ is minimal. Let $`c^{},d^{}`$ be the neighbours of $`e`$ in the path from $`c`$ to $`d`$, with $`c^{}`$ between $`e`$ and $`c`$ and $`d^{}`$ between $`e`$ and $`d`$. Then $`c^{}C[b..c]`$, and since $`m_{cd}=2`$ it follows that $`m_{c^{}d}=2`$. Now since $`d^{}C[b..d]`$ and $`m_{c^{}d}=2`$ it follows that $`m_{c^{}d^{}}=2`$. Thus the set $`L\mathrm{\Pi }`$ consisting of $`c^{}`$ and $`d^{}`$ and the vertices on the path from $`a`$ to $`e`$ form a system of type $`D`$ (or $`A_3`$ if $`e=a`$). So $`L`$ is spherical, and since $`b\mathrm{FC}(r_a)`$ it follows that $`L\{b\}`$ is spherical also. But this is impossible since $`L\{b\}`$ is connected, has an edge labelled 4 (namely, $`\{b,a\}`$), and has a vertex of valency 3 (namely $`e`$). ∎
The situation of alternative (3) of Proposition 27 is very similar to that of alternative (2), and in this case it turns out that $`\mathrm{Odd}(a)`$ has a half-focus in $`\mathrm{\Pi }`$.
###### Proposition 29.
Suppose that $`aJ\mathrm{\Pi }`$ with $`W_J=\mathrm{FC}(r_a)`$ and $`\{a\}`$ a component of $`J`$, and suppose that $`J\mathrm{Odd}(a)\{a\}`$. Then either $`\mathrm{Odd}(a)`$ is a spherical component of $`\mathrm{EOdd}(a)`$, or else there exists an element $`b\mathrm{Odd}(a)`$ such that $`\{a,b\}`$ is a half focus of $`\mathrm{Odd}(a)`$ in $`\mathrm{\Pi }`$.
###### Proof.
Let $`b(J\mathrm{Odd}(a))\{a\}`$, and let $`wW`$ with $`wa=b`$. Then $`wN_W(W_J)`$, by Proposition 25, and so
$$\mathrm{FC}(r_b)=\mathrm{FC}(wr_aw^1)=w\mathrm{FC}(r_a)w^1=wW_Jw^1=W_J.$$
Moreover, $`w\mathrm{\Phi }_J=\mathrm{\Phi }_J`$, and since $`ac=0`$ for all $`c\mathrm{\Phi }_J\{a\}`$, it follows that $`wad=0`$ for all $`d\mathrm{\Phi }_J\{wa\}`$. So $`\{b\}`$ is a component of $`J`$. Note that $`m_{ab}=2`$, since $`a`$ and $`b`$ are in different components of $`J`$.
Let $`c\mathrm{\Pi }\{a,b\}`$, and suppose first of all that $`2<m_{bc}<\mathrm{}`$. Since $`\{b,c\}`$ is spherical and $`W_J=\mathrm{FC}(r_b)`$ it follows that $`J\{c\}`$ is spherical. Let $`L=J\{c\}`$ and let $`L_0`$ be the component of $`L`$ containing $`a`$. By Proposition 27, every component of $`J`$ that is not contained in $`L_0`$ is a component of $`L`$. But $`b`$ is adjacent to $`c`$ in $`L`$; so $`\{b\}`$ is not a component of $`L`$, and it follows that $`bL_0`$. Now $`\{a\}JL_0`$, since $`bJL_0`$, and $`JL_0L_0`$, since $`cL_0`$ and $`cJ`$ (since $`\{b\}`$ is a component of $`J`$). Furthermore, the conditions of alternative (2) of Proposition 27 are not satisfied, since $`a`$ and $`b`$ are not adjacent in $`J`$. So either alternative (1) or alternative (3) must hold, and since $`c`$ is the only element of $`L`$ not in $`J`$ it follows that $`L_0=\{a,c,b\}`$, with $`m_{ac}=3`$. But a symmetrical argument, with the roles of $`a`$ and $`b`$ reversed, shows that every $`d\mathrm{\Pi }`$ with $`2<m_{ad}<\mathrm{}`$ has the property that $`m_{bd}=3`$. So $`m_{ac}=m_{bc}=3`$, and $`\{a,c,b\}`$ is of type $`A_3`$.
Now suppose that $`m_{bc}=2`$. Again since $`\{b,c\}`$ is spherical it follows that $`J\{c\}`$ is spherical, and so $`m_{ac}<\mathrm{}`$. If $`m_{ac}>2`$ then, as we have just observed, it follows that $`m_{bc}=3`$, contrary to our assumption that $`m_{bc}=2`$. So $`m_{ac}=m_{bc}=2`$, and we have now shown that whenever $`m_{bc}<\mathrm{}`$ we have $`m_{ac}=m_{bc}\{2,3\}`$. Since a symmetrical argument gives the same conclusion whenever $`m_{ac}<\mathrm{}`$, we conclude also that $`m_{ac}=\mathrm{}`$ if and only if $`m_{bc}=\mathrm{}`$.
We now use induction on $`k`$ to prove that for all $`k3`$, if $`b=c_1,a=c_2,c_3,\mathrm{},c_k`$ are simple roots satisfying
1. $`2<m_{c_ic_{i+1}}<\mathrm{}`$ for all $`i\{2,3,\mathrm{},k1\}`$, and
2. $`c_1,c_2,\mathrm{},c_k`$ are distinct from each other,
then $`\{c_1,c_2,\mathrm{},c_k\}`$ forms a system of type $`D_k`$ or $`A_3`$. The case $`k=3`$ follows from what we have proved above.
Suppose that $`k>3`$. The inductive hypothesis tells us that $`\{c_1,c_2,\mathrm{},c_{k1}\}`$ is of type $`D_{k1}`$ (or $`A_3`$ if $`k=4`$). The element $`w=v[c_{k1},\{c_{k2}\}]\mathrm{}v[c_4,\{c_3\}]v[c_3,\{c_2\}]`$ has the property that $`wa=wc_2=c_{k1}`$, and so if we write $`d=wb`$ then
$$r_d=wr_{c_1}w^1w\mathrm{FC}(r_a)w^1=\mathrm{FC}(c_{k1}),$$
since it is given that $`b\mathrm{FC}(r_a)`$. But $`W_{\{c_{k1},c_k\}}`$ is finite, and so it follows that $`\{r_d,r_{c_{k1}},r_{c_k}\}`$ generates a finite group. Now $`dc_{k1}=ba=0`$ and $`c_{k1}c_k=\mathrm{cos}(\pi /m)`$ for some $`m>2`$. If $`c_kc_i0`$ for some $`i\{1,2,\mathrm{},k2\}`$ then
$$c_kd=c_k\left(c_1+c_2+c_{k1}+2\underset{i=3}{\overset{k2}{}}c_i\right)(\frac{1}{2}+\mathrm{cos}\frac{\pi }{m})1,$$
whence the reflection subgroup $`W_{\{r_d,r_{c_k}\}}`$ is infinite (by Theorem 10), a contradiction. So $`c_kd=c_{k1}c_k=\mathrm{cos}\frac{\pi }{m}`$. Since the reflection subgroup generated by $`\{r_d,r_{c_{k1}},r_{c_k}\}`$ is finite it follows that $`m=3`$. So we have shown that $`m_{c_kc_{k1}}=3`$ and $`m_{c_kc_i}=2`$ for $`i<k1`$, and since $`\{c_1,c_2,\mathrm{},c_{k1}\}`$ is a system of type $`D_{k1}`$ it follows that $`\{c_1,c_2,\mathrm{},c_k\}`$ is a system of type $`D_k`$, as claimed.
Note that $`\mathrm{Odd}(a)\{b\}`$ and $`\mathrm{Odd}(a)\{a\}`$ are both connected, since each element $`c\mathrm{Odd}(a)`$ that is is adjacent to $`a`$ is also adjacent to $`b`$, and vice versa. If there were $`c,d\mathrm{Odd}(a)\{b\}`$ with $`3<m_{cd}<\mathrm{}`$ then $`b`$ together with a minimal length odd-labelled path from $`a`$ to $`\{c,d\}`$ would yield $`c_1,c_2,\mathrm{},c_k\mathrm{\Pi }`$ satisfying (1) and (2) above and not forming a system of type $`D_k`$, contradicting the result proved above. The same argument yields a contradiction whenever $`c\mathrm{Odd}(a)\{a,b\}`$ and $`d\mathrm{\Pi }\mathrm{Odd}(a)`$ with $`3<m_{cd}<\mathrm{}`$. So all edge labels in $`\mathrm{Odd}(a)`$ are $`3`$, if $`c,d\mathrm{Odd}(a)`$ are not adjacent in $`\mathrm{Odd}(a)`$ then $`m_{cd}\{2,\mathrm{}\}`$, and if $`c\mathrm{Odd}(a)`$ and $`d\mathrm{\Pi }\mathrm{Odd}(a)`$ then $`m_{cd}\{2,\mathrm{}\}`$. Furthermore, any circuit in $`\mathrm{Odd}(a)\{b\}`$ would similarly yield a contradiction (by combining the circuit with a minimal finite-labelled path connecting it to $`b`$). So $`\mathrm{Odd}(a)\{b\}`$ is tree. Of course, $`\mathrm{Odd}(b)\{a\}`$ is also a tree, by the same argument.
For each $`c\mathrm{Odd}(a)\{a,b\}`$ let $`D[a,b..c]\mathrm{\Pi }`$ consist of $`b`$ and the unique path from $`a`$ to $`c`$ in $`\mathrm{Odd}(a)\{b\}`$. The discussion above shows that $`D[a,b..c]`$ is of type $`D`$. Now suppose that $`c\mathrm{Odd}(a)\{a,b\}`$ and $`e\mathrm{\Pi }D[a,b..c]`$ with $`m_{ce}=2`$. Write $`D[a,b..c]=\{c_1,c_2,\mathrm{},c_k\}`$, with $`c_1=b`$, $`c_2=a`$ and $`c_k=c`$, and let $`d=c_1+c_2+c_k+2_{i=3}^{k1}c_i`$. An argument similar to one used above shows that $`r_d\mathrm{FC}(c)`$, and hence $`W_{\{d,c,e\}}`$ is finite. So $`de>1`$. If $`c_ie0`$ then $`c_ie1/2`$; so it follows that $`\{ic_ie0\}`$ is a subset of $`\{1,2,k\}`$ with at most one element. But $`c_ke=0`$ since $`m_{ce}=2`$, and $`c_1e=c_2e`$ since $`m_{af}=m_{bf}`$ for all $`f\mathrm{\Pi }`$. So $`c_ie=0`$ for all $`i\{1,2,\mathrm{},k\}`$. In particular, if $`e\mathrm{\Pi }\mathrm{Odd}(a)`$ and $`m_{ce}\mathrm{}`$ for some $`c\mathrm{Odd}(a)`$ then $`m_{ce}=2`$, as shown above, and so $`m_{de}=2`$ for all $`dD[a,b..c]`$.
All that remains to prove now is that if $`c,d\mathrm{Odd}(a)\{a,b\}`$ with $`cD[a,b..d]`$ and $`dD[a,b..c]`$, then $`m_{cd}=\mathrm{}`$. Since $`c`$ and $`d`$ are not adjacent in $`\mathrm{Odd}(a)`$ the only alternative is that $`m_{cd}=2`$; so suppose, for a contradiction, that this holds. Choose the vertex $`e\mathrm{Odd}(a)\{b\}`$ on the (unique) path from $`c`$ to $`d`$ such that the distance from $`e`$ to $`a`$ is minimal. Let $`c^{},d^{}`$ be the neighbours of $`e`$ in the path from $`c`$ to $`d`$, with $`c^{}`$ between $`e`$ and $`c`$ and $`d^{}`$ between $`e`$ and $`d`$. Then $`c^{}C[b..c]`$, and since $`m_{cd}=2`$ it follows that $`m_{c^{}d}=2`$. Now since $`d^{}C[b..d]`$ and $`m_{c^{}d}=2`$ it follows that $`m_{c^{}d^{}}=2`$. Thus the set $`L\mathrm{\Pi }`$ consisting of $`c^{}`$ and $`d^{}`$ and the vertices on the path from $`a`$ to $`e`$ form a system of type $`D_k`$, or $`A_3`$ if $`e=a`$. So $`L`$ is spherical, and since $`b\mathrm{FC}(r_a)`$ it follows that $`L\{b\}`$ is spherical also. If $`L=A_3`$ then $`m_{ac^{}}=m_{ad^{}}=3`$, and since $`m_{bc^{}}=m_{ac^{}}`$ and $`m_{bd^{}}=m_{ad^{}}`$ we see that $`L\{b\}`$ is of type $`\stackrel{~}{A}_3`$, contradicting the fact that $`L\{b\}`$ is spherical. Similarly if $`L`$ is of type $`D_k`$ then $`L\{b\}`$ is of type $`\stackrel{~}{D}_k`$, again giving a contradiction. ∎
We also need to obtain further information about the situation of alternative (1) of Proposition 27. So for the next three lemmas we assume that $`aJL\mathrm{\Pi }`$ with $`L`$ spherical and $`W_J=\mathrm{FC}(r_a)`$, and there exist $`bJ`$ and $`cLJ`$ such that $`L_0=\{a,c,b\}`$ is a component of $`L`$ of type $`C_3`$, with $`m_{ac}=3`$ and $`m_{cb}=4`$.
###### Lemma 30.
For all $`e\mathrm{\Pi }\{a,c,b\}`$, either $`m_{ce}=m_{ae}=m_{be}=2`$ or $`m_{ce}=\mathrm{}`$. Moreover, $`J\mathrm{Odd}(a)=\{a\}`$.
###### Proof.
If $`J\mathrm{Odd}(a)\{a\}`$ then, since $`\{a\}`$ is a component of $`J`$, Proposition 29 applies, and it follows in particular that no vertex in $`\mathrm{Odd}(a)`$ lies on an edge with finite label different from 3. This contradicts $`m_{bc}=4`$. So $`J\mathrm{Odd}(a)=\{a\}`$.
Suppose that $`e\mathrm{\Pi }\{a,c,b\}`$ with $`m_{ce}<\mathrm{}`$. The group $`r_cr_aW_{\{c,e\}}r_ar_c`$ is finite and contains $`r_cr_aW_{\{c\}}r_ar_c=W_{\{a\}}`$; so there exists a maximal finite subgroup $`G`$ of $`W`$ containing $`r_a`$ and the reflection along $`(r_cr_a)e`$. Since $`r_b\mathrm{FC}(r_a)G`$ it follows that $`W_{\{b,(r_cr_a)e\}}`$ is finite, and hence so is $`W_{\{(r_ar_c)b,e\}}=r_ar_cW_{\{b,(r_cr_a)e\}}r_cr_a`$. Hence
(4.1)
$$(b+\sqrt{2}c+\sqrt{2}a)e=(r_ar_c)be>1.$$
Assume, for a contradiction, that $`m_{ce}2`$. Then $`ce1/2<1/2\sqrt{2}`$, and so $`(b+\sqrt{2}a)e>1/2`$, giving a contradiction if either $`m_{be}2`$ or $`m_{ae}2`$. So $`be=ae=0`$, and the inequality 4.1 above gives $`ce>1/\sqrt{2}`$. So $`m_{ce}=3`$. But now $`W_{\{a,c,e\}}`$ is of type $`A_3`$, hence finite, and hence contained in a maximal finite subgroup that also contains $`\mathrm{FC}(r_a)=W_J`$. Since $`bJ`$ it follows that $`\{a,c,e,b\}`$ is spherical, which is false since it is of type $`\stackrel{~}{B}_3`$. So $`m_{ce}=2`$. and it remains to show that $`m_{ae}=m_{be}=2`$.
Since $`ce=0`$ we deduce from 4.1 that $`(b+\sqrt{2}a)e>1`$, and in particular it follows that $`m_{ae}`$ is 2 or 3. In either case $`\{e,a,c\}`$ is spherical (of type $`A_3`$ or $`A_1+A_2`$), and so $`\{e,a,c,b\}`$ is also spherical (since $`r_b\mathrm{FC}(r_a)`$). If either $`m_{ae}2`$ or $`m_{be}2`$ then applying Proposition 27 with $`L^{\prime \prime }=J\{e,c\}`$ in place of $`L`$ yields a contradiction, since if $`L_0^{\prime \prime }`$ is the the component of $`L^{\prime \prime }`$ containing $`a`$ then $`\{a,b\}L_0^{\prime \prime }JL_0^{\prime \prime }`$ (since $`cJ`$), but $`L_0^{\prime \prime }`$ is not of type $`C_3`$ or $`D_n`$ since it contains $`\{a,c,b,e\}`$. So $`m_{ae}=m_{be}=2`$, as required. ∎
###### Lemma 31.
Let $`J^{}=J\{a\}`$ and let $`d\mathrm{Odd}(a)`$. Then $`m_{cd}2`$. If $`m_{db^{}}2`$ for some $`b^{}J^{}`$ then $`m_{db^{}}=4`$, and there is a unique $`a^{}`$ adjacent to $`d`$ in $`\mathrm{Odd}(a)`$; moreover, $`\{a^{},d,b^{}\}`$ is of type $`C_3`$, and $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$. On the other hand, if $`m_{db^{}}=2`$ for all $`b^{}J^{}`$ then $`\mathrm{FC}(r_d)=W_{J^{}\{d\}}`$.
###### Proof.
We use induction on the distance from $`d`$ to $`a`$ in $`\mathrm{Odd}(a)`$. Observe that if $`d=a`$ then $`m_{db^{}}=2`$ for all $`b^{}J^{}`$, since $`\{a\}`$ is a component of $`J`$, and we have $`\mathrm{FC}(r_d)=W_J=W_{J^{}\{d\}}`$ and $`m_{cd}=32`$, as required.
Suppose now that $`da`$, and let $`a=d_1,d_2,\mathrm{},d_k=d`$ be a minimal length path from $`a`$ to $`d`$ in $`\mathrm{Odd}(a)`$. If $`2ik1`$ then $`d_i`$ does not have valency 1 in $`\mathrm{Odd}(a)`$, and so $`m_{d_ib^{}}=2`$ for all $`b^{}J^{}`$, by the inductive hypothesis. The same is true for $`i=1`$, since $`\{a\}`$ is a component of $`J`$.
We prove first that $`m_{cd}2`$. Assuming, for a contradiction, that $`m_{cd}=2`$, then clearly $`d\{a,c,b\}`$, and Lemma 30 tells us that $`m_{ad}=2`$ and $`m_{bd}=2`$. It follows that $`m_{bf}=2`$ for all $`f`$ in the set $`M=\{d_1,d_2,\mathrm{},d_k\}`$, since $`\{d_i\}`$ is a component of $`J^{}\{d_i\}`$ when $`1ik1`$. So $`wb=b`$ for all $`wW_M`$. Furthermore, since $`d`$ and $`a`$ lie in the same connected component of $`\mathrm{\Omega }(\mathrm{\Pi })`$, we can choose $`wW_M`$ such that $`wd=a`$. Now since $`wca=cd=2`$ we see that the reflection $`r_{wc}`$ centralizes $`r_a`$, and hence normalizes $`\mathrm{FC}(r_a)=W_J`$. By Lemma 11 it follows that either $`wc\mathrm{\Phi }_J`$ or $`wce=0`$ for all $`eJ`$. But $`wcb=cw^1b=cb0`$; so we must have $`wc\mathrm{\Phi }_J`$, and hence $`cw^1\mathrm{\Phi }_J\mathrm{\Phi }_{JM}`$. So $`cM`$, contradicting $`m_{bf}=2`$ for all $`fM`$. So $`m_{cd}2`$.
Write $`a^{}=d_{k1}`$ and $`\stackrel{~}{J}=J^{}\{a^{}\}`$. Note that $`\{a^{}\}`$ is a component of $`\stackrel{~}{J}`$, and $`\mathrm{FC}(r_a^{})=W_{\stackrel{~}{J}}`$ (by the inductive hypothesis). Now since $`\{d,a^{}\}`$ is spherical, $`\stackrel{~}{L}=\stackrel{~}{J}\{d\}`$ is spherical also. Let $`\stackrel{~}{L}_0`$ be the component of $`\stackrel{~}{L}`$ containing $`a^{}`$.
Consider first the case that $`m_{db^{}}=2`$ for all $`b^{}J^{}`$. Since also $`m_{a^{}b^{}}=2`$ for all $`b^{}J^{}`$, it follows that $`r_d`$ and $`r_a^{}`$ both fix all elements of $`J^{}`$. Since $`v=v[d,\{a^{}\}]W_{\{a^{},d\}}`$ satisfies $`va^{}=d`$, we conclude that
$$\mathrm{FC}(r_d)=v\mathrm{FC}(r_a^{})v^1=vW_{J^{}\{a^{}\}}v^1=W_{vJ^{}\{va^{}\}}=W_{J^{}\{d\}},$$
as required.
Now suppose that $`m_{db^{}}2`$ for some $`b^{}J^{}`$. Then $`b^{}\stackrel{~}{J}\stackrel{~}{L}_0`$, and so $`\{a^{}\}\stackrel{~}{J}\stackrel{~}{L}_0\stackrel{~}{L}_0`$. Applying Proposition 27, we see that the situation of alternative (1) must hold: alternative (2) is ruled out since $`a^{}`$ is a component of $`\stackrel{~}{J}`$, and alternative (3) is ruled out since $`J^{}\mathrm{Odd}(a)=\mathrm{}`$. Hence $`\stackrel{~}{L}_0=\{a^{},d,b^{}\}`$ is of type $`C_3`$, with $`m_{db^{}}=4`$ and $`m_{da^{}}=3`$. Furthermore, Lemma 30 tells us that $`m_{de}\{2,\mathrm{}\}`$ for all $`e\mathrm{\Pi }\{a^{},d,b^{}\}`$; so $`a^{}`$ is the unique neighbour of $`d`$ in $`\mathrm{Odd}(a)`$, as required. ∎
###### Lemma 32.
Let $`e\mathrm{EOdd}(a)\mathrm{Odd}(a)`$ with $`eb`$. Then $`m_{be}=2`$.
###### Proof.
If $`eJ`$ then it is clear that $`m_{be}=2`$, since $`b`$ is a component of $`J`$. So we may assume that $`eJ`$.
As above, write $`J^{}=J\{a\}`$. Since $`e\mathrm{EOdd}(a)`$ there exists a $`d\mathrm{Odd}(a)`$ with $`m_{de}`$ even. If $`d`$ is adjacent to some $`b^{}J^{}`$ then, by Lemma 31, there is a unique $`a^{}\mathrm{Odd}(a)`$ adjacent to $`d`$; furthermore, $`\{a^{},d,b^{}\}`$ is of type $`C_3`$, and $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$. By Lemma 30, since $`m_{de}\mathrm{}`$ it follows that $`m_{a^{}e}=m_{de}=2`$. On the other hand, if $`d`$ is not adjacent to any element of $`J^{}`$ then Lemma 31 tells us that $`\mathrm{FC}(r_d)=W_{J^{}\{d\}}`$. So in either case there is an $`a^{}\mathrm{Odd}(a)`$ with $`m_{a^{}e}`$ even and $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$.
Choose such an $`a^{}`$. Since $`m_{a^{}e}`$ is even, $`v[e,\{a^{}\}]a^{}=a^{}`$; moreover $`v[e,\{a^{}\}]`$ is the reflection along some root $`f=\lambda e+\mu a^{}`$. Note that $`fa^{}=0`$, and hence $`\lambda 0`$. Since $`eJ`$ it follows that $`f\mathrm{\Phi }_J`$. But $`r_f`$ centralizes $`r_a^{}`$, and hence normalizes $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$. By Lemma 11 it follows that $`fb=0`$. But also $`a^{}b=0`$, since $`\{a^{}\}`$ and $`\{b\}`$ are distinct components of $`J^{}\{a^{}\}`$; so it follows that $`eb=0`$. Thus $`m_{be}=2`$, as required. ∎
Lemmas 30, 31 and 32 combine to yield the following result.
###### Proposition 33.
Suppose that $`aJL\mathrm{\Pi }`$, with $`L`$ spherical and $`\mathrm{FC}(r_a)=W_J`$, and let $`L_0`$ be the component of $`L`$ containing $`a`$. Suppose that $`L_0=\{a,c,b\}`$ is of type $`C_3`$, with $`m_{ac}=3`$ and $`m_{bc}=4`$, and $`JL_0=\{a,b\}`$. Then $`b`$ is a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$. Furthermore, $`J\mathrm{Odd}(a)=\{a\}`$, and if $`a^{}\mathrm{Odd}(a)`$ is not adjacent to any $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$ then $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$, where $`J^{}=J\{a\}`$.
###### Proof.
If $`m_{bd}2`$ for some $`d\mathrm{Odd}(a)`$, then $`m_{db}=4`$, by Lemma 31. There is at least one $`d\mathrm{Odd}(a)`$ such that $`m_{bd}=4`$, namely $`d=c`$. Lemma 31 tells us that for each $`d\mathrm{Odd}(a)`$ with $`m_{bd}=4`$ there is an $`a^{}\mathrm{Odd}(a)`$ such that $`\{a^{},d,b^{}\}`$ is a system of type $`C_3`$. Moreover, by Lemma 30, if $`e\mathrm{\Pi }(\mathrm{Odd}(a)\{b\})`$ then either $`m_{de}=\mathrm{}`$ or $`m_{ae}=m_{be}=m_{de}=2`$, while if $`e\mathrm{Odd}(a)\{a,c\}`$ then $`m_{de}=\mathrm{}`$, since $`m_{de}2`$ by Lemma 31. And if $`e\mathrm{EOdd}(a)(\mathrm{Odd}(a)\{b\})`$ then $`m_{be}=2`$, by Lemma 32. So $`b`$ satisfies all the requirements of a $`C_3`$-neighbour of $`M=\mathrm{Odd}(a)`$, as specified in Definition 4.
It now follows from Lemma 31 that if $`a^{}\mathrm{Odd}(a)`$ is adjacent to some $`b^{}J^{}`$ then $`b^{}`$ is a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$, and if $`a^{}`$ is not adjacent to any such $`b^{}`$ then $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$. Finally, $`J\mathrm{Odd}(a)=\{a\}`$, by Lemma 30. ∎
Let $`a,a^{}\mathrm{\Pi }`$, and suppose that $`wW`$ has the property that $`wa=a^{}`$. By Proposition 12 there exist $`a_i\mathrm{Odd}(a)`$ and $`c_i\mathrm{\Pi }`$ such that
* $`a_1=a`$ and $`a_{k+1}=a^{}`$,
* $`m_{c_ia_i}\mathrm{}`$ and $`v[c_i,\{a_i\}]a_i=a_{i+1}`$, for all $`i\{1,2,\mathrm{},k\}`$,
* $`w=v[c_k,\{a_k\}]\mathrm{}v[c_2,\{a_2\}]v[c_1,\{a_1\}]`$.
Now let $`b`$ be a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$. For each $`c\mathrm{Odd}(a)`$ that is adjacent to $`b`$, define $`X(c)=b+\sqrt{2}c+\sqrt{2}\stackrel{~}{a}`$, where $`\stackrel{~}{a}`$ is the unique neighbour of $`c`$ in $`\mathrm{Odd}(a)`$, and for each $`c\mathrm{Odd}(a)`$ that is not adjacent to $`b`$, define $`X(c)=b`$. We show that $`v[c_i,\{a_i\}]X(a_i)=X(a_{i+1})`$ for all $`i\{1,2,\mathrm{}k\}`$.
Suppose first that neither $`c_i`$ nor $`a_i`$ is adjacent to $`b`$. Then $`X(a_i)=b`$, and since $`a_{i+1}\{a_i,c_i\}`$ we have that $`X(a_{i+1})=b`$ also. Since $`r_{a_i}`$ and $`r_{c_i}`$ both fix $`b`$, and $`v[c_i,\{a_i\}]W_{\{a_i,c_i\}}`$, it follows that
$$v[c_i,\{a_i\}]X(a_i)=v[c_i,\{a_i\}]b=b=X(a_{i+1}),$$
as required.
Next, suppose that $`c_i`$ is adjacent to $`b`$, but $`a_i`$ is not adjacent to $`b`$. Since $`m_{c_ia_i}\mathrm{}`$ and $`a\mathrm{Odd}(a)`$ it follows that $`c_i\mathrm{EOdd}(a)`$. Since $`b`$ is a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$, it is not adjacent to any element of $`\mathrm{EOdd}(a)\mathrm{Odd}(a)`$; so $`c_i\mathrm{Odd}(a)`$, and, moreover, $`m_{c_ie}=\mathrm{}`$ for all $`e\mathrm{Odd}(a)\{c_i\}`$ apart from the unique neighbour of $`c_i`$ in $`\mathrm{Odd}(a)`$. So $`a_i`$ is this unique neighbour, $`m_{c_ia_i}=3`$, and $`a_{i+1}=v[c_i,\{a_i\}]a_i=c_i`$. Moreover, $`m_{c_ib}=4`$ and $`m_{a_ib}=2`$. So
$$v[c_i,\{a_i\}]X(a_i)=r_{a_i}r_{c_i}b=b+\sqrt{2}c_i+\sqrt{2}a_i=X(c_i)=X(a_{i+1})$$
as required.
Now suppose that $`a_i`$ is adjacent to $`b`$, and let $`\stackrel{~}{a}`$ be the unique neighbour of $`a_i`$ in $`\mathrm{Odd}(a)`$. Since $`m_{a_ie}=\mathrm{}`$ for all $`e\mathrm{Odd}(a)\{a_i,\stackrel{~}{a}\}`$, if $`c_i\mathrm{Odd}(a)`$ then $`c_i=\stackrel{~}{a}`$. In this case we see that
$$v[c_i,\{a_i\}]X(a_i)=r_{a_i}r_{c_i}(b+\sqrt{2}c_i+\sqrt{2}a_i)=b=X(c_i)=X(a_{i+1}),$$
since $`a_{i+1}=v[c_i,\{a_i\}]a_i=c_i`$. If $`c_i=b`$ then $`v[c_i,\{a_i\}]=r_br_{a_i}r_b`$, which fixes both $`a_i`$ and $`X(a_i)=b+\sqrt{2}a_i+\sqrt{2}\stackrel{~}{a}`$. So $`v[c_i,\{a_i\}]X(a_i)=X(a_{i+1})`$ in this case too. Finally, suppose that $`c_i\mathrm{Odd}(a)\{b\}`$. Since $`m_{c_ia_i}\mathrm{}`$ we must have $`m_{c_i\stackrel{~}{a}}=m_{c_ia_i}=m_{c_ib}=2`$, (by the definition of a $`C_3`$-neighbour). So
$$v[c_i,\{a_i\}]X(a_i)=r_{c_i}(b+\sqrt{2}c_i+\sqrt{2}\stackrel{~}{a})=b+\sqrt{2}c_i+\sqrt{2}\stackrel{~}{a}=X(a_{i+1})$$
since $`a_{i+1}=r_{c_i}a_i=a_i`$.
We have now now covered all cases, and shown that $`v[c_i,\{a_i\}]X(a_i)=X(a_{i+1})`$ for all $`i\{1,2,\mathrm{}k\}`$. By a trivial induction it follows that $`X(a_{k+1})=wX(a_1)`$.
Thus we have established the following result.
###### Lemma 34.
Let $`a\mathrm{\Pi }`$ and $`wW`$ such that $`wa\mathrm{\Pi }`$. Suppose that $`b`$ is a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$ that is not adjacent to $`a`$. Then
$$wb=\{\begin{array}{cc}b\hfill & \text{if }wa\text{ is not adjacent to }b\hfill \\ b+\sqrt{2}wa+\sqrt{2}\stackrel{~}{a}\hfill & \text{if }wa\text{ is adjacent to }b\hfill \end{array}$$
where $`\stackrel{~}{a}`$ is adjacent to $`wa`$ in $`\mathrm{Odd}(a)`$.
We are now able to give a detailed description of the components of $`J`$ whenever $`W_J`$ is the finite continuation of a simple reflection.
###### Proposition 35.
Suppose that $`aJ\mathrm{\Pi }`$ with $`W_J=\mathrm{FC}(r_a)`$, and suppose that $`K`$ is a component of $`J`$. Then one of the following alternatives holds.
* $`K=\{a\}=J\mathrm{Odd}(a)`$.
* $`K=\{a,b\}`$ is of type $`C_2`$, and $`J\mathrm{Odd}(a)=\{a\}`$.
* $`K=\{a\}`$ or $`K=\{b\}`$, where $`\{a,b\}=J\mathrm{Odd}(a)`$ is of type $`A_1+A_1`$.
* $`K=\{b\}\mathrm{Odd}(a)`$, and $`b`$ is a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$.
* $`\mathrm{Odd}(a)K`$, and $`K`$ is a component of $`\mathrm{EOdd}(a)`$.
* $`K\mathrm{Odd}(a)=\mathrm{}`$, and $`K`$ is a component of $`\mathrm{EOdd}(a)`$.
###### Proof.
We consider first the case that $`K\mathrm{Odd}(a)\mathrm{}`$, and start by supposing that there exists a spherical $`L\mathrm{\Pi }`$ with $`JL`$ and $`K`$ not a component of $`L`$.
Choose such an $`L`$, and let $`L_0`$ be the component of $`L`$ containing $`a`$. By Proposition 27, since $`K`$ is not a component of $`L`$ we must have $`KL_0`$. So either $`K=\{a\}`$, in which case (a) above holds, or else $`\{a\}\{a\}KJL_0`$. Furthermore, $`JL_0L_0`$, since $`KL_0`$. So if (a) does not hold then $`\{a\}JL_0L_0`$, and so one of the alternatives (1), (2) or (3) of Proposition 27 must hold.
Suppose that alternative (2) holds, so that $`K=\{a,b\}=JL_0`$ for some $`b`$, and $`\{a,b\}`$ is of type $`C_2`$. By Proposition 28 we see that each $`c\mathrm{Odd}(a)\{a\}`$ lies in a type $`C`$ spherical subset $`L^{}`$ of $`\mathrm{\Pi }`$ containing $`\{a,b\}`$. Since $`JL^{}=\{a,b\}`$ (by Proposition 27) it follows that $`cJ`$. So $`J\mathrm{Odd}(a)=\{a\}`$, and (b) above is satisfied.
Suppose that alternative (3) of Proposition 27 holds, so that $`JL_0=\{a,b\}`$ is of type $`A_1+A_1`$, and $`b\mathrm{Odd}(a)`$. Proposition 29 immmediately yields that $`J\mathrm{Odd}(a)=\{a,b\}`$, and so (c) above is satisfied.
Suppose that alternative (1) of Proposition 27 holds, so that $`L_0=\{a,c,b\}`$ with $`m_{ac}=3`$ and $`m_{cb}=4`$, and $`JL_0=\{a,b\}`$. By Lemma 30 we know that $`b\mathrm{Odd}(a)`$, and since we have assumed that $`K\mathrm{Odd}(a)\mathrm{}`$, it follows that $`K=\{a\}=J\mathrm{Odd}(a)`$. So (a) holds.
We have now dealt with all cases that arise if there is a spherical $`L\mathrm{\Pi }`$ with $`JL`$ and $`K`$ not a component of $`L`$. So assume that $`K`$ is a component of every spherical $`L`$ containing $`J`$. We show that in this case $`\mathrm{Odd}(a)K`$, and $`K`$ is a component of $`\mathrm{EOdd}(a)`$; that is, (e) above holds.
To show that $`\mathrm{Odd}(a)K`$ it is clearly sufficient to show that if $`a^{}K\mathrm{Odd}(a)`$ and $`b`$ is adjacent to $`a^{}`$ in $`\mathrm{Odd}(a)`$ then $`bK`$. Note that since $`a^{}\mathrm{Odd}(a)`$ there exists $`wW`$ with $`a^{}=wa`$, and since Proposition 25 yields that $`wN_W(W_J)`$ it follows that $`\mathrm{FC}(r_a^{})=W_J`$. Now the assumption that $`b`$ and $`a^{}`$ are adjacent in $`\mathrm{Odd}(a)`$ implies that $`\{a^{},b\}`$ is spherical, and therefore $`J\{b\}`$ is spherical. But $`K`$ is a component of every spherical subset of $`\mathrm{\Pi }`$ containing $`J`$; so it is a component of $`J\{b\}`$. But $`a^{}K`$ and $`b`$ is adjacent to $`a^{}`$; so $`bK`$, as required.
Since $`KJ\mathrm{EOdd}(a)`$ and $`K`$ is connected, saying that $`K`$ is a component of $`\mathrm{EOdd}(a)`$ is equivalent to saying that $`m_{bc}=2`$ whenever $`bK`$ and $`c\mathrm{EOdd}(a)K`$. So suppose that $`c\mathrm{EOdd}(a)K`$. Then there exists an $`a^{}\mathrm{Odd}(a)`$ such that $`m_{a^{}c}`$ is even. Thus $`\{a^{},c\}`$ is spherical, and as above it follows that $`J\{c\}`$ is spherical. So $`K`$ must be a component of $`J\{c\}`$, and since $`cK`$ it follows that $`m_{bc}=2`$ for all $`bK`$, as required.
It remains to consider the case that $`K\mathrm{Odd}(a)=\mathrm{}`$; we must show that either (f) or (d) holds. We start by supposing that there exists a spherical $`L\mathrm{\Pi }`$ and a $`wW`$ with $`wJL`$ and $`wK`$ not a component of $`L`$.
Choose such $`L`$ and $`w`$, and let $`L_0`$ be the component of $`L`$ containing $`wa`$. By Proposition 27, since $`wK`$ is not a component of $`L`$ we must have $`wKL_0`$. Now $`wJL_0L_0`$ since $`wKL_0`$, and $`\{wa\}wJL_0`$ since $`wawK`$. So one of the alternatives (1), (2) or (3) of Proposition 27 must hold. Alternative (3) can be ruled out, since in that case $`wJL_0\mathrm{Odd}(wa)`$, which is impossible since $`K\mathrm{Odd}(a)=\mathrm{}`$. If alternative (2) holds then $`wK=wJL_0`$ contains $`wa`$ and is of type $`C_2`$, whence $`K`$ contains $`a`$ and is of type $`C_2`$, and (b) is satisfied. If alternative (1) holds then since $`wK\{wa\}`$ it follows from Proposition 32 that $`wK=\{b\}`$, with $`b`$ a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$. Since $`wa`$ is not adjacent to $`b`$, it follows from Lemma 34 that $`w^1b=b`$, unless $`a`$ is adjacent to $`b`$, in which case $`w^1b=b+\sqrt{2}a+\sqrt{2}\stackrel{~}{a}`$ for some $`\stackrel{~}{a}`$ in $`\mathrm{Odd}(a)`$. But this latter case cannot occur, since $`w^1bK\mathrm{\Pi }`$. So $`K=wK=\{b\}`$, with $`b`$ a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$, and (d) holds.
Finally, suppose that $`wK`$ is a component of every spherical $`L\mathrm{\Pi }`$ such that $`wJL`$ for some $`wW`$. For each $`c\mathrm{EOdd}(a)K`$ there is then a sequence $`a=a_0,a_1,\mathrm{},a_k=c`$ in $`\mathrm{\Pi }`$ such that $`m_{a_{i1}a_i}`$ finite for all $`i\{1,2,\mathrm{},k\}`$ and odd for all $`i\{1,2,\mathrm{},k1\}`$. We shall show that, for every such sequence, $`m_{ba_i}=2`$ for all $`bK`$ and $`i\{0,1,\mathrm{},k\}`$; in particular, this will show that $`m_{bc}=2`$ whenever $`bK`$ and $`c\mathrm{EOdd}(a)K`$, enabling us to conclude that $`K`$ is a component of $`\mathrm{EOdd}(a)`$.
The case $`k=0`$ is clear, since $`aJK`$ and $`K`$ is a component of $`J`$. Proceeding by induction, we may assume that $`k>0`$ and $`m_{ba_i}=2`$ for all $`i\{1,2,\mathrm{},k1\}`$ and all $`bK`$. We see that the element $`u=v[a_{k1},\{a_{k2}\}]v[a_{k2},\{a_{k3}\}]\mathrm{}v[a_1,\{a_0\}]`$ centralizes $`W_K`$ and has the property that $`ua=a_{k1}`$, since the labels in the path from $`a`$ to $`a_{k1}`$ are all odd. The group $`u^1W_{\{a_{k1},a_k\}}u`$ is finite and contains $`u^1r_{a_{k1}}u=r_a`$, and so there is a maximal finite subgroup $`G`$ of $`W`$ containing this group and also containing $`W_J`$.
Note that $`W_J\{u^1r_cu\}G=w^1W_Lw`$, for some $`wW`$ and spherical $`L\mathrm{\Pi }`$, the element $`u`$ being in the centralizer of $`W_K`$. We may choose $`w`$ to be the minimal length element of $`W_Lw=W_LwW_J`$, and it then follows from Lemma 15 that $`wJL`$. Hence $`wK`$ is a component of $`L`$. Furthermore, since $`wu^1r_cuw^1W_L`$ we see that the root $`wu^1c`$ is in $`\mathrm{\Phi }_L`$ and not in $`\mathrm{\Phi }_{wK}=wu^1\mathrm{\Phi }_K`$ (since $`c\mathrm{\Phi }_K`$). So $`wu^1cwu^1b=0`$ for all $`bK`$. So $`cb=0`$, or (equivalently) $`m_{bc}=2`$ for all $`bK`$, as required. ∎
To complement the results we have obtained so far, our next task is to find conditions that ensure that a visible subgroup $`W_K`$ is contained in $`\mathrm{FC}(r_a)`$.
###### Lemma 36.
Let $`a\mathrm{\Pi }`$ and $`K`$ a component of $`\mathrm{EOdd}(a)`$ such that $`W_K`$ is finite. Then $`W_K\mathrm{FC}(r_a)`$.
###### Proof.
Let $`F`$ be a maximal finite subgroup of $`W`$ with $`r_aF`$, and choose $`wW`$ such that $`wFw^1=W_L`$ for some $`L\mathrm{\Pi }`$. We may replace $`w`$ by the minimal length element in the double coset $`W_LwW_{\{a\}}`$, since this does not affect the condition $`wFw^1=W_L`$. So we have that $`w^1L\mathrm{\Phi }^+`$, and, moreover, $`r_aw^1LwW_{\{a\}}=W_{w^1L\{a\}}`$ by Lemma 15. So $`waL\mathrm{\Pi }`$, and by Lemma 12 we see that $`w`$ is a product of factors of the form $`v[d,\{c\}]`$, with $`c,d\mathrm{EOdd}(a)`$. Since $`K`$ is a component of $`\mathrm{EOdd}(a)`$ it follows that each $`v[d,\{c\}]`$ normalizes $`W_K`$, and therefore $`w`$ normalizes $`W_K`$. Moreover, since $`waL`$ and $`L`$ is spherical, it follows that $`L\mathrm{EOdd}(wa)=\mathrm{EOdd}(a)`$. So $`W_L`$ normalizes $`W_K`$. But $`W_K`$ is finite, by hypothesis, and $`W_L`$ is a maximal finite subgroup of $`W`$. So $`W_KW_L`$, and $`W_K=w^1W_Kww^1W_Lw=F`$. Thus $`W_K`$ is contained in all maximal finite subgroups of $`W`$ containing $`r_a`$, as required. ∎
###### Lemma 37.
Let $`a\mathrm{\Pi }`$ and let $`b`$ be a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$. If $`a`$ and $`b`$ are not adjacent in $`\mathrm{\Pi }`$ then $`r_b\mathrm{FC}(r_a)`$.
###### Proof.
Let $`F`$ be a maximal finite subgroup of $`W`$ with $`r_aF`$. As in the proof of Lemma 36 there exist a $`wW`$ and a maximal spherical $`L\mathrm{\Pi }`$ with $`wa=a^{}L`$ and $`F=w^1W_Lw`$. Since $`L`$ is spherical, $`L\mathrm{EOdd}(a)`$.
Suppose first that $`a^{}`$ is not adjacent to $`b`$. Then $`m_{ca^{}}=\mathrm{}`$ for every $`c\mathrm{Odd}(a)`$ that is adjacent to $`b`$, and since $`a^{}L`$ it follows that no such $`c`$ is in $`L`$. Thus $`m_{be}=2`$ for all $`eL\mathrm{Odd}(a)`$. But since also $`m_{be}=2`$ for all $`e\mathrm{EOdd}(a)(\mathrm{Odd}(a)\{b\})`$, it follows that $`m_{be}=2`$ for all $`eL\{b\}`$. Thus $`\{b\}`$ is a component of $`L\{b\}`$, and since $`L`$ is spherical it follows that $`L\{b\}`$ is spherical. Maximality of $`L`$ tells us that $`bL`$. Moreover, Lemma 34 gives $`wb=b`$, and so $`r_b=w^1r_bww^1W_Lw=F`$.
On the other hand, suppose that $`a^{}`$ is adjacent to $`b`$. In this case Lemma 34 gives $`wb=b+\sqrt{2}a^{}+\sqrt{2}\stackrel{~}{a}`$, where $`\stackrel{~}{a}`$ is the unique neighbour of $`a^{}`$ in $`\mathrm{Odd}(a)`$. Furthermore, since $`m_{a^{}e}\{2.\mathrm{}\}`$ for all $`e\mathrm{\Pi }\{\stackrel{~}{a},a^{},b\}`$, we see that $`m_{a^{}e}=2`$ for all $`eL\{\stackrel{~}{a},a^{},b\}`$ (since $`L`$ is spherical). But the definition of a $`C_3`$-vertex also requires that $`m_{\stackrel{~}{a}e}=m_{be}=2`$ whenever $`m_{a^{}e}=2`$; so it follows that $`\{\stackrel{~}{a},a^{},b\}`$ is a component of $`L\{\stackrel{~}{a},a^{},b\}`$, which is therefore spherical since $`L`$ and $`\{\stackrel{~}{a},a^{},b\}`$ are both spherical. Maximality of $`L`$ tells us that $`\{\stackrel{~}{a},a^{},b\}L`$; so $`wb=b+\sqrt{2}a^{}+\sqrt{2}\stackrel{~}{a}\mathrm{\Phi }_L`$, and $`r_b=w^1r_{wb}ww^1W_Lw=F`$.
So $`r_bF`$ in all cases, and so $`r_b`$ is contained in all maximal finite subgroups of $`W`$ containing $`r_a`$, as required. ∎
We now prove the converse to Proposition 28.
###### Proposition 38.
Let $`a\mathrm{\Pi }`$ and $`b\mathrm{\Pi }\mathrm{Odd}(a)`$, and suppose that $`(a,b)`$ is a focus of $`\mathrm{Odd}(a)`$ in $`\mathrm{\Pi }`$. Then $`\mathrm{FC}(r_a)=W_J`$, where $`J`$ is the union of $`\{a,b\}`$ and the spherical components of $`\mathrm{EOdd}(a)`$. Moreover, $`\mathrm{FC}(r_a^{})`$ is not visible for any $`a^{}\mathrm{Odd}(a)\{a\}`$.
###### Proof.
For each $`c\mathrm{Odd}(a)`$ let $`X(c)=b+\sqrt{2}_{i=1}^mc_i`$ and $`Y(c)=b+\sqrt{2}_{i=1}^{m1}c_i`$, where $`c_1=a,c_2,\mathrm{},c_m=c`$ is the unique path from $`a`$ to $`c`$ in $`\mathrm{Odd}(a)`$, noting that $`X(c)`$ and $`Y(c)`$ are roots in $`\mathrm{\Phi }_{C[b..c]}`$. We remark, for later use, that $`X(c)`$ and $`Y(c)`$ are fixed by the reflections $`r_b,r_{c_1},\mathrm{},r_{c_{m2}}`$.
Let $`F=w^1W_Lw`$ be a maximal finite subgroup of $`W`$ containing $`r_a`$, with $`L\mathrm{\Pi }`$ and $`w`$ of minimal length in $`W_LwW_a`$. Then $`wa=a^{}L`$, by Lemma 15. Put $`L_0=L\mathrm{Odd}(a)`$.
Choose $`cL_0`$ with $`C[b..c]`$ of maximal cardinality. If $`dL_0`$ then $`m_{cd}\mathrm{}`$ (since $`L_0`$ is spherical), whence $`dC[b..c]`$ by condition (3) of Definition 5. So $`L_0C[b..c]`$. Now if $`eLL_0`$ is arbitrary then $`e\mathrm{Odd}(a)`$ (since $`eL\mathrm{Odd}(a)`$) and $`m_{ce}<\mathrm{}`$ (since $`c,eL`$ and $`L`$ is spherical). By condition (4) of Definition 5 it follows that $`m_{de}=2`$ for all $`dC[b..c]`$. Since this holds for all $`eLL_0`$, and $`C[b..c]`$ and $`LL_0`$ are both spherical, it follows that $`C[b..c](LL_0)`$ is spherical. But this set contains $`L`$ (since $`L_0C[b..c]`$) and since $`L`$ is a maximal spherical subset of $`\mathrm{\Pi }`$ we conclude that $`L=C[b..c](LL_0)`$.
By Proposition 12 and Lemma 9 there exist simple roots $`e_1,e_2,\mathrm{},e_k`$ and $`d_1=a,d_2,\mathrm{},d_{k+1}=a^{}\mathrm{Odd}(a)`$ with $`w=v[e_k,\{d_k\}]\mathrm{}v[e_2,\{d_2\}]v[e_1,\{d_1\}]`$ and $`v[e_i,\{d_i\}]d_i=d_{i+1}`$ for all $`i\{1,2,\mathrm{},k\}`$. Moreover, $`m_{e_id_i}<\mathrm{}`$ for all $`i`$. Let $`w_0=1`$ and $`w_i=v[e_i,\{d_i\}]w_{i1}`$; we will show that
$$\{w_ib,w_ib,w_i(b+\sqrt{2}a),w_i(b+\sqrt{2}a)\}=\{X(d_{i+1}),X(d_{i+1}),Y(d_{i+1}),Y(d_{i+1})\}$$
for all $`i\{0,1,\mathrm{},k\}`$. The case $`i=0`$ is trivial.
Proceeding inductively, suppose that $`i>1`$ and
$$\{\pm w_{i1}b,\pm w_{i1}(b+\sqrt{2}a)\}=\{\pm X(d_i),\pm Y(d_i)\}.$$
It will be sufficient to show that $`v[e_i,\{d_i\}]X(d_i)`$ and $`v[e_i,\{d_i\}]Y(d_i)`$ both lie in the set $`\{\pm X(d_{i+1}),\pm Y(d_{i+1})\}`$.
Suppose first that $`d_i=a`$. Then $`X(d_i)=b+\sqrt{2}a`$ and $`Y(d_i)=b`$. If $`e_i\mathrm{Odd}(a)`$ then $`m_{e_id_i}`$ is even, and $`d_{i+1}=d_i=a`$. Furthermore, by condition (4) of Definition 5we have either $`\{e_i,d_i\}=\{b,a\}`$ or $`m_{e_ib}=m_{e_ia}=2`$. In the former case $`v[e_i,\{d_i\}]=v[b,\{a\}]=r_br_ar_b`$, giving $`v[e_i,\{d_i\}]b=b\sqrt{2}a=X(a)`$ and $`v[e_i,\{d_i\}](b+\sqrt{2}a)=b=Y(a)`$; in the latter case $`v[e_i,\{d_i\}]=v[e_i\{a\}]=r_{e_i}`$, giving $`v[e_i,\{d_i\}]b=b=Y(a)`$ and $`v[e_i,\{d_i\}](b+\sqrt{2}a)=b+\sqrt{2}a=X(a)`$. If $`e_i\mathrm{Odd}(a)`$ then $`aC[b..e_i]`$, and by condition (2) of Definition 5 we have $`m_{e_ib}=2`$ and either $`m_{e_ia}=2`$ or $`m_{e_ia}=3`$. If $`m_{e_ia}=2`$ then $`d_{i+1}=d_i=a`$, while if $`m_{e_ia}=3`$ then $`d_{i+1}=e_i`$. Furthermore, in former case we find that $`v[e_i,\{d_i\}]b=r_{e_i}b=b=Y(a)`$ and $`v[e_i,\{d_i\}](b+\sqrt{2}a)=b+\sqrt{2}a=X(a)`$, while in the latter case we find that $`v[e_i,\{d_i\}]b=r_ar_{e_1}b=b+\sqrt{2}a=Y(e_i)`$ and $`v[e_i,\{d_i\}](b+\sqrt{2}a+\sqrt{2}e_i)=X(e_i)`$.
Now suppose that $`d_ia`$. If $`e_i\mathrm{Odd}(a)\{b\}`$ then $`m_{e_id_i}=2`$ and $`d_{i+1}=d_i`$. Moreover, $`m_{e_id}=2`$ for all $`dC[b..d_i]`$, and so $`v[e_i,\{d_i\}]=r_{e_i}`$ fixes all the roots in $`\mathrm{\Phi }_{C[b..d_i]}`$, including $`X(d_i)=X(d_{i+1})`$ and $`Y(d_i)=Y(d_{i+1})`$. If $`e_i\mathrm{Odd}(a)\{b\}`$ and $`\{e_i,d_i\}`$ is not an edge of $`\mathrm{Odd}(a)`$ then we again have $`d_{i+1}=d_i`$ and $`v[e_i,\{d_i\}]=r_{e_i}`$. By condition (3) of Definition 5 we either have $`d_iC[b..e_i]`$ or $`e_iC[b..d_i]`$. In the former case we have $`m_{e_id}=2`$ for all $`dC[b..d_i]`$, and as above we see that $`r_{e_i}`$ fixes $`X(d_i)`$ and $`Y(d_i)`$. In the latter case the remark made at the start of the proof implies that it is still true that $`r_{e_i}`$ fixes $`X(d_i)`$ and $`Y(d_i)`$. So we have shown that when $`m_{e_id_i}=2`$ it is true that $`v[e_i,\{d_i\}]X(d_i)`$ and $`v[e_i,\{d_i\}]Y(d_i)`$ both lie in the set $`\{\pm X(d_{i+1}),\pm Y(d_{i+1})\}`$, and it remains only to consider the case that $`e_i`$ and $`d_i`$ are adjacent in $`\mathrm{Odd}(a)`$. Note that in this case $`d_{i+1}=e_i`$.
Let $`C[b..d_i]=\{b,c_1,\mathrm{},c_m\}`$ with $`c_1=a`$ and $`c_m=d_i`$, and suppose that $`e_i=c_{m1}`$ is the vertex adjacent to $`d_i`$ in $`C[b..d_i]`$. Then
$$\begin{array}{c}\text{ }v[e_i,\{d_i\}]X(d_i)=r_{c_m}r_{c_{m1}}\left(b+\sqrt{2}\underset{j=1}{\overset{m}{}}c_j\right)=r_{c_m}\left(b+\sqrt{2}\underset{j=1}{\overset{m}{}}c_j\right)\text{ }\hfill \\ \text{ }=b+\sqrt{2}\underset{j=1}{\overset{m1}{}}c_j=X(e_i),\text{ }\hfill \end{array}$$
and similarly
$$\begin{array}{c}\text{ }v[e_i,\{d_i\}]Y(d_i)=r_{c_m}r_{c_{m1}}\left(b+\sqrt{2}\underset{j=1}{\overset{m1}{}}c_j\right)=r_{c_m}\left(b+\sqrt{2}\underset{j=1}{\overset{m2}{}}c_j\right)\text{ }\hfill \\ \text{ }=b+\sqrt{2}\underset{j=1}{\overset{m2}{}}c_j=Y(e_i).\text{ }\hfill \end{array}$$
The alternative possibility is that $`d_i`$ is adjacent to $`e_i`$ in $`C[b..e_i]`$. Exactly the same calculations show that $`v[e_i,\{d_i\}]X(d_i)=X(e_i)`$ and $`v[e_i,\{d_i\}]Y(d_i)=Y(e_i)`$ in this case also.
The induction is now complete, and it follows in particular that $`wb=w_kb`$ is one of $`\pm X(a^{})`$ or $`\pm Y(a^{})`$. Hence
$$wb\mathrm{\Phi }_{C[b..a^{}]}\mathrm{\Phi }_{C[b..c]}\mathrm{\Phi }_L.$$
Thus $`wr_bw^1W_L`$, and so $`r_bw^1W_Lw=F`$. Since $`F`$ was an arbitrary maximal finite subgroup of $`W`$ containing $`r_a`$, this shows that $`r_b\mathrm{FC}(r_a)`$.
Let $`\stackrel{~}{M}`$ be the component of $`\mathrm{EOdd}(a)`$ containing $`\mathrm{Odd}(a)`$, and suppose, for a contradiction, that $`\stackrel{~}{M}`$ is spherical. Clearly $`b\stackrel{~}{M}`$, since $`m_{ba}=4`$, but $`\stackrel{~}{M}=\mathrm{Odd}(a)\{b\}`$ is not permitted, in view of condition (5) of Definition 5. So $`\stackrel{~}{M}(\mathrm{Odd}(a)\{b\}\mathrm{}`$. But for $`e\stackrel{~}{M}(\mathrm{Odd}(a)\{b\})`$ and $`c\mathrm{Odd}(a)`$ we have $`m_{ce}\mathrm{}`$, since $`\stackrel{~}{M}`$ is spherical, and by condition (4) of Definition 5 it follows that $`m_{be}=m_{ce}=2`$ for all $`c\mathrm{Odd}(a)`$. This contradicts the fact that $`\stackrel{~}{M}`$ is connected.
Now suppose that $`a^{}\mathrm{Odd}(a)`$ is such that $`\mathrm{FC}(r_a^{})=W_J`$ for some $`J\mathrm{\Pi }`$, and let $`J_0`$ be the component of $`J`$ containing $`a^{}`$. Since $`J_0\stackrel{~}{M}`$ it follows from Proposition 35 that $`J_0`$ has rank at most 2. Now since there exists $`wW_{C[b..a^{}]}`$ such that $`wa=a^{}`$ and $`wb=X(a^{})`$, and since $`r_b\mathrm{FC}(r_a)`$, it follows that $`r_{wb}=wr_bw^1\mathrm{FC}(wr_aw^1)=\mathrm{FC}(r_a^{})`$. Thus $`X(a^{})\mathrm{\Phi }_J`$, and so $`C[b..a^{}]J`$. Since $`J_0`$ has rank at most 2, this means that $`a^{}=a`$ and $`J_0=\{a,b\}`$.
It remains to prove that $`J`$ is the union of $`J_0`$ and the spherical components of $`\mathrm{EOdd}(a)`$. By Lemma 36 we know that all these components are contained in $`J`$. But if $`K`$ is any other component of $`J`$ such that $`K\mathrm{Odd}(a)=\mathrm{}`$, then by Proposition 35 we see that $`K=\{b^{}\}`$, with $`b^{}`$ a $`C_3`$-neighbour of $`\mathrm{Odd}(a)`$. Since $`b`$ is the only element of $`\mathrm{\Pi }`$ such that $`m_{bc}\{2,4\}`$ for all $`c\mathrm{Odd}(a)`$, we must have $`b^{}=b`$, contradicting the fact that the component of $`J`$ containing $`b`$ is $`J_0=\{a,b\}`$. ∎
Next, we have the converse to Proposition 29.
###### Proposition 39.
Let $`a\mathrm{\Pi }`$ and suppose that there exists a $`b\mathrm{Odd}(a)`$ such that $`\{a,b\}`$ is a half-focus of $`\mathrm{Odd}(a)`$ in $`\mathrm{\Pi }`$. Suppose also that the vertices $`\mathrm{Odd}(a)`$ do not comprise a spherical subset of $`\mathrm{\Pi }`$. Then $`\mathrm{FC}(r_a)=W_J`$, where $`J`$ is the union of $`\{a,b\}`$ and the spherical components of $`\mathrm{EOdd}(a)`$. Moreover, $`\mathrm{FC}(r_a^{})`$ is not visible for any $`a^{}\mathrm{Odd}(a)\{a,b\}`$.
###### Proof.
For each $`c\mathrm{Odd}(a)\{a,b\}`$, define
$$X(c)=b+a+c+2\underset{i=2}{\overset{m1}{}}c_i$$
where $`c_1=a,c_2,\mathrm{},c_m=c`$ is the unique path from $`a`$ to $`c`$ in $`\mathrm{Odd}(a)\{b\}`$. Then $`X(c)`$ is a root in $`\mathrm{\Phi }_{D[a,b..c]}`$ and is fixed by the reflections $`r_b,r_{c_1},\mathrm{},r_{c_{m2}}`$ and $`r_{c_m}`$. Define also $`X(a)=b`$ and $`X(b)=a`$.
Let $`F=w^1W_Lw`$ be a maximal finite subgroup of $`W`$ containing $`r_a`$, with $`L\mathrm{\Pi }`$ and $`w`$ of minimal length in $`W_LwW_a`$. Then $`wa=a^{}L`$, by Lemma 15. Put $`L_0=L\mathrm{Odd}(a)`$.
Choose $`cL_0`$ with $`D[a,b..c]`$ of maximal cardinality. If $`dL_0`$ then $`m_{cd}\mathrm{}`$ (since $`L_0`$ is spherical), whence $`dD[a,b..c]`$ by condition (4) of Definition 6. So $`L_0D[a,b..c]`$. Now if $`eLL_0`$ is arbitrary then $`e\mathrm{Odd}(a)`$ (since $`eL\mathrm{Odd}(a)`$) and $`m_{ce}<\mathrm{}`$ (since $`c,eL`$ and $`L`$ is spherical). By condition (5) of Definition 6 it follows that $`m_{de}=2`$ for all $`dD[a,b..c]`$. Since this holds for all $`eLL_0`$, and $`D[a,b..c]`$ and $`LL_0`$ are both spherical, it follows that $`D[a,b..c](LL_0)`$ is spherical. But this set contains $`L`$ (since $`L_0D[a,b..c]`$) and since $`L`$ is a maximal spherical subset of $`\mathrm{\Pi }`$ we conclude that $`L=D[a,b..c](LL_0)`$.
By Proposition 12 and Lemma 9 there exist simple roots $`e_1,e_2,\mathrm{},e_k`$ and $`d_1=a,d_2,\mathrm{},d_{k+1}=a^{}\mathrm{Odd}(a)`$ with $`w=v[e_k,\{d_k\}]\mathrm{}v[e_2,\{d_2\}]v[e_1,\{d_1\}]`$ and $`v[e_i,\{d_i\}]d_i=d_{i+1}`$ for all $`i\{1,2,\mathrm{},k\}`$. Furthermore, we have $`m_{e_id_i}<\mathrm{}`$ for all $`i`$. Let $`w_0=1`$, and $`w_i=v[e_i,\{d_i\}]w_{i1}`$ for $`i1`$. We will show that
$$\{w_ib,w_ib\}=\{X(d_{i+1}),X(d_{i+1})\}$$
for all $`i\{0,1,\mathrm{},k\}`$.
The case $`i=0`$ is trivial. Proceeding inductively, suppose that $`i>1`$ and $`w_{i1}b=\pm X(d_i)`$. It will be sufficient to show that $`v[e_i,\{d_i\}]X(d_i)=\pm X(d_{i+1})`$.
Suppose first that $`d_i=a`$, so that $`X(d_i)=b`$. If $`e_ib`$ then $`m_{e_ib}=m_{e_ia}\{2,3\}`$, since $`m_{e_ia}=m_{e_id_i}\mathrm{}`$. We also have $`m_{e_ia}=2`$ if $`e_i=b`$. In the case $`m_{e_ia}=3`$ we have $`v[e_i,\{d_i\}]=r_ar_{e_i}`$, and $`d_{i+1}=r_ar_{e_i}a=e_i`$. Furthermore,
$$v[e_i,\{d_i\}]X(d_i)=r_ar_{e_i}b=a+b+e_i=X(e_i)=X(d_{i+1}),$$
as required. In the case $`m_{e_ia}=2`$ we have $`v[e_i,\{d_i\}]=r_{e_i}`$, giving $`d_{i+1}=r_{e_i}a=a`$, and
$$v[e_i,\{d_i\}]X(d_i)=r_{e_i}b=\pm b=\pm X(d_{i+1}),$$
since either $`e_i=b`$ or $`m_{e_ib}=2`$.
The case $`d_i=b`$ is the same as the case $`d_i=a`$ with $`a`$ and $`b`$ interchanged; so suppose that $`d_i\{a,b\}`$. If $`e_i\mathrm{Odd}(a)`$ then $`m_{e_id_i}=2`$ and $`d_{i+1}=d_i`$. Moreover, $`m_{e_id}=2`$ for all $`dD[a,b..d_i]`$, and so $`v[e_i,\{d_i\}]=r_{e_i}`$ fixes all the roots in $`\mathrm{\Phi }_{D[a,b..d_i]}`$, including $`X(d_i)=X(d_{i+1})`$. If $`e_i\mathrm{Odd}(a)`$ and $`\{e_i,d_i\}`$ is not an edge of $`\mathrm{Odd}(a)`$ then we again have $`d_{i+1}=d_i`$ and $`v[e_i,\{d_i\}]=r_{e_i}`$. By condition (4) of Definition 6 we either have $`d_iD[a,b..e_i]`$ or $`e_iD[a,b..d_i]`$. In the former case we have $`m_{e_id}=2`$ for all $`dC[b..d_i]`$, and as above we see that $`r_{e_i}`$ fixes $`X(d_i)`$. In the latter case it is still true that $`r_{e_i}`$ fixes $`X(d_i)`$, since the only simple reflection of $`D[a,b..d_i]`$ that does not fix $`X(d_i)`$ is the one corresponding to the vertex adjacent to $`d_i`$. So we have shown that when $`m_{e_id_i}=2`$ it is true that $`v[e_i,\{d_i\}]X(d_i)`$ and $`v[e_i,\{d_i\}]Y(d_i)`$ both lie in the set $`\{\pm X(d_{i+1}),\pm Y(d_{i+1})\}`$, and it remains to consider the case that $`e_i`$ and $`d_i`$ are adjacent in $`\mathrm{Odd}(a)`$. Note that in this case $`d_{i+1}=e_i`$.
Let $`D[a,b..d_i]=\{b,c_1,\mathrm{},c_m\}`$ with $`c_1=a`$ and $`c_m=d_i`$. Suppose first that $`m>2`$, and suppose that $`e_i=c_{m1}`$ is the vertex adjacent to $`d_i`$ in $`D[a,b..d_i]`$. Then
$`v[e_i,\{d_i\}]X(d_i)`$ $`=r_{c_m}r_{c_{m1}}\left(b+a+c_m+2{\displaystyle \underset{j=2}{\overset{m1}{}}}c_j\right)`$
$`=r_{c_m}\left(b+a+c_m+c_{m1}+2{\displaystyle \underset{j=1}{\overset{m2}{}}}c_j\right)`$
$`=b+a+c_{m1}+2{\displaystyle \underset{j=1}{\overset{m2}{}}}c_j=X(e_i).`$
If $`m=2`$ and $`e_i=b`$ then
$$v[e_i,\{d_i\}]X(d_i)=r_br_{d_i}(a+b+d_i)=a=X(b)=X(e_i),$$
and if $`e_i=a`$ then similarly
$$v[e_i,\{d_i\}]X(d_i)=r_ar_{d_i}(a+b+d_i)=b=X(a)=X(e_i).$$
The alternative possibility is that $`d_i=c_{m1}`$ is the vertex adjacent to $`e_i=c_m`$ in $`D[a,b..e_i]=\{b,c_1,\mathrm{},c_m\}`$. We calculate that
$`v[e_i,\{d_i\}]X(d_i)`$ $`=r_{c_{m1}}r_{c_m}\left(b+a+c_{m1}+2{\displaystyle \underset{j=2}{\overset{m2}{}}}c_j\right)`$
$`=r_{c_{m1}}\left(b+a+c_m+c_{m1}+2{\displaystyle \underset{j=1}{\overset{m2}{}}}c_j\right)`$
$`=b+a+c_m+2{\displaystyle \underset{j=1}{\overset{m1}{}}}c_j=X(e_i),`$
as required.
The induction is now complete, and it follows that $`wb=w_kb=\pm X(a^{})`$. Hence
$$wb\mathrm{\Phi }_{D[a,b..a^{}]}\mathrm{\Phi }_{D[a,b..c]}\mathrm{\Phi }_L.$$
Thus $`wr_bw^1W_L`$, and so $`r_bw^1W_Lw=F`$. Since $`F`$ was an arbitrary maximal finite subgroup of $`W`$ containing $`r_a`$, this shows that $`r_b\mathrm{FC}(r_a)`$.
Note that since $`W`$ has a graph automorphism that swaps $`r_a`$ and $`r_b`$ and fixes all the other simple reflections, it must also be true that $`r_a\mathrm{FC}(r_b)`$.
Let $`\stackrel{~}{M}`$ be the component of $`\mathrm{EOdd}(a)`$ containing $`\mathrm{Odd}(a)`$, and suppose, for a contradiction, that $`\stackrel{~}{M}`$ is spherical. Note that $`\stackrel{~}{M}\mathrm{Odd}(a)`$, in view of condition (5) of Definition 5. So $`\stackrel{~}{M}\mathrm{Odd}(a)\mathrm{}`$. But for all $`e\stackrel{~}{M}\mathrm{Odd}(a)`$ and $`c\mathrm{Odd}(a)`$ we have $`m_{ce}\mathrm{}`$, since $`\stackrel{~}{M}`$ is spherical, and by conditions (1) and (5) of Definition 5 it follows that $`m_{ce}=2`$ for all $`c\mathrm{Odd}(a)`$. This contradicts the fact that $`\stackrel{~}{M}`$ is connected.
Suppose that $`a^{}\mathrm{Odd}(a)`$ is such that $`\mathrm{FC}(r_a^{})=W_J`$ for some $`J\mathrm{\Pi }`$, and let $`J_0`$ be the component of $`J`$ containing $`a^{}`$. Since $`J_0\stackrel{~}{M}`$ it follows from Proposition 35 that $`J\mathrm{Odd}(a)`$ has rank at most 2. Now suppose, for a contradiction, that $`a^{}\{a,b\}`$. Since there exists an element $`wW_{D[a,b..a^{}]}`$ such that $`wa=a^{}`$ and $`wb=X(a^{})`$, it follows that $`r_{wb}=wr_bw^1\mathrm{FC}(wr_aw^1)=\mathrm{FC}(r_a^{})`$. Thus $`X(a^{})\mathrm{\Phi }_J`$, and so $`D[a,b..a^{}]J`$, contradicting the fact that the rank of $`J\mathrm{Odd}(a)`$ is at most 2. So we deduce that $`a^{}=b`$ or $`a^{}=a`$. Moreover, in either case we know that $`\{a,b\}J\mathrm{Odd}(a)`$, and since $`J\mathrm{Odd}(a)`$ has rank at most 2 it follows that $`J\mathrm{Odd}(a)=\{a,b\}`$.
By Lemma 36 we know that all spherical components of $`\mathrm{EOdd}(a)`$ are components of $`J`$, and by Proposition 35 all other components of $`J`$ that intersect $`\mathrm{Odd}(a)`$ trivially correspond to $`C_3`$-neighbours of $`\mathrm{Odd}(a)`$. But clearly the conditions of Definition 6 imply that $`\mathrm{Odd}(a)`$ has no $`C_3`$-neighbours. So we we conclude that $`J`$ is the union of $`\{a,b\}`$ and the spherical components of $`\mathrm{EOdd}(a)`$, as required. ∎
### Proof of Theorem 7
.
Let $`M`$ be a connected component of $`\mathrm{\Omega }(\mathrm{\Pi })`$, and write $`\stackrel{~}{M}`$ for the component of $`\mathrm{E}(M)`$ containing $`M`$.
Suppose first that $`\stackrel{~}{M}`$ is spherical, so that the conditions of Case A of Theorem 7 are satisfied, and let $`aM`$ be arbitrary. Observe that all $`C_3`$-neighbours of $`M`$ are contained in $`\stackrel{~}{M}`$. Choose $`a^{}M`$ such that $`\mathrm{FC}(r_a^{})`$ is visible, and let $`\mathrm{FC}(r_a^{})=W_J`$. By Lemma 36 we know that $`\stackrel{~}{M}`$ is contained in $`J`$, and hence $`aJ`$. So by Proposition 25 it follows that $`\mathrm{FC}(r_a)=W_J`$ also. By Proposition 35 the only possible components of $`J`$ apart from $`J_0`$ are the other spherical components of $`\mathrm{E}(M)`$, and by Lemma 36 all of these are indeed components of $`J`$. So $`J`$ consists of the spherical components of $`\mathrm{E}(M)`$, as required.
Now suppose that $`\stackrel{~}{M}`$ is not spherical. If there exists a $`b\mathrm{\Pi }M`$ such that $`(a,b)`$ is a focus of $`M`$ then it follows from Proposition 38 that $`\mathrm{FC}(r_a)=W_J`$, where $`J`$ is composed of $`\{a,b\}`$ and the spherical components of $`\mathrm{E}(M)`$, and $`\mathrm{FC}(r_a^{})`$ is not visible for any $`a^{}\mathrm{Odd}(a)\{a\}`$. Similarly, if there exists a $`bM`$ such that $`\{a,b\}`$ is a half-focus of $`M`$, then it follows from Proposition 39 that $`\mathrm{FC}(r_a)=\mathrm{FC}(r_b)=W_J`$, where $`J`$ is composed of $`\{a,b\}`$ and the spherical components of $`\mathrm{E}(M)`$, and $`\mathrm{FC}(r_a^{})`$ is not visible for any $`a^{}\mathrm{Odd}(a)\{a,b\}`$.
Finally, suppose that $`\stackrel{~}{M}`$ is not spherical and $`M`$ does not have a focus or a half focus. Suppose that $`aM`$ is such that $`\mathrm{FC}(r_a)=W_J`$ for some $`J\mathrm{\Pi }`$, and let $`K`$ be the component of $`J`$ containing $`a`$.
Suppose first that alternative (b) of Proposition 35 holds, so that $`K=\{a,b\}`$ is of type $`C_2`$, and $`JM=\{a\}`$. Since $`M`$ does not have any focus in $`\mathrm{\Pi }`$, it follows from Proposition 28 that $`M\{b\}`$ is a spherical component of $`\mathrm{E}(M)`$. But the component of $`\mathrm{E}(M)`$ containing $`M`$ is $`\stackrel{~}{M}`$, which, by our assumptions, is not spherical. So this case does not arise.
Alternative (c) of Proposition 35 is similarly impossible, by Proposition 29, and alternative (e) is also incompatible with the assumption that $`\stackrel{~}{M}`$ is not spherical. So we conclude that alternative (a) holds: $`K=\{a\}=J\mathrm{Odd}(a)`$. Note that all spherical components of $`\mathrm{E}(M)`$ are components of $`J`$, and by Proposition 35 the only other possible components are the sets $`\{b\}`$ such that $`b`$ is a $`C_3`$-neighbour of $`M`$.
Suppose that $`b`$ is a $`C_3`$-neighbour of $`M`$ that is adjacent to $`a`$. Let $`\stackrel{~}{a}`$ be the unique neighbour of $`a`$ in $`M`$. By Lemma 37 we know that $`r_b\mathrm{FC}(r_{\stackrel{~}{a}})`$, and so it follows that $`\mathrm{FC}(r_a)=r_{\stackrel{~}{a}}r_a\mathrm{FC}(r_{\stackrel{~}{a}})r_ar_{\stackrel{~}{a}}`$ contains the reflection along the root $`r_{\stackrel{~}{a}}r_ab=b+\sqrt{2}a+\sqrt{2}\stackrel{~}{a}`$. Since $`\mathrm{FC}(r_a)=W_J`$, it follows that both $`b`$ and $`\stackrel{~}{a}`$ are in $`J`$. But this contradicts the fact that the component of $`J`$ containing $`a`$ is just $`\{a\}`$.
This reasoning has shown that $`aM`$ is adjacent in $`\mathrm{\Pi }`$ to a $`C_3`$-neighbour of $`M`$ then $`\mathrm{FC}(r_a)`$ is not visible. On the other hand, we know that there is at least one $`aM`$ such that $`\mathrm{FC}(r_a)`$ is visible. So we may choose an $`aM`$ such that $`\mathrm{FC}(r_a)=W_J`$ for some $`J\mathrm{\Pi }`$. Since $`a`$ is not adjacent to any $`C_3`$-neighbour of $`M`$ it follows by Lemma 37 that all $`C_3`$-neighbours of $`M`$ are in $`J`$. So we conclude that $`J=J^{}\{a\}`$, where $`J^{}`$ is the union of the spherical components of $`\mathrm{E}(M)`$ and the $`C_3`$-neighbours of $`M`$.
It remains to prove that is $`a^{}`$ is any other element of $`M`$ that is not adjacent to any $`C_3`$-neighbour of $`M`$ then $`\mathrm{FC}(r_a^{})=W_{J^{}\{a^{}\}}`$. Given such an $`a^{}`$, since $`a^{}`$ lies in $`M=\mathrm{Odd}(a)`$, we may choose $`wW`$ such that $`wa=a^{}`$. By Proposition 12 we see that $`wW_{\stackrel{~}{M}}`$, and so $`w`$ fixes all other components of $`\mathrm{E}(M)`$. And $`w`$ fixes all $`C_3`$-neighbours of $`M`$, by Lemma 34. So $`w`$ fixes $`J^{}`$, and it follows that
$$\mathrm{FC}(r_a^{})=w\mathrm{FC}(r_a)w^1=wW_{J^{}\{a\}}w^1=W_{wJ^{}\{wa\}}=W_{J^{}\{a^{}\}},$$
as required. This completes the proof of Theorem 7.
### Proof of Theorem 1
Let $`a\mathrm{\Pi }`$ and $`M=\mathrm{Odd}(a)`$. As $`\mathrm{\Pi }`$ is 2-spherical it follows that $`\mathrm{\Pi }=E(M)`$, and, as $`\mathrm{\Pi }`$ is non-spherical, it follows that Case A of Theorem 7 does not hold for $`M`$. As there are no $`\mathrm{}`$-labels in the Coxeter graph of $`\mathrm{\Pi }`$, Cases C and D do not hold either. Hence we are in Case B. Since there are no $`\mathrm{}`$-labels in the Coxeter graph of $`\mathrm{\Pi }`$, there are no $`C_3`$ neighbors of $`M`$. As $`E(M)=\mathrm{\Pi }`$ is irreducible, there are no spherical components of $`E(M)`$. It follows now from Theorem 7 that there is an $`a^{}\mathrm{Odd}(a)`$ such that $`\mathrm{FC}(r_a^{})=r_a^{}`$. As $`r_a`$ and $`r_a^{}`$ are $`W`$-conjugate we have $`\mathrm{FC}(r_a)=r_a`$ as well, and this completes the proof of Part a) of Theorem 1.
Let $`SW`$ be such that $`(W,S)`$ is a Coxeter system. It follows from Part a) and Corollary 24 that $`r_aS^W`$ for each $`a\mathrm{\Pi }`$, and hence $`\{r_aa\mathrm{\Pi }\}S^W`$. As $`\mathrm{\Pi }`$ is assumed to be non-spherical, irreducible and 2-spherical, it follows now from the main result of that there is an element $`wW`$ such that $`\{r_aa\mathrm{\Pi }\}=S^w`$. This completes the proof of Part b) of Theorem 1. As Part c) is an immediate consequence of Part b) we are done.
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# Periodicity and Growth in a Lattice Gas With Dynamical Geometry
## Abstract
We study a one-dimensional lattice gas “dynamical geometry model” in which local reversible interactions of counter-rotating groups of particles on a ring can create or destroy lattice sites. We exhibit many periodic orbits and show that all other solutions have asymptotically growing lattice length in both directions of time. We explain why the length grows as $`\sqrt{t}`$ in all cases examined. We completely solve the dynamics for small numbers of particles with arbitrary initial conditions.
## 1. Introduction
Lattice models are ubiquitous in physics, whether as regularizations for continuum theories (quantum field theory, quantum gravity), scaffolding for numerical methods (classical field theories, continuum mechanics), or because the lattice is physically real (condensed matter physics). In virtually all applications, however, the lattice structure and size are fixed, or at least not dynamical. (In numerical computation the lattice may be refined to maintain precision, but the evolution of the lattice is not part of the physical dynamics under study.) Among the exceptions known to us are the causal dynamical triangulation approach to quantum gravity \[AJL05\] and the variable-length lattice models of recent relevance to the AdS/CFT correspondence \[BCV05\].
In \[HM98\], Hasslacher and Meyer constructed a lattice gas model with dynamical geometry and a reversible evolution rule. It can be viewed as a toy model for general relativity in that the geometry (length) of the one-dimensional lattice changes in response to the motion (scattering) of the matter particles on it. The model is classical, but its quantization should present no problems. In this paper, we extend the previous analyses of the classical dynamics. The central issue is the long-time behavior of the lattice length. We will completely solve the dynamics for systems of a few particles, and explain the typical $`\sqrt{t}`$ growth of the length which has been seen previously in simulations.
The model consists of a one-dimensional lattice of $`L`$ sites with periodic boundary conditions (a ring), where $`L`$ may change with time. The initial state contains $`N_R`$ right-moving particles and $`N_L`$ left-moving ones, which may be placed arbitrarily on the sites subject to an exclusion principle: two particles moving in the same direction may not occupy the same site. The numbers of left- and right-movers are each conserved during evolution. Time proceeds in discrete steps. At each time step, the particles first advect: each particle moves one site in its own direction of motion. Then the particles scatter according to the rules (see Figure 1):
(1) If a right- and left-mover occupy the same site, this site is replaced by two sites, with the right-mover on the rightmost site and the left-mover on the leftmost site.
(2) If two adjacent sites are occupied, with a right-mover on the rightmost site and a left-mover on the leftmost, these sites are replaced by a single site occupied by both particles.
(3) Rules (1) and (2) both apply when two adjacent sites are occupied by three particles (or by four), with the singleton moving away from the doubly-occupied site. Since the application of rule (2) would violate the Exclusion Principle, in such cases only rule (1) is applied.
These rules define a reversible dynamics in that every state has not only a unique successor but also a unique predecessor.
Rule (3), the exclusion rule, is the major complication in analyzing the dynamics, as will be seen in section 2 below. The situation it describes can arise following advection when a single particle moving in one direction approaches a pair of particles on adjacent sites moving in the other. Such a pair of particles moving in the same direction on adjacent sites will be called simply a pair, and we will see that such “bound states” can play the role of quasiparticles in the system.
After initial explorations in \[HM98\], numerical simulations of the evolution of initial states rather densely populated with particles (25 particles on a 50 site lattice) were carried out in \[LBM04\]. Despite the reversibility of the dynamics, most initial states result in growth of the lattice size, empirically as $`L(t)\sqrt{t}`$ at late times. On small lattices, some “rogue states” were also found with $`L(t)`$ periodic, but the proportion of these dropped off rapidly with initial lattice size. Two versions of mean-field theory were proposed to explain the observed growth rate, one of which predicted $`\sqrt{t}`$ growth as observed while the other predicted $`t^{1/3}`$. In this paper we will identify many periodic solutions on lattices of arbitrary size, and also propose an alternative explanation for the typical $`\sqrt{t}`$ growth.
The reversibility of the dynamics leads to the following simple but fundamental
Evolution Theorem: Every solution of the model is either periodic, or grows without bound in both directions of time.
Proof: Consider a solution for which the lattice length remains bounded in one direction of time, say $`t\mathrm{}`$. Since there are only finitely many distinct states of this system on a lattice of given size, the evolution must eventually return to some previous state. The evolution is then periodic from this time on. By reversibility and uniqueness, it is periodic in backward time also.
The organization of this paper is as follows. In section 2 we give a general analysis of the evolution of initial states, assuming that no pairs are present or form later. We point out the crucial role of parity in the problem (first noted in \[LBM04\]) and establish the existence of many periodic solutions. We explain why growing solutions of this kind must grow as $`\sqrt{t}`$. In section 3 we exhaustively analyze the evolution of all initial states containing at most four particles, possibly including pairs. We see cases in which permanent pairs form and behave as quasiparticles or bound states in the system. The major obstacle to a complete solution of this model is the lack of a general framework for describing these quasiparticles and their effects. Section 4 contains conclusions and open problems.
We would like to acknowledge discussions at the early stages of this work with the authors of \[NLB05\], who have independently obtained similar results.
## 2. General Analysis of States Without Nearest-Neighbor Pairs
Although simulations of the time evolution of “random” initial states look quite complicated, there are several important general principles governing the dynamics which can be formulated. First, the translational (rotational) symmetry of the lattice allows the dynamics to be viewed in various reference frames. We have thus far used a frame fixed with respect to the lattice, but we can transform to the rest frame of either the right- or the left-moving particles. The rest frame of the left-movers, for example, is defined as follows. At each advection step, the left-movers move one site to the left. We can follow this advection step with a symmetry (gauge) transformation which rotates every particle one step to the right, thus undoing the advection for the left-movers. This is followed by the scattering step as usual. In this frame, or gauge, the left-movers do not advect, while the right-movers advect two sites per time step.
The use of the rest frame for one group of particles makes it clear that parity (mod 2) plays a crucial role in the dynamics. The size of the gap between two left-movers, or two right-movers, is trivially preserved by advection, but so is the parity of the (shrinking) interval between a left- and a right-mover. If this interval contains no other particles, its parity determines the character of the eventual interaction between these particles, assuming neither belongs to a nearest-neighbor pair: a site will be created (resp., destroyed) if this parity is even (resp., odd). In turn, either type of interaction will reverse the parity of a gap containing the newly created or destroyed site. Thus, we can reduce the dynamics mod 2, and study the evolution of the gap parities according to these simple rules. The presence of pairs can invalidate the analysis by preventing destruction events which would otherwise occur. The corresponding parity reversals then do not take place.
Thus, our plan for analyzing the dynamics is to assume first that no pairs are present, and no pairs form during evolution. Analysis via parity then leads to very simple and general results, which are valid if they are self-consistent, that is, if they do not lead to the formation of pairs. We then consider the excluded cases in which pairs are present initially or form later. The pairs can violate the parity “selection rules”, causing transitions between the types of behavior observed when those rules hold. At present, we can only analyze the effects of pairs in a laborious case-by-case manner with small numbers of particles. An important task for the future is to develop appropriate concepts for a general analysis of the effects of pairs, which might be thought of as quasiparticles in the system.
To introduce the technique of parity analysis, consider the case of one right-mover versus a string of $`NN_L`$ left-movers, without pairs. We work in the rest frame of the left-movers, so the right-moving particle advects by two sites per time step, preserving parities. We label the left-movers from 1 to $`N`$ in the order that the right-mover will encounter them, and describe the system by the gaps $`g_{12},g_{23},\mathrm{},g_{N1}`$ between successive left-movers as well as the parity $`p`$ of the separation between the right-mover and the first left-mover. (The term “gap” always refers to the distance between consecutive left-movers, or consecutive right-movers; we use the term “separation” for the shrinking distance between a left-mover and right-mover before they scatter.) This data is preserved by advection, and changes as follows under scattering:
$$[p;g_{12},g_{23},\mathrm{},g_{N1}][|p+g_{12}|;g_{23},\mathrm{},g_{N1},g_{12}+(1)^p],$$
where $`|x|`$ is the parity of $`x`$. The gaps are always listed in the order that the right-mover will encounter them.
Let the right-mover make a complete circuit, passing all $`N`$ left-movers and returning to face the first again. (Note: this complete circuit in the rest frame is only a half-circuit in the lattice-fixed frame. In that frame, the right-mover passes each left-mover twice as every particle makes a complete circuit of the lattice.) Each gap parity $`g_{i,i+1}`$ has been reversed, and $`p`$ has changed by the sum of the original gap parities, which is the parity of the original lattice length $`L`$. (It may be surprising that the parity of $`L`$ prior to the complete circuit determines whether the right-mover returns to its original position relative to the first left-mover or is displaced from it by one site. The point is that sites which may be created or destroyed during this circuit do not increase or decrease the distance the right-mover must advect to finish the circuit. This is because the right-mover is carried forward or backward with the newly created or destroyed site as part of the scattering step. Changes in $`L`$ during this circuit take effect at the next circuit.) The parity of the lattice length itself has changed by $`|N|`$ due to the $`N`$ scattering events. Now let the right-mover make another circuit. This restores all the original gap parities $`g_{i,i+1}`$, but not necessarily $`p`$, which differs from its original value by $`|2L+N|=|N|`$.
Supposing first that $`N`$ is even, the system is parity-periodic with period two circuits, that is, $`2N`$ interactions. This implies that further circuits will repeat the same pattern of creation and destruction interactions as the first two circuits. However, the actual gap lengths are not periodic, because the interactions of the second circuit do not undo the effects of the first. If, for example, $`L`$ is odd, it is easy to see that $`g_{12}`$ is unchanged after two circuits while $`g_{23}`$ has changed by $`\pm 2`$. Indeed, regardless of $`|L|`$, alternate gaps have changed by $`0,\pm 2,0,\pm 2,\mathrm{}`$ after two circuits. The system either grows forever (if all signs are $`+`$) or a pair eventually forms and invalidates the analysis.
Suppose next that $`N`$ is odd. After two circuits $`p`$ has reversed parity while all other gaps have their original parities. It follows that the interactions of circuits three and four will exactly undo the effects of circuits one and two respectively, and the evolution will be truly periodic with period four circuits, or $`4N`$ interactions. During these four circuits, each gap can change by at most $`\pm 2`$ from its initial size. Therefore, $`g_{i,i+1}>3`$ is sufficient to prevent the formation of pairs and render this behavior self-consistent. We have thus established the existence of a large class of periodic solutions.
Note that there are $`2^{N+1}`$ possible parity states, given by the parities of the $`N`$ gaps and of $`p`$. Since this number is finite, periodic evolution of the parities is inevitable. Since the period, $`2N`$ or $`4N`$, only grows linearly with $`N`$, the $`2^{N+1}`$ states clearly belong to a large number of disjoint parity orbits when $`N`$ is large. Some examples with $`N`$ small are given in the next section.
With a mild additional assumption this analysis can be extended to the general case of $`N_R`$ right-movers versus $`N_L`$ left-movers, of course assuming the absence of pairs. Provided that no right-mover is very near a left-mover in the initial state, we can again consider a complete circuit in which every right-mover passes every left-mover exactly once. Each right-mover interacts $`N_L`$ times, making $`N_RN_L`$ interactions in all. In the rest frame of the left-movers, it is clear that each gap between consecutive left-movers changes parity by $`|N_R|`$ during one circuit of the right-movers. Similarly, each gap between consecutive right-movers changes parity by $`|N_L|`$. There is also a parity change in the separation between a chosen left-mover and a chosen right-mover after one circuit. It suffices to compute this for one such choice, because the others are then determined by the known gap changes. However, this “offset parity” depends on the exact configuration of the particles around the ring. We will always work in the rest frame of the left-movers, and compute the offset between a chosen right-mover and the first left-mover it will encounter.
For example, suppose that initially one arc of the ring contains all the left-movers and no right-movers, and another contains all the right-movers. Consider the leading right-mover. As it passes the left-movers, it is carried along with the creation and destruction events, and its offset parity after one cycle is simply $`|L|`$. Contrast this with an initial configuration in which $`N_L=N_R=N`$ and the left- and right-movers alternate around the ring. Choose a “leading” right-mover arbitrarily. Number the left-movers in the order that this right-mover will encounter them. It passes the first and traverses the gap $`g_{12}`$, but the next gap has already been changed to $`g_{23}\pm 1`$ by the interaction with the right-mover ahead of the chosen one. Similarly the next gap will be $`g_{34}\pm 1\pm 1`$ (independent signs) when the chosen particle gets there, and the offset has the parity of $`|L+1+2+\mathrm{}+(N1)|=|L+\frac{1}{2}N(N1)|`$. In general, the offset is $`L`$ plus the number of interactions with left-movers which occur before the chosen right-mover reaches them. The latter contribution cancels out (mod 2) after two circuits.
Now we analyze the various parity combinations in detail. Suppose first that both $`N_R`$ and $`N_L`$ are even. Then the parity of every gap is unchanged after one circuit. If the offset parity for some particle is also even, then it is even for every particle, and the pattern of interactions at every successive circuit is identical. The system either grows indefinitely or eventually forms a pair and invalidates the analysis. If, however, the offset for some (hence every) particle was odd, then the interactions of the second circuit undo the effects of the first, and the system is truly periodic with period two circuits or $`2N_RN_L`$ interactions. No gap changes by more than $`\mathrm{max}(N_R,N_L)`$ during the evolution.
Next suppose that $`N_R`$ and $`N_L`$ are both odd. Now every gap changes parity in a circuit, as does the lattice length. After two circuits, the gap parities have their original values, but the offset is odd. As in the case of $`1`$ versus odd $`N`$ above, circuits three and four undo the effects of circuits one and two, and we get truly periodic solutions with period four circuits or $`4N_RN_L`$ interactions.
Finally, suppose $`N_R`$ is odd and $`N_L`$ is even (the opposite case being the same by symmetry). After one cycle the left-left gap parities are reversed, while the parities of the right-right gaps and the lattice length are unchanged. As in the case of 1 versus even $`N`$ above, the interactions of the next cycle cancel those of the first for every other left-left gap, but augment those of the first for the remaining left-left gaps. The result is either net growth or pair formation.
To summarize, in the absence of pairs periodic solutions are quite generic in the cases $`N_R,N_L`$ both odd, and both even with offset parity odd.
Now consider any solution which grows indefinitely. Because the gap and offset parities repeat after a certain number of circuits, and determine the pattern of interactions and thus the net number of sites created, the growth is characterized by some average number of sites created per circuit. Let $`k`$ be this number, and $`L(t)`$ be the lattice length at time $`t`$. Since a circuit takes $`L/2`$ time steps, in the limit of continuous time we have the differential equation
$$\frac{dL}{dt}=\frac{2k}{L},$$
with asymptotic solution $`L2\sqrt{kt}`$. The $`\sqrt{t}`$ growth observed in simulations thus simply reflects constant average growth per circuit.
It is tempting to claim that this reasoning is completely general, applying even if nearest-neighbor pairs form. The argument would be that the pattern of interactions is still determined by a finite set of data, namely the parities of all gaps and an offset, and a list of which particles are paired. As this finite set of data changes, it must eventually return to a former state, from which point the pattern of interactions will be periodic. There will be a net number of sites created per period, leading again to the $`\sqrt{t}`$ growth, on a sufficiently long time scale. However, this finite set of data is in fact insufficient to determine the sequence of interactions, because one needs the actual gap sizes, not just their parities, to predict when a new pair will form. Thus, at present we cannot prove that every nonperiodic solution grows according to the $`\sqrt{t}`$ law.
## 3. Few-Particle Systems
In this section we will completely solve the dynamics for all initial states containing at most four particles. By symmetry we may assume at most two are right-movers; for the moment we consider a single right-mover with up to three left-movers. We describe the states of the system at time $`t`$ by the parity of the separation between the right-mover and the left-mover it will encounter next, and by the gaps $`g_{i,i+1}`$ between left-movers $`i`$ and $`i+1`$ (where the indices $`i,i+1`$ are taken modulo $`N_L`$). After every interaction we give the new state.
### 3.1. One Against $`N_L`$
###### Example 3.1.
One against one
For completeness and to establish notation we begin with the trivial case $`N_R=N_L=1`$. We keep track of the parity $`p`$ of the separation $`d(r,l)`$ and the parity of the gap $`g_{11}`$, in other words of the length of the lattice. The state is described by $`[p,|L|]`$. To the right we give the value of the growth $`\mathrm{\Delta }L`$ following each interaction.
$$\begin{array}{cc}[1;1]& \\ & 1\\ [0;0]& \\ & +1\\ [0;1]& \\ & +1\\ [1;0]& \\ & 1\\ [1;1]& \end{array}$$
Any state in this cycle can be viewed as the “initial” state; all possible parities for one against one belong to the same orbit. This orbit has net growth zero, and is truly periodic.
###### Example 3.2.
One against two
The first interesting case is one right-mover against two left-movers. Here, the exclusion rule can apply (nearest-neighbor pairs suppressing the destruction of sites).
We describe the possible states as follows: $`[p;g_{i,i+1},g_{i+1,i+2};|L|]`$, where $`p`$ is the parity of the separation $`d(r,l_i)`$ between $`r`$ and the closest left mover $`l_i`$ ($`i=1,2`$), and $`|L|`$ is the parity of the lattice length.
With this data we can keep track of the position of $`r`$: we can always see which particles interact next. Note that $`|L|=|g_{12}+g_{21}|`$. For convenience we display parity by using bold letters $`𝐠_{𝐢,𝐢+\mathrm{𝟏}}`$ if and only if the gap $`g_{i,i+1}`$ is odd.
There are eight possible combinations of the parities of $`d(r,l_i)`$, $`g_{i,i+1}`$, $`g_{i+1,i+2}`$. We start with the two states (1) and (i), assuming at first that no pairs occur (see below).
$$\begin{array}{cccccc}(1)& [1;𝐠_{\mathrm{𝟏𝟐}},g_{21};1]\hfill & & (i)& [1;g_{12},g_{21};0]\hfill & \\ & & \hfill 1& & & \hfill 1\\ (2)& [0;g_{21},g_{12}1;0]\hfill & & (ii)& [1;g_{21},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};1]\hfill & \\ & & \hfill +1& & & \hfill 1\\ (3)& [0;g_{12}1,𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};1]\hfill & & (iii)& [1;𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏};0]\hfill & \\ & & \hfill +1& & & \hfill 1\\ (4)& [0;𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}};0]\hfill & & (iv)& [0;𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏},g_{12}2;1]\hfill & \\ & & \hfill +1& & & \hfill +1\\ (5)& [1;𝐠_{\mathrm{𝟏𝟐}},g_{21}+2;1]\hfill & & (v)& [1;g_{12}2,g_{21};0]\hfill & \end{array}$$
In this chart, the quantities $`g_{12}`$ and $`g_{21}`$ denote the initial values of these gaps; current values are indicated by position within the brackets. For example, line (3) indicates that the right-mover is about to encounter $`l_1`$, the gap from $`l_1`$ rightward to $`l_2`$ is currently $`g_{12}1`$, and the gap from $`l_2`$ rightward to $`l_1`$ is $`g_{21}+1`$. Observe that the eight states (1)-(4) and (i)-(iv) cover all the parity configurations. Thus if no pairs occur, these eight states describe all possible behaviors of one particle against two.
Consider the first example, states (1) - (5). Note that the parities of state (5) are a repetition of the parities of state (1). There is one gap that has grown by two, $`g_{21}g_{21}+2`$, while the other gap remained constant. So this is a parity-periodic growing orbit as long as no pairs appear. The only places where a pair can form are in (1), (5), (9), and so on. The restriction $`g_{12}3`$ prevents the formation of pairs. In particular, we have a growing orbit if and only if $`g_{12}3`$. The growth rate is as $`\sqrt{t}`$ as discussed previously.
The second example describes a shrinking system. After four interactions we get back to the same configuration of parities but with a shorter lattice. Two sites have been eliminated. It is clear that eventually pairs will form and alter the evolution in the states (iii), (vii), (xi), etc. where we have an odd distance of $`r`$ to $`l_1`$. Although we will consider the effects of these pairs below, the eventual fate of this system can be determined by the following time-reversal argument. Running time backward from the initial state, we would obviously see this system grow, with no pair formation. By the Evolution Theorem, this system must eventually grow in the future as well.
To complete the analysis of one right-mover against two left-movers, we now describe what happens if pairs form. In that case $`g_{12}=1`$, and we can assume that the distance $`d(r,l_1)`$ is odd, since otherwise the existence of the pair does not affect the evolution. Since we have the choice of the parity of $`g_{21}`$, there are two types of states with a pair:
$$\begin{array}{cccc}(A)\hfill & & (B)\hfill & \\ & & & \\ \left[1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1\right]\hfill & & \left[1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0\right]\hfill & \\ & \hfill +\mathrm{𝟏}& & \hfill +\mathrm{𝟏}\\ \left[0;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0\right]\hfill & & \left[1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}+1;1\right]\hfill & \end{array}$$
In the above evolution, the boldface values of $`\mathrm{\Delta }L=+\mathrm{𝟏}`$ indicate steps at which the evolution was altered by the presence of the pair. Note that (A) has evolved into step (4). Since in (A) we started with $`g_{21}2`$, we get $`g_{21}+13`$ and remain in the growing orbit. Case (B) has evolved into (A) after one step, so both (A) and (B) evolve to the growing orbit. The pair breaks, and the gap between these neighbors subsequently grows.
Now let $`r`$ face a sequence of left-movers, $`l_1,\mathrm{},l_{N_L}`$, the initial state being $`[p;g_{12},g_{23},\mathrm{},g_{N_L1}]`$. If particles 1 and 2 do not form a pair, then the state following the first interaction is as given in the previous section, namely
$$[|p+g_{12}|;g_{23},\mathrm{},g_{N_L1},g_{12}+(1)^p].$$
If particles 1 and 2 do form a pair, this changes the outcome iff $`p=1`$. In that case, $`r`$ does not interact with $`l_1`$, and following the interaction with $`l_2`$ the state will be
$$[|g_{23}|;g_{34},\mathrm{},g_{N_L1},g_{12},g_{23}+1].$$
The interaction between $`r`$ and a single pair can be conceptualized as follows. If the gap $`d(r,l_1)`$ to the leading member of the pair is odd, the right-mover lands on the trailing member of the pair. The pair remains intact and a site is created behind it, so that $`\mathrm{\Delta }L=1`$. In this case the pair behaves as a unit like a single left-mover. If the gap $`d(r,l_1)`$ is even, the right-mover lands on the leading member of the pair. A site is created between the left-movers, breaking the pair, and another site is destroyed behind the pair, resulting in $`\mathrm{\Delta }L=0`$. More globally, under the right conditions a pair may persist indefinitely, behaving like a quasiparticle in the system, as we will see later. Or a pair might break and re-form repeatedly. Notice that, in contrast, three left-movers on adjacent sites cannot form a stable “triple”; it will necessarily be broken by the next interaction with a right-mover.
###### Example 3.3.
One against three
Let us turn to the case of one right-mover against three left-movers.
Here we describe a state by the parity $`p`$ of $`d(r,l_1)`$ (assuming $`r`$ faces $`l_1`$ next), by $`g_{12}`$, $`g_{23}`$, $`g_{31}`$ and (for convenience, although it is redundant) the parity $`|L|`$ of the length of the lattice. There are sixteen possible parity configurations for the gaps and separation $`p`$.
We obtain two periodic orbits of length $`12`$ under the assumption that no pairs form. We give them here:
$$\begin{array}{cccccc}\hfill (1)& [1;g_{12},g_{23},g_{31};0]\hfill & & \hfill (a)& [1;g_{12},𝐠_{\mathrm{𝟐𝟑}},g_{31};1]\hfill & \\ & & 1& & & 1\\ \hfill (2)& [1;g_{23},g_{31},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};1]\hfill & & \hfill (b)& [1;𝐠_{\mathrm{𝟐𝟑}},g_{31},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};0]\hfill & \\ & & 1& & & 1\\ \hfill (3)& [1;g_{31},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏};0]\hfill & & \hfill (c)& [0;g_{31},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},g_{23}1;1]\hfill & \\ & & 1& & & +1\\ \hfill (4)& [1;𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏};1]\hfill & & \hfill (d)& [0;𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},g_{23}1,𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏};0]\hfill & \\ & & 1& & & +1\\ \hfill (5)& [0;𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},g_{12}2;0]\hfill & & \hfill (e)& [1;g_{23}1,𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏},g_{12};1]\hfill & \\ & & +1& & & 1\\ \hfill (6)& [1;𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},g_{12}2,g_{23};1]\hfill & & \hfill (f)& [1;𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏},g_{12},𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟐};0]\hfill & \\ & & 1& & & 1\\ \hfill (7)& [0;g_{12}2,g_{23},g_{31}2;0]\hfill & & \hfill (g)& [0;g_{12},𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟐},g_{31};1]\hfill & \\ & & +1& & & +1\\ \hfill (8)& [0;g_{23},g_{31}2,𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};1]\hfill & & \hfill (h)& [0;𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟐},g_{31},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏};0]\hfill & \\ & & +1& & & +1\\ \hfill (9)& [0;g_{31}2,𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟏};0]\hfill & & \hfill (i)& [1;g_{31},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏},g_{23}1;1]\hfill & \\ & & +1& & & 1\\ \hfill (10)& [0;𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏};1]\hfill & & \hfill (j)& [1;𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏},g_{23}1,𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏};0]\hfill & \\ & & +1& & & 1\\ \hfill (11)& [1;𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},g_{12};0]\hfill & & \hfill (k)& [0;g_{23}1,𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},g_{12};1]\hfill & \\ & & 1& & & +1\\ \hfill (12)& [0;𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},g_{12},g_{23};1]\hfill & & \hfill (l)& [0;𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},g_{12},𝐠_{\mathrm{𝟐𝟑}};0]\hfill & \\ & & +1& & & +1\\ \hfill (1)& [1;g_{12},g_{23},g_{31};0]\hfill & & \hfill (a)& [1;g_{12},𝐠_{\mathrm{𝟐𝟑}},g_{31};1]\hfill & \end{array}$$
The first system is free of pairs if and only if $`g_{12}4`$, $`g_{23}2`$ and $`g_{31}4`$: the only cases where the exclusion rule can apply are in steps (4) if $`g_{12}1=1`$, in (6) if $`g_{31}1=1`$ and in (11) if $`g_{23}+1=1`$. (Of course, this would literally imply $`g_{23}=0`$, which is not possible by the Exclusion Principle. What is meant is that an initial state with the parities of line (11) would contain a pair if the entry $`g_{23}+1`$ were $`1`$ instead.)
The second system is free of pairs if and only if $`g_{12}2`$, $`g_{23}3`$, $`g_{31}2`$. The only cases where the exclusion rule can apply are in steps (b) if $`g_{23}=1`$, in (f) if $`g_{31}+1=1`$ and in (j) if $`g_{12}+1=1`$.
Note that all sixteen parity configurations appear in the two orbits. The first orbit covers $`12`$ parity configurations. In the second orbit, the parity configurations repeat after $`4`$ steps, but the actual gap sizes have period $`12`$.
What is left is to understand the cases where pairs form (and thus the exclusion rule applies). In other words we have to study the systems with $`d(r,l_1)`$ odd, $`g_{12}=1`$ and all possible parities of $`g_{23}`$, $`g_{31}`$. We give them labels as follows:
$$\begin{array}{cc}\text{parities of}(g_{23},g_{31})& \text{label}\\ & \\ [0,0]& (A)\\ [0,1]& (B)\\ [1,0]& (C)\\ [1,1]& (D)\end{array}$$
States of types (B) and (C) belong to a single orbit. The pair stays intact and acts as a permanent quasiparticle, and the lattice grows exactly as in the 1 vs. 2 case:
$$\begin{array}{ccc}\hfill (C)& [1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}},g_{31};0]\hfill & \\ & & +\mathrm{𝟏}\\ & [1;g_{31},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23}+1;1]\hfill & \\ & & 1\\ \hfill (B)& [1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23}+1,𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏};0]\hfill & \\ & & +\mathrm{𝟏}\\ & [0;𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟐};1]\hfill & \\ & & +1\\ \hfill (C)& [1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟐},g_{31};0]\hfill & \end{array}$$
The pattern of type (A) is given below. Note that after eight interactions, we have a state of type (D), so type (D) is contained in that orbit. As long as no other pairs appear, type (A) forms a parity-periodic orbit of length sixteen. $`\mathrm{\Delta }L=2`$ after one such orbit: the gap between particles $`l_2`$, $`l_3`$ decreases by two. Pairs can appear in state (d) if $`g_{23}+1=1`$, in state (j) if $`g_{31}+1=1`$ and in state (o) if $`g_{23}1=1`$.
Suppose $`g_{31}+1=1`$ in (j). This is a pair as in type (C), so from here on, the orbit is growing.
Let $`g_{23}2k+1=1`$. The change in the pattern occurs in the $`k`$-th run through the orbit (a)-(p). For $`k=0`$ state (d) contains a pair of type (C). For $`k>0`$ state (o) contains a pair of type (D), $`[1;1,1,1]`$ with first and third gap length equal to one. This state evolves to state (j) with first gap of length one: a type (C) pair has formed. We already know that type (C) belongs to a growing orbit.
So systems of type (A) and (D) produce a parity-periodic orbit of decreasing length. As soon as a second pair forms, we observe a transition via (D) to the growing orbit of type (B).
Type (A):
$$\begin{array}{ccc}(a)& [1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23},g_{31};1]\hfill & \\ & & +\mathrm{𝟏}\\ (b)& [0;g_{31},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟏};0]\hfill & \\ & & +1\\ (c)& [0;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏};1]\hfill & \\ & & +1\\ (d)& [1;𝐠_{\mathrm{𝟐𝟑}}+\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏},g_{12}+1;0]\hfill & \\ & & 1\\ (e)& [0;𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏},g_{12}+1,g_{23};1]\hfill & \\ & & +1\\ (f)& [1;g_{12}+1,g_{23},g_{31}+2;0]\hfill & \\ & & 1\\ (g)& [1;g_{23},g_{31}+2,𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};1]\hfill & \\ & & 1\\ (h)& [1;g_{31}+2,𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏};0]\hfill & \\ & & 1\\ (i)& [1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏};1]\hfill & \\ & & +\mathrm{𝟏}\\ (j)& [1;𝐠_{\mathrm{𝟑𝟏}}+\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23};0]\hfill & \\ & & 1\\ (k)& [0;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23},g_{31};1]\hfill & \\ & & +1\\ (l)& [1;g_{23},g_{31},g_{12}+1;0]\hfill & \\ & & 1\\ (m)& [1;g_{31},g_{12}+1,𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏};1]\hfill & \\ & & 1\\ (n)& [1;g_{12}+1,𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏};0]\hfill & \\ & & 1\\ (o)& [1;𝐠_{\mathrm{𝟐𝟑}}\mathrm{𝟏},𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};1]\hfill & \\ & & 1\\ (p)& [0;𝐠_{\mathrm{𝟑𝟏}}\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23}2;0]\hfill & \\ & & +1\\ (a^{})& [1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{23}2,g_{31};1]\hfill & \end{array}$$
### 3.2. Two Against Two
This case needs more information: now there are two right-movers and they interact in a row. We have to adapt the notation and keep track of the positions of the right-movers.
Suppose $`r_1`$ interacts first with $`l_1`$. Then there are several possibilities:
* Particle $`r_1`$ interacts next with $`l_2`$, then $`r_2`$ interacts with $`l_1`$.
* There is an interlaced pattern: particle $`r_2`$ interacts with $`l_1`$, then $`r_1`$ interacts with $`l_2`$, then $`r_2`$ with $`l_2`$.
* Particles $`r_1`$ and $`l_2`$ interact at the same time step as $`r_2`$ with $`l_1`$.
It can be checked that it does not matter which interaction is first (as long as no pairs occur). This is a reflection of the locality of the scattering.
We start by looking at the cases where no pairs appear. In order to describe the evolution of the system, we give pairs of $`4`$-tuples: $`[d_{11};g_{12},g_{21};|L|][d_{21};g_{12},g_{21};|L|].`$
In contrast to the previous notation, we give the actual separation between the right-mover $`r_i`$ and the next left-mover $`l_j`$ that $`r_i`$ faces. Let $`h_{21}:=d(r_2,r_1)`$ be the gap between the two right-movers. As before, bold font is used to denote odd length. Furthermore, we add an arrow $`\pm 1`$ to show the change of the length of the lattice. The arrows also indicate which pair is interacting.
Figure 4 shows a “noninterlaced” state in which there is no left-mover between the two right-movers, and vice-versa. We can always assume such an initial state, unless the four particles are located symmetrically around the ring: right-movers diametrically opposite one another, left-movers likewise along a perpendicular diameter. It is easy to see that these exceptional initial states evolve as periodic orbits, and we will not consider them further.
Note that $`d_{21}=d_{11}+h_{21}`$. We know that the order of the interactions does not matter. Thus in the example below we can fix the order of interactions by assuming $`h_{21}g_{12}2`$ in (2) and $`h_{21}+1g_{21}2`$ in (4).
We first give the two periodic orbits. Then we will describe the growing orbits and those with decreasing length. In the end we discuss the cases with pairs.
###### Example 3.4.
We start with a system where $`g_{12}`$ is even, $`g_{21}`$ odd (i.e. the lattice has odd length). Also let $`h_{21}`$ and $`d_{11}`$ be even. Recall that notation such as $`g_{12}`$ below denotes this gap in the initial state (1); as the gaps change during evolution their current values are indicated by their positions.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [d_{11};g_{12},𝐠_{\mathrm{𝟐𝟏}};1]\hfill & & [h_{21}+d_{11};g_{12},𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \\ & & +1& & \\ (2)& [g_{12};𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏};0]\hfill & & [h_{21};𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & & & +1\\ (3)& [g_{12}h_{21};𝐠_{\mathrm{𝟐𝟏}},g_{12}+2;1]\hfill & & [𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏};g_{12}+2,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \\ & & +1& & \\ (4)& [𝐠_{\mathrm{𝟐𝟏}};g_{12}+2,g_{21}+1;0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};g_{21}+1,g_{12}+2;0]\hfill & \\ & & & & 1\\ (5)& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟐};g_{12}+2,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & & [𝐠_{\mathrm{𝟐𝟏}};g_{12}+2,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \\ & & 1& & \\ (6)& [𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & & & 1\\ (7)& [𝐠_{\mathrm{𝟏𝟐}}𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}},g_{12};1]\hfill & & [g_{12};𝐠_{\mathrm{𝟐𝟏}},g_{12};1]\hfill & \\ & & 1& & \\ (8)& [g_{21}1;g_{12},g_{21}1;0]\hfill & & [h_{21};g_{21}1,g_{12};0]\hfill & \\ & & & & +1\\ (1)^{}& [g_{21}h_{21}1;g_{12},𝐠_{\mathrm{𝟐𝟏}};1]\hfill & & [g_{21}1;g_{12},𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \end{array}$$
The assumption $`g_{21}3`$ (i.e. the gap between $`l_2`$ and $`l_1`$ is at least three) ensures that the exclusion rule does not apply in step (7).
There are 16 possible parity combinations for $`h_{21},d_{11},g_{12},g_{21}`$. In that language, the parities that occur in Example 3.4 above are
$`[0,0,0,1]`$ in (1), $`[1,0,1,0]`$ in (3), $`[0,1,0,1]`$ in (5) and $`[1,1,1,0]`$ in (7).
###### Example 3.5.
Next we assume $`d_{11}`$ to be even, $`h_{21}`$, $`g_{21}`$ even and $`g_{12}`$ odd.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [d_{11};𝐠_{\mathrm{𝟏𝟐}},g_{21};1]\hfill & & [h_{21}+d_{11};𝐠_{\mathrm{𝟏𝟐}},g_{21};1]\hfill & \\ & & +1& & \\ (2)& [𝐠_{\mathrm{𝟏𝟐}};g_{21},g_{12}+1;0]\hfill & & [h_{21};g_{12}+1,g_{21};0]\hfill & \\ & & & & +1\\ (3)& [𝐠_{\mathrm{𝟏𝟐}}𝐡_{\mathrm{𝟐𝟏}};g_{21},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐};1]\hfill & & [g_{12}+1;g_{21},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐};1]\hfill & \\ & & 1& & \\ (4)& [g_{21}1;𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐},𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏};0]\hfill & & [h_{21};𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐};0]\hfill & \\ & & & & +1\\ (5)& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐},g_{21};1]\hfill & & [𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐},g_{21};1]\hfill & \\ & & 1& & \\ (6)& [g_{12}+1;g_{21},g_{12}+1;0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{12}+1,g_{21};0]\hfill & \\ & & & & 1\\ (7)& [g_{12}h_{21}+1;g_{21},𝐠_{\mathrm{𝟏𝟐}};1]\hfill & & [𝐠_{\mathrm{𝟏𝟐}};g_{21},𝐠_{\mathrm{𝟏𝟐}};1]\hfill & \\ & & +1& & \\ (8)& [g_{21};𝐠_{\mathrm{𝟏𝟐}},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}};0]\hfill & \\ & & & & 1\\ (1)^{}& [g_{21}h_{21};𝐠_{\mathrm{𝟏𝟐}},g_{21};1]\hfill & & [g_{21};𝐠_{\mathrm{𝟏𝟐}},g_{21};1]\hfill & \end{array}$$
This is again a periodic orbit. The only situations where the exclusion rule can apply are in state (5) if $`g_{12}+2=1`$ (i.e. $`d(l_1,l_2)=1`$) or in state (3) if $`g_{12}+1(g_{12}h_{21})=h_{21}+1=1`$ (i.e. $`d(r_2,r_1)=1`$). Note that the odd numbered states are the interlaced ones. The parity configurations of the distances $`d(r_{i+1},r_i),d(r_i,l_j),d(l_j,l_{j+1})`$, $`d(l_{j+1},l_{j+2})`$ are
$`[0,0,1,0]`$ in (1), $`[1,1,0,1]`$ in (3), $`[0,1,1,0]`$ in (5) and $`[1,1,0,1]`$ in (7).
Now we describe the systems where the length of the lattice is increasing.
###### Example 3.6.
Let $`h_{21}`$ and $`d_{11}`$ be even, $`g_{12}`$ and $`g_{21}`$ odd.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [d_{11};𝐠_{\mathrm{𝟏𝟐}},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [h_{21}+d_{11};𝐠_{\mathrm{𝟏𝟐}},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & +1& & \\ (2)& [𝐠_{\mathrm{𝟏𝟐}};𝐠_{\mathrm{𝟐𝟏}},g_{12}+1;1]\hfill & & [h_{21};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \\ & & & & +1\\ (3)& [𝐠_{\mathrm{𝟏𝟐}}𝐡_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐};0]\hfill & & [g_{12}+1;𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐};0]\hfill & \\ & & 1& & \\ (4)& [g_{21}1;𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐},g_{21}1;1]\hfill & & [h_{21};g_{21}1,𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐};1]\hfill & \\ & & & & +1\\ (1)^{}& [g_{21}h_{21}1;𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [g_{21}1;𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟐},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \end{array}$$
After four interactions we return to the initial parity configuration, with $`\mathrm{\Delta }L=2`$ (the gap $`d(l_1,l_2)`$ grows by two). In the noninterlaced states we have the following parity configurations of the distances $`d(r_{i+1},r_i),d(r_i,l_j),d(l_j,l_{j+1})`$, $`d(l_{j+1},l_{j+2})`$:
$`[0,0,1,1]`$ in (1) and $`[1,1,1,1]`$ in (3).
###### Example 3.7.
In this case, $`h_{21}`$ is odd, $`d_{11}`$ even, $`g_{12}`$ and $`g_{21}`$ even.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [d_{11};g_{12},g_{21};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+𝐝_{\mathrm{𝟏𝟏}};g_{12},g_{21};0]\hfill & \\ & & +1& & \\ (2)& [g_{12};g_{21},𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏};1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}+\mathrm{𝟏},g_{21};1]\hfill & \\ & & & & 1\\ (3)& [g_{12}h_{21}1;g_{21},g_{12};0]\hfill & & [g_{12};g_{21},g_{12};0]\hfill & \\ & & +1& & \\ (4)& [g_{21};g_{12},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};1]\hfill & & [h_{21}+1;𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏},g_{12};1]\hfill & \\ & & & & +1\\ (1)^{}& [g_{21}h_{21}1;g_{12},g_{21}+2;0]\hfill & & [𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};g_{12},g_{21}+2;0]\hfill & \end{array}$$
This is another growing lattice. After four interactions, the gap between $`l_2`$ and $`l_1`$ has grown by two, with the other gaps unchanged. The parities are $`[1,0,0,0]`$ in (1) and $`[0,0,0,0]`$ in (3).
There are four remaining parity combinations. They belong to two systems with decreasing length. We describe them in the two examples below:
###### Example 3.8.
Let $`h_{21}`$, $`g_{12}`$ and $`g_{21}`$ be even, $`d_{11}`$ odd.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [𝐝_{\mathrm{𝟏𝟏}};g_{12},g_{21};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+𝐝_{\mathrm{𝟏𝟏}};g_{12},g_{21};0]\hfill & \\ & & 1& & \\ (2)& [𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};g_{21},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏},g_{21};1]\hfill & \\ & & & & 1\\ (3)& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{21},g_{12}2;0]\hfill & & [g_{12}2;g_{21},g_{12}2;0]\hfill & \\ & & 1& & \\ (4)& [𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{12}2,𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏};1]\hfill & & [h_{21}2;𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏},g_{12}2;1]\hfill & \\ & & & & +1\\ (1)^{}& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};g_{12}2,g_{21};0]\hfill & & [𝐠_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{12}2,g_{21};0]\hfill & \end{array}$$
After four interactions we return to the same parities but with $`\mathrm{\Delta }L=2`$ (gap $`g_{12}g_{12}2`$). In other words, this system is shrinking as long as no pairs occur. A pair can form in state (3). There the distance between particles $`r_2`$ and $`r_1`$ after $`k`$ cycles (1)-(4) is $`g_{12}2k(g_{21}h_{21}1)`$, which will eventually shrink to $`1`$. In this case we obtain pairs as in the system with label b) of Subsection 3.3 by switching the roles of the left-movers and the right-movers. The parities of this example are $`[0,1,0,0]`$ in (1) and $`[1,1,0,0]`$ in (3).
###### Example 3.9.
Finally, let $`d_{11}`$, $`g_{12}`$, $`g_{21}`$ be odd and $`h_{21}`$ be even.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & 1& & \\ (2)& [g_{12}1;𝐠_{\mathrm{𝟐𝟏}},g_{12}1;1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{12}1,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \\ & & & & 1\\ (3)& [g_{12}h_{21}1;𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟐};0]\hfill & & [𝐠_{\mathrm{𝟏𝟐}};𝐠_{\mathrm{𝟐𝟏}},g_{12}2;0]\hfill & \\ & & +1& & \\ (4)& [𝐠_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟐},g_{21}+1;1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{21}+1,𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟐};1]\hfill & \\ & & & & 1\\ (1)^{}& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟐},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [𝐠_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟐},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \end{array}$$
After four interactions, the length of the lattice has decreased by two (gap $`g_{12}g_{12}2`$). This is a shrinking system as long as no pairs appear. A pair appears in state (1), as soon as $`g_{12}2k=1`$ \[i.e. after $`k`$ cycles (1)-(4)\], when we obtain a pair as in the system with label b), cf. Subsection 3.3. The parities of this example are $`[0,1,1,1]`$ in (1) and $`[1,0,1,1]`$ in (3).
We now describe the patterns when pairs are present initially or form during the evolution.
### 3.3. Two Against Two With Pairs
There are six different configurations with pairs. The gap $`d_{11}`$ has to be odd if the pair is to affect the evolution of the system, and we assume $`g_{12}=1`$. Then we have the following parity combinations for the remaining distances:
$$\begin{array}{cccc}\text{Label}& h_{21}& & g_{21}\\ & & & \\ a)& 0& & 0\\ b)& 0& & 1\\ c)& 1& \text{with }h_{21}>1\hfill & 0\\ d)& 1& \text{with }h_{21}>1\hfill & 1\\ e)& 1& \text{with }h_{21}=1\hfill & 0\\ f)& 1& \text{with }h_{21}=1\hfill & 1\end{array}$$
###### Example 3.10 (Cases a),d) and c)).
We start with $`h_{21}`$ and $`g_{21}`$ even (i.e. a lattice of odd length). Assume that no other pair forms. The corresponding conditions are $`g_{21}(h_{21}1)3`$ in state (2), and $`g_{21}(g_{21}h_{21}1)=h_{21}+13`$ in state (7). We will discuss these other cases below.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1]\hfill & \\ & & +\mathrm{𝟏}& & \\ (2)& [g_{21};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0]\hfill & \\ & & & & +\mathrm{𝟏}\\ (3)& [g_{21}h_{21};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}+2;1]\hfill & & [𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}+2;1]\hfill & \\ & & +1& & \\ (4)& [𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};g_{21}+2,g_{12}+1;0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};g_{12}+1,g_{21}+2;0]\hfill & \\ & & 1& & \\ (5)& [𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};1]\hfill & \\ & & & & 1\\ (6)& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0]\hfill & & [𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};0]\hfill & \\ & & & & 1\\ (7)& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1]\hfill & & [g_{21};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1]\hfill & \\ & & +\mathrm{𝟏}& & \\ (8)& [g_{21};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0]\hfill & & [h_{21};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};0]\hfill & \\ & & & & +1\\ (9)& [g_{21}h_{21};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};1]\hfill & & [𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏},g_{12}+1;1]\hfill & \\ & & & & 1\\ (10)& [g_{21}h_{21}g_{12}1;g_{12}+1,g_{21};0]\hfill & & [g_{21};g_{12}+1,g_{21};0]\hfill & \end{array}$$
The first observation is that state (2) is as case d) (with particles $`r_1`$, $`r_2`$ switched) by the assumption that the gap between $`r_1`$ and $`r_2`$ is at least $`3`$.
Similarly, state (7) is as case c) The distance between $`r_2`$ and $`r_1`$ is odd and at least $`3`$ (by the assumptions), the distance between $`l_2`$ and $`l_1`$ is even.
Finally note that state (10) is as state (3) in Example 3.7. This means that the system evolves into a growing lattice as in Example 3.7.
Consider now other pairs of nearest neighbors.
Let $`g_{21}=h_{21}`$ (i.e. $`g_{21}(h_{21}1)=1`$ in (2)). Then the state (2) has the same parities as case f): the right-movers form a pair, as do the left-movers, and the gap between the particles $`r_2`$ and $`l_1`$ is odd. So for $`g_{21}=h_{21}`$ the system will evolve as f).
Let $`h_{21}+1=1`$. Then state (7) has the same parities as case e).
###### Example 3.11 (Case b)).
We start with $`h_{21}`$ even and $`g_{21}`$ odd (i.e. a lattice of even length).
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}+𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & +\mathrm{𝟏}& & \\ (2)& [𝐠_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟐𝟏}}=\mathrm{𝟏},g_{21}+1;1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}+1;1]\hfill & \\ & & & & +\mathrm{𝟏}\end{array}$$
Note that state (2) is as in a).
It remains to discuss cases e) and f). These are the cases where two pairs of nearest neighbors face each other with an odd distance between them. That means we have $`g_{12}=h_{21}=1`$, and necessarily $`g_{21}3`$.
###### Example 3.12.
Let particles $`r_2`$, $`r_1`$ be paired as well as particles $`l_1`$, $`l_2`$. Assume that $`g_{21}`$ is even (i.e. lattice of odd length). So $`g_{21}`$ is at least equal to four.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1]\hfill & & [h_{21}+d_{11};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21};1]\hfill & \\ & & +\mathrm{𝟏}& & +\mathrm{𝟏}\\ (2)& [g_{21};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏};1]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}=\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}}+\mathrm{𝟏},g_{12}+1;1]\hfill & \\ & & & & 1\\ (3)& [g_{21}h_{21}1;g_{12}+1,g_{21};0]\hfill & & [g_{21};g_{12}+1,g_{21};0]\hfill & \end{array}$$
Note that state (3) has the same parities as state (3) in Example 3.7. So the system evolves into a growing orbit as in Example 3.7.
###### Example 3.13.
Let $`r_2`$, $`r_1`$ and $`l_1`$, $`l_2`$ each be pairs, with $`g_{21}`$ odd ($`g_{21}3`$). The length of the lattice is then even.
$$\begin{array}{ccccc}& \text{Particle }r_1\hfill & & \text{Particle }r_2\hfill & \\ & & & & \\ (1)& [𝐝_{\mathrm{𝟏𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [h_{21}+d_{11};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & +\mathrm{𝟏}& & +\mathrm{𝟏}\\ (2)& [𝐠_{\mathrm{𝟐𝟏}};g_{12}+1,g_{21}+1;0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}}=\mathrm{𝟏};g_{21}+1,g_{12}+1;0]\hfill & \\ & & & & 1\\ (3)& [𝐠_{\mathrm{𝟐𝟏}}𝐡_{\mathrm{𝟐𝟏}}\mathrm{𝟏};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & & [𝐠_{\mathrm{𝟐𝟏}};g_{12}+1,𝐠_{\mathrm{𝟐𝟏}};1]\hfill & \\ & & 1& & \\ (4)& [𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & \\ & & & & +\mathrm{𝟏}\\ (5)& [g_{21}h_{21};g_{21}+1,𝐠_{\mathrm{𝟏𝟐}},1]\hfill & & [𝐠_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}+1;1]\hfill & \\ & & +1& & \\ (6)& [𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};g_{21}+1,g_{12}+1;0]\hfill & & [𝐡_{\mathrm{𝟐𝟏}};g_{12}+1,g_{21}+1;0]\hfill & \\ & & 1& & 1\\ (7)& [𝐠_{\mathrm{𝟐𝟏}};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},𝐠_{\mathrm{𝟐𝟏}};0]\hfill & & [𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};𝐠_{\mathrm{𝟐𝟏}},𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏};0]\hfill & \\ & & & & 1\\ (8)& [𝐠_{\mathrm{𝟐𝟏}}𝐠_{\mathrm{𝟏𝟐}}\mathrm{𝟏};𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}1;1]\hfill & & [g_{21}1;𝐠_{\mathrm{𝟏𝟐}}=\mathrm{𝟏},g_{21}1;1]\hfill & \end{array}$$
Note that state (8) has the same parities as state (1) in Example 3.12. So the system will switch to that example and then to a growing orbit (as in Example 3.7).
## 4. Conclusions
In this paper we have studied the lattice gas model with dynamical geometry introduced by Hasslacher and Meyer \[HM98\]. We first gave a general discussion of the pair-free evolution, establishing the importance of parity and the existence of many periodic orbits. The $`\sqrt{t}`$ growth of the lattice length observed in simulations is expected whenever the length grows by a constant amount, on average, as the particles complete one circuit. This is the case without pairs, and it is plausible that it is the general asymptotic behavior, but we have no proof of this. At present we can include the effects of pairs only by an exhaustive case-by-case analysis, which we carried out for systems of at most four particles. The pairs can form permanent “bound states”, and it seems promising to view them as quasiparticles. With an effective description of these quasiparticles, it should be possible to solve the dynamics of this model completely.
Although the microscopic dynamics of this model is reversible, it exhibits macroscopic “irreversibility” as the length grows for “most” states. It might be interesting to quantify what fraction of initial states having given $`N_L,N_R,L`$ ultimately grow.
The natural conjecture regarding the asymptotics of the model is that all nonperiodic solutions grow as $`\sqrt{t}`$, reflecting a constant average growth per circuit. Is this really true? Perhaps there are solutions with a characteristic time scale much longer than one circuit. For example, imagine a solution in which the length initially has average growth zero per circuit. A randomly chosen right-right gap would have size of order $`L/N_R`$ and could shrink to form a pair in a time $`L^2/N_RN_L`$. If this pair alters the evolution to produce constant average growth on this time scale, then $`dL/dt1/L^2`$ leads to $`Lt^{1/3}`$. Are there solutions with this behavior? Do the fluctuations in a solution growing as $`\sqrt{t}`$ include intervals when the growth is as $`t^{1/3}`$, correlated with the formation and destruction of pairs?
Finally, it should be straightforward to quantize this model lattice gas, generalizing \[M96, M97\]. The Hilbert space would be the direct sum $`H=_{L=1}^{\mathrm{}}H_L`$ of the Hilbert spaces for lattices of all fixed lengths $`L`$. Advection occurs within each $`H_L`$, but scattering causes transitions between them. The quantized model would be similar to that of \[BCV05\], but with evolution in discrete rather than continuous time. It is possible that the model, at least with few particles, is solvable by Bethe ansatz or other methods.
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# A Unified Framework for Tree Search Decoding: Rediscovering the Sequential Decoder
## 1 Introduction
Recent years have witnessed a growing interest in the closest lattice point search (CLPS) problem. This interest was primarily sparked by the connection between CLPS and maximum likelihood (ML) decoding in multiple-input multiple-output (MIMO) channels . On the positive side, MIMO channels offer significant advantages in terms of increased throughput and reliability. The price entailed by these gains, however, is a more challenging decoding task for the receiver. For example, naive implementations of the ML decoder have complexity that grows exponentially with the number of transmit antennas. This observation inspired several approaches for sub-optimal decoding that offer different performance-complexity tradeoffs (e.g., ).
Reduced complexity decoders are typically obtained by exploiting the codebook structure. The scenario considered in our work is no exception. In principle, the decoders considered here exploit the underlying lattice structure of the received signal to cast the decoding problem as a CLPS. Some variants of such decoders are known in the literature as sphere decoders (e.g., ). These decoders typically exploit number-theoretic ideas to efficiently span the space of allowed codewords (e.g., ). The complexity of such decoders were shown, via simulation and numerical analysis, to be significantly smaller than the naive ML decoder in many scenarios of practical interest (e.g., ). The complexity of the state of the art sphere decoder, however, remains prohibitive for problems characterized by a large dimensionality . This observation is one of the main motivations for our work.
The overriding goal of our work is to establish a general framework for the design and analysis of tree search algorithms for joint detection and decoding. Towards this goal, we first divide the decoding task into two interrelated stages; namely, 1) preprocessing and 2) tree search. The preprocessing stage is primarily concerned with exposing the underlying tree structure from the noisy received signal. Here, we discuss the integral roles of minimum mean square error decision feedback (MMSE-DFE) filtering, lattice reduction techniques, and relaxing the boundary control (i.e., lattice decoding) in tree search decoding. We then proceed to the search stage where a general framework based on the branch and bound (BB) algorithm is presented. This framework establishes, rigorously, the equivalence in terms of performance and complexity between different sphere and sequential decoders. We further use the proposed framework to classify the different search algorithms and identify their advantages/disadvantages. The MMSE-Fano decoder emerges as a special case of our general framework that enjoys a favorable performance-complexity tradeoff. We establish the superiority of the proposed decoder via numerical results and analytical arguments in several relevant scenarios corresponding to coded as well as uncoded transmission over MIMO and inter-symbol-interference (ISI) channels. More specifically, in our simulation experiments, we apply the tree search decoding framework to uncoded V-BLAST , linear dispersion space-time codes , algebraic space-time codes , and trellis codes over ISI channels . In all these cases, our results show that the MMSE-Fano decoder achieves near-ML performance with a much smaller complexity.
The rest of the paper is organized as follows. Section 2 introduces our system model and notation. In Section 3, we consider the design of the preprocessing stage and discuss the interplay between this stage and the tree search stage. In Section 4, we present a general framework for designing tree search decoders based on the branch and bound (BB) algorithm. In Section 5, we establish the superior performance-complexity tradeoff achieved by the proposed MMSE-Fano decoder, using analytical arguments and numerical results, in several interesting scenarios. Finally, we offer some concluding remarks in Section 6
## 2 System Model
We consider the transmission of lattice codes over linear channels with white Gaussian additive noise (AWGN). The importance of this problem stems from the fact that several very relevant applications arising in digital communications fall in this class, as it will be illustrated by some examples at the end of this section. Let $`\mathrm{\Lambda }\text{}^m`$ be an $`m`$-dimensional lattice, i.e., the set of points
$$\mathrm{\Lambda }=\{𝝀=\mathrm{𝐆𝐱}:𝐱\text{}^m\}$$
(1)
where $`𝐆\text{}^{m\times m}`$ is the lattice generator matrix. Let $`𝐯\text{}^m`$ be a vector and $``$ a measurable region in $`\text{}^m`$. A lattice code $`𝒞(\mathrm{\Lambda },𝐯,)`$ is defined as the set of points of the lattice translate $`\mathrm{\Lambda }+𝐯`$ inside the shaping region $``$, i.e.,
$$𝒞(\mathrm{\Lambda },𝐯,)=\{\mathrm{\Lambda }+𝐯\}.$$
(2)
Without loss of generality, we can also see $`𝒞(\mathrm{\Lambda },𝐯,)`$ as the set of points $`𝐜+𝐯`$, such that the codewords $`𝐜`$ are given by
$$𝐜=\mathrm{𝐆𝐱},\text{for}𝐱𝒰$$
(3)
where $`𝒰\text{}^m`$ is the code information set.
The linear additive noise channel is described, in general, by the input-output relation
$$𝐫=𝐇(𝐜+𝐯)+𝐳$$
(4)
where $`𝐫\text{}^n`$ denotes the received signal vector, $`𝐳𝒩(\mathrm{𝟎},𝐈)`$ is the AWGN vector, and $`𝐇^{n\times m}`$ is a matrix that defines the channel linear mapping between the input and the output.
Consider the following communication problem: a vector of information symbols $`𝐱`$ is generated with uniform probability over $`𝒰`$, the corresponding codeword $`𝐜=\mathrm{𝐆𝐱}`$ is produced by the encoder and the signal $`𝐜+𝐯`$ is transmitted over the channel (4). Assuming $`𝐇`$ and $`𝐯`$ known to the receiver, the ML decoding rule is given by
$$\widehat{𝐱}=\text{arg}\underset{𝐱𝒰}{\mathrm{min}}|𝐫\mathrm{𝐇𝐯}\mathrm{𝐇𝐆𝐱}|^2$$
(5)
The constraint $`𝒰^m`$ implies that the optimization problem in (5) can be viewed as a constrained version of the CLPS with lattice generator matrix given by $`\mathrm{𝐇𝐆}`$ and constraint set $`𝒰`$.
A few remarkable examples of the above framework are:
1. MIMO flat fading channels: One of the simplest and most widely studied examples is a MIMO V-BLAST system with squared QAM modulation, $`M`$ transmit and $`N`$ receive antennas, operating over a flat Rayleigh fading channel. The baseband complex received signal<sup>1</sup><sup>1</sup>1We use the superscript <sup>c</sup> to denote complex variables. in this case can be expressed as
$$𝐫^c=\sqrt{\frac{\rho }{M}}𝐇^c𝐜^c+𝐳^c$$
(6)
where the complex channel matrix $`𝐇^c^{N\times M}`$ is composed of i.i.d elements $`h_{i,j}^c𝒩_𝒞(0,1)`$, the input complex signal $`𝐜^c`$ has components $`𝐜_i^c`$ chosen from a unit-energy $`Q^2`$-QAM constellation, the noise has i.i.d. components $`z_i^c𝒩_𝒞(0,1)`$ and $`\rho `$ denotes the signal to noise ratio (SNR) observed at any receive antenna. The system model in (6) can be expressed in the form of (4) by appropriate scaling and by separating the real and imaginary parts using the vector and the matrix transformations defined by
$$𝐮^c𝐮=[\mathrm{Re}\{𝐮^c\}^^T,\mathrm{Im}\{𝐮^c\}^^T]^^T,$$
$$𝐌^c𝐌=\left[\begin{array}{cc}\mathrm{Re}\{𝐌^c\}& \mathrm{Im}\{𝐌^c\}\\ \mathrm{Im}\{𝐌^c\}& \mathrm{Re}\{𝐌^c\}\end{array}\right].$$
The resulting real model is given by (4) where $`n=2N`$, $`m=2M`$ and the constraint set is given by $`𝒰=\text{}_Q^m`$, with $`\text{}_Q=\{0,\mathrm{},Q1\}`$ denoting the set of integers residues modulo $`Q`$.
In the case of V-BLAST, the lattice code generator matrix $`𝐆=\kappa 𝐈`$, where $`\kappa `$ is a normalizing constant, function of $`Q`$, that makes the (complex) transmitted signal of unit energy per symbol. This formulation extends naturally to MIMO channels with more general lattice coded inputs . In general, a space-time code of block length $`T`$ is defined by a set of matrices $`𝐂^c=[𝐜_1^c,\mathrm{},𝐜_T^c]`$ in $`\text{}^{M\times T}`$. The columns of the codeword $`𝐂^c`$ are transmited in parallel on the $`M`$ transmit antennas in $`T`$ channel uses. The received signal is given by the sequence of vectors
$$𝐫_t^c=\sqrt{\frac{\rho }{M}}𝐇^c𝐜_t^c+𝐳_t^c,t=1,\mathrm{},T$$
(7)
Lattice space-time codes are obtained by taking a lattice code $`𝒞(\mathrm{\Lambda },𝐯,)`$ in $`\text{}^{2MT}`$, and mapping each codeword $`𝐜`$ into a complex matrix $`𝐂^c`$ according to some linear one-to-one mapping $`\text{}^{2MT}\text{}^{M\times T}`$. It is easy to see that a lattice-coded MIMO system can be again expressed by (4) where the channel matrix $`𝐇`$ is proportional (through an appropriate scaling factor) to the block-diagonal matrix
$$𝐈_T\left[\begin{array}{cc}\mathrm{Re}\{𝐇^c\}& \mathrm{Im}\{𝐇^c\}\\ \mathrm{Im}\{𝐇^c\}& \mathrm{Re}\{𝐇^c\}\end{array}\right]$$
(8)
In this case, we have $`n=2NT`$ and $`m=2MT`$. It is interesting to notice that for a wide class of linear dispersion (LD) codes , the information set $`𝒰`$ is still given by $`\text{}_Q^m`$, as in the simple V-BLAST case, although the generator matrix $`𝐆`$ is generally not proportional to $`𝐈`$. For other classes of lattice codes , with more involved shaping regions $``$, the information set $`𝒰`$ does not take on the simple form of an “hypercube”. For example, consider $`\mathrm{\Lambda }`$ obtained by construction A , i.e., $`\mathrm{\Lambda }=C+Q\text{}^m`$, where $`C\text{}_Q^m`$ is a linear code over $`\text{}_Q`$ with generator matrix in systematic form $`[𝐈,𝐏^^T]^^T`$. A generator matrix of $`\mathrm{\Lambda }`$ is given by
$$𝐆=\left[\begin{array}{cc}𝐈& \mathrm{𝟎}\\ 𝐏& Q𝐈\end{array}\right].$$
(9)
Typically, the shaping region $``$ of the lattice code $`𝒞(\mathrm{\Lambda },𝐯,)`$ can be an $`m`$-dimensional sphere, the fundamental Voronoi region of a sublattice $`\mathrm{\Lambda }^{}\mathrm{\Lambda }`$, or the $`m`$-dimensional hypercube. In all these cases, the information set $`𝒰`$ may be difficult to describe.
2. ISI Channels: For simplicity, we consider a baseband real single-input single-output (SISO) inter-symbol-interference (ISI) channel with the input and output sequences related by
$$r_i=\underset{\mathrm{}=0}{\overset{L}{}}h_{\mathrm{}}c_i\mathrm{}+z_i,$$
where $`(h_0,\mathrm{},h_L)`$ denotes the discrete-time channel impulse response, assumed of finite length $`L+1`$. The extension to the complex baseband model is immediate. Assuming that the transmitted signal is padded by $`L`$ zeroes, the channel can be written in the form (4) where the channel matrix takes on the tall banded Toeplitz form
$$𝐇=\left[\begin{array}{cccc}h_0& & & \\ h_1& h_0& & \\ \mathrm{}& \mathrm{}& \mathrm{}& \\ h_L& \mathrm{}& \mathrm{}& h_0\\ & h_L& \mathrm{}& h_1\\ & & \mathrm{}& \mathrm{}\\ & & & h_L\end{array}\right].$$
A wide family of trellis codes obtained as coset-codes , including binary linear codes, can be formulated as lattice codes where $`\mathrm{\Lambda }`$ is a Construction A lattice and the shaping region $``$ is chosen appropriately. In particular, coded modulation schemes based on the $`Q`$-PAM constellation obtained by mapping group codes over $`\text{}_Q`$ onto the $`Q`$-PAM constellation can be seen as lattice codes with hypercubic shaping $``$. The important case of binary convolutional codes falls in this class for $`Q=2`$. Again, the information set $`𝒰`$ corresponding to $``$ may, in general, be very complicated.
## 3 The Preprocessing Stage
In our framework, we divide the CLPS into two stages; namely, 1) preprocessing and 2) tree search. The complexity and performance of CLPS algorithms depend critically on the efficiency of the preprocessing stage. Loosely, the goal of preprocessing is to transform the original constrained CLPS problem, described by the lattice generator matrix $`\mathrm{𝐇𝐆}`$ and by the constraint set $`𝒰`$, into a form which is friendly to the search algorithm used in the subsequent stage. In the following, we discuss the different tasks performed in the preprocessing stage. In general, a friendly tree structure can be exposed through three steps: left preprocessing, right preprocessing, and forming the tree.
Some options for these three steps are illustrated in the following subsections. However, before entering the algorithmic details, it is worthwhile to point out some general considerations. The classical sphere decoding approach to the solution of the original constrained CLPS problem (5) consists of applying QR decomposition on the combined channel and code matrix, i.e., letting $`\mathrm{𝐇𝐆}=\mathrm{𝐐𝐑}`$ where $`𝐐\text{}^{n\times m}`$ has orthonormal columns and $`𝐑\text{}^{m\times m}`$ is upper triangular. Using the fact that $`𝐐^^T`$ is an isometry with respect to the Euclidean distance, (5) can be written as
$$\widehat{𝐱}=\text{arg}\underset{𝐱𝒰}{\mathrm{min}}\left|𝐲^{}\mathrm{𝐑𝐱}\right|^2$$
(10)
where $`𝐲^{}=𝐐^^T(𝐫\mathrm{𝐇𝐯})`$. If rank$`(\mathrm{𝐇𝐆})=m`$, $`𝐑`$ has non-zero diagonal elements and its triangular form can be exploited to search for all the points $`𝐱𝒰`$ such that $`\mathrm{𝐑𝐱}`$ is in a sphere of a given search radius centered in $`𝐲^{}`$. If the sphere is non-empty, the ML solution is guaranteed to be found inside the sphere, otherwise, the search radius is increased and the search is restarted. Different variations on this main theme have been proposed in the literature, and will be reviewed in Section 4 as special cases of a general BB algorithm. Nevertheless, it is useful to point out here the two main sources of inefficiency of the above approach: 1) It does not apply to the case rank$`(\mathrm{𝐇𝐆})<m`$ and, even when rank$`(\mathrm{𝐇𝐆})=m`$ but $`\mathrm{𝐇𝐆}`$ is ill-conditioned, the spread (or dynamic range) of the diagonal elements of $`𝐑`$ is large. This entails large complexity of the tree search . Intuitively, when $`\mathrm{𝐇𝐆}`$ is ill-conditioned, the lattice generated by $`\mathrm{𝐇𝐆}`$ has a very skewed fundamental cell such that there are directions in which it is very difficult to distinguish the points $`\{\mathrm{𝐇𝐆𝐱}:𝐱𝒰\}`$; 2) Enforcing the condition $`𝐱𝒰`$, can be very difficult because a lattice code $`𝒞(\mathrm{\Lambda },𝐯,)`$ with non-trivial shaping region $``$ might have an information set $`𝒰`$ with a complicated shape. Hence, just checking the condition $`𝐱𝒰`$ during the search may entail a significant complexity.
Left preprocessing can be seen as an effort to tackle the first problem: it modifies the channel matrix and the noise vector such that the resulting CLPS problem is non-equivalent to ML (therefore, it is suboptimal), but it has a much better conditioned “channel” matrix. The second problem can be tackled by relaxing the constraint set $`𝒰`$ to the whole $`\text{}^m`$, i.e., searching over the whole lattice $`\mathrm{\Lambda }`$ instead of only the lattice code $`𝒞`$ (or lattice decoding). In general, lattice decoding is another source of suboptimality. Nevertheless, once the boundary region is removed, we have the freedom of choosing the lattice basis which is more convenient for the search algorithm. This change of lattice basis is accomplished by right preprocessing. Finally, the tree structure is obtained by factorizing the resulting combined channel-lattice matrix in upper triangular form, as in classical sphere decoding. Overall, left and right preprocessing combined with lattice decoding are a way to reduce complexity at the expense of optimality. Fortunately, it turns out that an appropriate combination of these elements yields very significant saving in complexity with very small degradation with respect to the ML performance. Thus, it yields a very attractive decoding solution. While the outstanding performance of appropriate preprocessing and lattice decoding can be motivated via rigorous information theoretic arguments , here we are more concerned with the algorithmic aspects of the decoder and we shall give some heuristic motivation based on “signal-processing” arguments.
Finally, we note that the notion of complexity adopted in this work does not capture the complexity of the preprocessing stage (mostly cubic in the lattice dimension). In practice, this assumption is justified in slowly varying channels where the complexity of the preprocessing stage will be shared by many transmission frames (e.g., a wired ISI channel or a wireless channel with stationary terminals). If the number of these frames is large enough, i.e., the channel is slow enough, the preprocessing complexity can be ignored compared to the complexity of the tree search stage which has to be independently performed in every frame. Optimizing the complexity of the preprocessing stage, however, is an important topic, especially for fast fading channels.
### 3.1 Taming the Channel: Left Preprocessing
In the case of uncoded transmission ($`𝐆=𝐈`$), QR decomposition of the channel matrix $`𝐇`$ (assuming rank$`(𝐇)=m`$) allows simple recursive detection of the information symbols $`𝐱`$. Indeed, $`𝐐`$ is the feedforward matrix of the zero-forcing decision feedback equalizer (ZF-DFE) . In general, sphere decoders can be seen as ZF-DFEs with some reprocessing capability of their tentative decisions.
It is well-known that ZF-DFE is outperformed by the MMSE-DFE in terms of signal-to-interference plus noise ratio (SINR) at the decision point, under the assumption of correct decision feedback . This observation motivates the proposed approach for left preprocessing . This new matrix can be obtained through the QR decomposition of the augmented channel matrix
$`\stackrel{~}{𝐇}`$ $`\stackrel{}{=}`$ $`\left[\begin{array}{c}𝐇\\ 𝐈\end{array}\right]=\stackrel{~}{𝐐}𝐑_1`$ (13)
where $`\stackrel{~}{𝐐}\text{}^{(n+m)\times m}`$ has orthonormal columns and $`𝐑_1`$ is upper triangular. Let $`𝐐_1`$ be the upper $`n\times m`$ part of $`\stackrel{~}{𝐐}`$. $`𝐐_1`$ and $`𝐑_1`$ are the MMSE-DFE forward and backward filters, respectively. Thus, the transformed channel matrix and the received sequence are given by $`𝐑_1`$ and $`𝐲^{}=(𝐐_1^^T𝐫𝐑_1𝐯)`$, respectively. The transformed CLPS
$$\underset{𝐱𝒰}{\mathrm{min}}\left|𝐲^{}𝐑_1\mathrm{𝐆𝐱}\right|^2$$
(14)
is not equivalent to (5) since, in general, $`𝐐_1`$ does not have orthonormal columns. The additive noise $`𝐰=𝐲^{}𝐑_1\mathrm{𝐆𝐱}`$ in (14) contains both a Gaussian component, given by $`𝐐_1^^T𝐳`$, and a non-Gaussian (signal-dependent) component, given by $`(𝐐_1^^T𝐇𝐑_1)(𝐜+𝐯)`$. Nevertheless, for lattice codes such that cov$`(𝐜+𝐯)=𝐈`$, it can be shown that cov$`(𝐰)=𝐈`$ . Hence, the additive noise component $`𝐰`$ in (14) is still white, although non-Gaussian and data dependent. Therefore, the minimum distance rule (14) is expected to be only slightly suboptimal.<sup>2</sup><sup>2</sup>2This argument can be made rigorous by considering certain classes of lattices of increasing dimension, Voronoi shaping and random uniformly distributed dithering common to both the transmitter and the receiver, as shown in . On the other hand, the augmented channel matrix $`\stackrel{~}{𝐇}`$ in (13) has always rank equal to $`m`$ and it is well conditioned, since $`𝐑_1^^T𝐑_1=𝐈+𝐇^^T𝐇`$. Therefore, in some sense we have tamed the channel at the (small) price of the non-Gaussianity of the noise. The better conditioning achieved by the MMSE-DFE preprocessing is illustrated in Fig.2 (b) and (c).
### 3.2 Inducing Sparsity: Right Preprocessing
In order to obtain the tree structure, one needs to put $`𝐑_1𝐆`$ in upper triangular form $`𝐑`$ via QR decomposition. The sparser the matrix $`𝐑`$, the smaller the complexity of the tree search algorithm. For example, a diagonal $`𝐑`$ means that symbol-by-symbol detection is optimal, i.e., the tree search reduces to exploring a single path in the tree. Loosely, if one adopts a depth first search strategy, then a sparse $`𝐑`$ will lead to a better quality of the first leaf node found by the algorithm.<sup>3</sup><sup>3</sup>3More details on the different search strategies are reported in Section 4. Consequently, the algorithm finds the closest point in a shorter time .
While we have no rigorous method for relating the “sparsity” of $`𝐑`$ to the complexity of the tree search, inspired by decision feedback equalization in ISI channel, we define the sparsity index of the upper triangular matrix $`𝐑`$ as follows
$`\text{S}(𝐑)`$ $`\stackrel{}{=}`$ $`\underset{i\{1,\mathrm{},m\}}{\mathrm{max}}{\displaystyle \frac{_{j=i+1}^mr_{i,j}^2}{r_{i,i}^2}}.`$ (15)
where $`r_{i,j}`$ denotes the $`(i,j)`$-th element of $`𝐑`$. One can argue that the smaller $`\text{S}(𝐑)`$ the sparser $`𝐑`$ (e.g., $`\text{S}(𝐑)=0`$ for $`𝐑`$ diagonal). The goal of right preprocessing is to find a change of basis of the lattice $`\{𝐑_1\mathrm{𝐆𝐱}:𝐱\text{}^m\}`$, such that the new lattice generator matrix, $`𝐒`$, satisfies $`𝐒=\mathrm{𝐐𝐑}`$ with $`\text{S}(𝐑)`$ as small as possible. This amounts to finding a unimodular matrix $`𝐓`$ (i.e., the entries of $`𝐓`$ and $`𝐓^1`$ are integers) such that $`𝐑_1𝐆=\mathrm{𝐐𝐑𝐓}`$ with $`𝐐`$ unitary and $`\text{S}(𝐑)`$ minimized over the group of unimodular matrices. This optimization problem appears very difficult to solve; however, there exist many heuristic approaches to find unimodular matrices that give small values of $`\text{S}(𝐑)`$. Examples of such methods, considered here, are lattice reduction, column permutation and a combination thereof.
Lattice reduction finds a reduced lattice basis, i.e., the columns of the reduced generator matrix $`𝐒`$ have “minimal” norms and are as orthogonal as possible.<sup>4</sup><sup>4</sup>4For more details on the different notions and methods of lattice reduction, the reader is referred to . The most widely used reduction algorithm is due to Lenstra, Lenstra and Lovász (LLL) and has a polynomial complexity in the lattice dimension. An enhanced version of the LLL algorithm, namely the deep insertion modification, was later proposed by Schnorr and Euchner . LLL with deep insertion gives a reduced basis with significantly shorter vectors . In practice, the complexity of the LLL with deep insertion is similar to the original one even though it is an exponential time algorithm in the worst case sense .
Another method for decreasing $`\text{S}(𝐑)`$ consists of ordering the columns of $`𝐑_1𝐆`$, i.e., by right-multiplication by a permutation matrix $`𝚺`$. In the sequel, we shall use the V-BLAST greedy ordering strategy proposed in . This algorithm finds a permutation matrix $`𝚺`$ such that $`𝐑_1𝐆=\mathrm{𝐐𝐑}𝚺`$ maximizes $`\mathrm{min}_ir_{i,i}^2`$. Since $`𝐑^^T𝐑=𝚺^{}{}_{}{}^{T}𝐆_{}^{^T}𝐑_1^^T𝐑_1𝐆𝚺^1`$, i.e., the set $`\{_jr_{i,j}^2:i=1,\mathrm{},m\}`$ depends only on $`𝐑_1𝐆`$ and not on $`𝚺`$, by maximizing the minimum $`r_{i,i}^2`$ this algorithm miminizes $`\text{S}(𝐑)`$ over the group of permutation matrices (a subgroup of the unimodular matrices).
Lattice reduction and column permutation can be combined. This yields an unimodular matrix $`𝐓=𝚺𝐓_1`$, where $`𝐓_1`$ is obtained by lattice-reducing $`𝐑_1𝐆`$ and $`𝚺`$ by applying the V-BLAST greedy algorithm on the resulting reduced matrix $`𝐑_1\mathrm{𝐆𝐓}_1^1`$.
As observed before, the unimodular right multiplication does not change the lattice but may significantly complicate the boundary control. In fact, we have
$`\underset{𝐱𝒰}{\mathrm{min}}\left|𝐲^{}𝐑_1\mathrm{𝐆𝐱}\right|^2`$ $`=`$ $`\underset{𝐱𝒰}{\mathrm{min}}\left|𝐲^{}\mathrm{𝐐𝐑𝐓𝐱}\right|^2`$ (16)
$`=`$ $`\underset{𝐱𝐓𝒰}{\mathrm{min}}\left|𝐐^^T𝐲^{}\mathrm{𝐑𝐱}\right|^2`$
The new constraint set $`𝐓𝒰`$ might be even more complicated to enforce than the original information set $`𝒰`$ (see Fig. 2(d)). However, it is clear that although modifying the boundary control may result in a significant complexity increase for ML decoding, lattice decoding is not affected at all, since $`𝐓\text{}^m=\text{}^m`$.
### 3.3 Forming the Tree
The final step in preprocessing is to expose the tree structure of the problem. In this step, QR decomposition is applied on the transformed combined channel and lattice matrix $`𝐐_1^^T\mathrm{𝐇𝐆𝐓}^1`$, after left and right preprocessing. The upper triangular nature of $`𝐑`$ means that a tree search can now be used to solve the CLPS problem. Fig. 1 illustrates an example of such a tree.
Here, we wish to stress that our approach for exposing the tree is fundamentally different from the one traditionally used for codes over finite alphabets (e.g., linear block codes, convolutional codes, trellis coset codes in AWGN channels). Here, we operate over the field of real numbers and consider the lattice corresponding to the joint effect of encoding and channel distortion. In the conventional approach, the tree is generated from the trellis structure of the code alone, and hence, does not allow for a natural tree search that handles jointly detection (the linear channel) and decoding. In fact, joint detection and decoding is achieved at the expenses of an increase of the overall system memory (joint trellis), or by neglecting some paths in the search (e.g., by per-survivor reduced state processing). Since operating on the full joint trellis is usually too complex, both the proposed and the conventional per-survivor (reduced state) approach are suboptimal, and the matter is to see which one achieves the best performance/complexity tradeoff.
For the sake of convenience, in the following we shall denote again by $`𝐲`$ the channel output after all transformations, i.e., the tree search is applied to the CLPS problem $`\mathrm{min}_{𝐱\text{}^m}|𝐲\mathrm{𝐑𝐱}|^2`$ with $`𝐑`$ in upper triangular form. The components of vectors and matrices are numbered in reverse order, so that the preprocessed received signal can finally be written as
$$\left(\begin{array}{c}y_m\\ \mathrm{}\\ y_1\end{array}\right)=\left(\begin{array}{cccc}r_{m,m}& \mathrm{}& \mathrm{}& r_{m,1}\\ 0& r_{m1,m1}& \mathrm{}& r_{m1,1}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& r_{1,1}\end{array}\right)\left(\begin{array}{c}x_m\\ \mathrm{}\\ x_1\end{array}\right)+\left(\begin{array}{c}w_m\\ \mathrm{}\\ w_1\end{array}\right).$$
(17)
Notice that after preprocessing the problem is always squared, of dimension $`m`$, even though the original problem has arbitrary $`m`$ and $`n`$. Throughout the paper, we consider a tree rooted at a fixed dummy node $`x_0`$. The node at level $`k`$ is denoted by the label $`𝐱_1^k=(x_1,x_2,\mathrm{},x_k)`$. Moreover, every node $`𝐱_1^k`$ is associated with the the squared distance
$$w_k(𝐱_1^k)=\left|y_k\underset{j=1}{\overset{k}{}}r_{k,j}x_j\right|^2.$$
(18)
The difference between the transmitted codeword $`\widehat{𝐱}`$ and any valid codeword $`𝐱`$ is denoted by $`\stackrel{~}{𝐱}`$, i.e., $`\stackrel{~}{𝐱}=\widehat{𝐱}𝐱`$.
We hasten to stress that the preprocessing steps highlighted in Sections 3.1-3.3 are for a general setting. In some special cases, some steps can be eliminated or alternative options can be used. Some of these cases are listed hereafter.
1. Upper Triangular Code Generator Matrix:
In this case, after taming the channel, $`𝐇𝐑_1`$, the new combined matrix $`𝐑_1𝐆`$ is also upper triangular and can be directly used to form the tree without any further preprocessing (if one decides against right preprocessing).
2. Uncoded V-BLAST:
For uncoded V-BLAST systems (i.e., $`𝐆=𝐈`$), applying the MMSE-DFE greedy ordering of may achieve better complexity of the tree search stage than applying MMSE-DFE left preprocessing, lattice reduction, and greedy ordering of the final QR decomposition. This is especially true for large dimensions, where lattice reduction is less effective .
3. The Hermite Normal Form Transformation:
Ultimately, any hardware implementation of the decoder requires finite arithmetics. In this case, all quantities are scaled and quantized such that they take on integer values. While all the preprocessing steps in Sections 3.1-3.3 can be easily adapted to finite arithmetics, there exist other efficient transformations for integral matrices that may yield smaller complexity over the ones mentioned above. For example, one can apply the Hermite normal form (HNF) directly on the scaled (and quantized) matrix $`\tau \mathrm{𝐇𝐆}`$, such $`\tau \mathrm{𝐇𝐆}=\mathrm{𝐑𝐓}`$, with $`𝐓`$ unimodular and $`𝐑`$ upper triangular with the property that each diagonal element dominates the rest of the entries on the same row (i.e., $`r_{i,i}>r_{i,j}0,i=1,\mathrm{},m,j=i+1,\mathrm{},m`$). Interestingly, the HNF transformation improves the sparsity index and reduces the preprocessing to a single step.
## 4 The Tree Search Stage
After proper preprocessing, the second stage of the CLPS corresponds to an instance of searching for the best path in a tree. In this setting, the tree has a maximum depth $`m`$, and the goal is to find the node(s) at level $`m`$ that has the least squared distance, where the squared distance for any node $`𝐱_1^m`$ at level $`m`$ (called leaf node) is given by
$`d^2(𝐱,𝐱_1^m)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{m}{}}}w_i(𝐱_1^i)`$ (19)
Visiting all leaf nodes to find the one with the least metric, is either prohibitively complex (exponential in $`m`$), or not possible, as with lattice decoding. The complexity of tree search can be reduced by the branch and bound (BB) algorithm which determines if an intermediate node $`𝐱_1^k`$, on extending, has any chance of yielding the desired leaf node. This decision is taken by comparing the cost function assigned to the node by the search algorithm, against a bounding function. In the following section, we propose a generic tree search stage, inspired by the BB algorithm, that encompasses many known algorithms for CLPS as its special cases. We further use this algorithm to classify various tree search algorithms and elucidate some of their structural properties.
### 4.1 Generic Branch and Bound Search Algorithm
Before describing the proposed algorithm, we first need to introduce some more notation.
* ACTIVE is an ordered list of nodes.
* $`f(𝐱_1^k)`$ is the cost function of any node $`𝐱_1^k`$ in the tree, and $`𝐭^{m\times 1}`$ is the bounding function.
* Any node $`𝐱_1^k`$ in the search space of the search algorithm is a valid node, if $`f(𝐱_1^k)<t_k`$.
* A node is generated by the search algorithm, if the node occupies any position in ACTIVE at some instant during the search.
* sort” is a rule for ordering the nodes in the list ACTIVE.
* gen” is a rule defining the order for generating the child nodes of the node being extended.
* $`g_1`$ and $`g_2`$ are rules for tightening the bounding function.
* At any instant, the leaf node with the least distance generated by the search algorithm so far in the search process is stored in $`\widehat{𝐱}`$.
* We define the search complexity of a tree search algorithm as the number of nodes generated by the algorithm.
* Two search algorithms are said to be equivalent if they generate the same set of nodes.
* A BB algorithm whose solution is guaranteed to be (one of) leaf node(s) with least distance is called an optimal BB algorithm. If the solution is not guaranteed to have the least distance to $`𝐲`$ among all leaf nodes, then the BB algorithm is a heuristic BB algorithm.
We are now ready to present our generic search algorithm.
GBB($`f`$, $`𝐭`$, sort, gen, $`g_1`$, $`g_2`$):
1. Create the empty list ACTIVE, and place the root node in ACTIVE. Set $`n_c1`$.
2. Let $`𝐱_1^k`$ be the top node of ACTIVE.
If $`𝐱_1^k`$ is a leaf node ($`k=m`$), then
$`𝐭g_1(𝐭,f(𝐱_1^m))`$ and $`\widehat{𝐱}\text{arg}\mathrm{min}(_{i=1}^mw_i(𝐱_1^i),_{i=1}^mw_i(\widehat{𝐱}_1^i))`$.
Remove $`𝐱_1^m`$ from ACTIVE.
Go to step 4.
If $`𝐱_1^k`$ is not a valid node, then remove $`𝐱_1^k`$ from ACTIVE. Go to step 4.
If all valid child nodes of $`𝐱_1^k`$ have already been generated, then remove $`𝐱_1^k`$ from ACTIVE. Go to step 4.
Generate a valid child node $`𝐱_1^{k+1}`$ of $`𝐱_1^k`$, not generated before, according to the order gen, and place it in ACTIVE. Set $`n_cn_c+1`$. Set $`𝐭g_2(𝐭,n_c,\text{ACTIVE})`$. Update $`f(𝐱_1^k),f(𝐱_1^{k+1})`$.
3. Sort the nodes in ACTIVE according to sort.
4. If ACTIVE is empty, then exit. Else, Go to step 2.
In GBB, $`g_1`$ allows one to tighten the bounding function when a leaf node reaches the top of ACTIVE, whereas $`g_2`$ allows for restricting the search space in heuristic BB algorithms. For example, setting
$$g_2(𝐭,n_{c,t},\text{ACTIVE})=[\mathrm{},\mathrm{},\mathrm{},\mathrm{}]^^T,$$
will force the search algorithm to terminate when the number of nodes generated increases beyond a tolerable limit on the complexity given by $`n_{c,t}`$. Whenever a leaf node reaches the top of ACTIVE, $`\widehat{𝐱}`$ is updated if appropriate. Now, we use GBB to classify various tree search algorithms in three broad categories. This classification highlights the structural properties and advantages/disadvantages of the different search algorithms.
#### 4.1.1 Breadth First Search
GBB becomes a Breadth First Search (BrFS) if $`g_1(𝐭,f(𝐱_1^m))=𝐭`$, and the cost function $`f`$ of any node, once determined, is never updated. Ultimately, all nodes $`𝐱_1^k`$ whose cost function along the path $`𝐱_1^k`$ does not rise above the bounding function, are generated before the algorithm terminates, unless the $`g_2`$ function removes their parent nodes from ACTIVE. Now, we can establish the equivalence between various sphere/sequential decoders and BrFS.
The first algorithm is the Pohst enumeration strategy reported in . In this strategy, the bounding function $`𝐭`$ consists of equal components $`C_0`$, where $`C_0`$ is a constant chosen before the start of search<sup>5</sup><sup>5</sup>5For the sake of simplicity, we assumed in the above classification that the bounding function is chosen such that at least one leaf node is found before the search terminates. If, however, no leaf node is found before the search terminates, the bounding function is relaxed and the search is started afresh., and the cost function of a node $`𝐱_1^k`$ is $`f(𝐱_1^k)=_{i=1}^kw_i(𝐱_1^i)`$. Therefore, all nodes $`𝐱_1^k`$ in the search space that satisfy
$$\underset{i=1}{\overset{k}{}}w_i(𝐱_1^i)C_0$$
(20)
are generated before termination. Generating the child nodes in this strategy is simplified by the following observation. For any parent node $`𝐱_1^k`$, the condition $`{\displaystyle \underset{i=1}{\overset{k+1}{}}}w_i(𝐱_1^i)C_0`$ for the set of generated child nodes implies that the $`(k+1)`$th component of the generated child nodes lies in some interval $`[a_0,a_1]`$. The second example is the statistical pruning (SP) decoder which is equivalent to a heuristic BrFS decoder. Two variations of SP are proposed in , the increasing radii (IR) and elliptical pruning (EP) algorithms. The IR algorithm is a BrFS with the bounding function $`𝐭=\{t_1,\mathrm{},t_m\}`$, where $`t_k,1km`$ are constants chosen before the start of search. The cost function for any node in IR is the same as in Pohst enumeration. The EP algorithm is given by the bounding function $`𝐭=\{1,\mathrm{},1\}`$, and the cost function for the node $`𝐱_1^k`$ given by $`f(𝐱_1^k)={\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \frac{w_i(𝐱_1^i)}{e_k}}`$, where $`e_k,1km`$ are constants. More generally, when $`g_2(.)=𝐭`$, i.e., $`g_2`$ is not used, the resulting BrFS algorithm is equivalent to the Wozencraft sequential decoder where, depending on the cost function, the decoder can be heuristic or optimal.
The $`M`$ algorithm and $`T`$-algorithm are also examples of heuristic BrFS. Here, however, $`g_2`$ serves an important role in restricting the search space. In both algorithms, sort is defined as follows. Any node in ACTIVE at level $`k`$ is placed above any node at level $`k+1`$, and nodes in the same level are sorted in ascending order of their cost functions. In the $`M`$-algorithm, after the first node at level $`k+1`$ is generated (indicating that all valid nodes at level $`k`$ have already been generated), $`g_2`$ sets $`t_k`$ to the cost function of the $`M`$-th node at level $`k`$ (where $`M`$ is an initial parameter of the $`M`$-algorithm). In the $`T`$-algorithm, $`g_2`$ sets $`t_k`$ to $`(f(\overline{𝐱}_1^k)+T)`$, where $`\overline{𝐱}_1^k`$ is the top node at level $`k`$ in ACTIVE, and $`T`$ is a parameter of the $`T`$-algorithm. After $`t_k`$ is tightened in this manner, all nodes in ACTIVE at level $`k`$, that satisfy $`f(𝐱_1^k)>t_k`$ are rendered invalid, and are subsequently removed from ACTIVE.
In general, BrFS algorithms are naturally suited for applications that require soft-outputs, as opposed to a hard decision on the transmitted frame. The reason is that such algorithms output an ordered<sup>6</sup><sup>6</sup>6The list is ordered based on the cost function of the different candidates list of candidate codewords. One can then compute the soft-outputs from this list using standard techniques (e.g., ,). Here, we note that in the proposed joint detection and decoding framework, soft outputs are generally not needed. Another advantage of BrFS is that the complexity of certain decoders inspired by this strategy is robust against variations in the SNR and channel conditions. For example, the $`M`$-algorithm has a constant complexity independent of the channel conditions. This property is appealing for some applications, especially those with hard limits on the maximum, rather than average, complexity. On the other hand, decoders inspired by the BrFS strategy usually offer poor results in terms of the average complexity, especially at high SNR. One would expect a reduced average complexity if the bounding function is varied during the search to exploit the additional information gained as we go on. This observation motivates the following category of tree search algorithms.
#### 4.1.2 Depth First Search
GBB becomes a depth first search (DFS) when the following conditions are satisfied. The sorting rule sort orders the nodes in ACTIVE in reverse order of generation, i.e., the last generated node occupies the top of ACTIVE, and
$$g_1(𝐭,f(𝐱_1^m))=[\mathrm{min}(t_1,f(𝐱_1^m)),\mathrm{},\mathrm{min}(t_m,f(𝐱_1^m))]^^T.$$
As in BrFS, the cost function of any node, once generated, remains constant. Even among algorithms within the class of DFS algorithms, other parameters, like gen and $`g_2`$, can significantly alter the search behavior. To illustrate this point, we contrast in the following several sphere decoders which are equivalent to DFS strategies.
The first example of such decoders is the modified Viterbo-Boutros (VB) decoder reported in . In this decoder, $`g_2(𝐭,n_c)=𝐭`$, and the cost function for any node $`𝐱_1^k`$ is $`f(𝐱_1^k)=_{i=1}^kw_i(𝐱_1^i)`$. For any node $`𝐱_1^k`$ and its corresponding interval $`[a_0,a_1]`$ for valid child nodes, the function gen generates the child node with $`a_0`$ as its $`(k+1)^{th}`$ component first. Our second example is the Schnorr Euchner (SE) search strategy first reported in . This decoder shares the same cost functions and $`g_2`$ with the modified VB decoder, but differs from it in the order of generating the child nodes. For any node $`𝐱_1^k`$ and its corresponding interval $`[a_0,a_1]`$ for the valid child nodes, let $`a_m\stackrel{}{=}\frac{a_0+a_1}{2}`$ and $`\delta \stackrel{}{=}\text{sign}(w_{k+1}(𝐱_1^{k+1}))`$. Then, the function gen in the SE decoder generates nodes according to the order $`\{a_m,a_m+\delta ,a_m\delta ,a_m+2\delta ,\mathrm{}\}`$.
Due to the adaptive tightening of the bounding function, DFS algorithms have a lower average complexity than the corresponding BrFS algorithms with the same cost functions, especially at high SNR. Another advantage of the DFS approach is that it allows for greater flexibility in the performance-complexity tradeoff through carefully constructed termination strategy. For example, if we terminate the search after finding the first leaf node, i.e., $`n_c=m`$, then we have the MMSE-Babai point decoder . This decoder corresponds to the MMSE-DFE solution aided with the right preprocessing stage. It was shown in that the performance of this decoder is within a fraction of a dB from the ML decoder in systems with small dimensions. The fundamental weakness of DFS algorithms is that the sorting rule is static and does not exploit the information gained thus far to speed up the search process.
#### 4.1.3 Best First Search
GBB becomes a best first search (BeFS) when the following conditions are satisfied. The nodes in ACTIVE are sorted in ascending order of their cost functions, and
$$g_1(𝐭,f(𝐱_1^m))=[\mathrm{min}(t_1,f(𝐱_1^m)),\mathrm{},\mathrm{min}(t_m,f(𝐱_1^m))]^^T.$$
Note that in BeFS, the search can be terminated once a leaf node reaches the top of the list, since this means that all intermediate nodes have cost functions higher than that of this leaf node. Thus, the bounding function is tightened just once in this case. The stack algorithm is an example of BeFS decoder obtained by setting $`g_2(𝐭,n_c)=𝐭`$, and the cost function of any node in ACTIVE at any instant defined as follows: If $`𝐱_1^k`$ is a leaf node, then $`f(𝐱_1^k)=\mathrm{}`$. Otherwise, we let $`𝐱_{1,g}^{k+1}`$ be the best child node of $`𝐱_1^k`$ not generated yet, and define $`f(𝐱_1^k)=_{i=1}^{k+1}w_i(𝐱_{1,g}^{k+1})b(k+1)`$, where we refer to $`b^+`$ as the bias. Because of the efficiency of the sorting rule, BeFS algorithms are generally more efficient than the corresponding BrFS and DFS algorithms. This fact is formalized in the following theorems. Theorem 1 establishes the efficiency of the stack decoder with $`b=0`$ among all known sphere decoders.
###### Theorem 1
The stack algorithm with $`b=0`$ generates the least number of nodes among all optimal tree search algorithms.
The following result compares the heuristic stack algorithm, i.e., $`b>0`$, with a special case of the IR algorithm , where the bounding function takes the form $`t_k=bk+\delta `$.
###### Theorem 2
The IR algorithm with cost function $`\{𝐭:t_k=bk+\delta \}`$, generates at least as many nodes as those generated by the stack algorithm when the same bias $`b`$ is used.
Proof: Appendix C
At this point, it is worth noting that in our definition of search complexity, we count only the number of generated nodes, i.e., nodes that occupy some position in ACTIVE at some instant. In general, this is a reasonable abstraction of the actual computational complexity involved. However, in the stack algorithm, for each node generated, the cost functions of two nodes are updated instead of one; one for the generated node, and one for the parent node. Thus, the comparisons in Theorems 1 and 2 are not completely fair.
Finally, we report the following two advantages offered by the the stack algorithm. First, it offers a natural solution for the problem of choosing the initial radius (or radii), which is commonly encountered in the design of sphere decoders (e.g., ). By setting all the components of $`𝐭`$ to $`\mathrm{}`$, it is easy to see that we are guaranteed to find the closest lattice point while generating the minimum number of nodes (among all search algorithms that guarantee finding the closest point). Second it allows for a systematic approach for trading-off performance for complexity. To illustrate this point, if we set $`b=0`$, we obtain the closest point lattice decoder (i.e., best performance but highest complexity). On the other extreme, when $`b\mathrm{}`$, the stack decoder reduces to the MMSE-Babai point decoder discussed in the DFS section (the number of nodes visited is always equal to $`m`$). In general, for systems with small $`m`$, one can obtain near-optimal performance with a relatively large values of $`b`$. As the number of dimensions increases, more complexity must be expended (i.e., smaller values of $`b`$) to approach the optimal performance.
### 4.2 Iterative Best First Search
In Section 4.1, our focus was primarily devoted to complexity, defined as the number of nodes visited by the tree search algorithm. Another important aspect is the memory requirement entailed by the search. Straightforward implementation of the GBB algorithm requires maintaining the list ACTIVE, which can have a prohibitively long length in certain application. This motivates the investigation of modified implementations of these search strategies that are more efficient in terms of storage requirements. The BrFS and DFS sphere decoders discussed in Sections 4.1.1 and 4.1.2 lend themselves naturally to storage efficient implementations. Such implementations have been reported in .
In order to exploit the complexity reduction offered by BeFS strategy in practice, it is therefore important to seek modified memory-efficient implementations of such algorithms. This can be realized by storing only one node at a time, and allowing nodes to be visited more than once. The search in this case progresses in contours of increasing bounding functions, thus allowing more and more nodes to be generated at each step, finally terminating once a leaf node is obtained. The Fano decoder is the iterative BeFS variation of the stack algorithm. Although the stack algorithm and the Fano decoder, with the same cost functions, generate essentially the same set of nodes , the Fano decoder visits some nodes more than once. However, the Fano decoder requires essentially no memory, unlike the stack algorithm. Appendix A provides an algorithmic description of the Fano decoder and a brief description of the relevant parameters. Overall, the proposed decoder consists of left preprocessing (MMSE-DFE) and right preprocessing (combined lattice reduction and greedy ordering), followed by the Fano (or stack) search stage for lattice, not ML, decoding.
## 5 Analytical and Numerical Results
To illustrate the efficiency and generality of the proposed framework, we utilize it in three distinct scenarios. First, we consider uncoded transmission over MIMO channels (i.e., V-BLAST). Here, we present analytical, as well as simulation, results that demonstrate the excellent performance-complexity tradeoff achieved by the proposed Stack and Fano decoders. Then, we proceed to coded MIMO systems and apply tree search decoding to two different classes of space-time codes. Finally, we conclude with trellis coded transmission over ISI channels.
### 5.1 The V-BLAST Configuration
Unfortunately, analytical characterization of the performance- complexity tradeoff for sequential/sphere decoders with arbitrary $`\mathrm{𝐇𝐆}`$ and $`𝒰`$ still appears intractable. To avoid this problem, we restrict ourselves in this section to uncoded transmission over flat Rayleigh MIMO channels. In our analysis, we further assume that ZF-DFE pre-processing is used. The complexity reductions offered by the proposed preprocessing stage are demonstrated by numerical results.
###### Theorem 3
The Stack algorithm and the Fano decoder with any finite bias $`b`$, achieve the same diversity as the ML decoder when applied to a V-BLAST configuration.
Proof : Appendix D.
The result shows that the Fano decoder, unlike other heuristic algorithms like nulling-and-canceling, does not lead to a lower diversity than the ML decoder.
###### Theorem 4
In a V-BLAST system with $`Q^2`$-QAM, the average complexity per dimension of the stack algorithm for a sufficiently large bias $`b`$ is linear in $`m`$ when the SNR $`\rho `$ grows linearly with $`m`$ and $`r=nm0`$.
Proof : Appendix E.
Thus, one can achieve linear complexity with the stack algorithm by allowing the SNR to increase linearly with the lattice dimension. To validate our theoretical claims, we further report numerical results in selected scenarios. In our simulations, we assume that the channel matrix is square and choose the SE enumeration as the reference sphere decoder for comparison purposes. In all the figures, the subscript $`Z`$ refers to ZF-DFE left preprocessing and the subscript $`M`$ denotes MMSE-DFE left preprocessing followed by LLL reduction and V-BLAST greedy ordering for right preprocessing. In Fig. 3, the average complexity per lattice dimension and frame error rate of Fano decoder with $`b=1`$ and the SE sphere decoder are shown for different values of SNR in a $`20\times 20`$ $`16`$QAM V-BLAST system. Thus, for $`m=40`$, the Fano decoder can offer a reduction in complexity up-to a factor of 100. Moreover, the performance of the the Fano decoder is seen to be only a fraction of a dB away from that of the SE decoder, which achieves ML performance. We also see that the frame error rate curves for both the Fano decoder and the SE (ML) decoder have the same slope in the high SNR region, as expected from our analysis. Fig. 4 compares the complexity and performance of the Fano decoder with ZF-DFE and MMSE-DFE based preprocessing, respectively, in a $`30\times 30`$ $`4`$QAM V-BLAST system (i.e., $`m=60`$). From the figures, we see that the MMSE-DFE based preprocessing plays a crucial role in lowering the search complexity of the Fano decoder, despite the apparent increase in search space due to lattice decoding. Fig. 5 reports the dependence of the complexity of the Fano decoder on the value of $`b`$. The complexity attains a local minimum for some $`b^{}>1`$, and for large values of $`b`$, the complexity of the Fano decoder decreases as $`b`$ is increased. The error rate, however, increases monotonically with $`b`$ and approaches that of the MMSE-DFE Babai decoder as $`b\mathrm{}`$. For small dimensions, the performance of the MMSE-DFE based Babai decoder is remarkable. This DFS decoder terminates after finding the first leaf node. Fig. 6 compares the performance of this decoder with the ML performance for a $`4\times 4`$, $`4`$QAM V-BLAST system. We also report the performance of the Yao-Wornell and Windpassinger-Fischer (YWWF) decoder which has the same complexity as the MMSE-DFE Babai decoder . It is shown that the performance of the proposed decoder is within a fraction of a dB from that of ML decoder, whereas the algorithm in exhibits a loss of more than $`3`$ dB.
### 5.2 Coded MIMO Systems
In this section, we consider two classes of space-time codes. The first class is the linear dispersion (LD) codes which are obtained by applying a linear transformation (over $``$) to a vector of PAM symbols. For convenience, we follow the set-up of Dayal and Varanasi where two variants of the threaded algebraic space-time (TAST) constellations are used in a $`3\times 1`$ MIMO channel. This setup also allows for demonstrating the efficiency of the MMSE-DFE frontend in solving under-determined systems. In , the rate-$`1`$ TAST constellation uses $`64`$-QAM inputs at a rate of one symbol per channel use. The rate-$`3`$ TAST constellation, on the other hand, uses $`4`$-QAM inputs to obtain the same throughput as the rate-$`1`$ constellation. As observed in , one obtains a sizable performance gain when using rate-$`3`$ TAST constellation under ML decoding. The main disadvantage, however, of the rate-3 code is that it corresponds to an under-determined system with $`6`$ excess unknowns which significantly complicates the decoding problem. Fig. 7 shows that the performance of the proposed MMSE-DFE lattice decoder is less than $`0.1`$ dB away from the ML decoder for both cases. In order to quantify the complexity reduction offered by our approach, compared with the generalized sphere decoder (GSD) used in , we measure the average complexity increase with the excess dimensions. If we define
$$\gamma \stackrel{}{=}\frac{\text{Average complexity of decoding rate-}3\text{ constellation}}{\text{Average complexity of decoding rate-}1\text{ constellation}},$$
(21)
then a straightforward implementation of the GSD, as outlined in for example, would result in $`\gamma =𝒪\left(4^6\right)`$. In fact, even with the modification proposed in , Dayal and Varanasi could only bring this number down to $`\gamma =460`$ at an SNR of $`30`$ dB. In Table 1, we report $`\gamma `$ for the proposed algorithm at different SNRs, where one can see the significant reduction in complexity (i.e., from $`460`$ to $`12`$ at an SNR of $`30`$ dB). Based on experimental observations, we also expect this gain in complexity reduction to increase with the excess dimension $`mn`$.
The second space-time coding class is the algebraic codes proposed in . This approach constructs linear codes, over the appropriate finite domain, and then the encoded symbols are mapped into QAM constellations. The QAM symbols are then parsed and appropriately distributed across the transmit antennas to obtain full diversity. It has been shown that the complexity of ML decoding of this class of codes grows exponentially with the number of transmit antennas and data rates. Here, we show that the proposed tree search framework allows for an efficient solution to this problem. Figure 8 shows the performance of MMSE-DFE lattice decoding for two such constructions of space-time codes i.e., Golay space-time code for two transmit antennas and the companion matrix code for three transmit antennas . In both case, the performance of the MMSE-DFE lattice decoder is seen to be essentially same as the ML performance. In the proposed decoder, we use the lattice $`𝚲`$ obtained from underlying algebraic code through construction A. The ML performance, obtained via exhaustive search in Figure 8, is not feasible for higher dimensions due to exponential complexity in the number of dimensions.
### 5.3 Coded Transmission over ISI Channels
In this section, we compare the performance of the MMSE-Fano decoder with the Per-Survivor-Processing (PSP) algorithm for convolutionally coded transmission over ISI channels. Our MMSE-Fano decoder uses the construction A lattice obtained from the convolutional code. For this scenario, it is known that PSP achieves near-ML frame error rate performance . Figure 9 compares the Frame and Bit Error Rates for a $`4`$state, rate $`1/2`$ convolutional code with generator polynomials given by $`(5,7)`$ and code length $`200`$, over a $`5`$tap ISI channel. The channel impulse response was chosen as $`(0.848,0.424,0.2545,0.1696,0.0848)`$. The Fano decoder with $`b=1`$ and stepsize $`1`$ is seen to achieve essentially the same performance as the PSP algorithm for this code, with reasonable search complexity over the entire SNR range. We again note that the loss in lattice decoding as opposed to finite search space is negligible, due to MMSE-DFE preprocessing of the channel prior to the search. Moreover, the complexity of PSP algorithm, although linear in frame length, increases exponentially with the constraint length of the convolutional code used, while that of the Fano decoder is essentially independent of the constraint length. Figure 9 also shows the performance of the Fano decoder for a rate $`1/2`$, $`1024`$-state convolutional code with generator polynomials $`(4672,7542)`$, with the same frame size. Due to the increased constraint length, the performance is significantly better (with almost no increase in complexity). The complexity of PSP algorithm, on the other hand, is significantly higher for this code.
## 6 Conclusions
A central goal of this paper was to introduce a unified framework for tree search decoding in wireless communication applications. Towards this end, we identified the roles of two different, but inter-related, components of the decoder, namely; 1) Preprocessing and 2) Tree Search. We presented a preprocessing stage composed of MMSE-DFE filtering for left preprocessing and lattice reduction with column ordering for right preprocessing. We argued that this preprocessor allows for ignoring the boundary control in the tree search stage while entailing only a marginal loss in performance. By relaxing the boundary control, we were able to build a generic framework for designing tree search strategies for joint detection and decoding. Within this framework, BeFS emerged as a very efficient solution that offers many valuable advantages. To limit the storage requirement of BeFS, we re-discovered the Fano decoder as our proposed tree search algorithm. Finally, we established the superior performance-complexity tradeoff of the Fano decoder analytically in a V-BLAST configuration and demonstrated its excellent performance and complexity in more general scenarios via simulation results.
## Appendix A The Fano Decoder
In this section, we obtain the cost function used in the proposed Fano/Stack decoder from the Fano metric defined for tree codes over general point-to-point channels, and give a brief description of the Fano decoder and its properties.
### A.1 Generic Cost Function of the Fano Decoder
For the transmitted sequence $`\widehat{𝐱}`$, let
$$𝐲=𝐑\widehat{𝐱}+𝐰$$
(22)
be the system model, as in Section 2. In (22), the noise sequence $`𝐰`$ is composed of i.i.d Gaussian noise components with zero mean and unit variance.
For a general point-to-point channel with continuous output, the Fano metric of the node $`𝐱_1^k`$ can be written as
$$\mu (𝐱_1^k)=\mathrm{log}\left(\frac{Pr((𝐱_1^k))p(𝐲_1^k|(𝐱_1^k))}{p(𝐲_1^k)}\right)$$
(23)
where $`(𝐱_1^k)`$ is the hypothesis that $`𝐱_1^k`$ form the first $`k`$ symbols of the transmitted sequence.
For $`1km`$, if $`Pr((𝐱_1^k))`$ is uniform over all nodes $`𝐱_1^k`$ that consist of the first $`k`$ components of any valid codeword in $`𝒞`$, from (23), the cost function for the Fano decoder for our system model (22) can be simplified as
$$f(𝐱_1^k)=\mu (𝐱_1^k)=\mathrm{log}\left(\underset{𝐱_1^k}{}e^{\frac{_{j=1}^kw_j(𝐱_1^j)}{2}}\right)+\frac{_{j=1}^kw_j(𝐱_1^j)}{2}.$$
(24)
Since summation over $`𝐱_1^k`$ in (24) is not feasible, we use the following approximations: first, $`\mathrm{log}(a_i)\mathrm{log}(\mathrm{max}(a_i))`$, so the sum can be approximated by the largest term. Second, for moderate to high SNRs, the transmitted sequence is actually the closest vector with a high probability, i.e., the largest term corresponds to the transmitted sequence. Thus, (24) can be approximated as
$$\mathrm{log}\left(\underset{𝐱_1^k}{}e^{\frac{_{j=1}^kw_j(𝐱_1^j)}{2}}\right)\frac{|𝐰_1^k|^2}{2}.$$
(25)
After averaging (25) over noise samples and scaling, we have,
$$f(𝐱_1^k)=\underset{j=1}{\overset{k}{}}w_j(𝐱_1^j)k$$
In general, the cost function for the Fano/Stack decoder can be written in terms of the parameter $`b`$, the bias, as
$$f(𝐱_1^k)=\underset{j=1}{\overset{k}{}}w_j(𝐱_1^j)bk.$$
### A.2 The Algorithm
The operation of the Fano decoder with no boundary control (lattice decoding) follows the following steps:
* Step 1: (Initialize) Set $`k0`$, $`T0`$, $`𝐱x_0`$.
* Step 2: (Look forward) $`𝐱_1^{k+1}(𝐱_1^k,x_{k+1})`$, where $`x_{k+1}`$ is the $`(k+1)^{th}`$ component of the best child node of $`𝐱_1^k`$.
* Step 3:
If $`f(𝐱_1^{k+1})T`$,
If $`k+1=m`$ (leaf node), then $`\widehat{𝐱}=𝐱_1^m`$; exit.
Else (move forward), $`kk+1`$.
If $`f(𝐱_1^{k1})>T\mathrm{\Delta }`$,
while $`f(𝐱_1^k)T\mathrm{\Delta }`$, $`TT\mathrm{\Delta }`$ (tighten threshold).
Go to step 2.
Else
If $`(k=0`$ or $`f(𝐱_1^{k1})>T)`$, $`TT+\mathrm{\Delta }`$ (cannot move back, so relax threshold).
Go to step 2.
Else (move back and look forward to the next best node)
$`𝐱_1^k\{𝐱_1^{k1},x_k\}`$, where $`x_k`$ is the last component of the next best child node of $`𝐱_1^{k1}`$.
$`kk1`$.
Go to Step 3. $`\mathrm{}`$
Note that $`T`$ (i.e., the threshold) is allowed to take values only in multiples of the step size $`\mathrm{\Delta }`$ (i.e., $`0,\pm \mathrm{\Delta },\pm 2\mathrm{\Delta },\mathrm{}`$). When a node is visited by the Fano decoder for the first time, the threshold $`T`$ is tightened to the least possible value while maintaining the validity of the node. If the current node does not have a valid child node, then the decoder moves back to the parent node (if the parent node is valid) and attempts moving forward to the next best node. However, if the parent node is not valid, the threshold is relaxed and attempt is made to move forward again, proceeding in this way until a leaf node is reached.
The determination of best and next best child nodes is simplified in CLPS problem; the child node generation order gen in SE enumeration (section 4.1.2) generates child nodes with cost functions in ascending order, given any node $`𝐱_1^k`$.
### A.3 Properties of the Fano Decoder
The main properties of the Fano decoder used in our analysis are :
1. A node $`𝐱_1^k`$ is generated by the Fano decoder only if its cost function is not greater than the bound $`T`$.
2. Let correct path be defined as the path corresponding to the transmitted codeword, and let $`f_M`$ be the maximum cost function along the correct path. The bound $`T`$ is always less than $`(f_M+\mathrm{\Delta })`$, where $`\mathrm{\Delta }`$ is the step-size of the Fano decoder; that is, $`\mathrm{max}\{T\}<T_M\stackrel{}{=}f_M+\mathrm{\Delta }`$.
All nodes that are generated by the Fano decoder are necessarily those with cost function less than the bound $`T`$, by Property $`(1)`$. However, even though the cost function of some node $`𝐱_1^k`$ may be smaller than the bound, the node itself might not be visited when bound takes the value $`T`$. If any of the cost functions along the path $`\{𝐱_1^r,r<k\}`$ increases above $`T`$, the node $`𝐱_1^r`$ is not generated and thus $`𝐱_1^k`$ is not visited. Hence, this is not a sufficient condition for a node to be generated.
Moreover, in Property $`(2)`$, the bound $`T`$ is always lesser than $`(f_M^{}+\mathrm{\Delta })`$, where $`f_M^{}`$ is the maximum cost function along any path of length $`m`$. A tight bound is obtained only when the maximum cost function corresponding to the path with the least $`f_M^{}`$ is chosen. However, $`f_M`$ along the transmitted path is usually easier to characterize statistically than $`f_M^{}`$.
## Appendix B Properties of the Stack Decoder
For any node $`𝐱_1^k`$ in the tree, let $`h(𝐱_1^k)\stackrel{}{=}_{i=1}^kw_i(𝐱_1^k)bk`$. For the stack algorithm, the cost function of any node in ACTIVE at any instant defined as follows: If $`𝐱_1^k`$ is a leaf node, then $`f(𝐱_1^k)=\mathrm{}`$. Otherwise, we let $`𝐱_{1,g}^{k+1}`$ be the best child node of $`𝐱_1^k`$ not generated yet, and define $`f(𝐱_1^k)=h(𝐱_{1,g}^k)`$. We note that $`h`$ of any node, once generated, remains constant throughout the algorithm, and $`f`$ of any node is non-decreasing as the algorithm progresses.
###### Proposition 1
Let $`\overline{𝐱}_1^m=(\overline{x}_1,\mathrm{}\overline{x}_m)`$ be the path chosen by the stack algorithm, and $`𝐱_1^m=(x_1,\mathrm{},x_m)`$ be any path in the tree. Then,
$$\underset{1jm}{\mathrm{max}}h(\overline{𝐱}_1^j)\underset{1jm}{\mathrm{max}}h(𝐱_1^j)$$
(26)
Proof : On the contrary, assume there exists a path $`(\overline{x}_1,\overline{x}_2,\mathrm{},\overline{x}_d,\stackrel{˘}{x}_{d+1},\mathrm{},\stackrel{˘}{x}_m)`$ that does not satisfy (26). Here, the path is assumed to share the same nodes with the chosen path until level $`d`$, and diverges from the chosen path from level $`d+1`$ onwards. Since this path does not satisfy (26),
$$\underset{d+1jm}{\mathrm{max}}h(\overline{𝐱}_1^j)>\underset{d+1jm}{\mathrm{max}}h(\stackrel{˘}{𝐱}_1^j)$$
(27)
Let $`\overline{𝐱}_1^k`$, $`k>d`$, be the node for which $`\mathrm{max}_{d+1jm}h(\overline{𝐱}_1^j)`$ occurs. Then, we have,
$$h(\overline{𝐱}_1^k)>h(\stackrel{˘}{𝐱}_1^j),d<jm$$
(28)
Since $`\overline{𝐱}_1^m`$ is the chosen path, the node $`\overline{𝐱}_1^k`$ is generated at some instant before the search terminates. Just before $`\overline{𝐱}_1^k`$ is generated, $`h(\overline{𝐱}_1^{k1})=g(\overline{𝐱}_1^k)`$, since $`\overline{𝐱}_1^k`$ is the best child node of $`\overline{𝐱}_1^{k1}`$ not generated yet. Moreover, since $`h(\overline{𝐱}_1^k)>h(\stackrel{˘}{𝐱}_1^{d+1})`$, the node $`\overline{𝐱}_1^d`$ with cost function $`f(\overline{𝐱}_1^d)=h(\stackrel{˘}{𝐱}_1^{d+1})`$ appears at the top of the stack at some instant before $`\overline{𝐱}_1^k`$ is generated. Therefore, $`\stackrel{˘}{𝐱}_1^{d+1}`$ is generated before $`\overline{𝐱}_1^k`$ is generated. Since the search does not terminate before $`\overline{𝐱}_1^k`$ is generated, applying the same argument, one sees that all the nodes $`\stackrel{˘}{𝐱}_1^{d+2},\mathrm{},\stackrel{˘}{𝐱}_1^m`$ are generated before $`\overline{𝐱}_1^k`$ is generated. However, once $`\stackrel{˘}{𝐱}_1^m`$ is generated by the stack algorithm, the search terminates, with $`(\overline{x}_1,\mathrm{},\overline{x}_d,\stackrel{˘}{x}_{d+1},\mathrm{},\stackrel{˘}{x}_m)`$ as the chosen path. Since $`1dm`$ can take any value, the inequality in (26) is satisfied by all paths.
###### Proposition 2
If
$$\underset{1jd}{\mathrm{max}}h(𝐱_1^d)>\underset{1jm}{\mathrm{max}}h(\overline{𝐱}_1^j),$$
(29)
then, the node $`𝐱_1^d`$ is not generated.
Proof : First, we show that if
$$h(𝐱_1^d)>\underset{1jm}{\mathrm{max}}h(\overline{𝐱}_1^j),$$
(30)
then $`𝐱_1^d`$ is not generated. Let (30) be true, and assume $`𝐱_1^d`$ is generated. Then, just before $`𝐱_1^d`$ is generated, its parent node $`𝐱_1^{d1}`$ is at the top of ACTIVE, with cost function $`f(𝐱_1^{d1})=h(𝐱_1^d)`$. However, since $`h(𝐱_1^d)>h(\overline{𝐱}_1^j),1jm`$, all nodes along the chosen path are generated before $`𝐱_1^d`$ is generated, and the hence the search terminates before $`𝐱_1^d`$ is generated. Noting that $`𝐱_1^d`$ can be generated only if all the nodes $`𝐱_1^1,\mathrm{},𝐱_1^{d1}`$ are generated, and applying the same argument for $`𝐱_1^{d1},\mathrm{},𝐱_1^1`$, we have (29).
## Appendix C Proof of Theorem 2
Let $`𝒜_{IR}`$ be the set of nodes generated by the IR algorithm, where the bounding function $`𝐭`$ has components given by $`t_k=bk+\delta `$. Let $`𝒜_s`$ be the set of nodes generated by the stack decoder with the bias $`b`$. The IR algorithm in Theorem 2 can be defined with bounding function given by $`\{t_k=bk+\delta ,1km\}`$, and the cost function for any node $`𝐱_1^k`$ given by $`_{i=1}^kw_i(𝐱_1^i)`$, or equivalently, with the bounding function $`t_k=\delta `$ and cost function $`\left(_{i=1}^kw_i(𝐱_1^i)bk\right)`$. If $`\delta `$ is the bound of the IR algorithm, then any node $`𝐱_1^k`$ is generated by the algorithm if and only if all the conditions $`\left\{_{i=1}^1w_i(𝐱_1^i)b<\delta ,_{i=1}^2w_i(𝐱_1^i)2b<\delta ,\mathrm{},_{i=1}^kw_i(𝐱_1^i)bk<\delta \right\}`$, are satisfied. Therefore,
$$𝒜_{IR}=\{𝐱_1^k:\underset{1jk}{\mathrm{max}}\left(\underset{i=1}{\overset{k}{}}w_i(𝐱_1^i)bk\right)<\delta \}.$$
(31)
Moreover, $`\delta `$ should be such that at least one sequence $`𝐱𝒰`$ is included within the search space.<sup>7</sup><sup>7</sup>7Otherwise, $`\delta `$ is increased and search is repeated afresh Let $`\widehat{𝐱}_{IR}`$ be a leaf node such that
$$\widehat{𝐱}_{IR}=\text{arg}\underset{𝐱𝒰}{\mathrm{min}}\left(\underset{1km}{\mathrm{max}}\left(w_i(𝐱_1^i)bk\right)\right).$$
(32)
i.e., $`\widehat{𝐱}_{IR}`$ has the least value of maximum cost function among all paths of length $`m`$. If
$$\delta <\underset{1jm}{\mathrm{max}}\left(w_i(\widehat{𝐱}_{IR,1}^j)bk\right),$$
then no $`𝐱𝒰`$ lies within the search space, and the search space is empty. If lattice decoding is used, then the minimum in (32) is taken over all $`𝐱^m`$. Therefore, $`\delta >\mathrm{max}_{1jm}\left(w_i(\widehat{𝐱}_{IR,1}^j)bk\right)`$. From Section B, Prop. 1, the path chosen by the stack algorithm, $`𝐱_1^m`$ satisfies
$$\left(\underset{1km}{\mathrm{max}}\left(w_i(\overline{𝐱}_1^i)bk\right)\right)\left(\underset{1km}{\mathrm{max}}\left(w_i(𝐱_1^i)bk\right)\right)$$
(33)
where $`𝐱_1^m`$ is any other path.
From (32) and (33),
$$\left(\underset{1km}{\mathrm{max}}\left(w_i(\overline{𝐱}_1^i)bk\right)\right)=\left(\underset{1km}{\mathrm{max}}\left(w_i(\widehat{𝐱}_{IR,1}^i)bk\right)\right)<\delta $$
(34)
From Proposition 2 and (33), $`𝒜_s𝒜_{IR}`$.
## Appendix D Proof of Theorem 3
In this section, we derive an upper bound to the frame error rate for a V-BLAST system with uncoded input (with $`Q`$-PAM constellation for the components), for the Fano decoder that visits paths in the regular $`Q`$-PAM signal space. The preprocessing assumed here is QR transformation of $`𝐇`$.
Let $`_f`$ be the event that the Fano decoder makes an erroneous detection, conditioned on $`T_M\stackrel{}{=}f_M+\mathrm{\Delta }`$. Then, $`P_e=E_{T_M}(Pr(_f))`$ is the frame error rate of the Fano decoder. In this section, we derive an upper bound on $`P_e`$. From property $`(2)`$ in Section A.3, $`T<(f_M+\mathrm{\Delta })`$, where $`\mathrm{\Delta }`$ is the step size of the Fano decoder. Any sequence $`𝐱\widehat{𝐱}`$ can be decoded as the closest point by the Fano decoder only if its cost function is lesser than $`T_M`$. One has
$`𝐲`$ $`=`$ $`\mathrm{𝐇𝐱}+𝐳=𝐐\left(\begin{array}{c}𝐑\\ \mathrm{𝟎}\end{array}\right)𝐱+𝐳,`$ (35)
and therefore
$`𝐲`$ $``$ $`𝐐^^T𝐲=\left(\begin{array}{c}𝐑\\ \mathrm{𝟎}\end{array}\right)𝐱+\left(\begin{array}{c}𝐰_{r+1}^n\\ 𝐰_1^r\end{array}\right)`$ (36)
where $`r=nm`$ is the excess degrees of freedom in the V-BLAST system. Since the cost function of a leaf node $`𝐱_1^m`$ is $`f(𝐱_1^m)={\displaystyle \underset{i=1}{\overset{m}{}}}w_i(𝐱_1^i)bm=|𝐑\stackrel{~}{𝐱}+𝐰_{r+1}^n|^2bm`$, $`P(_f)`$ can be upper bounded as
$`P(_f)`$ $``$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}({\displaystyle \underset{j=1}{\overset{m}{}}}w_j(𝐱_1^j)bm<T_M)`$ (37)
$`=`$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}(|𝐑\stackrel{~}{𝐱}+𝐰_{r+1}^n|^2<bm+f_M+\mathrm{\Delta })`$ (38)
where
$$f_M=\mathrm{max}\{0,|𝐰_{r+1}^{r+1}|^2b,|𝐰_{r+1}^{r+2}|^22b,\mathrm{},|𝐰_{r+1}^n|^2mb\}$$
is the maximum cost function along the transmitted sequence path. The upper bound in (37) follows from the union bound, and due to the fact that in general, $`f(𝐱_1^m)<T_M`$ is only a necessary condition for $`𝐱_1^m`$ to be decoded by the Fano decoder.
The bound in (38) can be rewritten as
$`P(_f)`$ $``$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}\left(\left|\left(\begin{array}{c}𝐑\\ \mathrm{𝟎}\end{array}\right)\stackrel{~}{𝐱}+𝐰_1^n\right|^2<bm+f_M+\mathrm{\Delta }+|𝐰_1^r|^2\right)`$ (39)
$`=`$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}\left(\left|\left(\begin{array}{c}𝐑\\ \mathrm{𝟎}\end{array}\right)\stackrel{~}{𝐱}+𝐰_1^n\right|^2|𝐰_1^n|^2<bm+f_M+\mathrm{\Delta }|𝐰_{r+1}^n|^2\right)`$ (40)
$`=`$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}(|𝐇\stackrel{~}{𝐱}|^2+2(𝐇\stackrel{~}{𝐱})^^T𝐳<bm+f_M+\mathrm{\Delta }|𝐰_{r+1}^n|^2)`$ (41)
$``$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}(|𝐇\stackrel{~}{𝐱}|^2+2(𝐇\stackrel{~}{𝐱})^^T𝐳<bm+\mathrm{\Delta }),`$ (42)
since $`f_M|𝐰_{r+1}^n|^2=\mathrm{max}\{|𝐰_{r+1}^n|^2,|𝐰_{r+2}^n|^2b,\mathrm{},mb\}0`$. The bound in (42) is now independent of the value of $`f_M`$, and hence represents a bound on the frame error rate. Note that the corresponding expression in (42) for ML decoding is $`\mathrm{Pr}(|𝐇\stackrel{~}{𝐱}|^2+2(𝐇\stackrel{~}{𝐱})^^T𝐳<0)`$. For any $`𝐱𝒰`$ and $`𝐱\widehat{𝐱}`$, let $`d^2(\widehat{𝐱},𝐱)=|𝐇\stackrel{~}{𝐱}|^2`$ represent the squared Euclidean distance between the lattice points $`\mathrm{𝐇𝐱}`$ and $`𝐇\widehat{𝐱}`$. Then,
$$\mathrm{Pr}(|𝐇\stackrel{~}{𝐱}+𝐳|^2|𝐳|^2mb+\mathrm{\Delta })\{\begin{array}{cc}e^{\left(\frac{1}{8}(d^2(\widehat{𝐱},𝐱)mb\mathrm{\Delta })^2/d^2(\widehat{𝐱},𝐱)\right)},\hfill & d^2(\widehat{𝐱},𝐱)>mb+\mathrm{\Delta }\hfill \\ 1\hfill & d^2(\widehat{𝐱},𝐱)mb+\mathrm{\Delta }\hfill \end{array}$$
(43)
by Chernoff bound.
For $`d^2(\widehat{𝐱},𝐱)>mb+\mathrm{\Delta }`$, equation (43) can be rewritten as
$`\mathrm{Pr}(|𝐇\stackrel{~}{𝐱}+𝐳|^2|𝐳|^2mb+\mathrm{\Delta })`$ $``$ $`e^{\frac{1}{8}\left(d^2(\widehat{𝐱},𝐱)+\frac{(mb+\mathrm{\Delta })^2}{d^2(\widehat{𝐱},𝐱)}2(mb+\mathrm{\Delta })\right)}`$ (44)
$``$ $`e^{\frac{1}{8}\left(d^2(\widehat{𝐱},𝐱)\right)}e^{(mb+\mathrm{\Delta })/4}`$ (45)
since $`e^{(mb+\mathrm{\Delta })^2/(8d^2(\widehat{𝐱},𝐱))}<1`$, for $`d^2(\widehat{𝐱},𝐱)>0`$. Let
$`q`$ $`\stackrel{}{=}`$ $`\underset{𝐱_i,𝐱_j𝒰,ij}{\mathrm{min}}(|𝐇(𝐱_i𝐱_j)|^2)`$ (46)
and let
$`g(q)`$ $`\stackrel{}{=}`$ $`\{\begin{array}{cc}_{𝐱𝒰,\stackrel{~}{𝐱}\mathrm{𝟎}}e^{(mb+\mathrm{\Delta })/4}e^{\left(q/8\right)}\hfill & q>mb+\mathrm{\Delta }\hfill \\ 1\hfill & q<mb+\mathrm{\Delta }\hfill \end{array}`$ (47)
Then, from (45) and (43), $`P_eE_q(g(q))`$. An upper bound on the probability density function (pdf) of $`q`$ is given by
$`p(q)<p_\chi (q){\displaystyle \underset{k=1}{\overset{m}{}}}\left(\begin{array}{c}m\\ k\end{array}\right){\displaystyle \frac{1}{k}}`$ (50)
where $`p_\chi (q)`$ is the pdf of a scaled chi-square random variable with $`n`$ degrees of freedom and mean $`\frac{n\rho }{m}`$ (i.e., a random variable that is the sum of squares of $`n`$ i.i.d zero-mean Gaussian variables with variance $`\frac{\rho }{m}`$). Then, (47) and (50) give
$`P_e`$ $``$ $`Q^me^{(mb+\mathrm{\Delta })/4}{\displaystyle _{mb+\mathrm{\Delta }}^{\mathrm{}}}e^{\left(q/8\right)}p(q)𝑑q+{\displaystyle _0^{mb+\mathrm{\Delta }}}p(q)𝑑q`$ (51)
$`=`$ $`AQ^me^{(mb+\mathrm{\Delta })/4}{\displaystyle _{mb+\mathrm{\Delta }}^{\mathrm{}}}e^{\left(q/8\right)}{\displaystyle \frac{q^{(n/21)}e^{q/(2\sigma ^2)}}{2^{n/2}\mathrm{\Gamma }(\frac{n}{2})\sigma ^n}}𝑑q+A\gamma ({\displaystyle \frac{mb+\mathrm{\Delta }}{2\sigma ^2}},{\displaystyle \frac{n}{2}})`$ (52)
$``$ $`AQ^me^{(mb+\mathrm{\Delta })/4}{\displaystyle _0^{\mathrm{}}}e^{\left(q/8\right)}{\displaystyle \frac{q^{(n/21)}e^{q/(2\sigma ^2)}}{2^{n/2}\mathrm{\Gamma }(\frac{n}{2})\sigma ^n}}𝑑q+A\gamma ({\displaystyle \frac{mb+\mathrm{\Delta }}{2\sigma ^2}},{\displaystyle \frac{n}{2}})`$ (53)
$`=`$ $`{\displaystyle \frac{AQ^me^{(mb+\mathrm{\Delta })/4}}{\left(1+\frac{\rho }{4m}\right)^{n/2}}}+A\gamma ({\displaystyle \frac{mb+\mathrm{\Delta }}{2\sigma ^2}},{\displaystyle \frac{n}{2}})`$ (54)
where $`\sigma ^2=\frac{\rho }{m}`$, $`A\stackrel{}{=}{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \frac{1}{k}}\left({\displaystyle \genfrac{}{}{0pt}{}{m}{k}}\right)`$ is a constant independent of $`q`$ or $`\rho `$, and $`\gamma (x,a)`$ is the incomplete gamma function. If $`b`$ is bounded (i.e., $`b<<\mathrm{}`$) $`\rho `$, then $`e^{(mb/4)}`$ is also bounded for all $`\rho `$ and finite $`m`$. The error performance of the Fano decoder can now be characterized by the sum of two terms. The dependence of the first error term on $`\rho `$ is of the form $`\rho ^{(n/2)}`$ for large values of SNR, and hence has the same diversity as the ML decoder. The second term can also be bounded as
$`\gamma ({\displaystyle \frac{mb+\mathrm{\Delta }}{2\sigma ^2}},{\displaystyle \frac{n}{2}})`$ $``$ $`\left(1e^{(mb+\mathrm{\Delta })/(2\sigma ^2)}\right)^{(n/2)}`$ (55)
$``$ $`\left({\displaystyle \frac{mb+\mathrm{\Delta }}{2\sigma ^2}}\right)^{(n/2)}`$ (56)
$`=`$ $`\left({\displaystyle \frac{m(mb+\mathrm{\Delta })}{2\rho }}\right)^{(n/2)}`$ (57)
where (55) follows from the inequality $`\gamma (x,a)(1e^x)^a`$ (Appendix E.2), and (56) from $`(1e^x)<x`$ for $`x>0`$. The second term also has the dependence $`\rho ^{(n/2)}`$, and hence the Fano decoder achieves the same diversity as that of the ML decoder for this system.
The above derivation also applies to the Stack algorithm, with minor modifications. Let $`_s`$ be the event that the stack algorithm makes an erroneous detection, conditioned on the value of $`f_M`$. Then, $`P_e=E_{f_M}(\mathrm{Pr}(_s))`$ is the word error rate of the stack algorithm. Since any path $`𝐱\widehat{𝐱}`$ is decoded as the closest point by the stack algorithm only if $`h(𝐱)=_{i=1}^mw_i(𝐱_1^i)bm`$ is not greater than $`f_M`$ (Prop. 1, Section B , $`P(_s)`$ can be written as
$`P(_s)`$ $``$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}\left\{{\displaystyle \underset{j=1}{\overset{m}{}}}w_j(𝐱_1^j)bm<f_M\right\}`$ (58)
$`=`$ $`{\displaystyle \underset{𝐱𝒰,𝐱\widehat{𝐱}}{}}\mathrm{Pr}\left\{|𝐇\stackrel{~}{𝐱}+𝐳|^2<bm+f_M\right\}`$ (59)
From (59) and (38), it is easy to see that the error probability expression for the stack algorithm is the same as that for the Fano decoder, when $`\mathrm{\Delta }=0`$. Thus, the stack algorithm too achieves the same diversity as the ML decoder for a V-BLAST system, for any finite value of $`b`$. $`\mathrm{}`$
## Appendix E Proof of Theorem 4
The following are required for the proof.
### E.1 Wald’s inequality
Let $`S_0=0,S_1,S_2,\mathrm{}`$ be a random walk, with $`S_j=_{i=1}^jX_i`$, where $`X_i`$s are i.i.d random variables such that $`\mathrm{Pr}(X_i>0)>0`$, $`\mathrm{Pr}(X_i<0)>0`$, and $`E(X_i)<0`$. Let $`g(\lambda )=E(e^{\lambda X_i})`$ be the moment generating function of $`X_i`$. Let $`\lambda _0>0`$ be a root of $`g(\lambda )=1`$. Then, from Wald’s identity ,
$$\mathrm{Pr}(S_{\mathrm{max}}>u)e^{\lambda _0u}$$
(60)
where $`S_{\mathrm{max}}=\mathrm{max}_j(S_j)`$.
For the random walk with $`X_i=w_i^2b`$, where $`w_i𝒩(0,1)`$, the above conditions are satisfied if $`b>1`$. The moment generating function for $`X_i=w_i^2b`$ is given by
$$g(\lambda )=\frac{e^{\lambda b}}{\sqrt{12\lambda }}.$$
(61)
From (61), $`\lambda _0>0`$ can be found as the positive root of the equation
$$2\lambda b=\mathrm{log}(12\lambda ).$$
Notice that since $`\mathrm{log}(12\lambda )`$ decreases from $`0`$ to $`\mathrm{}`$ as $`\lambda `$ increases from $`0`$ to $`\frac{1}{2}`$, $`\lambda _0`$ satisfies $`\lambda _0(0,0.5)`$. Since $`\mathrm{max}_{0jm}S_j\mathrm{max}_{j0}S_j`$, the bound in (60) is also valid for any stopped random walk.
### E.2 Upper bounds on $`\gamma (\beta ,k)`$
For a scaled chi-square random variable $`X`$ with $`k`$ degrees of freedom and mean $`k\sigma ^2`$,
$$\mathrm{Pr}(X\beta )=\gamma (\frac{\beta }{2\sigma ^2},\frac{k}{2})$$
where $`\gamma `$ is known as the incomplete gamma function. From Chernoff bound, we have
$$\gamma (\frac{\beta }{2\sigma ^2},\frac{k}{2})=\mathrm{Pr}(X\beta )\{\begin{array}{cc}\left(\frac{\beta }{\sigma ^2k}\right)^{k/2}e^{(\frac{k}{2}\frac{\beta }{2\sigma ^2})}\hfill & \frac{\beta }{2\sigma ^2}<\frac{k}{2}\hfill \\ 1\hfill & \frac{\beta }{2\sigma ^2}\frac{k}{2}\hfill \end{array}$$
(62)
A simpler, though looser, upper bound is given in :
$$\gamma (x,a)(1e^x)^a$$
(63)
### E.3 Proof
Let $`𝐱_1^k`$ be any path in the tree, and $`h(𝐱_1^k)={\displaystyle \underset{i=1}{\overset{k}{}}}w_i(𝐱_1^i)bk`$, as in Section B. Let $`f_M=\mathrm{max}_{1im}h(\widehat{𝐱}_1^i)`$ be the maximum cost function along the transmitted path. From Section B, Prop. 1, $`f_M`$ is not lesser than the maximum of the cost functions along the path chosen by the stack decoder. From Prop. 2, it is easy to see that any node $`𝐱_1^k`$ is generated, only if the maximum of the cost functions along the path $`𝐱_1^k`$ does not increase above $`f_M`$.
Let $`𝒜_{s,b}`$ be the set of generated nodes. In the proof, we upper bound the number of all the paths visited by the algorithm that are different from the correct path, and then we add the complexity of finding the correct path (i.e., $`m`$). Then, $`𝒜_{s,b}`$ is a subset of the set of nodes that satisfy $`f(𝐱_1^k)<f_M`$. Let $`𝐑_{k,k}`$ be the lower $`k\times k`$ part of the $`𝐑`$ matrix, i.e.,
$$𝐑_{k,k}=\left(\begin{array}{ccc}r_{k,k}& \mathrm{}& r_{k,1}\\ & \mathrm{}& \mathrm{}\\ \mathrm{𝟎}& & r_{1,1}\end{array}\right)$$
Then, we have,
$`P(𝐱_1^k𝒜_{s,b})`$ $``$ $`P(|𝐑_{k,k}\stackrel{~}{𝐱}_1^k+𝐰_{r+1}^{r+k}|^2bk<f_M)`$ (64)
$`=`$ $`P\left(\left|\left(\begin{array}{c}𝐑_{k,k}\\ \mathrm{𝟎}\end{array}\right)\stackrel{~}{𝐱}_1^k+𝐰_1^{r+k}\right|^2bk<f_M+|𝐰_1^r|^2\right)`$ (65)
$`=`$ $`P\left(\left|\left(\begin{array}{c}𝐑_{k,k}\\ \mathrm{𝟎}\end{array}\right)\stackrel{~}{𝐱}_1^k\right|^2+|𝐰_1^{r+k}|^2+2(𝐰_1^{r+k})^^T\left(\begin{array}{c}𝐑_{k,k}\\ \mathrm{𝟎}\end{array}\right)\stackrel{~}{𝐱}_1^kbk<f_M+|𝐰_1^r|^2\right)`$ (66)
where $`r=nm`$ is the excess degrees of freedom in the V-BLAST system. From , for each $`km`$, one can find an $`(r+k)\times k`$ matrix $`\overline{𝐇}_{r+k,k}`$ that has the same distribution as the lower $`(r+k)\times k`$ part of $`𝐇`$, and an $`(r+k)\times (r+k)`$ unitary matrix $`𝚯^{(r+k)}`$ whose distribution is independent of $`𝐑_{k,k}`$, such that $`\overline{𝐇}_{r+k,k}=𝚯^{(r+k)}\left(\begin{array}{c}𝐑_{k,k}\\ \mathrm{𝟎}\end{array}\right)`$.
Let $`\overline{𝐳}_1^{r+k}=𝚯^{(r+k)}𝐰_1^{r+k}`$. The bound in (66) can now be rewritten as
$`P(𝐱_1^k𝒜_{s,b})`$ $``$ $`P(|\overline{𝐇}_{r+k,k}\stackrel{~}{𝐱}_1^k|^2+2(\overline{𝐳}_1^{r+k})^^T\overline{𝐇}_{r+k,k}\stackrel{~}{𝐱}_1^k<f_M+bk|𝐰_{r+1}^{r+k}|^2)`$ (67)
In (67), $`f_M+bk|𝐰_{r+1}^{r+k}|^2`$ can be bounded as
$`f_M+bk|𝐰_{r+1}^{r+k}|^2`$ $`=`$ $`\mathrm{max}\{bk|𝐰_{r+1}^{r+k}|^2,b(k1)|𝐰_{r+2}^{r+k}|^2,\mathrm{},f_M^{}\}`$ (68)
$``$ $`\mathrm{max}\{bk,f_M^{}\}`$ (69)
where $`f_M^{}=\mathrm{max}\{0,|𝐰_{r+k+1}^{r+k+1}|^2b,\mathrm{},|𝐰_{r+k+1}^n|^2b(mk)\}`$. Let $`\beta =\mathrm{max}\{bk,f_M^{}\}`$. (67) can now be rewritten as
$$P(𝐱_1^k𝒜_{s,b})P(|\overline{𝐇}_{r+k,k}\stackrel{~}{𝐱}_1^k|^2+2(\overline{𝐳}_1^{r+k})^^T\overline{𝐇}_{r+k,k}\stackrel{~}{𝐱}_1^k<\beta )$$
(70)
Using Chernoff bound, (70) can be written as
$$P(𝐱_1^k_{s,b}|q_k,\beta )\{\begin{array}{cc}e^{(q_k\beta )^2/(8q_k)},\hfill & q_k>\beta \hfill \\ 1,\hfill & q_k<\beta \hfill \end{array}$$
(71)
where $`q_k=|\overline{𝐇}_{r+k,k}\stackrel{~}{𝐱}_1^k|^2`$. Let $`\eta =\frac{\rho }{m}`$. Then, in (71), $`\frac{1}{\eta |\stackrel{~}{𝐱}_1^k|^2}q_k`$ is a chi-square random variable with $`(r+k)`$ degrees of freedom. Since the three random variables, $`\overline{𝐇}_{r+k,k}\stackrel{~}{𝐱}_1^k`$, $`𝐳_1^k`$ and $`\beta `$ are independent, averaging over $`q_k`$ and $`\beta `$ gives
$`P(𝐱_1^k𝒜_{s,b})`$ $``$ $`E_\beta \left({\displaystyle _0^\beta }f_{q_k}(q_k)𝑑q_k+{\displaystyle _\beta ^{\mathrm{}}}e^{(q_k\beta )^2/(8q_k)}f_{q_k}(q_k)𝑑q_k\right)`$ (72)
$``$ $`P(\beta =bk)\left({\displaystyle _0^{bk}}f_{q_k}(q_k)𝑑q_k+{\displaystyle _{bk}^{\mathrm{}}}e^{(q_kbk)^2/(8q_k)}f_{q_k}(q_k)𝑑q_k\right)+P(\beta >bk)`$ (73)
In (73), $`P(\beta >bk)=P(f_M^{}>bk)e^{bk\lambda _0}`$ for $`b>1`$ (see Section E.1). Note that one requires $`b>1`$ because the distribution of the maximum of the cost functions along the transmitted path will depend on $`m`$ otherwise. <sup>8</sup><sup>8</sup>8Later, we will require a stronger condition on $`b`$ to guarantee the convergence of the sums in (82). The bound in (73) amounts to counting all the nodes $`𝐱_1^k`$ in the search space when $`\beta >bk`$. Since $`P(\beta >bk)`$ decreases sufficiently fast as $`k`$ increases, this upper bound is still tight for our purposes. Now, (73) can be further simplified as,
$`P(𝐱_1^k𝒜_{s,b})`$ $``$ $`{\displaystyle _0^{bk}}f_{q_k}(q_k)𝑑q_k+{\displaystyle _{bk}^{\mathrm{}}}e^{(q_kbk)^2/(8q_k)}f_{q_k}(q_k)𝑑q_k+e^{bk\lambda _0}`$ (74)
$``$ $`{\displaystyle _0^{bk}}f_{q_k}(q_k)𝑑q_k+{\displaystyle _0^{\mathrm{}}}e^{(q_kbk)^2/(8q_k)}f_{q_k}(q_k)𝑑q_k+e^{bk\lambda _0}`$ (75)
$``$ $`\gamma ({\displaystyle \frac{bk}{2\eta |\stackrel{~}{𝐱}_1^k|^2}},{\displaystyle \frac{r+k}{2}})+e^{bk/4}{\displaystyle _0^{\mathrm{}}}e^{q_k/8}f_{q_k}(q_k)𝑑q_k+e^{bk\lambda _0}`$ (76)
$``$ $`\gamma ({\displaystyle \frac{bk}{2\eta }},{\displaystyle \frac{r+k}{2}})+{\displaystyle \frac{e^{bk/4}}{\left(1+\frac{\eta }{4}\right)^{(r+k)/2}}}+e^{bk\lambda _0}`$ (77)
with $`\gamma (,)`$ as the incomplete gamma function. Assuming $`r0`$, (77) can be bounded as
$`P(𝐱_1^k𝒜_{s,b})`$ $``$ $`\left({\displaystyle \frac{b}{\eta }}e^{1\frac{b}{\eta }}\right)^{\frac{k}{2}}+\left({\displaystyle \frac{e^{b/2}}{(1+\frac{\eta }{4})}}\right)^{\frac{k}{2}}+e^{bk\lambda _0}`$ (78)
for $`\eta >b`$. The inequality in (78) follows from an upper bound on the incomplete gamma function (see Section E.2). For a node $`𝐱_1^k`$, let $`G(𝐱_1^k)=1`$ if the node is generated and $`0`$ otherwise. Then, the expected number of nodes generated by the algorithm (i.e., complexity) is $`{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{𝐱_1^k}{}}E[G(𝐱_1^k)]`$, where the expectation is over all channel realizations. Let $`C_m`$ be the expected complexity per dimension. Then, assuming a bounded $`r`$, $`C_m`$ is written as <sup>9</sup><sup>9</sup>9The first term in the RHS of (79) comes from counting the complexity of the finding correct path, i.e., $`\stackrel{~}{𝐱}=\mathrm{𝟎}`$.
$`mC_mm+{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{𝐱_1^k}{}}\left(\left({\displaystyle \frac{b}{\eta }}e^{1\frac{b}{\eta }}\right)^{\frac{k}{2}}+\left({\displaystyle \frac{e^{b/2}}{(1+\frac{\eta }{4})}}\right)^{\frac{k}{2}}+e^{bk\lambda _0}\right).`$ (79)
The complexity per dimension, $`C_m`$, can now be upper bounded as
$`C_m`$ $``$ $`1+{\displaystyle \frac{1}{m}}{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{𝐱_1^k}{}}\left(\left({\displaystyle \frac{b}{\eta }}e^{1\frac{b}{\eta }}\right)^{\frac{k}{2}}+\left({\displaystyle \frac{e^{b/2}}{(1+\frac{\eta }{4})}}\right)^{\frac{k}{2}}+e^{bk\lambda _0}\right)`$ (80)
$``$ $`1+{\displaystyle \frac{1}{m}}{\displaystyle \underset{k=1}{\overset{m}{}}}Q^k\left(\left({\displaystyle \frac{b}{\eta }}e^{1\frac{b}{\eta }}\right)^{\frac{k}{2}}+\left({\displaystyle \frac{e^{b/2}}{(1+\frac{\eta }{4})}}\right)^{\frac{k}{2}}+e^{bk\lambda _0}\right)`$ (81)
$``$ $`1+{\displaystyle \frac{1}{m}}\left({\displaystyle \frac{1}{1Q^2\frac{b}{\eta }e^{1\frac{b}{\eta }}}}+{\displaystyle \frac{1}{1Q^2\frac{e^{b/2}}{(1+\frac{\eta }{4})}}}+{\displaystyle \frac{1}{1Qe^{b\lambda _0}}}\right),`$ (82)
when $`b`$ and $`\eta `$ are sufficiently large, so that all the three sums converge. The inequality in (81) is true, since the number of nodes at level $`k`$ is $`Q^k`$. Since the terms inside the parenthesis in (82) are all independent of $`m`$, the number of nodes visited by the stack algorithm scales at most linearly, when $`\eta >\eta _0`$, where $`\eta _0`$ is the minimum $`\frac{\rho }{m}`$ ratio required for convergence of the sums in (82). $`\mathrm{}`$
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# Untitled Document
Brownian motion and magnetism
Supurna Sinha
Centre for Theoretical Studies, Indian Institute of Science,
Bangalore, India 560 012
Joseph Samuel
Raman Research Institute, Bangalore, India 560 080
We present an interesting connection between Brownian motion and magnetism. We use this to determine the distribution of areas enclosed by the path of a particle diffusing on a sphere. In addition, we find a bound on the free energy of an arbitrary system of spinless bosons in a magnetic field. The work presented here is expected to shed light on polymer entanglement, depolarized light scattering, and magnetic behavior of spinless bosons.
Consider a particle diffusing on a sphere. If the diffusing particle returns to its starting point at time $`\beta `$ its path subtends a solid angle $`\mathrm{\Omega }`$ at the center of the sphere. We ask: what is the probability distribution of $`\mathrm{\Omega }`$ ? This problem comes up if one considers a spin-$`\frac{1}{2}`$ system in a random magnetic field. As is well known, the state (up to a phase) of a spin-$`\frac{1}{2}`$ system can be represented as a point on the Poincarè sphere. Under the influence of a random Hamiltonian, the state of the system diffuses on the Poincarè sphere. From the work of Berry<sup>1</sup> and others, it is known that the system picks up a geometric phase $`\gamma `$ equal to half the solid angle swept out on the Poincarè sphere. To compute the distribution of geometric phases one is led to the question posed above. A closely related problem has already been studied in the context of polymer entanglement:<sup>2</sup> given that a Brownian path on the plane is closed at time $`\beta `$, what is the probability that it encloses a given area $`A`$?
In this paper we present a general method of solving these problems by using a connection between Brownian motion and magnetism. The qualitative idea is to use a magnetic field as a “counter,” to measure the area enclosed in a Brownian motion. We derive a relation between the distribution of areas in a Brownian motion and the partition function of a magnetic system, which can be used to cast light on both subjects. Despite its apparent simplicity, this relation does not seem to have been noticed or exploited so far. Our main purpose here is to illustrate its usefulness. We first discuss the planar problem solved earlier. We then go on to solve the (as yet unsolved) problem of diffusion on the sphere. We also exploit the relation to learn about the magnetic properties of bosonic systems. Here we recover previously known results and arrive at some others. We conclude the paper with a few remarks.
Let a diffusing particle start from a point on a plane at time $`\tau =O`$. Given that the path is closed at time $`\beta `$ (not necessarily for the first time), what is the conditional probability that it encloses a given area $`A`$? By “area” we mean the algebraic area, including sign. The area enclosed to the left of the diffusing particle counts as positive and the area to the right as negative. This problem has been posed and solved<sup>2</sup> by polymer physicists, since it provides an idealized model for the entanglement of polymers. We present a method of solving this simple problem.
Let $`\{\stackrel{}{x}(\tau ),0\tau \beta ,\stackrel{}{x}(0)=\stackrel{}{x}(\beta )\}`$ be any realization of a closed Brownian path on the plane. As is well known, Brownian paths are distributed according to the Wiener measure:<sup>3</sup> if $`f[\stackrel{}{x}(\tau )]`$ is any functional on paths, the expectation value of $`f`$ is given by
$`f[\stackrel{}{x}(\tau )]_𝒲`$
$`{\displaystyle \frac{𝒟[\stackrel{}{x}(\tau )]f[\stackrel{}{x}(\tau )]\mathrm{exp}\left[_0^\beta \left\{\frac{1}{2}\frac{d\stackrel{}{x}}{d\tau }\frac{d\stackrel{}{x}}{d\tau }d\tau \right\}\right]}{𝒟[\stackrel{}{x}(\tau )]\mathrm{exp}\left[_0^\beta \left\{\frac{1}{2}\frac{d\stackrel{}{x}}{d\tau }\frac{d\stackrel{}{x}}{d\tau }d\tau \right\}\right]}}.`$ (1)
In Eq. (1) the functional integrals<sup>4</sup> are over all closed paths (the starting point is also integrated over). (We set the diffusion constant equal to half throughout this paper.) Let $`𝒜[\stackrel{}{x}(\tau )]`$ be the algebraic area enclosed by the path $`\stackrel{}{x}(\tau )`$. Clearly, the normalized probability distribution of areas P(A) is given by
$$P(A)\delta (𝒜[\stackrel{}{x}(\tau )]A)_𝒲.$$
(2)
The expectation value $`\stackrel{~}{\varphi }`$ of any function $`\varphi (A)`$ of the area is given by $`P(A)\varphi (A)d(A)`$. As is usual in probability theory we focus on the generating function $`\stackrel{~}{P}(B)`$ of the distribution $`P(A)`$:
$$\stackrel{~}{P}(B)\overline{e^{ieBA}}=P(A)e^{ieBA}𝑑A,$$
(3)
which is simply the Fourier transform of $`P(A)`$. For future convenience we write the Fourier transform variable as $`eB`$. The distribution $`P(A)`$ can be recovered from its generating function by an inverse Fourier transform. From Eqs. (2) and (3) above we find
$$\stackrel{~}{P}(B)=e^{ieB𝒜}_𝒲.$$
(4)
Notice that $`B𝒜`$ can be expressed as
$$B𝒜=_0^\beta \stackrel{}{A}(\stackrel{}{x})\frac{d\stackrel{}{x}}{d\tau }𝑑\tau ,$$
(5)
where $`\stackrel{}{A}(\stackrel{}{x})`$ is any vector potential whose curl is a homogeneous magnetic field $`B`$. Equations (1), (4), and (5) yield
$$\stackrel{~}{P}(B)=\frac{𝒟[\stackrel{}{x}(\tau )]\mathrm{exp}\left[_0^\beta \left\{\frac{1}{2}\frac{d\stackrel{}{x}}{d\tau }\frac{d\stackrel{}{x}}{d\tau }d\tau \right\}+ie_0^\beta \left\{\stackrel{}{A}\frac{d\stackrel{}{x}}{d\tau }d\tau \right\}\right]}{𝒟[\stackrel{}{x}(\tau )]\mathrm{exp}\left[_0^\beta \left\{\frac{1}{2}\frac{d\stackrel{}{x}}{d\tau }\frac{d\stackrel{}{x}}{d\tau }d\tau \right\}\right]}.$$
(6)
By inspection of Eq.(6) we arrive at
$$\stackrel{~}{P}(B)=Z(B)/Z(0),$$
(7)
where $`Z(B)`$ is the partition function $`(Z(B)=Tr\{\mathrm{exp}[\beta H(B)]\})`$ for a quantum particle of charge e in a homogeneous magnetic field $`B`$ at an inverse temperature $`\beta `$. This is the central result of this paper and it relates Brownian motion and magnetism. As the reader can easily verify, the relation (7) holds even if there is an arbitrary biasing potential. The plane can also be replaced by a sphere or $`(^3)^N`$, the configuration space of $`N`$ particles in $`(^3)`$. In the last case, the area of interest is the sum of the weighted areas of the projections of the closed Brownian paths onto the $`x`$y plane. Now we demonstrate the utility of Eq. (7) by computing the distribution of areas for diffusion on a plane. The partition function $`Z(B)`$ for a particle of unit mass in a constant magnetic field, is easily computed from the energies $`E_n=(n+\frac{1}{2})eB`$ and degeneracy (or the number of states per unit area) $`(eB/2\pi )`$ of Landau levels<sup>5</sup> (throughout this paper we set $`\mathrm{}=c=1)`$:
$$Z(B)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{eB}{2\pi }\mathrm{exp}[(n+\frac{1}{2})\beta eB]=\frac{eB/4\pi }{\mathrm{sinh}(\beta eB/2)}.$$
From (7) we find $`\stackrel{~}{P}(B)=(\beta eB/2)[\mathrm{sinh}(\beta eB/2)]^1`$. Taking the Fourier transform of $`\stackrel{~}{P}(B)`$ by contour integration we get the result $`P(A)=(\pi /2\beta )[\mathrm{cosh}(\pi A/\beta )]^2`$ derived in Ref. 2. This provides a check on Eq. (7) and illustrates its use.
Let us now address the problem posed at the beginning of this paper: what is the distribution $`P(\mathrm{\Omega })`$ of solid angles enclosed by a diffusing particle on a unit sphere? Unlike the planar case, $`P(\mathrm{\Omega })`$ is a periodic<sup>6</sup> function with period $`4\pi `$. The generating function $`P_g`$ of the distribution of solid angles is given by
$$\stackrel{~}{P}_g=_0^{4\pi }𝑑\mathrm{\Omega }P(\mathrm{\Omega })e^{i\frac{g\mathrm{\Omega }}{2}}$$
(8)
with $`g`$ an integer. $`P(\mathrm{\Omega })`$ is expressed in terms of $`\stackrel{~}{P}_g`$ by a Fourier series
$$P(\mathrm{\Omega })=\frac{1}{4\pi }\underset{g=\mathrm{}}{\overset{\mathrm{}}{}}e^{i\frac{g\mathrm{\Omega }}{2}}\stackrel{~}{P}_g$$
(9)
rather than an integral (3). Relation (7) now takes the form
$$\stackrel{~}{P}_g=Z_g/Z_0,$$
(10)
where $`Z_g`$ is the partition function for a particle of charge $`e`$ on a sphere subject to a magnetic field created by a monopole of quantized strength $`G=g/e`$ (Ref. 7) at the center of the sphere. The energy levels of this system are easily computed:<sup>8</sup>
$$E_j=[j(+1)g^2]/2,$$
where $`j`$, the total angular momentum quantum number ranges from $`|g|`$ to infinity, and the $`j\mathrm{th}`$ level is $`(2j+1)`$-fold degenerate. The partition function is consequently given by
$$Z_g=\underset{j=|g|}{\overset{\mathrm{}}{}}(2j+1)e^{\frac{\beta \{j(j+1)g^2\}}{2}}.$$
(11)
Combining (9),(10), and (11) and rearranging the summations we arrive at
$`P(\mathrm{\Omega })=\mathrm{Re}{\displaystyle \frac{1}{2\pi Z_0}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\{(2l+1){\displaystyle \frac{1+\zeta }{2(1\zeta )}}`$
$`+\left[{\displaystyle \frac{2\zeta }{(1\zeta )^2}}\right]\}e^{\frac{\beta }{2}l(l+1)},`$ (12)
where $`\zeta (l,\beta ,\mathrm{\Omega })=\mathrm{exp}[1/2\{\beta (2l+1)+i\mathrm{\Omega }\}].`$ The function (12) is plotted numerically for various values of $`\beta `$ in Fig. 1. The qualitative nature of these plots is easily understood. For small values of $`\beta `$ the particle tends to make small excursions and its path encloses solid angles close to 0 or $`4\pi `$ and consequently the plots are peaked around these two values. As the available time $`\beta `$ increases, other values of $`\mathrm{\Omega }`$ are also probable and the peaks tend to spread and the curves to flatten out. Finally in the limit of $`\beta \mathrm{}`$ the particle has enough time to enclose all possible solid angles with equal probability. These plots give the answer to the question that was raised in the beginning of the paper.
Now we turn to the magnetic properties of spinless bosons. An $`N`$ particle system in three dimensions placed in a homogeneous external magnetic field which is along the $`z`$ direction has the Hamiltonian
FIG.1. The probability distribution $`P(\mathrm{\Omega })`$ of solid angles for closed randon walks lasting a time $`\beta `$. $`P(\mathrm{\Omega })`$ is plotted above for four values of $`\beta `$:0.5, 2, 5, and 10.
$$H(\stackrel{}{x^a},\stackrel{}{p_a})=\underset{a=1}{\overset{N}{}}\frac{[\stackrel{}{p_a}e_a\stackrel{}{A}(\stackrel{}{x^a})]^2}{2m_a}+V(\stackrel{}{x^a}),$$
where $`\stackrel{}{A}(\stackrel{}{x^a})`$ is the vector potential of the external magnetic field. $`V(\stackrel{}{x_a)}`$ includes an arbitrary interaction between the particles as well as an external potential, $`m_a`$ and $`e_a`$ are the masses and charges of the particles. $`\{\stackrel{}{x^a},a=1,2,\mathrm{},N\}`$ are the position vectors of the $`N`$ particles. The configuration space of the system is given by $`Q=(^3)^N/`$, where $``$ means that we identify points in $`(^3)^N`$ which differ by an exchange of identical particles. For simplicity we give the argument for $`N`$ identical particles with unit mass and charge $`e_a=e`$. The argument is easily adapted to several species of particles of arbitrary charge and mass.
Now consider a diffusion on $`Q`$ biased by the potential $`V(\stackrel{}{x^a})`$. The Wiener measure is now appropriately modified:
$$f[\stackrel{}{x}(\tau )]_{𝒲(V)}\frac{𝒟[\stackrel{}{x}(\tau )]f[\stackrel{}{x}(\tau )]\mathrm{exp}\left[_0^\beta \left\{\frac{1}{2}\left(_a\frac{d\stackrel{}{x_a}}{d\tau }\frac{d\stackrel{}{x_a}}{d\tau }\right)+V(\stackrel{}{x_a})\right\}𝑑\tau \right]}{𝒟[\stackrel{}{x}(\tau )]\mathrm{exp}\left[_0^\beta \left\{\frac{1}{2}\left(_a\frac{d\stackrel{}{x_a}}{d\tau }\frac{d\stackrel{}{x_a}}{d\tau }\right)+V(\stackrel{}{x^a})\right\}𝑑\tau \right]}.$$
The area whose distribution we are interested in is defined as follows: Let $`q(\tau )`$ be a closed curve in $`Qq(\tau )`$ determines trajectories of particles $`\{\stackrel{}{x^a}(\tau ),a=1,2,\mathrm{},N\}`$ in $`^3`$. The area functional of interest is $`𝒜[q(\tau )]_a\stackrel{}{A}(\stackrel{}{x^a})𝑑\stackrel{}{x^a}`$. The area functional has the following interpretation. If the final positions of the $`N`$ particles are the same as the initial ones (direct processes), $`𝒜[q(\tau )]`$ is simply the sum of the areas enclosed by the projection of the particle trajectories on the $`(xy)`$ plane. If the final positions differ from the initial ones by a permutation (exchange processes), the projections of the particle trajectories still define closed curves on the $`(xy)`$ plane. $`𝒜[q(\tau )]`$ is defined as the sum of areas enclosed by these closed curves.
As before we find that $`\stackrel{~}{P}(B)`$, which is the Fourier transform of the distribution $`P(A)\delta (𝒜[\stackrel{}{x}(\tau )]A)_{𝒲(V)}`$ of areas, is given by Eq. (7). It is crucial for our argument that the particles obey Bose statistics.<sup>9</sup> Since $`P(A)`$ is a probability distribution, $`\stackrel{~}{P}(B)=Z(B)/Z(O)`$ is the Fourier transform of a positive function. This places strong restrictions on the partition function $`Z(B)`$. Let $`u_i,i=1,\mathrm{},n`$ be $`n`$ real numbers. If one defines the $`n\times n`$ matrix $`D_{ij}^{(n)}=\stackrel{~}{P}(u_iu_j)`$, the necessary and sufficient condition for $`\stackrel{~}{P}(B)`$ to be the Fourier transform of a positive function is<sup>10,11</sup>
$$\mathrm{\Delta }^{(n)}\mathrm{Det}D^{(n)}0\mathrm{for}\mathrm{all}n.$$
(13)
This imposes restrictions on the free energy $`F(B)=(1/\beta )\mathrm{ln}Z(B)`$ of the system in the presence of a magnetic field $`B`$.
For the simplest nontrivial case $`n=2`$, the inequality (13) with $`B=u_1u_2`$ leads to
$$Z(B)Z(0)$$
(14)
or equivalently, $`F(B)F(0)`$. Since the free energy of the system increases in the presence of a magnetic field, the material is diamagnetic. This universal diamagnetic behavior of spinless bosons at all temperatures is known in the mathematical physics literature.<sup>11</sup> However, our approach may be accessible to a wider community of physicists. Our approach relating Brownian motion to magnetism enriches both fields and provides each field with intuition derived from the other. For instance, the zero-field susceptibility $`\chi =^2F(B)/B^2|_{B=0}`$ of the magnetic system is related to the variance of the distribution of areas in the diffusion problem:
$$\chi =\frac{1}{\beta }[\mathrm{ln}\stackrel{~}{P}(B)]^{\prime \prime }|_{B=0}=\frac{e^2}{\beta }\overline{(A\overline{A})^2}=\frac{e^2}{\beta }\mathrm{Var}A0.$$
(15)
It is curious that the zero-field susceptibility can be interpreted as the variance of the distribution of areas. Since the variance cannot be negative, it follows that $`\chi `$, the zero-field susceptibility cannot be positive and so these systems are diamagnetic.
Next consider the case $`n=3`$. The $`3\times 3`$ matrix $`D_{ij}^{(3)}`$ will then be a function of $`u=u_1u_2`$ and $`v=u_2u_3(u_1u_3`$, being expressible in terms of $`u`$ and $`v)`$. If we set $`u=0`$ (i.e., set $`u_1=u_2=0,u_3=v)`$, we find that $`\mathrm{\Delta }^{(3)}(u,v)|_{u=0}=0`$. It then follows from the inequality (13) that $`\mathrm{\Delta }^{(3)}(u,v)`$ has a minimum at $`u=0`$ for all $`v`$. This implies that $`^2\mathrm{\Delta }^{(3)}/u^2|_{u=0,v=B}0`$. Defining the function $`𝒰(B)=\stackrel{~}{P}^2[1\stackrel{~}{P}^2]^1`$, where the prime means derivative with respect to the magnetic field, we find
$$𝒰(B)𝒰(0),$$
(16)
As can be seen by taking the limit $`B0`$, $`𝒰(0)=\stackrel{~}{P^{\prime \prime }}(0)=\beta \chi (0)`$. We define a critical field $`B_c=\pi /[2\sqrt{\beta \chi (0)}]`$. The inequality (16) implies a bound on the partition function. Notice that $`\stackrel{~}{P}^{}`$ lies in a cone defined by the lines of slope $`(\pi /2B_c)\sqrt{(1\stackrel{~}{P}^2)}`$ and $`(\pi /2B_c)\sqrt{(1\stackrel{~}{P}^2)}`$. It follows that
$$\stackrel{~}{P}(B)\mathrm{cos}(\pi B/2B_c)\mathrm{for}|B|B_c.$$
(17)
The diamagnetic inequality due to Simon and Nelson<sup>11</sup> gives an upper bound on the partition function $`Z(B)`$ of a system of spinless bosons. The new inequality stated in (17) gives us a lower bound on $`Z(B)`$ (see Fig. 2) (or equivalently, an upper bound on the free energy).
Fig. 2. The region forbidden by the bounds \[inequalities (14) and (18) on the partition function. These bounds are shown as solid lines. The dotted curve is the partitition for a charged simple harmonic oscillator in an external magnetic field. Notice that the dotted curve lies outside the forbidden region.
As an explicit check on this new bound on the free energy we considered a simple system–a charged particle in a magnetic field subject to a harmonic oscillator potential. The calculated partition function of this system is close to, but above the lower bound set by (17). Needless to say, our bound is derived for an arbitrary interacting system of spinless bosons. The new bound presented here along with the earlier (14) diamagnetic inequality<sup>11</sup> places strong restrictions on the partition function of a bosonic system in the presence of a magnetic field. We find a curious and immediate consequence of these restrictions: if the zero-field susceptibility of the system vanishes, then Eqs. (14) and (17) imply that $`Z(B)=Z(O)`$, i.e., the system is nonmagnetic at all fields.
The key result of this paper is a connection between two apparently distinct classes of problems–Brownian motion and magnetism. This allows us to compute the distribution of solid angles enclosed in Brownian motion on a sphere. As mentioned earlier, this problem comes up when computing the distribution of Berry phases in a random magnetic field. A more classical context is depolarized light scattering. As is well known, a light ray following a space curve picks up a geometric phase,<sup>12,13</sup> equal to the solid angle swept out by the direction vector. If a light ray inelastically scatters off a random medium, its direction vector does a random walk on the unit sphere of directions. The distribution $`P(\mathrm{\Omega })`$ computed here is relevant to the extent of depolarization in such an experiment.<sup>14</sup>
In the domain of magnetism we find an independent way of arriving at the diamagnetic inequality<sup>11</sup> which states that the free energy of a system of spinless bosons always increases in the presence of a magnetic field. Spinless charged bosonic systems occur in the context of superconductors (which are perfect diamagnets) and neutron stars.<sup>15</sup> We believe that the community of physicists working in these areas may not be aware of the general results available in the mathematical literature. For instance, the diamagnetism of bosons may be relevant<sup>16</sup> to the interpretation of recent experiments<sup>17</sup> on high-$`T_c`$ superconductivity.
Throughout this paper we have only discussed homogeneous magnetic fields. It is easy to generalize our discussion to take into account arbitrary inhomogeneous fields: all one does is consider the distribution of weighted areas. An obvious application of this is the computation of the probability of entanglement of a polymer with a background lattice of polymers. We expect the new method outlined here to shed light on open problems in polymer entanglement involving more complicated configurations of polymers than the simplest one solved so far. One can also use the relation (7) to compute the distribution of winding numbers in diffusion in a multiply connected space.
It is a pleasure to thank N. Kumar for bringing up the problem of diffusion on a sphere and several discussions on this work; Barry Simon for his help in finding Ref. 11; Diptiman Sen for discussions and for giving us Ref. 8, and R. Nityananda for discussions and drawing our attention to Ref. 10.
$``$
* See Geometric Phases in Physics, edited by Alfred Shapere and Frank Wilczek (World Scientific, Singapore, 1989).
* M. G. Brereton and C. Butler, J. Phys. (London) A 20, 3955, (1987); D. C. Khandekar and F. W. Wiegel, J. Phys. (London) A 21, L563 (1988); J. Stat. Phys. 53, 1073 (1988).
* L. S. Schulman, Techniques and Applications of Path Integration, (Wiley, New York, 1981).
* R. P. Feynman and A. R. Hibbs, Quantum Mechanics and Path Integrals (McGraw-Hill, New York, 1965).
* L. D. Landau and E. M. Lifshitz, Quantum Mechanics (Pergamon, Oxford, 1977), p. 457.
* A path that encloses a solid angle $`\mathrm{\Omega }`$ to its left also encloses a solid angle of $`4\pi \mathrm{\Omega }`$ to its right. Using the symmetry $`P(\mathrm{\Omega })=P(\mathrm{\Omega })`$ we conclude that $`P(\mathrm{\Omega })`$ is periodic with period $`4\pi `$.
* P. A. M. Dirac, Proc. R. Soc. London A133, 60 (1931); M. N. Saha, Ind. J. Phys. 10, 141 (1936).
* S. Coleman, in The Unity of Fundamental Interactions, edited by A. Zichichi (Plenum, New York, 1983), p. 21-117.
* The identity (10) is not true for Fermi particles since in this case exchange processes are weighted negatively relative to direct ones in the expression for the partition function. This ruins the probability interpretation that we give to the partition function.
* M. Komesaroff et al., Astron. Astrophys. 93, 269 (1981).
* B. Simon, Phys. Rev. Lett. 36, 1083 (1976); Indiana Univ. Math. J. 26, 1067 (1977); J. Functional Anal. 32, 97 (1979); Functional Integration and Quantum Physics (Academic, New York, 1979), and references therein.
* L. D. Landau and E. M. Lifshitz, Electrodynamics of Continuous Media (Pergamon, Oxford, 1960), p. 271.
* A. Tomita and R. Chiao, Phy. Rev. Lett. 57, 937 (1986).
* S. Sanyal et al., Phys. Rev. Lett. 72, 2963 (1994).
* J. Daicic et al., Phys. Rep. 237, 65 (1994); J. Daicic et al., Phys. Rev. Lett. 71, 1779 (1993).
* In these experiments, diamagnetism is observed well above $`T_c`$. A possible explanation may be that the diamagnetism is due to the persistence of Cooper pairs (spinless bosons) well above $`T_c`$. We thank N. Kumar for this observation. See also A. S. Alexandrov and N. F. Mott, Phys. Rev. Lett. 71, 1075 (1993); D. C. Bardos et al., Phys. Rev. B 49, 4082 (1994).
* W. C. Lee et al., Phy. Rev. Lett. 63, 1012 (1989).
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# Feedback stabilization of quantum ensembles: a global convergence analysis on complex flag manifolds
## I Introduction
The use of feedback in quantum mechanics is limited by the phenomenon of wavefunction collapse following a measurement. In this work the problem is bypassed by considering density operators of quantum ensembles and completely noninvasive measurements. This allows also to relax the requirement of commutativity of the measured observables and in fact we shall assume to have a complete knowledge of the density operator for all times. Although physically this set up is realistic only for some applications (typically nuclear spin ensembles Cory1 ; Havel1 ), it is of widespread use for the purposes of model-based quantum control (often under the name “tracking control” Brown1 ; Zhu1 ), as it allows to generate control fields in spite of the high complexity of open loop control Boscain2 ; DAlessandro1 ; Schirmer5 . Furthermore, while the formulation comes from quantum control, our motivations for this work are mostly mathematical, namely feedback design and convergence analysis for a class of bilinear control systems living on a particular family of compact manifolds and evolving isospectrally. As the system has a drift term which cannot be canceled without incurring in singularities of the control law, the most natural problem formulation is to seek for a stabilizer to the periodic orbit drawn by the drift. Rather than studying this problem like an orbital stabilization problem Bacciotti1 , we reformulate and solve it as a state tracking problem, thus avoiding the obstruction to semiglobal convergence of a periodic orbit, see Wilson1 , Corollary 1.6 (where it is called stability in the large). In fact, with our feedback design the state will converge to the orbit, but the entire orbit is not an invariant set, only a point moving along it is invariant. As a matter of fact, by passing to a suitable rotating frame, our time-dependent trajectory tracking problem can be reformulated completely in terms of time-varying feedback law for the fixed point of a nonautonomous system. The Lyapunov design is essentially of the Jurdjevic-Quinn type Jurdjevic5 , for which the usual LaSalle invariance principle is applicable in spite of the time-dependence of the closed loop, and does not differ much from what has already been proposed in the literature for wavefunctions Ferrante1 ; Vettori1 ; Grivopoulos1 ; Mirrahimi1 .
What is nontrival is to ascertain the convergence of certain initial conditions and to provide a global description of the region of attraction. In fact, the sufficient condition used in Jurdjevic5 to prove asymptotic convergence and based on the so-called $`\mathrm{ad}`$-condition or Jurdjevic-Quinn condition Bacciotti2 , is never verified globally for $`N>2`$. This is due to the presence of an abelian subalgebra (Cartan subalgebra) that can never be fully spanned by $`\mathrm{ad}`$-commutators alone. It will be shown, however, that the undesired critical points are not only unstable but also repulsive, meaning that the Jurdjevic-Quinn condition or, equivalently, the controllability of the linearization along the desired orbit Mirrahimi1 , guarantees convergence for all initial conditions outside the set of equilibrium points. To attain a complete and time-independent description of the critical set and thus of the domain of attraction, a thorough geometric and topologic characterization of the state manifold of mixed density operators is required. A unitary evolution like that appearing in a Liouville equation is isospectral, as the eigenvalues of the density operator form a complete set of invariants. Unlike for a wavefunction, the state space has dimension and structure which depend on the multiplicities of such eigenvalues. Since these form a flag in dimension $`N`$, all complex flag manifolds obtained as homogeneous spaces of $`U(N)`$ (or $`SU(N)`$) by the Cartesian products of subgroups of dimensions given by the multiplicities are admissible state spaces Bengtsson1 ; Picken1 ; Zyczkowski2 . Since the set of density operators up to the imaginary unit “overlaps” with the Lie algebra $`𝔲(N)`$ (or $`𝔰𝔲(N)`$, excluding the constant trace), these complex flag manifolds can also be intended as the orbits of the (co)adjoint action of $`U(N)`$ (or $`SU(N)`$) on its Lie algebra. This is a well-studied action and the structure of its orbits is well-known Marsden2 ; Frankel1 : for example a fundamental topological invariant like the Euler characteristic acquires the meaning of number of nontrivial possible permutations of the eigenvalues of the density operator, see also Bengtsson1 ; Chaturvedi1 ; Ercolessi1 ; Zyczkowski2 . For the purposes of stabilizability, this is an important feature, because it will be shown that each complex flag manifold has a number of “antipodal” points equal to the Euler characteristic, and that these points must be equilibria of the closed loop system. In order to give a complete description of the region of attraction, we use the resemblance between the set of density operators and the Lie algebra $`𝔰𝔲(N)`$, and a few tools deriving from the root space decomposition of a compact Lie algebra, namely its orthogonal decomposition into Cartan subalgebra plus root spaces and the invariance properties of the root spaces under certain commutators (like the ad-commutators) Cla-contr-root1 . This “graph-like” approach yields simple, time-independent characterizations of all converging initial conditions for a given reference orbit and Hamiltonian. Also the Kalman controllability of the linearization admits an intrinsic formulation in these terms. The formalism used gives insight into the problem of choosing reference orbits having a large domain of attraction.
It is known Bhat1 ; Koditschek3 , that compact manifolds do not admit a global asymptotically stable equilibrium because they are not contractible. This is a topological property and corresponds to a set being homotopy equivalent to a point Guillemin1 . The region of convergence of an asymptotically stable attractor must be in such a homotopy class Bhat1 ; Wilson1 . For our complex flag manifolds, it will be shown that the antipodal points represent topological obstructions to global stabilizability.
In order to simplify the treatment, an equivalent real representation of density operators is used throughout, given by the so-called coherence vector and corresponding to the vector of expectation values with respect to a complete orthonormal set of Hermitian matrices Alicki1 ; Bengtsson1 ; Schirmer7 . It provides a linear representation of the adjoint action occurring in a Liouville equation Cla-spin-tens1 , and it allows to formulate the control system in terms of standard bilinear systems on smooth manifolds which are real representations of the complex flag manifolds.
## II Driven Liouville-von Neumann equation
With a given Hamiltonian $`H=H_A+uH_B`$, $`iH_A,iH_B𝔰𝔲(N)`$, $`uC^{\mathrm{}}()`$ a control field, one can form a Schrödinger equation for the wavefunction $`|\psi `$ (in atomic units, $`\mathrm{}=1`$)
$$|\dot{\psi }=i\left(H_A+uH_B\right)|\psi ,|\psi (t)𝕊^{2N1},$$
(1)
or a Liouville-von Neumann equation for the density operator $`\rho `$
$$\dot{\rho }=i[H_A+uH_B,\rho ],\rho =\rho ^{}0,\mathrm{tr}(\rho )=1,\mathrm{tr}(\rho ^2)1.$$
(2)
Eq. (2) holds for a quantum ensemble, hence it is more general than (1); the two being equivalent only when $`\mathrm{tr}(\rho ^2)=1`$ (i.e., $`\rho `$ is a rank-one operator: $`\rho =|\psi \psi |`$).
### II.1 Gell-Mann basis and adjoint representation
The left hand side of (2) contains a conjugation action on state matrices. In order to deal with it, we can use one of the features of the adjoint representation, namely its providing a linear representation of one-parameter groups of automorphisms of $`𝔰𝔲(N)`$ (see Appendix A), and reformulate (2) as a standard bilinear control systems on a suitable manifold. To do that, we use the coherence vector representation of $`\rho `$, whose key property is that both $`H`$ and $`\rho `$ are expressed in terms of the same complete orthonormal set of Hermitian matrices. Let $`\lambda _0=\frac{1}{\sqrt{N}}𝟙_{}`$ and call $`𝝀`$ the $`n`$-dimensional vector of $`N\times N`$ Gell-Mann matrices Georgi1 , $`n=N^21`$. Since $`𝔰𝔲(N)`$ contains traceless skew-hermitian matrices, $`\mathrm{span}\{i𝝀\}=𝔰𝔲(N)`$. Denote with $`𝔥`$ the Cartan subalgebra of $`𝔰𝔲(N)`$, i.e., the maximally abelian subalgebra in $`𝔰𝔲(N)`$, $`\mathrm{dim}(𝔥)=N1`$. In the Gell-Mann basis $`𝝀`$, $`𝔥`$ corresponds to the $`N1`$ diagonal matrices. Denote $`𝔨`$ the vector space such that $`𝔰𝔲(N)=𝔥𝔨`$, with $`𝔥𝔨`$ in a standard biinvariant $`𝔰𝔲(N)`$ metric (e.g. the Killing metric). While $`𝔥`$ is an abelian subalgebra, $`𝔨`$ is only a vector space. In correspondence of this direct sum, we have the decomposition of $`𝝀`$ into $`𝝀_𝔥`$ and $`𝝀_𝔨`$ so that $`𝔥=\mathrm{span}\{i𝝀_𝔥\}`$ and $`𝔨=\mathrm{span}\{i𝝀_𝔨\}`$. If $`E_j\mathrm{}`$ is the elementary $`N\times N`$ matrix having 1 in the $`(j\mathrm{})`$ slot and 0 elsewhere, then the matrices $`𝝀`$ are given by
$$\left\{\lambda _{𝔥,j},1jN1\right\}=\left\{(E_{11}+\mathrm{}+E_{jj}jE_{j+1,j+1})/\sqrt{j(j+1)},1jN1\right\}$$
(3)
for the diagonal part, and
$`\left\{\lambda _{𝔨,\mathrm{},j\mathrm{}},1j<\mathrm{}N\right\}`$ $`=`$ $`\left\{(E_j\mathrm{}+E_\mathrm{}j)/\sqrt{2},1j<\mathrm{}N\right\}`$ (4)
$`\left\{\lambda _{𝔨,\mathrm{},j\mathrm{}},1j<\mathrm{}N\right\}`$ $`=`$ $`\left\{i(E_j\mathrm{}+E_\mathrm{}j)/\sqrt{2},1j<\mathrm{}N\right\}`$ (5)
for the off-diagonal part. Calling $`𝔨_j\mathrm{}=\mathrm{span}\{i\lambda _{𝔨,\mathrm{},j\mathrm{}},i\lambda _{𝔨,\mathrm{},j\mathrm{}}\}`$, then we have the further splitting of $`𝔨`$ into “root spaces”
$$𝔨=\underset{1j<\mathrm{}N}{}𝔨_j\mathrm{}$$
(6)
with the following commutation relations (see for instance Cla-contr-root1 for the details):
$$[𝔥,𝔨_j\mathrm{}]=𝔨_j\mathrm{},$$
(7)
$$[𝔨_j\mathrm{},𝔨_{pq}]=\{\begin{array}{cc}\mathrm{}\hfill & \text{ if }\mathrm{}p\text{ and }jq\hfill \\ 𝔨_{jq}\hfill & \text{ if }\mathrm{}=p\hfill \\ 𝔨_p\mathrm{}\hfill & \text{ if }j=q\hfill \\ 𝔥\hfill & \text{ if }j=p\text{ and }\mathrm{}=q.\hfill \end{array}$$
(8)
Assume $`H_A`$ is diagonal
$$H_A=\left[\begin{array}{c}_1\\ & \mathrm{}& \\ & & _N\end{array}\right],_1+\mathrm{}+_N=0,$$
and nondegenerate, i.e., $`_j_{\mathrm{}}`$, $`j\mathrm{}`$, where the $`_j`$ are supposed ordered: $`_1<_2<\mathrm{}<_N`$. The $`_j`$ are the energy levels of the (unforced) system (1), i.e., the eigenvalues of the stationary Schrödinger equation $`H_A|\psi _j=_j|\psi _j`$ of eigenstates $`|\psi _j=\mathrm{e}_j`$ ($`\mathrm{e}_j`$ is the elementary basis vector), $`j=1,\mathrm{},N`$. Assume also that the transition frequencies are nondegenerate, i.e., that the levels are not equispaced $`_j_{\mathrm{}}_p_q`$, $`(j\mathrm{})(pq)`$ $`j\mathrm{}`$, $`pq`$. Further standard assumptions are that $`H_B`$ is off-diagonal and $`\mathrm{Graph}(H_B)`$ connected. A stronger assumption we shall need is that $`H_B`$ enables all transitions among adjacent energy levels: $`\mathrm{tr}\left(H_B𝔨_{j,j+1}\right)0`$ $`j=1,\mathrm{},N1`$. Beside connectivity of $`\mathrm{Graph}(H_B)`$, it guarantees that all “fundamental root spaces” (see Cla-contr-root1 ) are excited by the dynamics. Since $`H_A`$ is diagonal and traceless, $`iH_A𝔰𝔲(N)`$. A nondegenerate element of the Cartan subalgebra (like $`H_A`$ above) is called regular. It is called strongly regular if in addition it has all nondegenerate transitions.
In terms of $`𝝀`$, $`H_A=𝒉_A𝝀`$ and $`H_B=𝒉_B𝝀`$, with $`𝒉_A,𝒉_B^n`$. Likewise $`\rho =\varrho _0\lambda _0+\mathit{\varrho }𝝀`$, where $`\varrho _0=\frac{1}{\sqrt{N}}`$ and the coherence vector $`\mathit{\varrho }`$ is composed of expectation values along the basis elements (3)-(5): $`\varrho _{𝔥,j}=\mathrm{tr}(\rho \lambda _{𝔥,j})`$, $`j=1,\mathrm{},N1`$, $`\varrho _{𝔨,\mathrm{},j\mathrm{}}=\mathrm{tr}(\rho \lambda _{𝔨,\mathrm{},j\mathrm{}})`$, $`\varrho _{𝔨,\mathrm{},j\mathrm{}}=\mathrm{tr}(\rho \lambda _{𝔨,\mathrm{},j\mathrm{}})`$, $`1j<\mathrm{}N`$. Any density $`\rho `$ can be split as $`\rho =\varrho _0\lambda _0+\rho _𝔥+\rho _𝔨`$, or in correspondence of (3)-(5), $`\rho =\varrho _0\lambda _0+\mathit{\varrho }_𝔥𝝀_𝔥+\mathit{\varrho }_𝔨𝝀_𝔨=\varrho _0\lambda _0+_{1j<N}\varrho _{𝔥,j}\lambda _{𝔥,j}+_{1j<\mathrm{}N}\left(\varrho _{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}+\varrho _{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}\right)`$. The nonzero components of the coherence vector $`\mathit{\varrho }`$ uniquely identify a subset of the $`𝔨_j\mathrm{}`$. Denote $`𝔣_𝔨(\rho )`$ the “support” of $`\rho `$ in $`𝔨`$, i.e., the set of root spaces “touched” by $`\rho `$: $`𝔣_𝔨(\rho )=\rho 𝔨`$. Also let $`_𝔨(\rho )=\{(j\mathrm{})\text{ s.t. }\mathrm{tr}\left(\rho 𝔨_j\mathrm{}\right)0,\mathrm{\hspace{0.33em}1}j<\mathrm{}N\}`$ be the corresponding set of index pairs. When $`(j\mathrm{})_𝔨(\rho )`$, then $`(\varrho _{𝔨,\mathrm{},j\mathrm{}},\varrho _{𝔨,\mathrm{},j\mathrm{}})(0,0)`$. Likewise $`𝔣_𝔥(\rho )=\rho 𝔥`$ and $`_𝔥(\rho )=\{(j)\text{ s.t. }\varrho _{𝔥,j}0,\mathrm{\hspace{0.33em}1}j<N\}`$. In the following we shall use both symbols $`\rho `$ and $`\mathit{\varrho }`$ for densities and we shall refer to “diagonal” and “off-diagonal” $`\mathit{\varrho }`$ with an obvious abuse of notation. The Hilbert-Schmidt norm on density operators induces for $`\mathit{\varrho }`$ the standard Euclidean norm $``$: $`\mathrm{tr}\left(\rho ^2\right)=\varrho _0^2+\mathit{\varrho },\mathit{\varrho }=\varrho _0^2+\mathit{\varrho }^2`$, where we indicate with $`,`$ the $`^n`$-Euclidean inner product. Due to the trace-class constraint, the notion of distance between the densities $`\rho _1`$ and $`\rho _2`$ having the same purity is $`d(\rho _1,\rho _2)=\mathrm{tr}\left(\rho _1^2\right)\mathrm{tr}\left(\rho _1\rho _2\right)`$, see e.g. Zyczkowski2 , or in terms of $`\mathit{\varrho }`$:
$$d(\mathit{\varrho }_1,\mathit{\varrho }_2)=\varrho _0^2+\mathit{\varrho }_1^2\varrho _0^2\mathit{\varrho }_1,\mathit{\varrho }_2=\mathit{\varrho }_1^2\mathit{\varrho }_1^T\mathit{\varrho }_2[0,\mathrm{tr}\left(\rho _1^2\right)].$$
(9)
Thanks to the use of the same basis for $`\rho `$ and the Hamiltonian, up to the imaginary unit the trajectories of (2) can be identified with the adjoint orbits of $`SU(N)`$ on its Lie algebra, see Section III for a thorough description. Following Appendix A, we can replace the matrix ODE (2) with the linear vector ODE
$$\dot{\mathit{\varrho }}=\left(𝑨+u𝑩\right)\mathit{\varrho },$$
(10)
where $`𝑨=i\mathrm{ad}_{H_A}=i𝒉_A\mathrm{ad}_𝝀`$, $`𝑩=i\mathrm{ad}_{H_B}=i𝒉_B\mathrm{ad}_𝝀`$ and $`𝑨,𝑩\mathrm{ad}_{𝔰𝔲(N)}𝔰𝔬(n)`$. The Lie algebra $`\mathrm{ad}_{𝔰𝔲(N)}`$ is the adjoint representation of $`𝔰𝔲(N)`$, hence $`\mathrm{dim}(\mathrm{ad}_{𝔰𝔲(N)})=\mathrm{dim}(𝔰𝔲(N))=n`$. Therefore, for $`N>2`$, $`\mathrm{ad}_{𝔰𝔲(N)}𝔰𝔬(n)`$.
The isomorphism $`𝔰𝔲(N)\mathrm{ad}_{𝔰𝔲(N)}`$ induces an orthogonal splitting also in the adjoint representation: $`\mathrm{ad}_{𝔰𝔲(N)}=\mathrm{ad}_𝔥\mathrm{ad}_𝔨`$. Obviously relations similar to (7)-(8) still hold and $`𝑨=i𝒉_{A,𝔥}\mathrm{ad}_{𝝀_𝔥}`$, $`𝑩=i𝒉_{B,𝔨}\mathrm{ad}_{𝝀_𝔨}`$. For $`C𝔰𝔲(N)`$, in components $`C=_{1j<N}c_{𝔥,j}\lambda _{𝔥,j}+_{1j<\mathrm{}N}\left(c_{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}+c_{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}\right)`$, we also shall indicate with $`𝔣_𝔥(C)`$, $`𝔣_𝔨(C)`$ the support of $`C`$ in, respectively, $`𝔥`$, $`𝔨`$, of indices $`_𝔥(C)`$, $`_𝔨(C)`$.
For later use, we need to compute some of the commutators of (7)-(8) more in detail.
$$[\lambda _{𝔨,\mathrm{},j\mathrm{}},\lambda _{𝔨,\mathrm{},j\mathrm{}}]=i(E_{jj}E_{\mathrm{}\mathrm{}})=\{\begin{array}{cc}\sqrt{\frac{j1}{j}}\lambda _{𝔥,j1}+_{p=j}^{\mathrm{}}\frac{1}{\sqrt{p(p+1)}}\lambda _{𝔥,p}+\sqrt{\frac{\mathrm{}}{\mathrm{}1}}\lambda _{𝔥,\mathrm{}1}\hfill & \text{ if }j>1\text{ and }\mathrm{}>2\hfill \\ _{p=j}^{\mathrm{}}\frac{1}{\sqrt{p(p+1)}}\lambda _{𝔥,p}+\sqrt{\frac{\mathrm{}}{\mathrm{}1}}\lambda _{𝔥,\mathrm{}1}\hfill & \text{ if }j=1\text{ and }\mathrm{}>2\hfill \\ \sqrt{\frac{\mathrm{}}{\mathrm{}1}}\lambda _{𝔥,\mathrm{}1}\hfill & \text{ if }\mathrm{}=2\hfill \end{array}$$
(11)
For $`(j\mathrm{})(pq)`$:
$$\begin{array}{cc}\hfill [\lambda _{𝔨,\mathrm{},j\mathrm{}},\lambda _{𝔨,\mathrm{},pq}]& =\frac{i}{\sqrt{2}}\left(\delta _\mathrm{}p\lambda _{𝔨,\mathrm{},jq}+\delta _{jp}\lambda _{𝔨,\mathrm{},\mathrm{}q}+\delta _{jq}\lambda _{𝔨,\mathrm{},\mathrm{}p}+\delta _{lq}\lambda _{𝔨,\mathrm{},jp}\right)\hfill \\ \hfill [\lambda _{𝔨,\mathrm{},j\mathrm{}},\lambda _{𝔨,\mathrm{},pq}]& =\frac{i}{\sqrt{2}}\left(\delta _\mathrm{}p\lambda _{𝔨,\mathrm{},jq}\delta _{jp}\lambda _{𝔨,\mathrm{},\mathrm{}q}+\delta _{jq}\lambda _{𝔨,\mathrm{},\mathrm{}p}+\delta _{lq}\lambda _{𝔨,\mathrm{},jp}\right)\hfill \\ \hfill [\lambda _{𝔨,\mathrm{},j\mathrm{}},\lambda _{𝔨,\mathrm{},pq}]& =\frac{i}{\sqrt{2}}\left(\delta _\mathrm{}p\lambda _{𝔨,\mathrm{},jq}+\delta _{jp}\lambda _{𝔨,\mathrm{},\mathrm{}q}\delta _{jq}\lambda _{𝔨,\mathrm{},\mathrm{}p}+\delta _{lq}\lambda _{𝔨,\mathrm{},jp}\right)\hfill \end{array}$$
(12)
### II.2 Unforced equation
For pure states in an orthonormal basis, the eigenvectors $`|\psi _j=\mathrm{e}_j`$ of the stationary Shrödinger equation are mapped into the diagonal density operator $`|\psi _j\psi _j|=E_{jj}`$. More generally, for quantum ensembles, after a suitable diagonalization, $`\stackrel{~}{\rho }=\mathrm{diag}(w_1,\mathrm{},w_N)`$, $`0w_j1`$, $`_{j=1}^Nw_j=1`$. The eigenvalues $`w_j`$ represent the populations of the various energy levels and provide a complete set of invariants for (2), call it $`𝒥=\{w_1,\mathrm{},w_N\}`$, since (2) is isospectral.
###### Proposition 1
Consider the system (10) with $`H_A`$ strongly regular. The state $`\mathbf{\varrho }`$ is an equilibrium point of (10) for $`u=0`$ if and only if $`\rho =\varrho _0\lambda _0+\mathbf{\varrho }_𝔥𝛌_𝔥`$. Furthermore, if $`\mathbf{\varrho }_𝔨0`$, then for $`u=0`$
1. $`𝔣_𝔨(\rho (0))=𝔣_𝔨(\rho (t))`$;
2. $`\varrho _{𝔨,\mathrm{},j\mathrm{}}^2+\varrho _{𝔨,\mathrm{},j\mathrm{}}^2=\mathrm{const}`$;
3. for $`\delta t`$ small, $`\varrho _{𝔨,\mathrm{},j\mathrm{}}(t)\varrho _{𝔨,\mathrm{},j\mathrm{}}(t+\delta t)`$ and $`\varrho _{𝔨,\mathrm{},j\mathrm{}}(t)\varrho _{𝔨,\mathrm{},j\mathrm{}}(t+\delta t)`$ $`(j\mathrm{})_𝔨(\rho )`$.
###### Proof.
When $`u=0`$, for a given $`\rho =\varrho _0\lambda _0+\rho _𝔥+\rho _𝔨`$,
$`i[H_A,\rho _𝔥]`$ $`=`$ $`0`$ (13)
$`i[H_A,\rho ]`$ $`=`$ $`i[H_A,\rho _𝔨]=i[H_A,{\displaystyle \underset{(j\mathrm{})_𝔨(\rho )}{}}\varrho _{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}+\varrho _{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}]`$ (14)
$`=`$ $`{\displaystyle \underset{(j\mathrm{})_𝔨(\rho )}{}}\left(_j_{\mathrm{}}\right)\left(\varrho _{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}\varrho _{𝔨,\mathrm{},j\mathrm{}}\lambda _{𝔨,\mathrm{},j\mathrm{}}\right).`$
In terms of the coherence vector $`\mathit{\varrho }`$ and using the isomorphism (30) (meaning $`𝑨\mathit{\varrho }(t)i[H_A,\rho ]`$), from (13) if $`\mathit{\varrho }_𝔨=0`$, $`𝑨\mathit{\varrho }=0`$, i.e., $`\rho =\varrho _0\lambda _0+\rho _𝔥`$ is a fixed point. To show the other direction, notice that in (14) $`_j_{\mathrm{}}0`$ $`(j\mathrm{})`$ $`1j<\mathrm{}N1`$, since $`H_A`$ is nondegenerate. Hence whenever $`_𝔨(\rho )0`$, $`𝑨\mathit{\varrho }=𝑨\mathit{\varrho }_𝔨0`$, because of the invariance of the $`𝔨_{pq}`$ subspaces under $`𝔥`$, see also (7). Therefore when $`\mathit{\varrho }_𝔨0`$ the unforced system flows along nontrivial periodic orbits. Condition 1 of the last part also follows from (7). Since $`\mathit{\varrho }_𝔥(t)=\mathit{\varrho }_𝔥(0)`$, it must be $`\mathit{\varrho }_𝔨(t)=\mathrm{const}`$ $`t`$. This, together with the invariance property (7) yields 2. Finally, Item 3 follows from $`\rho _𝔨`$ never being fixed under the flow of the drift. To see it, consider a small time increment $`\delta t`$. In the first order approximation, one can write
$$\mathit{\varrho }(t+\delta t)=e^{\delta t𝑨}\mathit{\varrho }(t)=\left(I+\delta t𝑨\right)\mathit{\varrho }(t)$$
i.e., the increment at $`\delta t`$ is given by (14) and the claim follows from the fact that $`H_A`$ is strongly regular, i.e., $`_j_{\mathrm{}}_p_q`$, $`(j\mathrm{})(pq)`$ $`j\mathrm{}`$, $`pq`$.∎
Since $`\mathrm{dim}(𝑨)=n>N`$, the stationary Liouville equation has more eigenvalues than those referable to the eigenvalues of the corresponding Schrödinger equation. From $`𝑨=i\mathrm{ad}_{H_A}`$, these are the roots of the Lie algebra $`𝔰𝔲(N)`$ computed at the element $`H_A`$ of the Cartan subalgebra $`𝔥`$, and, from (14), they correspond to the transition frequencies $`_j_{\mathrm{}}`$. The regularity of $`H_A`$ guarantees that these extra eigenvalues of $`𝑨`$ are all nonzero: $`dim\left(\mathrm{ker}(𝑨)\right)=N1=dim(𝔥)`$, thus providing an alternative proof of the first part of Proposition 1.
For $`\mathrm{tr}(\rho ^2)=1`$, the mapping $`|\psi |\psi \psi |`$ eliminates the ambiguity in the (unobservable) global phase: $`|\psi \psi |=e^{i\phi }|\psi \psi |e^{i\phi }`$ $`\phi `$ and, as before, the same property holds also for mixed states. Proposition 1 affirms that, consequently, the corresponding one-parameter orbit passing through each eigenstate $`|\psi _j`$ (due to the global phase) collapses into a fixed point of the unforced Liouville equation. Rephrasing in terms of density operators (part of) Proposition 1, we have the following.
###### Corollary 1
Any diagonal density operator is a fixed point of (2) when $`u=0`$. More generally, for any density operator both the diagonal part and the trace square norm of the off-diagonal part are integrals of motion of (2) when $`u=0`$.
## III Structure of the state space: complex flag manifolds
It is possible to give a more thorough interpretation of Proposition 1 by studying the structure of the manifold in which $`\rho `$ is living, call it $`𝒮`$. $`𝒮`$ is a connected, simply connected submanifold of $`𝕊^{n1}`$ (the $`(n1)`$-dimensional sphere of radius $`\mathit{\varrho }`$) whose dimension depends on the multiplicities of the eigenvalues of $`\rho `$. For $`N>2`$, $`𝒮𝕊^{n1}`$, since the Lie group $`\mathrm{exp}\left(\mathrm{ad}_{𝔰𝔲(N)}\right)`$ is not acting transitively on the entire $`𝕊^{n1}`$. $`𝒮`$, instead, is a homogeneous space of $`\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})`$, the action being left matrix multiplication, and can be described as a (co)adjoint orbit of $`SU(N)`$ on its Lie algebra as follows. Consider a diagonal density $`\stackrel{~}{\mathit{\varrho }}_𝔥𝒮`$. Call $`C_{\stackrel{~}{\mathit{\varrho }}_𝔥}`$ the stabilizer of $`\stackrel{~}{\mathit{\varrho }}_𝔥`$, $`C_{\stackrel{~}{\mathit{\varrho }}_𝔥}=\{g\mathrm{exp}\left(\mathrm{ad}_{𝔰𝔲(N)}\right)\text{ s.t. }g\stackrel{~}{\mathit{\varrho }}_𝔥=\stackrel{~}{\mathit{\varrho }}_𝔥\}`$. Because of the identification (up to the imaginary unit) of the density operators with (a convex set in) $`𝔰𝔲(N)`$, the coset space $`\mathrm{exp}\left(\mathrm{ad}_{𝔰𝔲(N)}\right)/C_{\stackrel{~}{\mathit{\varrho }}_𝔥}`$ is the adjoint orbit of $`SU(N)`$ on its Lie algebra passing through $`\stackrel{~}{\mathit{\varrho }}_𝔥`$. Because of transitivity, this orbit can be identified with $`𝒮`$: $`𝒮=\mathrm{exp}\left(\mathrm{ad}_{𝔰𝔲(N)}\right)/C_{\stackrel{~}{\mathit{\varrho }}_𝔥}\stackrel{~}{\mathit{\varrho }}_𝔥`$. The $`dim(𝒮)`$ is always even (each (co)adjoint orbit has a symplectic structure as is well-known). The orbit $`𝒮`$ is transverse to $`𝔥`$ and meets $`𝔥`$ in a number of disjoint points equal to the number of distinct permutations of the entries of $`\stackrel{~}{\mathit{\varrho }}_𝔥`$. Such number is equal to the cardinality of the Weyl group as well as to the Euler characteristic $`\chi (𝒮)`$ of the orbit, see Ercolessi1 ; Zyczkowski2 and Theorem E.2 of Frankel1 . These points form the vertices of a polygon in the $`N1`$-dimensional eigenensemble sitting in $`𝔥`$ and are sometimes denoted Weyl chambers. Inspired by the $`𝕊^2`$ case (see Example 1 below), we shall call them antipodal. If $`\stackrel{~}{\rho }_𝔥=\mathrm{diag}(\stackrel{~}{w}_1,\mathrm{},\stackrel{~}{w}_N)`$, $`_{j=1}^N\stackrel{~}{w}_j=1`$, $`0\stackrel{~}{w}_j1`$, then the $`\chi (𝒮)1`$ antipodal points are given by $`\mathrm{diag}(\stackrel{~}{w}_{\sigma (1)},\mathrm{},\stackrel{~}{w}_{\sigma (N)})`$ with $`\sigma (1),\mathrm{},\sigma (N)`$ a permutation of $`1,\mathrm{},N`$ such that $`\mathrm{diag}(\stackrel{~}{w}_{\sigma (1)},\mathrm{},\stackrel{~}{w}_{\sigma (N)})\stackrel{~}{\rho }_𝔥`$. While the topology of the diagonal coset representatives is particularly easy to visualize, the entire orbit enjoys the same topological structure of $`\stackrel{~}{\mathit{\varrho }}_𝔥`$. To see it, simply notice that applying a rotation in $`\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})`$ to two or more diagonal antipodal points they remain antipodal. Since $`\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})`$ acts transitively on $`𝒮`$, this is true on the entire orbit: each $`\mathit{\varrho }𝒮`$ has $`\chi (𝒮)1`$ antipodal states in $`𝒮`$. What is not known a priori is the isotropy subgroup, which depends on $`𝒥`$. In fact, the convex set of $`N`$-level density operators is foliated into leaves of different dimensions, depending on the number of distinct $`w_j`$ and on their multiplicities. For example, for a pure/pseudopure state $`𝒥=\{w_1,w_2,\mathrm{},w_2\}`$, $`w_1w_2`$, $`w_1+(N1)w_2=1`$, $`𝒮=\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})/\left(𝕊^1\times \mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N1)})\right)`$ and $`dim(𝒮)=2N2`$. In the pure state case $`𝒥=\{1,0,\mathrm{},0\}`$, it is well-known that the map $`|\psi |\psi \psi |`$ can be seen as a Hopf fibration: $`𝕊^{2N1}\stackrel{𝕊^1}{}𝒮=P^{N1}`$, with fibers representing the global phase. At the other extreme, if $`𝒥=\{w_1,w_2,\mathrm{},w_N\}`$, $`w_jw_{\mathrm{}}`$, $`_{j=1}^Nw_j=1`$, then $`𝒮=\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})/\left(𝕊^1\right)^N`$ of dimension $`N^2N`$. Hence if $`m=dim(𝒮)`$, $`2N2mN^2N`$, $`m`$ even. In between lies the flag manifolds with flag determined by the multiplicities of $`w_j`$. If such multiplicities are given by $`j_1,\mathrm{},j_{\mathrm{}}`$, $`j_1+\mathrm{}+j_{\mathrm{}}=N`$, $`2\mathrm{}N`$,
$$𝒮=\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})/\left(\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(j_1)})\times \mathrm{}\times \mathrm{exp}(\mathrm{ad}_{𝔰𝔲(j_{\mathrm{}})})\times \left(𝕊^1\right)^\mathrm{}1\right).$$
When $`\mathrm{}=2`$ we have Grassmannian manifolds. Normally, in the literature these are known as complex flag manifolds and are given directly in terms of unitary group actions, see Adelman1 ; Boya1 ; Zyczkowski2 :
$`𝒮`$ $`=`$ $`U(N)/\left(U(j_1)\times \mathrm{}\times U(j_{\mathrm{}})\right),j_1+\mathrm{}+j_{\mathrm{}}=N,2\mathrm{}N.`$
$`=`$ $`SU(N)/\left(SU(j_1)\times \mathrm{}\times SU(j_{\mathrm{}})\times \left(𝕊^1\right)^\mathrm{}1\right)`$
In terms of unitary actions, the two extreme cases of pure states and all different eigenvalues are, respectively, $`𝒮=U(N)/\left(U(N1)\times U(1)\right)=SU(N)/\left(SU(N1)\times 𝕊^1\right)`$ and $`𝒮=U(N)/\left(U(1)\right)^N=SU(N)/\left(𝕊^1\right)^{N1}`$. The description adopted here is just an isomorphic real representation of such complex flag manifolds deriving from the use of the adjoint representation.
###### Example 1
$`N=2`$, $`𝒥=\{1,0\}`$. The case $`N=2`$ is the only easy one, as $`𝒮=𝕊^2P^1`$. On the great horizontal circle of $`𝕊^2`$, $`\mathit{\varrho }=\mathit{\varrho }_𝔨`$. In terms of the Bloch vector, the diagonal antipodal states become the north and south poles of the Bloch sphere,
$`\rho _{v_1}=|00|=\mathrm{diag}(1,\mathrm{\hspace{0.17em}0})`$ $``$ $`\mathit{\varrho }_{v_1}=\left[\begin{array}{ccc}0& 0& \frac{1}{\sqrt{2}}\end{array}\right]^T`$
$`\rho _{v_2}=|11|=\mathrm{diag}(0,\mathrm{\hspace{0.17em}1})`$ $``$ $`\mathit{\varrho }_{v_2}=\left[\begin{array}{ccc}0& 0& \frac{1}{\sqrt{2}}\end{array}\right]^T`$
and $`𝔥`$, $`\mathrm{dim}(𝔥)=1`$, corresponds to the vertical line passing through $`\mathit{\varrho }_{v_1}`$, $`\mathit{\varrho }_{v_2}`$. Everything extends unchanged to mixed states, since $`𝒮`$ is still equal to $`𝕊^2`$ regardless of the purity. Since each $`𝒮`$ crosses $`𝔥`$ exactly twice, $`\chi (𝒮)=2`$. For any $`\mathit{\varrho }𝒮`$ the antipodal state is $`\mathit{\varrho }`$.
###### Example 2
$`N=3`$, $`𝒥=\{1,0,0\}`$. Since the isotropy subgroup in this case is $`SO(3)\times 𝕊^1`$ of dimension 4 (recall that $`dim(\mathrm{ad}_{𝔰𝔲(3)})=8`$), $`dim(𝒮)=4`$ and $`\chi (𝒮)=3`$. Following the standard ordering convention, the 3-level Gell-Mann basis (see e.g. Georgi1 , p. 99) is
$$\{\lambda _{𝔨,\mathrm{},12},\lambda _{𝔨,\mathrm{},12},\lambda _{𝔥,1},\lambda _{𝔨,\mathrm{},13},\lambda _{𝔨,\mathrm{},13},\lambda _{𝔨,\mathrm{},23},\lambda _{𝔨,\mathrm{},23},\lambda _{𝔥,2}\}.$$
The three diagonal antipodal states are
$`\rho _{v_1}=\mathrm{diag}(1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0})`$ $``$ $`\mathit{\varrho }_{v_1}=\left[\begin{array}{cccccccc}0& 0& \frac{1}{\sqrt{2}}& 0& 0& 0& 0& \frac{1}{\sqrt{6}}\end{array}\right]^T`$
$`\rho _{v_2}=\mathrm{diag}(0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}0})`$ $``$ $`\mathit{\varrho }_{v_2}=\left[\begin{array}{cccccccc}0& 0& \frac{1}{\sqrt{2}}& 0& 0& 0& 0& \frac{1}{\sqrt{6}}\end{array}\right]^T`$
$`\rho _{v_3}=\mathrm{diag}(0,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1})`$ $``$ $`\mathit{\varrho }_{v_3}=\left[\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& \frac{2}{\sqrt{6}}\end{array}\right]^T.`$
The structure of $`𝒮𝕊^7`$ is studied in detail in Ercolessi1 ; Chaturvedi1 ; Kimura1 ; Schirmer7 . In terms of the coherence vector, one has that each state of a triplet of antipodal states is at an angle of $`\frac{2\pi }{3}`$ from the other two, see Fig. 1. In particular, if $`\mathit{\varrho }𝒮`$, then $`\mathit{\varrho }𝒮`$: the single antipodal point of the case $`N=2`$ is replaced by two symmetrically distributed and equidistant antipodal points. As expected, only $`\lambda _{𝔥,1}`$ and $`\lambda _{𝔥,2}`$ are of concern when $`\rho `$ is diagonal. In an attempt to visualize the entire $`𝒮`$ manifold, one should replace each edge connecting to vertices with a sphere $`𝕊^2`$. The same numbers occur for pseudopure states, which have the same twofold degeneracy and are obtained by rescaling down the coherence vector $`\mathit{\varrho }_{v_j}`$ by a constant factor. For the all different eigenvalue case $`\rho _𝔥=\mathrm{diag}(w_1,w_2,w_3)`$, $`w_jw_k`$, which is the generic case, the stabilizer is the torus $`𝕊^1\times 𝕊^1`$, $`𝒮=\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})/(𝕊^1\times 𝕊^1)`$ and $`dim(𝒮)=6`$. The only diagonal matrices that are conjugate with $`\rho _𝔥`$ are its five element permutations, i.e., $`\chi (𝒮)=6`$ in this case. The six vertices are given by
$$\begin{array}{cc}\hfill \mathrm{diag}(a,b,c)& \mathit{\varrho }_{g_1}=\left[\begin{array}{cccccccc}0& 0& \frac{ab}{\sqrt{2}}& 0& 0& 0& 0& \frac{a+b2c}{\sqrt{6}}\end{array}\right]^T\hfill \\ \hfill \mathrm{diag}(b,a,c)& \mathit{\varrho }_{g_2}=\left[\begin{array}{cccccccc}0& 0& \frac{ba}{\sqrt{2}}& 0& 0& 0& 0& \frac{a+b2c}{\sqrt{6}}\end{array}\right]^T\hfill \\ \hfill \mathrm{diag}(c,a,b)& \mathit{\varrho }_{g_3}=\left[\begin{array}{cccccccc}0& 0& \frac{ca}{\sqrt{2}}& 0& 0& 0& 0& \frac{a+c2b}{\sqrt{6}}\end{array}\right]^T\hfill \\ \hfill \mathrm{diag}(c,b,a)& \mathit{\varrho }_{g_4}=\left[\begin{array}{cccccccc}0& 0& \frac{cb}{\sqrt{2}}& 0& 0& 0& 0& \frac{b+c2a}{\sqrt{6}}\end{array}\right]^T\hfill \\ \hfill \mathrm{diag}(b,c,a)& \mathit{\varrho }_{g_5}=\left[\begin{array}{cccccccc}0& 0& \frac{bc}{\sqrt{2}}& 0& 0& 0& 0& \frac{b+c2a}{\sqrt{6}}\end{array}\right]^T\hfill \\ \hfill \mathrm{diag}(a,c,b)& \mathit{\varrho }_{g_6}=\left[\begin{array}{cccccccc}0& 0& \frac{ac}{\sqrt{2}}& 0& 0& 0& 0& \frac{a+c2b}{\sqrt{6}}\end{array}\right]^T.\hfill \end{array}$$
(15)
The 6 vertices correspond to $`\frac{2\pi }{3}`$ rotations, plus their reflections around the three axes bisecting the triangle of pure states, see Fig. 1. Represented in the $`(\lambda _{𝔥,1},\lambda _{𝔥,2})`$ plane, the eigenensemble of $`\rho `$ has the shape of a hexagon inscribed in the triangle having as vertices the pure states, see Fig. 1.
## IV Feedback stabilization for $`N`$-level quantum ensembles
### IV.1 Problem formulation
For the system (10), we are interested in the problem of tracking a periodic orbit. More precisely, the stabilization problem is the following.
> Given $`\rho _d𝒮`$, find $`u=u(\mathbf{\varrho }_d,\mathbf{\varrho })`$ such that, for $`t\mathrm{}`$, $`\mathbf{\varrho }\mathbf{\varrho }_d`$, where
>
> $$\dot{\mathit{\varrho }}_d=𝑨_d\mathit{\varrho }_d,𝑨_d=𝒉_{A_d}\mathrm{ad}_{𝝀_𝔥}.$$
> (16)
This is a full state tracking problem which, from Proposition 1, reduces to stabilization to an equilibrium point when $`\rho _d=\varrho _0\lambda _0+\mathit{\varrho }_{d,𝔥}𝝀_𝔥`$.
### IV.2 A modified Jurjevic-Quinn condition and antipodal points
The algorithm for the feedback design resembles the one used for $`|\psi `$ discussed in Ferrante1 ; Vettori1 ; Grivopoulos1 ; Mirrahimi1 ; Mirrahimi2 and indeed the standard Jurdjevic-Quinn method for bilinear systems Jurdjevic5 . It consists in choosing a distance-like candidate Lyapunov function $`V=V(\mathit{\varrho }_d,\mathit{\varrho })`$. From (9), consider <sup>1</sup><sup>1</sup>1Notice that rather than the “distance inherited from the $`𝕊^{n1}`$ sphere” $`V`$ used here, one could use the $`^n`$ distance $`\delta \mathit{\varrho }^2=\mathit{\varrho }_d\mathit{\varrho }^2`$ as candidate Lyapunov function. Up to a scalar factor, the two give the same gradient, hence the same control design. $`V=\mathit{\varrho }^2\mathit{\varrho }_d,\mathit{\varrho }`$. If $`\mathit{\varrho }_d`$ obeys to (16),
$$\begin{array}{cc}\hfill \dot{V}=\dot{V}(\mathit{\varrho }_d,\mathit{\varrho })& =𝑨_d\mathit{\varrho }_d,\mathit{\varrho }\mathit{\varrho }_d,𝑨\mathit{\varrho }u\mathit{\varrho }_d,𝑩\mathit{\varrho }\hfill \\ & =\mathit{\varrho }_d,(𝑨_d𝑨)\mathit{\varrho }u\mathit{\varrho }_d,𝑩\mathit{\varrho }.\hfill \end{array}$$
(17)
This expression has a drift term which is in general sign indefinite and for which there is no global smooth feedback compensation. Whenever $`\dot{V}`$ can be rendered homogeneous in $`u`$ (namely when $`H_{A_d}=H_A`$, which is assumed thereafter), there is an obvious choice of feedback guaranteeing positive semidefiniteness. For example,
$$u=\mathit{\varrho }_d,𝑩\mathit{\varrho }$$
(18)
makes $`\dot{V}=u^20`$. In order to apply LaSalle invariance principle, one can try to adapt the “$`\mathrm{ad}`$-condition” of Jurdjevic-Quinn Jurdjevic5 to the case at hand. The largest invariant set $``$ in $`𝒩=\{\mathit{\varrho }\text{ s.t. }\dot{V}=0\}`$ can be computed imposing $`u=\frac{du}{dt}=\mathrm{}=\frac{d^{\mathrm{}}u}{dt^{\mathrm{}}}=0`$, where, in correspondence of $`u=0`$:
$$\frac{du}{dt}=\mathit{\varrho }_d^T[𝑨,𝑩]\mathit{\varrho }=0$$
and, similarly,
$$\frac{d^{\mathrm{}}u}{dt^{\mathrm{}}}=(1)^{\mathrm{}}\mathit{\varrho }_d^T\underset{\text{ }\mathrm{}\text{ times}}{\underset{}{[𝑨,\mathrm{},[𝑨,}}𝑩]\mathrm{}]\mathit{\varrho }=0.$$
(19)
For $`|\psi `$ which is an eigenfunction, this condition, used implicitly in Vettori1 , is made explicit in Mirrahimi1 as a Kalman rank condition on the linearized tangent system. It can be reformulated for densities as follows. For a given $`\mathit{\varrho }_d𝒮`$, write $`\mathit{\varrho }=\mathit{\varrho }_d+\delta \mathit{\varrho }`$. The linearization of (10) around $`\mathit{\varrho }_d`$
$$\frac{d\delta \mathit{\varrho }}{dt}=𝑨\delta \mathit{\varrho }+𝒃u$$
(20)
where
$$𝒃=𝑩\mathit{\varrho }_d,$$
(21)
is living on $`T_{\mathit{\varrho }_d}𝒮`$. Since $`dim(T_{\mathit{\varrho }_d}𝒮)=m`$, if the Kalman rank condition
$$\mathrm{rank}\left[𝒃𝑨𝒃\mathrm{}𝑨^{m1}𝒃\right]=m$$
(22)
is satisfied, then in the same spirit of the original Jurdjevic-Quinn work, this implies that $``$ contains no other trajectory than $`\mathit{\varrho }_d`$, at least locally. Of course if $`\rho _d=\varrho _0\lambda _0+\mathit{\varrho }_{d,𝔥}𝝀`$, then under (16) $`\rho _d(t)=\rho _d(0)`$ and $`𝒃(t)=𝒃(0)`$ in (21). However, when $`\rho _d(0)`$ is not diagonal, $`\rho _d(t)`$ is time-varying and so is $`𝒃(t)`$, complicating the verification of (22) considerably. Furthermore, in this case the condition (22) does not have a global character because of the topological structure of $`𝒮`$. To see this, consider the case of diagonal density operators in which the linearization (20) is time-invariant. Call $`\mathit{\varrho }_p`$ an antipodal state of $`\mathit{\varrho }_d`$. Then also $`\mathit{\varrho }_p`$ is diagonal and so is $`\delta \mathit{\varrho }_p=\mathit{\varrho }_p\mathit{\varrho }_d`$. Hence $`𝑨\delta \mathit{\varrho }_p=0`$, i.e., (20) has vanishing drift in correspondence of the $`\chi (𝒮)1`$ antipodal states of $`\mathit{\varrho }_d`$. Therefore, in spite of the Kalman condition (22), in this case when checking LaSalle invariance principle $`\frac{d\delta \mathit{\varrho }}{dt}|_{𝒩,\mathit{\varrho }=\mathit{\varrho }_p}=0`$ for (20) since $`u=0`$ in $`𝒩`$. This argument can be generalized as follows.
###### Proposition 2
Given $`\mathbf{\varrho }_d𝒮`$, any of the other $`\chi (𝒮)1`$ antipodal states $`\mathbf{\varrho }_p𝒮`$ is an equilibrium point for the system (10) with the feedback (18).
###### Proof.
We make use of the isomorphism (30). For any given $`\rho _p`$, $`\rho _d`$ diagonal, $`[iH_B,\rho _p]𝔨/i`$ is off-diagonal i.e., $`𝑩\mathit{\varrho }_p=\mathit{\varrho }_{p,𝔨}^{}`$. But if $`\rho _d`$ diagonal $`\rho _d=\varrho _0\lambda _0+\mathit{\varrho }_{𝔥_d}𝝀`$ then $`u=\mathit{\varrho }_d,𝑩\mathit{\varrho }_p=\mathit{\varrho }_d,\mathit{\varrho }_{p,𝔨}^{}=0`$, since $`𝔥𝔨`$. Hence no feedback is produced. When instead $`\mathit{\varrho }_p,\mathit{\varrho }_d𝒮`$ are antipodal but not diagonal, then (by construction) $``$ $`𝑹\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})`$ such that $`\stackrel{~}{\mathit{\varrho }}_p=𝑹\mathit{\varrho }_p`$ and $`\stackrel{~}{\mathit{\varrho }}_d=𝑹\mathit{\varrho }_d`$ are both diagonal. The skew symmetric matrix $`𝑹^T(i\mathrm{ad}_{H_B})𝑹`$ belongs to the adjoint orbit of $`\mathrm{exp}(\mathrm{ad}_{𝔰𝔲(N)})`$ in $`\mathrm{ad}_{𝔰𝔲(N)}`$, hence $``$ $`C𝔰𝔲(N)`$ such that $`𝑹^T𝑩𝑹=\mathrm{ad}_C`$. Therefore $`u=\mathit{\varrho }_d,𝑩\mathit{\varrho }_p=\stackrel{~}{\mathit{\varrho }}_d,𝑹^T𝑩𝑹\stackrel{~}{\mathit{\varrho }}_p=\stackrel{~}{\mathit{\varrho }}_d,\mathrm{ad}_C\stackrel{~}{\mathit{\varrho }}_p=0`$ because $`\mathrm{ad}_C\stackrel{~}{\mathit{\varrho }}_p[C,\stackrel{~}{\rho }_p]𝔨/i`$, while $`\stackrel{~}{\rho }_d𝔥`$. ∎
###### Corollary 2
For pure or pseudopure states, the $`N`$ antipodal points of $`𝒮`$ are all equidistant. For pure states they are also maximally distant.
###### Proof.
It is enough to notice that for any triple of antipodal points $`\rho _{p_1}`$, $`\rho _{p_2}`$ and $`\rho _{p_3}`$ ($`\rho _d`$ included), $`\mathrm{tr}\left(\rho _{p_1}\rho _{p_2}\right)=\mathrm{tr}\left(\rho _{p_1}\rho _{p_3}\right)=\mathrm{tr}\left(\rho _{p_2}\rho _{p_3}\right)`$, hence $`V(\mathit{\varrho }_{p_j},\mathit{\varrho }_p_{\mathrm{}})=\mathrm{const}`$ $`j,\mathrm{}=1,\mathrm{},N`$, $`j\mathrm{}`$. For pure states, in addition, $`\mathrm{tr}\left(\rho _{p_j}\rho _p_{\mathrm{}}\right)=0`$, hence $`V(\mathit{\varrho }_{p_j},\mathit{\varrho }_p_{\mathrm{}})=\mathrm{tr}(\rho _{p_j}^2)`$ are maximally distant. ∎
As will be shown in next Section, the antipodal points are not the only states lacking attractivity, and the linearization alone is not enough to investigate the domain of attraction of the feedback stabilizer.
###### Remark 1
The trajectory tracking problem presented above admits a reformulation as a point stabilization for a nonautonomous system. Consider a frame rotating with $`𝑨`$. Call $`\widehat{\mathit{\varrho }}_d`$ and $`\widehat{\mathit{\varrho }}`$ the new reference and state vectors. Then $`\widehat{\mathit{\varrho }}(t)=e^{t𝑨}\mathit{\varrho }(t)`$ and $`\widehat{\mathit{\varrho }}_d(t)=e^{t𝑨}\mathit{\varrho }_d(t)=\mathit{\varrho }_d(0)`$, i.e., the reference trajectory becomes a fixed point. Using a variation of constants formula, we obtain for (10)
$$\{\begin{array}{cc}\dot{\widehat{\mathit{\varrho }}}\hfill & =ue^{t𝑨}𝑩e^{t𝑨}\widehat{\mathit{\varrho }}\hfill \\ \widehat{\mathit{\varrho }}(0)\hfill & =\mathit{\varrho }(0).\hfill \end{array}$$
(23)
The Lyapunov distance is $`V=\widehat{\mathit{\varrho }}_d^2\widehat{\mathit{\varrho }}_d,\widehat{\mathit{\varrho }}`$ and its derivative $`\dot{V}=u\widehat{\mathit{\varrho }}_d,\dot{\widehat{\mathit{\varrho }}}`$. The uniformity of the asymptotic stability for the nonautonomous system (23) with the same feedback stabilizer as (18) follows directly.
### IV.3 Time-independent convergence conditions
To formulate a convergence condition in a more geometric manner, rewrite (19) in terms of bilinear forms of skew-symmetric matrices as follows. Call
$$𝔚^\alpha =\mathrm{span}\{𝑩,[𝑨,𝑩],\mathrm{},\underset{\text{ }\alpha \text{ times }}{\underset{}{[𝑨,\mathrm{},[𝑨,}}𝑩]\mathrm{}]\}$$
and $`𝔚_A^\alpha =\mathrm{span}\{𝑨,𝔚^\alpha \}`$. $`𝔚^\alpha `$ contains skew-symmetric matrices, and the conditions $`\frac{d^{\mathrm{}}u}{dt^{\mathrm{}}}`$, $`\mathrm{}=0,\mathrm{},\alpha `$, written compactly as $`\mathit{\varrho }_d^T𝔚^\alpha \mathit{\varrho }`$, are bilinear forms. If $`\mathit{\varrho }//\mathit{\varrho }_d`$ then $`\mathit{\varrho }_d^T𝔚^\alpha \mathit{\varrho }=0`$. On a sphere $`\mathit{\varrho }//\mathit{\varrho }_d`$ means $`\mathit{\varrho }=\pm \mathit{\varrho }_d`$. However $`𝒮`$ is only a submanifold of $`𝕊^{n1}`$, and $`\mathit{\varrho }_d`$ may not belong to it at all. $`𝔚^\alpha `$ is invariant to the so-called “$`\mathrm{ad}`$-brackets” but not necessarily a Lie subalgebra. Of course if for the system (10) the $`\mathrm{ad}`$-brackets are generating, i.e., if, for some $`\alpha `$, $`𝔚_A^\alpha =\mathrm{ad}_{𝔰𝔲(N)}`$, then almost global convergence on $`𝒮`$ is always verified. However we have the following negative result.
###### Lemma 1
If $`N3`$, for the system (10) with $`H_A`$ strongly regular and $`\mathrm{Graph}(H_B)`$ connected, $`\mathrm{Lie}(𝔚_A^\alpha )=\mathrm{ad}_{𝔰𝔲(N)}`$ but $`𝔚_A^\alpha \mathrm{ad}_{𝔰𝔲(N)}`$ $`\alpha >0`$.
###### Proof.
Since $`H_A`$ is strongly regular and $`\mathrm{Graph}(H_B)`$ is connected, it follows from Theorem 2 of Cla-contr-root1 that the smallest subalgebra containing $`iH_A`$, $`iH_B`$ is $`𝔰𝔲(N)`$. Hence the same holds for the adjoint representation. For the second part, recall that $`dim(𝔥)=N1`$. From the Lie bracket relations (7), $`𝑨\mathrm{ad}_𝔥`$, $`𝑩\mathrm{ad}_𝔨`$ implies $`[𝑨,\mathrm{},[𝑨,𝑩]\mathrm{}]\mathrm{ad}_𝔨`$. Even adding $`𝑨`$, $`𝔚_A^\alpha `$ alone cannot fully generate $`\mathrm{ad}_𝔥`$ for any $`\alpha `$. ∎
The first part of Lemma 1 is also known as the strong accessibility condition Nijmeijer1 . Since $`\mathrm{ad}_{𝔰𝔲(N)}`$ is compact, it suffices for controllability. The second part is the Jurdjevic-Quinn condition mentioned above. If $`H_B`$ is not off-diagonal, then the statement of Lemma 1 should be reformulated as “$`N>3`$”.
The following Theorem provides a time-independent condition for asymptotic stabilizability to any $`\rho _d𝒮`$, and a global description of the region of attraction of the controller.
###### Theorem 1
Consider the system (10) with the feedback (18), where $`\rho _d𝒮`$ obeys to (16). Assume that $`H_A`$ is strongly regular and that $`H_B`$ is such that $`(h_{B,\mathrm{},jj+1},h_{B,\mathrm{},jj+1})(0,\mathrm{\hspace{0.17em}0})`$. An initial condition $`\rho (0)𝒮`$ is asymptotically converging to $`\rho _d(t)`$ if
1. $`\rho (0)`$ is not an antipodal point of $`\rho _d(0)`$,
2. $`\left([H_B,\rho _d]\right)\left(\rho (0)\right)0`$,
3. $`\mathrm{Card}_𝔨\left([H_B,\rho _d]\right)m/2`$
where $`\mathrm{Card}`$ denotes the number of pairs of indexes in $`_𝔨`$.
In order to prove the Theorem we need a few preliminary results.
###### Lemma 2
Under the assumption of strong regularity of $`H_A`$, the following three conditions are equivalent:
1. the Kalman rank condition (22) is satisfied;
2. $`\mathrm{rank}\left(𝔚^{m1}\mathit{\varrho }_d\right)=m`$;
3. $`\mathrm{Card}_𝔨([H_B,\rho _d])m/2`$.
###### Proof.
Given $`C𝔰𝔲(N)`$, strong regularity of $`H_A`$ implies that $`C,[H_A,C],\mathrm{},[H_A,\mathrm{},[H_A,C]\mathrm{}]`$ are all linearly independent up to a number $`\alpha 1`$ , $`\alpha =2\mathrm{Card}_𝔨(C)`$, of nested $`H_A`$ commutators, see Theorem 2 in Cla-contr-root1 . Using $`C=𝒄𝝀`$ and the isomorphism given by the adjoint representation, the vectors $`𝒄,𝑨𝒄,\mathrm{}𝑨^{\alpha 1}𝒄`$ are all linearly independent. If $`𝒄=𝒃=𝑩\mathit{\varrho }_d`$ as in (21), then this is the Kalman controllability condition provided $`\alpha m`$. If $`𝒃=𝒃_𝔥+𝒃_𝔨`$, $`𝑨𝒃_𝔥=0`$, hence only the off-diagonal part of $`[H_B,\rho _d]`$ matters. The support $`𝔣_𝔨([H_B,\rho _d])`$ intersects a number of “root spaces” $`𝔨_j\mathrm{}`$ (each has real dimension 2) equal to $`\mathrm{Card}_𝔨([H_B,\rho _d])`$. Furthermore, since $`iH_A𝔰𝔲(N)`$, the invariance property (7) applies. Written in terms of the original commutators
$$_𝔨([H_B,\rho _d])=_𝔨([H_A,[H_B,\rho _d]])=\mathrm{}=_𝔨([H_A,\mathrm{},[H_A,[H_B,\rho _d]]\mathrm{}]).$$
(24)
For the $`\mathrm{}`$-th order commutator of $`𝔚^{\alpha 1}`$, one has the binomial-like expansion:
$$\begin{array}{cc}& \underset{\mathrm{}\text{ times}}{\underset{}{[𝑨,\mathrm{},[𝑨}},𝑩]\mathrm{}]\hfill \\ & =𝑨^{\mathrm{}}𝑩+(1)^1\left(\genfrac{}{}{0pt}{}{\mathrm{}}{1}\right)𝑨^\mathrm{}1𝑩𝑨+\mathrm{}\hfill \\ & +(1)^\mathrm{}1\left(\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}1}\right)𝑨𝑩𝑨^\mathrm{}1+(1)^{\mathrm{}}𝑩𝑨^{\mathrm{}}.\hfill \end{array}$$
(25)
After the linearization around $`\rho _d`$, only the first term of this expression is retained. Since $`iH_A𝔥`$, from (7), $`𝔣_𝔨\left(𝑨^{\mathrm{}}\mathit{\varrho }_d\right)=𝔣_𝔨(\mathit{\varrho }_d)`$, while $`𝔣_𝔥\left(𝑨^{\mathrm{}}\mathit{\varrho }_d\right)=0`$, $`\mathrm{}1`$. From (24), $`_𝔨`$ is the same for all terms in (25), and similarly,
$$_𝔨\left(𝑩\mathit{\varrho }_d\right)=_𝔨\left([𝑨,𝑩]\mathit{\varrho }_d\right)=\mathrm{}=_𝔨\left([𝑨,\mathrm{},[𝑨,𝑩]\mathrm{}]\mathit{\varrho }_d\right).$$
(26)
In summary, strong regularity of $`H_A`$ guarantees the full spanning of a linear space whose dimension is determined uniquely by $`\mathrm{Card}_𝔨([H_B,\rho _d])`$. This space is identifiable with the tangent space $`T_{\mathit{\varrho }_d}𝒮`$ as well as with $`𝔚^{m1}\mathit{\varrho }_d`$. The equivalence of the three conditions follows consequently. ∎
###### Remark 2
Lemma 1 and Condition 2 of Lemma 2 imply that although the vector space $`𝔚^{m1}`$ is never the entire Lie algebra $`\mathrm{ad}_{𝔰𝔲(N)}`$ acting transitively on $`𝒮`$, it may nevertheless span the entire tangent space at a point. The same holds for the Kalman controllability.
###### Remark 3
In general $`\mathrm{Card}_𝔨(H_B)\mathrm{Card}_𝔨([H_B,\rho _d])`$, hence the controllability of the linearization depends from the reference trajectory $`\rho _d`$ chosen. The meaning of Lemma 2 is that in order to have linear controllability the reference trajectory (16) must be “rich enough” along $`H_B`$.
###### Remark 4
While the Kalman condition (22) seems time-varying as soon as $`\rho _{d,𝔨}0`$, the equivalent condition 3 of Lemma 2 is always time-independent since $`\left([H_B,\rho _d(0)]\right)=\left([H_B,\rho _d(t)]\right)`$.
###### Remark 5
The conditions of Lemma 2 depend on $`\rho _d`$, $`H_A`$ and $`H_B`$ but not on the state $`\rho `$, meaning that alone they are not enough to guarantee convergence of a given $`\rho (0)`$.
The Lyapunov derivative in (17) is made homogeneous in $`u`$ by the cancellation of the drift term and therefore the notion of attractivity provided by $`\dot{V}`$ must be rendered invariant under such flow (in a way similarly to the orbital stabilization problem, see Bacciotti1 ). The following Lemma gives an alternative attractivity condition which is fully invariant under the drift and generically (i.e., almost always under $`e^{t𝑨}`$) equivalent to the usual Lyapunov convergence property. This last in fact may fail in isolated points: certain critical points of $`V`$ are not invariant under the flow of the drift (see Section V).
###### Lemma 3
Consider the system (10) with the feedback (18), where $`\rho _d`$ obeys to (16). If $`H_A`$ is strongly regular, the following conditions are generically equivalent under the flow of the drift term:
1. $`\left([H_B,\rho _d]\right)\left(\rho \right)0`$;
2. $`\dot{V}(\mathit{\varrho }_d,\mathit{\varrho })<0`$;
3. $`\mathit{\varrho }_d^T𝔚^\alpha \mathit{\varrho }=\{z_0,z_1,\mathrm{}z_\alpha \}`$, $`z_j0`$.
###### Proof.
Clearly $`\dot{V}=\mathit{\varrho }_d,𝑩\mathit{\varrho }^2<0`$ implies $`\left([H_B,\rho _d]\right)\left(\rho \right)0`$. To prove that also the contrary is generically true, it is enough to show that when $`\left([H_B,\rho _d]\right)\left(\rho \right)0`$ the zero crossing of the inner product can occur only at isolated points along the trajectories of the closed loop system. Assume Item 1 holds and, at time $`t`$, $`\mathit{\varrho }_d,𝑩\mathit{\varrho }=0`$. If $`\delta t`$ is a small time increment, then from Item 1 of Proposition 1, $`\left([H_B,\rho _d]\right)`$ and $`\left(\rho \right)`$ remain the same, while, from Item 3 of Proposition 1
$`(\mathit{\varrho }_{𝔨,\mathrm{},j\mathrm{}}(t+\delta t),\mathit{\varrho }_{𝔨,\mathrm{},j\mathrm{}}(t+\delta t))`$ $``$ $`(\mathit{\varrho }_{𝔨,\mathrm{},j\mathrm{}}(t),\mathit{\varrho }_{𝔨,\mathrm{},j\mathrm{}}(t))`$
$`\mathit{\varrho }_{𝔥,j}(t+\delta t)`$ $`=`$ $`\mathit{\varrho }_{𝔥,j}(t)`$
$`(\mathit{\varrho }_{d,𝔨,\mathrm{},j\mathrm{}}(t+\delta t),\mathit{\varrho }_{d,𝔨,\mathrm{},j\mathrm{}}(t+\delta t))`$ $``$ $`(\mathit{\varrho }_{d,𝔨,\mathrm{},j\mathrm{}}(t),\mathit{\varrho }_{d,𝔨,\mathrm{},j\mathrm{}}(t))`$
$`\mathit{\varrho }_{d,𝔥,j}(t+\delta t)`$ $`=`$ $`\mathit{\varrho }_{d,𝔥,j}(t).`$
If $`_𝔥\left([H_B,\rho _d]\right)_𝔥\left(\rho \right)0`$, then from the last row of (8) only $`\mathit{\varrho }_{d,𝔨}`$ matters in the computation of $`_𝔥\left([H_B,\rho _d]\right)`$, and $`\mathit{\varrho }_{d,𝔨}(t+\delta t),𝑩\mathit{\varrho }_𝔥(t+\delta t)0`$ since $`\mathit{\varrho }_{d,𝔨}(t+\delta t)\mathit{\varrho }_{d,𝔨}(t)`$, while $`\mathit{\varrho }_𝔥(t+\delta t)=\mathit{\varrho }_𝔥(t)`$. If, instead, $`_𝔨\left([H_B,\rho _d]\right)_𝔨\left(\rho \right)0`$, then we have two possible contributions to consider: $`_𝔨\left([H_B,\rho _{d,𝔥}]\right)`$ and $`_𝔨\left([H_B,\rho _{d,𝔨}]\right)`$. In the first case the conclusion follows from the same argument used above since now $`\mathit{\varrho }_{d,𝔥}(t+\delta t)=\mathit{\varrho }_{d,𝔥}(t)`$ while $`\mathit{\varrho }_𝔨(t+\delta t)\mathit{\varrho }_𝔨(t)`$. In the second case it follows from the observation that $`_𝔨\left([H_B,\rho _{d,𝔨}]\right)_𝔨\left(\rho \right)0`$ implies $`_𝔨(\rho _{d,𝔨})_𝔨\left(\rho \right)`$ (see the explicit computations of the commutators in (12)). The general case $`\left([H_B,\rho _d]\right)\left(\rho \right)0`$ is the sum of the two situations just described. Concerning Item 3, it is enough to notice that generically $`\mathit{\varrho }_d,𝑩\mathit{\varrho }0`$ if and only if $`\mathit{\varrho }_d,[𝑨,\mathrm{},[𝑨,𝑩]\mathrm{}]\mathit{\varrho }0`$. The argument is of the same type used in the proof of Lemma 2. For example if $`_𝔨\left([H_B,\rho _{d,𝔨}]\right)_𝔨\left(\rho \right)0`$ then just apply (26). If, instead, we are in the case $`_𝔥\left([H_B,\rho _{d,𝔨}]\right)_𝔥\left(\rho \right)0`$, then the only useful term in the expansion (25) is the last one, but this is enough to prove the claim. The genericity of the argument can be shown as above. ∎
###### Proof.
(of Theorem 1) Consider the set $`𝒩`$. We want to show that the largest invariant set $``$ in $`𝒩`$ is given by $`\rho _d`$ only. Condition 2 guarantees that locally around $`\rho _d(t)`$ there is no other equilibrium point in $`𝒩`$, as, from Lemma 2, the linearization at $`\rho _d`$ is controllable. Hence $`\rho _d`$ is a locally asymptotically stable equilibrium for the closed loop system and $`\rho _d`$ is isolated in $`𝒩`$. Consider $`\rho _e𝒩`$, $`\rho _e\rho _d`$. This implies $`\rho _e`$ disjoint from $`\rho _d`$ and $`V(\mathit{\varrho }_d,\mathit{\varrho }_e)>0`$. We need to show that $`\rho _e`$ must be a repulsive equilibrium for the closed loop system <sup>2</sup><sup>2</sup>2Since $`\mathit{\varrho }_e`$ may not be isolated in $`𝒩`$, the term repulsive has to be intended as “semi-repulsive”.. For $`\mathit{\varrho }_e`$ which is an antipodal point of a pure state $`\mathit{\varrho }_d`$, this is follows from Corollary 2, since $`V(\mathit{\varrho }_d,\mathit{\varrho }_e)`$ is maximal in $`𝒮`$ while $`\dot{V}0`$. For any other $`\rho _e𝒩`$, it is enough to perturb $`\rho _e`$ to $`\stackrel{~}{\rho }_e𝒮`$ so that $`\left([H_B,\rho _d]\right)\left(\stackrel{~}{\rho }_e\right)0`$. It is always possible to do this in a neighborhood of $`\rho _e`$ since $`\left([H_B,\rho _d]\right)`$ has cardinality at least $`m/2`$ and $`(h_{B,\mathrm{},jj+1},h_{B,\mathrm{},jj+1})(0,\mathrm{\hspace{0.17em}0})`$ implies that $`\mathrm{Graph}(H_B)`$ is connected and that there is no subspace $`𝔨_j\mathrm{}`$ invariant under $`𝑩`$. But then, from Lemma 3, $`\mathit{\varrho }_d,𝑩\stackrel{~}{\mathit{\varrho }}_e0`$ and $`V(\mathit{\varrho }_d,\stackrel{~}{\mathit{\varrho }}_e)<V(\mathit{\varrho }_d,\mathit{\varrho }_e)`$, i.e., $`\stackrel{~}{\rho }_e`$ is attracted to $`\rho _d`$. To show that $`V(\mathit{\varrho }_e,\stackrel{~}{\mathit{\varrho }}_e)`$ increases, assume by contradiction that
$$\dot{V}(\mathit{\varrho }_e,\stackrel{~}{\mathit{\varrho }}_e)=\dot{\mathit{\varrho }}_e,\stackrel{~}{\mathit{\varrho }}_e+\mathit{\varrho }_e,\dot{\stackrel{~}{\mathit{\varrho }}}_e=\mathit{\varrho }_d,𝑩\stackrel{~}{\mathit{\varrho }}_e\mathit{\varrho }_e,𝑩\stackrel{~}{\mathit{\varrho }}_e<0.$$
(27)
Consider the geodesic line in $`𝒮`$ connecting $`\mathit{\varrho }_d`$ with $`\mathit{\varrho }_e`$: $`\mathit{\varrho }_\varphi (s)=\mathit{\varrho }_d+\mathit{\varphi }(s)`$ such that $`\mathit{\varphi }(0)=0`$ and $`\mathit{\varphi }(s_e)=\mathit{\varrho }_e\mathit{\varrho }_d`$. Along this line,
$$\dot{V}(\mathit{\varrho }_\varphi (s),\stackrel{~}{\mathit{\varrho }}_e)=\mathit{\varrho }_d,𝑩\stackrel{~}{\mathit{\varrho }}_e^2\mathit{\varrho }_d,𝑩\stackrel{~}{\mathit{\varrho }}_e\mathit{\varphi }(s),𝑩\stackrel{~}{\mathit{\varrho }}_e,s[0,s_e]$$
is a function linear in $`\mathit{\varphi }(s)`$ and such that, by the assumption (27),
$`\dot{V}(\mathit{\varrho }_\varphi (0),\stackrel{~}{\mathit{\varrho }}_e)`$ $`=`$ $`\dot{V}(\mathit{\varrho }_d,\stackrel{~}{\mathit{\varrho }}_e)=\mathit{\varrho }_d,𝑩\stackrel{~}{\mathit{\varrho }}_e^2<0`$
$`\dot{V}(\mathit{\varrho }_\varphi (s_e),\stackrel{~}{\mathit{\varrho }}_e)`$ $`=`$ $`\dot{V}(\mathit{\varrho }_e,\stackrel{~}{\mathit{\varrho }}_e)<0.`$
But then $`\dot{V}(\mathit{\varrho }_\varphi (s),\stackrel{~}{\mathit{\varrho }}_e)<0`$ $`s[0,s_e]`$, $`\dot{V}(\mathit{\varrho }_\varphi (s),\mathit{\varrho }_\varphi (s))=0`$, meaning that $`\stackrel{~}{\mathit{\varrho }}_e`$ is attracted to the entire geodesic segment $`\mathit{\varrho }_\varphi (s)`$, $`s[0,s_e]`$, which is a contradiction, since $`\mathit{\varrho }_d`$ is an isolated equilibrium point. Hence it must be $`\dot{V}(\mathit{\varrho }_e,\stackrel{~}{\mathit{\varrho }}_e)0`$ i.e., $`\mathit{\varrho }_e`$ is a repulsive equilibrium point. Therefore $`\mathit{\varrho }_e`$ cannot belong to $``$. From Lemma 3, all conditions (19) are satisfied or violated simultaneously respectively when $`\dot{V}=0`$ or $`\dot{V}<0`$, i.e., when $`\left([H_B,\rho _d]\right)\left(\stackrel{~}{\rho }_e\right)=0`$ or $`0`$. Hence outside $`𝒩`$ the Jurdjevic-Quinn condition applies and $`\rho (0)`$ must converge to $`\rho _d(t)`$ since any other $`\rho _e𝒩`$ is repulsive. ∎
###### Remark 6
Condition 2 of Theorem 1 is obviously a necessary condition for convergence. Condition 3 instead is sufficient but not necessary, see Example 1 in Section V.
While, from Lemma 2, the linear span at $`\mathit{\varrho }_d`$ of the linearized system and of the $`𝔚^\alpha `$ yield a space of the same dimension, Item 3 of Lemma 3 holds for the bilinear forms but it is in general not true for the linearization.
###### Corollary 3
For $`H_A`$ strongly regular:
1. $`\mathrm{rank}𝔚^{\alpha 1}\mathit{\varrho }_d=\mathrm{rank}\left[𝒃𝑨𝒃\mathrm{}𝑨^{\alpha 1}𝒃\right],\alpha =0,\mathrm{},m1`$;
2. $`\mathit{\varrho }_d,[\underset{\text{ }\alpha \text{ times }}{\underset{}{𝑨,\mathrm{}[𝑨}},𝑩]\mathrm{}]\mathit{\varrho }0\overline{)}\mathit{\varrho },𝑨^\alpha 𝒃0`$
###### Proof.
The first point follows from the strong regularity of $`H_A`$ and from (26), which implies that the maximum number of independent vectors in the two sequences above is the same for all $`\alpha `$. The second from (25) and $`𝑨^\alpha \mathit{\varrho }_𝔥=0`$. ∎
The consequence is that the linearization alone is inconclusive about the region of attraction of the reference trajectory in the closed loop system, while instead the $`\mathrm{ad}`$-commutators completely specify it.
###### Corollary 4
When $`\mathrm{Card}_𝔨\left([H_B,\rho _d]\right)m/2`$, the region of attraction of the system (10) with the feedback (18) is given by $`=𝒮𝒩`$.
### IV.4 Global stabilization and topological obstructions
The notions from differential topology used in this Section are recalled in Appendix B. A compact manifold like $`𝒮`$ cannot be globally asymptotically stabilized because it lacks the contractivity property, i.e., it is not homotopy equivalent to a point, see Bhat1 , Proposition 1 and Theorem 1, and Wilson1 . Proposition 2 suggests that for $`𝒮`$ this is due to the antipodal points.
###### Proposition 3
For $`𝒮`$, the $`\chi (𝒮)1`$ antipodal points are irremovable topological obstruction to global stabilizability by smooth feedback.
###### Proof.
Contractivity is a necessary condition for global asymptotic stabilizability. For example, that $`𝒮=𝕊^2`$ with a point removed is homeomorphic (and hence homotopy equivalent) to $`^2`$ is well-known through the stereographic projection (see e.g. Armstrong1 , p. 34). Since it is known that the domain of attraction of an asymptotically stable point must be homotopy equivalent to $`^m`$ for some $`m`$ Wilson1 , then for $`N=2`$ this is enough to affirm that convergence can be rendered global up to the antipodal point. For $`N>2`$, in order to show that the antipodal points are all obstructions to contractivity of $`𝒮`$, consider the equilibrium $`\rho _d`$ and one of its antipodal points $`\rho _{p_1}`$. By suitable change of basis, $`\rho _d`$ and $`\rho _{p_1}`$ can be rendered diagonal simultaneously. By the transversality of the coadjoint orbit on $`𝔥`$, it is possible to determine a submanifold of $`𝒮`$ connecting $`\rho _d`$ and $`\rho _{p_1}`$ and not passing through any other of the antipodal points. To see it, notice that for $`N>2`$ it is always possible to adjust the basis (3)-(5) so as to attain a 3-dimensional simple subalgebra of $`𝔰𝔲(N)`$, which, as described in Section III, draws an $`𝕊^2`$ orbit under the adjoint action (see Sanchez1 for the details of this construction). This is a well-defined compact submanifold of $`𝒮`$ and it is not contractible for what said above for $`N=2`$.
## V A few cases of physical interest
The methods developed above yield considerable insight into the stabilizability and convergence properties of a quantum density operator. A few interesting cases for $`N`$-level systems are now described. It is followed by a more detailed description for systems with $`N=2,\mathrm{\hspace{0.17em}3}`$.
* Since $`mN^2N`$, and $`\mathrm{Card}_𝔨\left([H_B,\rho _d]\right)(N^2N)/2`$, (i.e., the maximal number of off-diagonal terms), each complex flag manifold $`𝒮`$ may admit a controllable linearization (depending on $`\rho _d`$).
* The assumption of direct coupling between nearest energy levels $`(h_{B,\mathrm{},jj+1},h_{B,\mathrm{},jj+1})(0,\mathrm{\hspace{0.17em}0})`$, is needed in order to exclude the existence of subset of $`𝒮`$ which remains invariant under the closed loop dynamics. It is a common assumption in most practical cases (dipole approximation Dahleh1 ). See also Example 2 below (last item).
* The full connectivity of $`\mathrm{Graph}(H_B)`$ is neither a sufficient nor a necessary condition for asymptotic stability.
* If $`\rho _d`$ is an eigenstate and $`\rho `$ another eigenstate then there is never convergence, not even if $`\mathrm{Graph}(H_B)`$ is fully connected, because $`\rho `$ is antipodal to $`\rho _d`$.
* For pure states and not fully connected $`\mathrm{Graph}(H_B)`$, certain eigenstates are easier to stabilize than others. The easiest is the one of energy $`_j`$ such that the index $`j`$ appears more often in $`_𝔨\left([H_B,\rho _d]\right)`$. In Example 2 below with $`H_B`$ in (28), it is easier to stabilize to the eigenstate of intermediate energy than to the ground state or to the most excited state. When $`\mathrm{Graph}(H_B)`$ is fully connected, there is no such difference. From Theorem 1, this does not mean that all initial conditions have the same convergence properties to a given $`\rho _d`$.
* If $`\rho _d`$ and $`\rho (0)`$ are both block diagonal and the blocks do not overlap
$$\rho _d=\left[\begin{array}{cc}\begin{array}{ccc}& \mathrm{}& \\ \mathrm{}& & \mathrm{}\\ & \mathrm{}& \end{array}& \end{array}\right],\rho (0)=\left[\begin{array}{cc}& \\ & \begin{array}{cc}& \\ & \mathrm{}& \\ \mathrm{}& & \mathrm{}\\ & \mathrm{}& \end{array}\end{array}\right],$$
then
$$[H_B,\rho _d]=\left[\begin{array}{cc}\begin{array}{ccc}& \mathrm{}& \\ \mathrm{}& & \\ & \mathrm{}& \end{array}& \begin{array}{ccc}& \mathrm{}& \\ & & \mathrm{}\\ & \mathrm{}& \end{array}\\ \begin{array}{ccc}& \mathrm{}& \\ \mathrm{}& & \mathrm{}\\ & \mathrm{}& \end{array}& \begin{array}{cc}& \\ 0& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0\end{array}\end{array}\right],$$
which implies $`\left([H_B,\rho _d]\right)\left(\rho (0)\right)=0`$ and $`\dot{V}=0`$, i.e., $`\rho (0)`$ is not attracted to $`\rho _d`$. Since $`\mathrm{tr}\left(\rho \rho _d\right)=0`$, $`\rho _d`$ and $`\rho `$ are as distant as antipodal states.
* Not all states in $`𝒩`$ are maximally distant from $`\rho _d`$. Assume $`\rho _d,\rho `$ such that $`_𝔥\left([H_B,\rho _d]\right)_𝔥\left(\rho (0)\right)0`$, $`_𝔨\left(\rho _d\right)=0`$, $`_𝔨\left(H_B\right)_𝔨\left(\rho \right)=0`$. Also in this case $`\left([H_B,\rho _d]\right)\left(\rho (0)\right)=0`$ and $`\rho `$ is not converging. However, since $`\mathrm{tr}\left(\rho _d\rho \right)0`$, $`\rho _d`$ and $`\rho `$ are not maximally distant.
* A typical example of initial condition such that $`\dot{V}(\mathit{\varrho }_d(0),\mathit{\varrho }(0))=0`$ but not invariant under the drift (see paragraph before Lemma 3) is attained when $`_𝔨\left([H_B,\rho _d]\right)_𝔨\left(\rho (0)\right)0`$ but $`[H_B,\rho _d(0)]`$, $`\rho (0)`$ both real or purely imaginary. This follows from Proposition 1.
Example 1 (cont’d) Assume
$$H_A=\frac{h_{A,𝔥,1}}{\sqrt{2}}\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right]=h_{A,𝔥,1}\lambda _{𝔥,1}\text{ and }H_B=\frac{h_{B,𝔨,\mathrm{},12}}{\sqrt{2}}\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]=h_{B,𝔨,\mathrm{},12}\lambda _{𝔨,\mathrm{},12}$$
Then
$$𝑨=i\sqrt{2}h_{A,𝔥,1}\mathrm{ad}_{\lambda _{𝔥,1}}=2h_{A,𝔥,1}\left[\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right]\text{ and }𝑩=i\sqrt{2}h_{B,𝔨,\mathrm{},12}\mathrm{ad}_{\lambda _{𝔨,\mathrm{},12}}=2h_{B,𝔨,\mathrm{},12}\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right]$$
From Proposition 1, both $`\mathit{\varrho }_𝔨`$ and $`\mathit{\varrho }_𝔥`$ are integrals of motion of the unforced dynamics, while the two components of $`\mathit{\varrho }_𝔨`$ evolve according to a sinusoidal law. When applying Theorem 1 to the system plus the feedback (18), we have the following for the closed loop system:
* any $`\rho _d`$ has a single antipodal point which also is an equilibrium;
* if $`\rho _d`$ diagonal, $`_𝔨\left([H_B,\rho _d]\right)=\{(12)\}`$, the linearization is controllable and any nondiagonal $`\rho `$ satisfies Theorem 1. Hence any $`\rho (0)`$ such that $`\rho _𝔨(0)0`$ is attracted to $`\rho _d`$ diagonal;
* if $`\rho _d`$ off-diagonal, $`\mathrm{Card}_𝔨\left([H_B,\rho _d]\right)=0`$, and the sufficient condition of Theorem 1 does not apply. However, $`_𝔥\left([H_B,\rho _d]\right)0`$ and as long as $`_𝔥\left([H_B,\rho _d]\right)_𝔥\left(\rho (0)\right)0`$, i.e., whenever $`\rho _𝔥0`$, $`\rho \rho _d`$. This is a special situation due to $`\mathrm{dim}(𝔥)=1`$, and has no counterpart for $`N>2`$.
In summary, there is always almost global convergence except when $`\mathit{\varrho }_{d,𝔥}=\mathit{\varrho }_𝔥=0`$, i.e., except when both $`\mathit{\varrho }_d`$ and $`\mathit{\varrho }`$ belong to great horizontal circles. ∎
Example 2 (cont’d) The drift of the system is given by
$$H_A=\frac{h_{A,𝔥,1}}{\sqrt{2}}\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right]+\frac{h_{A,𝔥,2}}{\sqrt{6}}\left[\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right]=h_{A,𝔥,1}\lambda _{𝔥,1}+h_{A,𝔥,2}\lambda _{𝔥,2}.$$
We shall consider the following control vector field
$$H_B=\frac{1}{\sqrt{2}}\left[\begin{array}{ccc}0& h_{B,𝔨,\mathrm{},12}& 0\\ h_{B,𝔨,\mathrm{},12}& 0& h_{B,𝔨,\mathrm{},23}\\ 0& h_{B,𝔨,\mathrm{},23}& 0\end{array}\right]=h_{B,𝔨,\mathrm{},12}\lambda _{𝔨,\mathrm{},12}+h_{B,𝔨,\mathrm{},23}\lambda _{𝔨,\mathrm{},23},$$
(28)
which has $`_𝔨(H_B)=\{(12),(23)\}`$ or, alternatively,
$$H_B=\frac{1}{\sqrt{2}}\left[\begin{array}{ccc}0& h_{B,𝔨,\mathrm{},12}& h_{B,𝔨,\mathrm{},13}\\ h_{B,𝔨,\mathrm{},12}& 0& h_{B,𝔨,\mathrm{},23}\\ h_{B,𝔨,\mathrm{},13}& h_{B,𝔨,\mathrm{},23}& 0\end{array}\right]=h_{B,𝔨,\mathrm{},12}\lambda _{𝔨,\mathrm{},12}+h_{B,𝔨,\mathrm{},13}\lambda _{𝔨,\mathrm{},13}+h_{B,𝔨,\mathrm{},23}\lambda _{𝔨,\mathrm{},23}$$
(29)
which has a “fully connected” graph, $`_𝔨(H_B)=\{(12),(13),(23)\}`$.
A list of interesting cases is the following:
* any of the (two for pure/psudopure, five for the generic case) antipodal points of any $`\rho _d`$ is also an equilibrium.
* $`\rho _d`$ diagonal: only the off-diagonal part of $`\rho `$ matters
+ $`\rho _d`$ pure (or pseudopure), e.g. $`\rho _d=\rho _{v_1}`$
- $`H_B`$ given in (28): $`_𝔨\left([H_B,\rho _d]\right)=\{(12)\}`$ $``$ the linearization is never controllable since $`2\mathrm{Card}_𝔨\left([H_B,\rho _d]\right)<4=m`$, hence Theorem 1 does not apply. Unlike the $`N=2`$ case, now in general $`\rho (0)\overline{)}\rho _d`$;
- $`H_B`$ given in (29): $`_𝔨\left([H_B,\rho _d]\right)=\{(12),(13)\}`$ $``$ the linearization is controllable. Any $`\rho (0)`$ such that $`_𝔨\left(\rho (0)\right)\{(12),(13)\}0`$ is converging. However, if one considers the pure state
$$\rho (0)=\frac{1}{2}\left[\begin{array}{ccc}0& 0& 0\\ 0& 1& 1\\ 0& 1& 1\end{array}\right],$$
then $`_𝔨\left(\rho \right)=\{(23)\}`$, implying $`\dot{V}(0)=u=\mathit{\varrho }_d(0),𝑩\mathit{\varrho }(0)=0`$, i.e., the system is not converging to $`\rho _d`$ in spite of the Kalman controllability condition on the linearization. Notice how for this example $`\mathrm{rank}\left(𝔚^3\mathit{\varrho }_d\right)=\mathrm{rank}\left(𝔚^3\mathit{\varrho }(0)\right)=4`$, while $`\mathit{\varrho }_d^T𝔚^3\mathit{\varrho }(0)=\{0,0,0,0\}`$.
+ $`\rho _d`$ pure (or pseudopure), but $`\rho _d=\rho _{v_2}`$
- $`H_B`$ either (28) or (29): $`_𝔨\left([H_B,\rho _d]\right)=\{(12),(23)\}`$ $``$ the linearization is always controllable. Any $`\rho (0)`$ such that $`_𝔨\left(\rho (0)\right)\{(12),(23)\}0`$ is converging.
+ $`\rho _d`$ with all different eigenvalues, e.g. $`\rho _d=\rho _{g_1}`$
- $`H_B`$ in (28): $`_𝔨\left([H_B,\rho _d]\right)=\{(12),(23)\}`$ $``$ the linearization is never controllable since now $`m=6`$;
- $`H_B`$ in (29): $`_𝔨\left([H_B,\rho _d]\right)=\{(12),(13),(23)\}`$ $``$ the linearization is always controllable. Any $`\rho (0)`$ such that $`_𝔨\left(\rho (0)\right)0`$ is converging, any $`\rho (0)`$ such that $`_𝔨\left(\rho (0)\right)=0`$ is antipodal.
* $`\rho _d\varrho _0\lambda _0`$ off-diagonal
+ $`H_B`$ in (28) and $`_𝔨\left(\rho _d\right)_𝔨\left(H_B\right)`$ $``$ linearization is never controllable, hence Theorem 1 does not apply and in general $`\rho (0)\overline{)}\rho _d(t)`$;
+ $`H_B`$ in (28) and $`_𝔨\left(\rho _d\right)\overline{)}_𝔨\left(H_B\right)`$ $``$ $`\mathrm{Card}_𝔨\left([H_B,\rho _d]\right)`$ is at least 2, implying that the linearization is controllable at least for pure/pseudopure states;
+ if $`_𝔨\left(\rho _d\right)_𝔨\left(H_B\right)0`$, then also $`_𝔥\left(\rho (0)\right)`$ matters for the convergence, see (11);
+ if $`_𝔨\left(\rho _d\right)_𝔨\left(H_B\right)=0`$ then convergence depends only on $`_𝔨\left(\rho (0)\right)`$ (plus controllability), see (12).
* If the control Hamiltonian is $`H_B=h_{𝔨,\mathrm{},12}\lambda _{𝔨,\mathrm{},12}+h_{𝔨,\mathrm{},13}\lambda _{𝔨,\mathrm{},13}`$, i.e., direct coupling between $`_2`$ and $`_3`$ is missing, then the sufficient condition of Theorem 1 does not apply. Assume for example
$$\rho _d=\left[\begin{array}{ccc}0& 0& 0\\ 0& & \\ 0& & \end{array}\right],\rho (0)=\left[\begin{array}{ccc}& & 0\\ & & 0\\ 0& 0& 0\end{array}\right].$$
Then $`_𝔨\left([H_B,\rho _d]\right)=\{(12),(13)\}`$ and $`_𝔨\left([H_B,\rho _d]\right)_𝔨\left(\rho (0)\right)=\{(12)\}`$. However, $`\rho (0)\overline{)}\rho _d(t)`$.
## VI Acknowledgments
The author would like to thank A. Agrachev and P. Rouchon for discussion on the topic of this work.
## Appendix A On the adjoint representation
A representation of a Lie algebra $`𝔤`$ on a vector space $`𝒳`$ is a mapping $`\mathrm{\Theta }:𝔤𝔤𝔩(𝒳)`$ which is a Lie algebra homomorphism, i.e., a map which
1. is linear $`\mathrm{\Theta }(\alpha _1A_1+\alpha _2A_2)=\alpha _1\mathrm{\Theta }(A_1)+\alpha _2\mathrm{\Theta }(A_2)`$, $`A_1,A_2𝔤`$ and $`\alpha _1,\alpha _2`$ in the field of $`𝒳`$;
2. preserves the Lie bracket $`\mathrm{\Theta }([A_1,A_2])=[\mathrm{\Theta }(A_1),\mathrm{\Theta }(A_2)]`$, $`A_1,A_2𝔤`$.
So a representation $`\mathrm{\Theta }`$ assigns to each $`A𝔤`$ a linear operator $`\mathrm{\Theta }(A)𝔤𝔩(𝒳)`$. A particularly useful representation is the adjoint representation on $`𝒳=^n`$. If for an $`n`$-dimensional Lie algebra $`𝔤`$ we choose the basis $`A_1,\mathrm{}A_n`$ then the Lie brackets of the basis elements are $`[A_j,A_k]=_{\mathrm{}=1}^nc_{jk}^{\mathrm{}}A_{\mathrm{}}`$. The components of the 3-tensor $`c_{jk}^{\mathrm{}}`$ are called structure constants of the Lie algebra with respect to the basis $`A_1,\mathrm{}A_n`$. The adjoint representation of $`𝔤`$, $`\mathrm{ad}_𝔤`$, with respect to the basis $`A_1,\mathrm{}A_n`$ is the representation having as basis elements the $`n\times n`$ matrices of structure constants $`𝑨_1,\mathrm{},𝑨_n`$, $`𝑨_j=\mathrm{ad}_{A_j}=[A_j,]`$ of entries $`\left(𝑨_j\right)_\mathrm{}k=c_{jk}^{\mathrm{}}`$. Notice how the two free indexes $`k`$ and $`\mathrm{}`$ identify respectively the columns and the rows of the new basis elements.
In general, the adjoint representation of a linear Lie algebra is a derivation of the algebra and corresponds to the infinitesimal representation of all the one-parameter groups of automorphisms. For a semisimple compact Lie algebra $`𝔤`$, the main features of $`\mathrm{ad}_𝔤`$ are (see e.g. Sattinger1 , p. 39 and 129):
* it is a real semisimple Lie algebra;
* it is isomorphic to $`𝔤`$;
* $`A,B𝔤`$: $`[\mathrm{ad}_A,\mathrm{ad}_B]=\mathrm{ad}_{[A,B]}`$.
Let us spend some more words on emphasizing how the “linearity” of the adjoint representation may be intended, which is one of the leitmotifs of the paper. If $`B𝔤`$ has the expression $`B=b_1A_1+\mathrm{}b_nA_n`$, then as long as we keep the basis fixed, $`B`$ is uniquely identified by its vector of components: $`B𝒃=\left[\begin{array}{ccc}b_1& \mathrm{}& b_n\end{array}\right]^T`$. Then
$$[A_j,B]=[A_j,\underset{k=1}{\overset{n}{}}b_kA_k]\underset{k,\mathrm{}=1}{\overset{n}{}}\mathrm{e}_{\mathrm{}}^T\left(\mathrm{ad}_{A_j}\right)_\mathrm{}kb_k=𝑨_j𝒃.$$
(30)
We will often make the double substitution $`\{A,B\}\{𝑨,𝒃\}`$ which will correspond to replacing the (bilinear) matrix commutator $`[,]:𝔤\times 𝔤𝔤`$ with the linear operation $`\mathrm{ad}_𝔤\times ^n^n`$, i.e., left matrix multiplication.
## Appendix B A few facts from topology
The material in this Appendix is taken from standard texts on (differential) topology e.g. Armstrong1 ; Guillemin1 . Let $`𝒳`$, $`𝒴`$ be topological spaces and $`f,g:𝒳𝒴`$ be continuous maps. The mapping $`f`$ is said homotopic to $`g`$ if “it can be continuously deformed to $`g`$”, i.e., if $``$ a continuous mapping $`h:𝒳\times [0,\mathrm{\hspace{0.17em}1}]𝒴`$ such that $`h(x,0)=f(x)`$, $`h(x,1)=g(x)`$ $`x𝒳`$. $`𝒳`$, $`𝒴`$ are said homotopy equivalent if $``$ maps $`f:𝒳𝒴`$ and $`g:𝒴𝒳`$ such that $`gf`$ and $`fg`$ are homotopic to the identity maps in $`𝒳`$ and $`𝒴`$ respectively. $`𝒴`$ is said contractible if the identity map on $`𝒴`$ is homotopic to the constant map $`𝒴x_{}`$ for any $`x_{}𝒴`$. A space is contractible if and only if it is homotopy equivalent to a point. No compact manifold is contractible. Homotopy equivalence is an equivalence relation on topological spaces and the classes of homotopy equivalent spaces are called homotopy types. Examples used in this paper are:
* $`𝕊^m\{p\}`$ is homotopy equivalent to $`^m`$ ($`p`$ is any point of $`𝕊^m`$);
* $`^m\{0\}`$ is homotopy equivalent to $`𝕊^{m1}`$
* $`^m`$ is homotopy equivalent to a point (and hence to any contractible space).
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# Determinant Formulas for Matrix Model Free Energy
## 1 Introduction
An interest to the multicut solutions to matrix models was inspired by the studies in $`𝒩=1`$ supersymmetric gauge theories due to Cachazo, Intrilligator and Vafa and Dijkgraaf, Vafa who proposed to calculate the nonperturbative superpotentials of $`𝒩=1`$ SUSY gauge theories in four dimensions using matrix models technique. This $`𝒩=1`$ theories contains the multiplet of $`𝒩=2`$ SUSY gauge theories but with nontrivial tree superpotential. The nonperturbative superpotential could be obtained from the partition functions of the one-matrix model (1MM) in the leading order in $`1/N`$, $`N`$ being the matrix size. Higher genus corrections are identified with certain holomorphic couplings of gauge theory to gravity.
The authors of proposed a new anzatz for $`_1`$ in the two-cut case (with absent double points) and made a perturbative check. Their formula in fact comes from the correspondence between the so called topological B-model on the local Calabi-Yau geometry $`\widehat{II}`$ and the cubic matrix model. Here we give complete proof of this formula and generalize it to the multi-cut case.
We start with definition of the matrix integral and introduce all relevant constructions. For a complete review of the subject, see and references there in. Consider the hermitian 1-matrix model:
$`{\displaystyle _{N\times N}}DXe^{\frac{1}{\mathrm{}}trV(X)}=e^{},`$ (1)
where $`V(X)=_{n1}t_nX^n`$, $`\mathrm{}=\frac{t_0}{N}`$ is a formal expansion parameter, the integration goes over the $`N\times N`$ matrices, $`DX_{ij}dX_{ij}`$
The topological expansion of the Feynman diagrams series is then equivalent to the expansion in even powers of $`\mathrm{}`$ for
$`(\mathrm{},t_0,t_1,t_2,\mathrm{})={\displaystyle \underset{h=0}{\overset{\mathrm{}}{}}}\mathrm{}^{2h2}_h,`$ (2)
Customarily $`t_0=\mathrm{}N`$ is the scaled number of eigenvalues. We assume the potential $`V(p)`$ to be a polynomial of the fixed degree $`m+1`$.
The averages, corresponding to the partition function (1) are defined as usual:
$`f(X)={\displaystyle \frac{1}{Z}}{\displaystyle _{N\times N}}DXf(X)\mathrm{exp}\left({\displaystyle \frac{1}{\mathrm{}}}trV(X)\right)`$ (3)
and it is convenient to use their generating functionals: the one-point resolvent
$`W(\lambda )`$ $`=\mathrm{}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{trX^k}{\lambda ^{k+1}}}.`$ (4)
as well as the $`s`$-point resolvents $`(s2)`$
$`W(\lambda _1,\mathrm{},\lambda _s)`$ $`=`$
$`\mathrm{}^{2s}{\displaystyle \underset{k_1,\mathrm{},k_s=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{trX^{k_1}\mathrm{}trX^{k_s}_{\mathrm{conn}}}{\lambda _1^{k_1+1}\mathrm{}\lambda _s^{k_s+1}}}`$ $`=`$
$`\mathrm{}^{2s}tr{\displaystyle \frac{1}{\lambda _1X}}\mathrm{}tr{\displaystyle \frac{1}{\lambda _sX}}_{\mathrm{conn}}`$ (5)
The genus expansion of the resolvent has the form
$`W(\lambda _1,\mathrm{},\lambda _s)={\displaystyle \underset{h=0}{\overset{\mathrm{}}{}}}\mathrm{}^{2h}W_h(\lambda _1,\mathrm{},\lambda _s),s1,`$ (6)
It satisfies the loop equation :
$`\left[V^{}(x)W(x)\right]_{}=W(x)^2+\mathrm{}^2W(x,x),`$ (7)
where $`\left[\mathrm{}\right]_{}`$ is the projector on the negative powers. In genus zero, loop equations have the solution
$`W_0(\lambda )`$ $`={\displaystyle \frac{1}{2}}(V^{}(\lambda )y)`$ (8)
$`y^2`$ $`=V^{}(\lambda )^2+4P_{m1}(\lambda ),`$ (9)
where $`P_{m1}`$ is an arbitrary polynomial of degree $`m1`$. If the curve (9) has $`n`$ cuts, it can be represented in terms of branching points $`\mu _\alpha `$
$`yM(\lambda )\stackrel{~}{y}M(\lambda )\sqrt{{\displaystyle _{\alpha =1}^{2n}}(\lambda \mu _\alpha )}.`$ (10)
In this article we concentrate on the case with $`m=n`$ (without double points, i.e. $`M(\lambda )`$ is a constant). Thus the full set of moduli is: $`t_I\{S_i,t_0,t_k\}`$, $`i=\overline{1,n1}`$, $`k=\overline{1,n}`$, where occupancy numbers $`S_i`$ are defined as integrals over A-cycles on the curve $`y`$,
$`S_i{\displaystyle \frac{1}{4\pi i}}{\displaystyle _{A_i}}y𝑑\lambda `$ (11)
To construct $`_1`$, we also define the polynomials $`H_I(\lambda )`$
$`{\displaystyle \frac{dy}{dt_I}}={\displaystyle \frac{H_I(\lambda )}{y(\lambda )}}`$ (12)
and matrix $`\sigma _{i,j}`$
$`\sigma _{j,i}{\displaystyle _{A_j}}{\displaystyle \frac{\lambda ^{i1}}{y(\lambda )}}𝑑\lambda ,i,j=\overline{1,n1}.`$ (13)
It can be shown that for polynomials $`H_k(\lambda )_{l=1}^{n1}H_{l,k}\lambda ^{l1}`$, $`k=\overline{1,n1}`$ corresponded $`S_k`$,
$`{\displaystyle \underset{l=1}{\overset{n1}{}}}\sigma _{j,l}H_{l,k}=\delta _{j,k}\text{for}j,k=\overline{1,n1}.`$ (14)
## 2 Two-cut case
According to paper the holomorphic part of the genus one B-model amplitude is, up to an additive constant,
$`_1={\displaystyle \frac{1}{2}}\mathrm{log}\left(det\left({\displaystyle \frac{\mu _i^{}}{S_j}}\right)\mathrm{\Delta }^{2/3}{\displaystyle \frac{2}{\mu _2^+\mu _1^+}}\right),`$ (15)
where $`\mu _1^{}=\frac{1}{2}(\mu _1\mu _2)`$, $`\mu _2^{}=\frac{1}{2}(\mu _3\mu _4)`$, $`\mu _1^+=\frac{1}{2}(\mu _1+\mu _2)`$ and $`\mu _2^+=\frac{1}{2}(\mu _3+\mu _4)`$. On the other hand there is an answer for $`_1`$ obtained directly from solving the loop equations (7) for matrix model , or using conformal field theory technique
$`_1={\displaystyle \frac{1}{24}}\mathrm{log}\left({\displaystyle \underset{\alpha =1}{\overset{2n}{}}}M(\mu _\alpha )\mathrm{\Delta }^4(\underset{i,j}{det}\sigma _{j,i})^{12}\right),`$ (16)
which, in the two-cut case without double points reads as
$`_1={\displaystyle \frac{1}{24}}\mathrm{log}\left(\mathrm{\Delta }^4\sigma ^{12}\right),`$ (17)
$`\sigma `$ (13) here is $`1\times 1`$ matrix. To obtain (15) from (17), one should prove the following formula
$`det\left({\displaystyle \frac{\mu _i^{}}{S_j}}\right)\mathrm{\Delta }{\displaystyle \frac{2}{\mu _3+\mu _4\mu _1\mu _2}}\sigma =1.`$ (18)
We can explicitly find the derivatives $`\frac{S_i}{\mu _j^{}}`$ (instead of $`\frac{\mu _i^{}}{S_j}`$), keeping times $`t_k`$ constant. To do so one should first write $`\frac{S_i}{\mu _j}`$ then make the change of variables from $`\{\mu _1`$, $`\mu _2`$, $`\mu _3`$, $`\mu _4\}`$ to $`\{t_1`$, $`t_2`$, $`\mu _1^{}`$, $`\mu _2^{}\}`$. Then
$`{\displaystyle \frac{S_i}{\mu _j^{}}}={\displaystyle \frac{S_i}{\mu _k}}{\displaystyle \frac{\mu _k}{\mu _j^{}}},i,j=\overline{1,n},k=\overline{1,2n}.`$ (19)
$`\frac{\mu _k}{\mu _j^{}}`$ here are obtained by inverting the matrix $`(\frac{\mu _j^{}}{\mu _k},\frac{S_j}{\mu _k})`$. After this, it is easy to rewrite (18) using the elliptic integrals:
$`\sqrt{{\displaystyle \frac{\mu _4\mu _2}{\mu _3\mu _1}}}({\displaystyle \frac{\mu _4\mu _1}{\mu _4\mu _2}}\mathrm{\Pi }({\displaystyle \frac{\mu _2\mu _1}{\mu _4\mu _2}},\kappa )+`$
$`{\displaystyle \frac{\mu _3\mu _2}{\mu _4\mu _2}}\mathrm{\Pi }({\displaystyle \frac{\mu _4\mu _3}{\mu _4\mu _2}},\kappa )+K(\kappa ))`$ $`={\displaystyle \frac{\pi }{2}}.`$ (20)
where $`\kappa =\sqrt{\frac{(\mu _2\mu _1)(\mu _4\mu _3)}{(\mu _4\mu _2)(\mu _3\mu _1)}}`$, $`\mathrm{\Pi }(\nu ,\kappa )`$ and $`K(\kappa )`$ are complete elliptic integrals of the third and first kinds respectively. To prove this statement, one can rewrite the elliptic integrals of the third kind via the complete and incomplete elliptic integrals of the first and the second kinds (these formulas can be found in (formulas 22, 24 from chapter 13.8); note, however, that in there is a misprint in these formulas)
$`k^2{\displaystyle \frac{\mathrm{sin}\theta \mathrm{cos}\theta }{\sqrt{1k^2\mathrm{sin}\theta ^2}}}[\mathrm{\Pi }(1k^2\mathrm{sin}\theta ^2,\kappa )K(\kappa )]`$ $`=`$
$`{\displaystyle \frac{\pi }{2}}(E(\kappa )K(\kappa ))F(\mathrm{sin}\theta ,k^{})K(\kappa )E(\mathrm{sin}\theta ,k^{})`$ (21)
$`{\displaystyle \frac{\sqrt{1k^2\mathrm{sin}\theta ^2}}{\mathrm{sin}\theta \mathrm{cos}\theta }}[\mathrm{\Pi }(k^2\mathrm{tan}\theta ^2,\kappa )K(\kappa )\mathrm{cos}\theta ^2]`$ $`=`$
$`(E(\kappa )K(\kappa ))F(\mathrm{sin}\theta ,k^{})K(\kappa )E(\mathrm{sin}\theta ,k^{})`$ (22)
where $`k^{}=\sqrt{1k^2}`$, $`\theta [0,\pi /2]`$. In this case one should put $`\mathrm{sin}^2\theta =\frac{\mu _3\mu _1}{\mu _4\mu _1}.`$
The same computation can be done for any other partition of $`\mu _i`$ into the two sets $`\mu _{1,2}^\pm `$ (without changing $`\sigma `$), say, for $`\mu _1^{}=\frac{1}{2}(\mu _1\mu _3)`$, $`\mu _2^{}=\frac{1}{2}(\mu _2\mu _4)`$, $`\mu _1^+=\frac{1}{2}(\mu _1+\mu _3)`$ and $`\mu _2^+=\frac{1}{2}(\mu _2+\mu _4)`$. It leads to the same result (15), however, the perturbative calculation in this case is irrelevant.
## 3 Generalization for n-cut solution
A natural generalisation of (15) is
$`_1={\displaystyle \frac{1}{2}}\mathrm{log}\left(det{\displaystyle \frac{\{\mu _j^{}\}}{\{S_i,S_n\}}}\mathrm{\Delta }^{2/3}\mathrm{\Delta }^1(\mu _j^+)\right),`$ (23)
where we divided all the branching points into two ordered sets $`\{\mu _j^{(1)}\}_{j=1}^n`$ and $`\{\mu _j^{(2)}\}_{j=1}^n`$ and performed a linear orthogonal transformation of $`\mu _j^{(1,2)}`$ to the quantities $`\{\mu _j^+\}`$ and $`\{\mu _j^{}\}`$, $`j=\overline{1,n}`$,
$`\mu _j^\pm =\mu _j^{(1)}\pm \mu _j^{(2)}.`$ (24)
To prove formula (23), one should calculate the derivative of the branching points $`\mu _j`$ with respect to the moduli $`\{t_K\}\{S_1..S_{n1},t_0..t_n\}`$ :
$`{\displaystyle \frac{\mu _\alpha }{t_K}}={\displaystyle \frac{H_K(\mu _\alpha )}{M(\mu _\alpha )_{\beta \alpha }(\mu _\alpha \mu _\beta )}}.`$ (25)
The polynomials $`H_I(\lambda )`$ corresponding to the variables $`t_k,k1`$ always have the coefficient $`k`$ at the highest term $`\lambda ^{n1+k}`$ and the polynomial corresponding to $`t_0`$ starts with $`\lambda ^{n1}`$. Therefore, one can find the determinant:
$`det{\displaystyle \frac{\{\mu _{\alpha _j}\}}{\{S_i,S_n,t_k\}}}={\displaystyle \frac{\mathrm{\Delta }(\mu _{\alpha _j})\left(\underset{l,m}{det}\sigma _{l,m}\right)^1}{\underset{i=1}{\overset{2n}{}}M(\mu _{\alpha _i})\underset{j=1}{\overset{2n}{}}\left(\underset{\beta \alpha _j}{\overset{2n}{}}(\mu _{\alpha _j}\mu _\beta )\right)}}`$ (26)
Indeed, consider the left hand side of (26).
$`det{\displaystyle \frac{\{\mu _{\alpha _j}\}}{\{S_i,S_n,t_k\}}}={\displaystyle \frac{\left(\underset{K,j}{det}H_K(\mu _{\alpha _j})\right)}{\underset{i=1}{\overset{2n}{}}M(\mu _{\alpha _i})\underset{j=1}{\overset{2n}{}}\left(\underset{\beta \alpha _j}{\overset{2n}{}}(\mu _{\alpha _j}\mu _\beta )\right)}}`$ (27)
The change of variables $`\{S_1,..,S_n\}\{S_1,\mathrm{},S_{n1},t_0\}`$ does not change the determinant. To obtain the Vandermonde determinant in the right hand side of (26), there should be, instead of the polynomials $`H_K`$, polynomials of degree $`2ni+1`$ where $`i`$ is the line number, with unit leading coefficients. To this end, one should multiply the matrix $`H_K(\mu _{\alpha _j})`$ with the block diagonal matrix
$`\stackrel{~}{\sigma }=\left(\begin{array}{cc}1& 0\\ 0& \sigma \end{array}\right).`$ (28)
This gives the factor $`(det\stackrel{~}{\sigma })^1=(\underset{l,m}{det}\sigma _{l,m})^1`$. Lines from $`1`$ to $`n+1`$ contribute to $`n!`$ which could be omitted from the free energy. The Vandermonde determinant $`\mathrm{\Delta }(\mu _{\alpha _j})`$ then combines with the rational factors in the denominator to produce $`(1)^{_{j=1}^n\alpha _j}\mathrm{\Delta }(\overline{\mu _{\alpha _j}})/\mathrm{\Delta }(\mu )`$, where $`\mathrm{\Delta }(\overline{\mu _{\alpha _j}})`$ is the Vandermonde determinant for the supplementary set of $`n`$ branching points not entering the set $`\{\mu _{\alpha _j}\}_{j=1}^n`$ whereas $`\mathrm{\Delta }(\mu )`$ is the total Vandermonde determinant. Now we should put $`M(\mu _\alpha )`$ constant independent of $`\alpha `$. Expanding the determinant in (23) by each line and neglecting the additive constant $`\frac{1}{2}\mathrm{log}2^n`$, one obtain (16).
Introducing the quantities
$`\varphi _I^\alpha {\displaystyle \frac{H_I(\mu _\alpha )}{M^{1/3}(\mu _\alpha )_{\beta \alpha }(\mu _\alpha \mu _\beta )^{2/3}}},`$ (29)
one can rewrite (16) in a more simple form:
$`_1={\displaystyle \frac{1}{2}}\mathrm{log}\left(\underset{I,\alpha }{det}\varphi _I^\alpha \right)`$ (30)
## 4 Perturbative Formula
We have also performed the perturbative check of (23) for the 3-cut case. It is easier to make the expansion not in the moduli $`S_i`$ but in the difference of the branching points $`\mu _j^{}`$. In order to calculate $`det\frac{\{S_i,S_n\}}{\{\mu _j^{}\}}`$, one should rewrite $`S_i`$ and $`\sigma _{i,j}`$ in terms of $`\mu _i^+`$, $`\mu _j^{}`$ and expand them in $`\mu _i^{}`$
$`S_l`$ $`={\displaystyle \frac{1}{2}}res_{\lambda =\mu _l^+}{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mu _i^{})^{2k}c_k}{(\lambda \mu _i^+)^{2k1}}}`$ (31)
$`\sigma _{l,j}`$ $`={\displaystyle \frac{1}{2}}res_{\lambda =\mu _l^+}\lambda ^{j1}{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\mu _i^{})^{2k}\stackrel{~}{c}_k}{(\lambda \mu _i^+)^{2k+1}}}`$ (32)
$`c_k`$ and $`\stackrel{~}{c}_k`$ are the Taylor coefficients for $`\sqrt{1x}`$ and $`\frac{1}{\sqrt{1x}}`$ respectively. It should be mentioned that derivatives $`\frac{\{S_i\}}{\{\mu _j^{}\}}`$ are taken at $`t_k`$ constant, while in (31) $`S_k`$ are functions of $`\mu ^+,\mu ^{}`$. This problem is solved by calculating the transition matrix from $`\{\mu _k^{},t_k\}`$ to $`\{\mu _k^{},\mu _k^+\}`$ and inverting it. We have done this calculation up to $`(\mu ^{})^3`$ and found it in perfect agreement with (16) (up to an additive constant mentioned).
## Acknowledgments
Our work is partly supported by Federal Program of the Russian Ministry of Industry, Science and Technology No 40.052.1.1.1112 and by the grant RFBR 04-02-16880
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# 1 Introduction
## 1 Introduction
The suggestion that, in analogy with $`K^o,\overline{K}^o`$ oscillations, there could be neutrino-antineutrino oscillations ( $`\nu \overline{\nu }`$), was considered by Pontecorvo in 1957. It was subsequently considered by Maki et al. and Pontecorvo that there could be mixings (and oscillations) of neutrinos of different flavors (i.e., $`\nu _e\nu _\mu `$ transitions).
In the general case there can be two schemes (types) of neutrino mixings (oscillations): mass mixing schemes and charge mixings scheme (as it takes place in the vector dominance model or vector boson mixings in the standard model of electroweak interactions) .
In the Standard theory of neutrino oscillations is supposed that physically observed neutrino states $`\nu _e,\nu _\mu ,\nu _\tau `$ have no definite masses and they are directly produced as mixture of the $`\nu _1,\nu _2,\nu _3`$ neutrino states. And if neutrino oscillations are generated by the neutrino mass matrix, then neutrino mixing parameters are expressed via elements of the neutrino mass matrix.
The mass lagrangian of two neutrinos ($`\nu _e,\nu _\mu `$) has the following form (for simplification the case of two neutrinos is considered):
$$\begin{array}{c}_M=\frac{1}{2}\left[m_{\nu _e}\overline{\nu }_e\nu _e+m_{\nu _\mu }\overline{\nu }_\mu \nu _\mu +m_{\nu _e\nu _\mu }\left(\overline{\nu }_e\nu _\mu +\overline{\nu }_\mu \nu _e\right)\right]\\ \frac{1}{2}(\overline{\nu }_e,\overline{\nu }_\mu )\left(\begin{array}{cc}m_{\nu _e}& m_{\nu _e\nu _\mu }\\ m_{\nu _\mu \nu _e}& m_{\nu _\mu }\end{array}\right)\left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)\end{array},$$
$`\left(1\right)`$
which is diagonalized by rotation on the angle $`\theta `$ and then this lagrangian (1) transforms into the following one (see ref. in ):
$$_M=\frac{1}{2}\left[m_1\overline{\nu }_1\nu _1+m_2\overline{\nu }_2\nu _2\right],$$
$`\left(2\right)`$
where
$$m_{1,2}=\frac{1}{2}\left[\left(m_{\nu _e}+m_{\nu _\mu }\right)\pm \left(\left(m_{\nu _e}m_{\nu _\mu }\right)^2+4m_{\nu _\mu \nu _e}^2\right)^{1/2}\right],$$
and angle $`\theta `$ is determined by the following expression:
$$tg\left(2\theta \right)=\frac{2m_{\nu _e\nu _\mu }}{\left(m_{\nu _\mu }m_{\nu _e}\right)},$$
$`\left(3\right)`$
$$\begin{array}{c}\nu _e=cos\theta \nu _1+sin\theta \nu _2,\\ \nu _\mu =sin\theta \nu _1+cos\theta \nu _2.\end{array}$$
$`\left(4\right)`$
Then $`\nu _e,\nu _\mu `$ masses are:
$$m_{\nu _e}=m_1cos^2\theta +m_2sin^2\theta ,$$
$$m_{\nu _\mu }=m_1sin^2\theta +m_2cos^2\theta ,$$
$`\left(5\right)`$
in contrast to the primary supposition that $`\nu _e,\nu _\mu ,\mu _\tau `$ neutrinos have no definite masses.
The probability of $`\nu _e\nu _e`$ is given by the following expression:
$$P\left(\nu _e\nu _e\right)=1\mathrm{sin}^2\left(2\theta \right)sin^2\left(\left(m_2^2m_1^2\right)/2p\right)t,$$
$`\left(6\right)`$
where
$$sin\theta =\frac{1}{\sqrt{2}}\left[1\frac{\left|m_{\nu _\mu }m_{\nu _e}\right|}{\sqrt{\left(m_{\nu _\mu }m_{\nu _e}\right)^2+\left(2m_{\nu _e\nu _\mu }\right)^2}}\right],$$
$`\left(7\right)`$
or
$$sin^2\left(2\theta \right)=\frac{\left(2m_{\nu _e\nu _\mu }\right)^2}{\left(m_{\nu _e}m_{\nu _\mu }\right)^2+\left(2m_{\nu _e\nu _\mu }\right)^2}.$$
$`\left(8\right)`$
Then the nondiagonal mass term $`m_{\nu _e\nu _\mu }`$ of the mass matrix in (1) can be interpreted as width of $`\nu _e\nu _\mu `$ transitions .
In this standard theory of neutrino oscillations neutrino oscillations are real even when neutrino masses are different therefore the law of energy momentum conservation is violated. In the corrected theory of neutrino oscillations the law of energy momentum conservation is fulfilled and neutrino oscillations are virtual if neutrino masses are different and real if neutrino masses are equal.
It is necessary to note that in physics all the processes are realized through dynamics. Unfortunately, in the above considered mass mixings scheme the dynamics is absent. Probably, this is an indication of the fact that this scheme is incomplete one, i.e., this scheme requires a physical substantiation. Below we consider neutrino oscillations which appear in the scheme of charge (couple constant) mixings, i.e. by using dynamics .
## 2 Theory of Neutrino Oscillations in the Framework of Charge Mixings Scheme
At first we consider a case of two neutrino mixings (oscillations) and then we consider the case of three neutrino oscillations.
### 2.1 The Case of Two Neutrino Mixings (Oscillations) in the Charge Mixings Scheme
In this scheme (or mechanism) the neutrino mixings or transitions can be realized by mixings of the neutrino fields in analogy with the vector dominance model ($`\gamma \rho ^o`$ and $`Z^o\gamma `$ mixings), the way it takes place in the particle physics. Then, in the case of two neutrinos, we have
$$\nu _1=\mathrm{cos}\theta \nu _e\mathrm{sin}\theta \nu _\mu ,$$
$`\left(9\right)`$
$$\nu _2=\mathrm{sin}\theta \nu _e+\mathrm{cos}\theta \nu _\mu .$$
The charged current in the standard model of weak interactions for two lepton families has the following form:
$$j^\alpha =\left(\begin{array}{cc}\overline{e}\overline{\mu }& \end{array}\right)_L\gamma ^\alpha V\left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)_L,$$
$$V=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),$$
$`\left(10\right)`$
and then the interaction Lagrangian is
$$=\frac{g}{\sqrt{2}}j^\alpha W_\alpha ^++\mathrm{h}.\mathrm{c}.$$
$`\left(11\right)`$
and
$$\begin{array}{c}\nu _e=\mathrm{cos}\theta \nu _1+\mathrm{sin}\theta \nu _2\\ \nu _\mu =\mathrm{sin}\theta \nu _1+\mathrm{cos}\theta \nu _2.\end{array}$$
$`\left(12\right)`$
The lagrangian (10)-(11) can be rewritten in the following form:
$$=\frac{g}{\sqrt{2}}j^\alpha W_\alpha ^++\mathrm{h}.\mathrm{c}.,$$
$`\left(13\right)`$
where $`j^\alpha `$ is
$$j^\alpha =\left(\begin{array}{cc}\overline{e}\overline{\mu }& \end{array}\right)_L\gamma ^\alpha \left(\begin{array}{c}\nu _1\\ \nu _2\end{array}\right)_L.$$
And the mass matrix is
$$\left(\begin{array}{cc}m_1& 0\\ 0& m_2\end{array}\right).$$
In this case the neutrino oscillations cannot take place, and even if neutrino oscillations take place, then there must be $`\nu _1\nu _2`$ neutrino oscillations but not $`\nu _e\nu _\mu `$ oscillations.
In this point some questions arise. Where have we taken $`\nu _e,\nu _\mu `$ neutrinos if in the weak interactions, given by expression (13), $`\nu _1,\nu _2`$ neutrinos are produced? From the all existent accelerator experiments very well know that in the weak interactions $`\nu _e,\nu _\mu `$ neutrinos are produced and that the $`l_{\nu _e},l_{\nu _\mu }`$ lepton numbers are well conserved ones. Obviously we must solve this problem. So, $`\nu _1,\nu _2`$ neutrinos are eigenstates of the weak interactions when we take mixing matrix $`V`$ into account and $`\nu _e,\nu _\mu `$ neutrinos are eigenstates of the weak interactions with $`W,Z^o`$ boson exchanges. Then we have to rewrite the lagrangian of the weak interaction in the correct form to describe neutrino productions and oscillations correctly. Then
$$=\frac{g}{\sqrt{2}}j^\alpha W_\alpha ^++\mathrm{h}.\mathrm{c}.$$
$`\left(14\right)`$
where $`j^\alpha `$ is
$$j^\alpha =\left(\begin{array}{cc}\overline{e}\overline{\mu }& \end{array}\right)_L\gamma ^\alpha \left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)_L,$$
How are the lepton numbers violated? It is necessary to suppose that after $`\nu _e,\nu _\mu `$ production the violation of lepton numbers takes place, i.e.,
$$\left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)_L=V\left(\begin{array}{c}\nu _1\\ \nu _2\end{array}\right)_L,V=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),$$
$`\left(15\right)`$
and then $`\nu _e,\nu _\mu `$ neutrinos become superpositions of $`\nu _1,\nu _2`$ neutrinos.
$$\begin{array}{c}\nu _e=cos\theta \nu _1+sin\theta \nu _2,\\ \nu _\mu =sin\theta \nu _1+cos\theta \nu _2.\end{array}$$
$`\left(16\right)`$
Taking into account that the charges of $`\nu _1,\nu _2`$ neutrinos are $`g_1,g_2`$, we get
$$g\mathrm{cos}\theta =g_1,g\mathrm{sin}\theta =g_2,$$
$`\left(17\right)`$
i.e.
$$\mathrm{cos}\theta =\frac{g_1}{g},\mathrm{sin}\theta =\frac{g_2}{g}.$$
$`\left(18\right)`$
Since $`\mathrm{sin}^2\theta +\mathrm{cos}^2\theta =1`$, then
$$g=\sqrt{g_1^2+g_2^2}$$
and
$$\mathrm{cos}\theta =\frac{g_1}{\sqrt{g_1^2+g_2^2}},\mathrm{sin}\theta =\frac{g_2}{\sqrt{g_1^2+g_2^2}}.$$
$`\left(19\right)`$
Since we suppose that $`g_1g_2\frac{g}{\sqrt{2}}`$, then
$$\mathrm{cos}\theta \mathrm{sin}\theta \frac{1}{\sqrt{2}}.$$
$`\left(20\right)`$
In the general case the couple constants $`g_1,g_2`$ and $`g`$ can have no connections and then we obtain only expressions (19).
What happens with the neutrino mass matrix in this case? The primary neutrino mass matrix has the following diagonal form:
$$\left(\begin{array}{cc}m_{\nu _e}& 0\\ 0& m_{\nu _\mu }\end{array}\right),$$
$`\left(21\right)`$
since in the weak interactions (with $`W,Z^o`$ bosons) the lepton numbers are conserved and then $`\nu _e,\nu _\mu `$ are eigenstates of these interactions.
It is interesting to note that the same situation takes place in the quark sector when we consider $`K^o,\overline{K}^o`$ oscillations. In the strong interactions only $`d,s,b`$ quarks are produced and the aroma numbers are well conserved in these interactions, i.e., these states are eigenstates of the strong interactions. Then oscillations appear while at violating the aroma numbers by the weak interactions with the Cabibbo-Kobayashi-Maskawa matrices.
Then due to the presence of terms violating the lepton numbers, the nondiagonal terms appear in this matrix and then this mass matrix is transformed into the following nondiagonal matrix (the case when $`CP`$ is conserved ):
$$\left(\begin{array}{cc}m_{\nu _e}& m_{\nu _e\nu _\mu }\\ m_{\nu _\mu \nu _e}& m_{\nu _\mu }\end{array}\right),$$
$`\left(22\right)`$
then the masses lagrangian of neutrinos takes the following form:
$$\begin{array}{c}_M=\frac{1}{2}\left[m_{\nu _e}\overline{\nu }_e\nu _e+m_{\nu _\mu }\overline{\nu }_\mu \nu _\mu +m_{\nu _e\nu _\mu }\left(\overline{\nu }_e\nu _\mu +\overline{\nu }_\mu \nu _e\right)\right]\\ \frac{1}{2}(\overline{\nu }_e,\overline{\nu }_\mu )\left(\begin{array}{cc}m_{\nu _e}& m_{\nu _e\nu _\mu }\\ m_{\nu _\mu \nu _e}& m_{\nu _\mu }\end{array}\right)\left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)\end{array}.$$
$`\left(23\right)`$
Masses lagrangian of the new states obtained by diagonalizing of this matrix while rotating on angle $`\theta `$, has the following form (these states are namely the same weak interactions states considered above):
$$_M=\frac{1}{2}(\overline{\nu }_e,\overline{\nu }_\mu )V^1\left(\begin{array}{cc}m_{\nu _1}& 0\\ 0& m_2\end{array}\right)V\left(\begin{array}{c}\nu _e\\ \nu _\mu \end{array}\right)=$$
$$\frac{1}{2}(\overline{\nu }_e,\overline{\nu }_\mu )\left)\right(\begin{array}{cc}m_1cos^2\theta +m_2sin^2\theta & \left(m_2m_1\right)cos\theta sin\theta \\ \left(m_2m_1\right)cos\theta sin\theta & m_1sin^2\theta +m_2cos^2\theta \end{array}\left)\right(\begin{array}{c}\nu _e\\ \nu _\mu \end{array})=$$
$$\frac{1}{2}\left[m_1\overline{\nu }_1\nu _1+m_2\overline{\nu }_2\nu _2\right],$$
$`\left(24\right)`$
where $`\nu _1,\nu _2`$ are eigenstates and $`m_1,m_2`$ are their eigenmasses and from expressions (23), (24) we obtain
$$\begin{array}{c}\nu _1=cos\theta \nu _esin\theta \nu _\mu ,\\ \nu _2=sin\theta \nu _e+cos\theta \nu _\mu ,\end{array}$$
$`\left(25\right)`$
$$m_{\nu _e}=m_1cos^2\theta +m_2sin^2\theta ,$$
$$m_{\nu _\mu }=m_1sin^2\theta +m_2cos^2\theta ,$$
$`\left(26\right)`$
$$m_{\nu _e\nu _\mu }=\left(m_2m_1\right)cos\theta sin\theta ,$$
or
$$m_1=\frac{\left(m_{\nu _e}cos^2\theta m_{\nu _\mu }sin^2\theta \right)}{\left(cos^2\theta sin^2\theta \right)},$$
$$m_2=\frac{\left(m_{\nu _e}sin^2\theta m_{\nu _\mu }cos^2\theta \right)}{\left(cos^2\theta sin^2\theta \right)},$$
$`\left(27\right)`$
where $`sin\theta ,cos\theta `$ are given by expressions (19)
Then $`\mathrm{\Delta }m^2=m_2^2m_1^2`$ is
$$\mathrm{\Delta }m^2=\frac{\left(m_{\nu _\mu }^2cos^4\theta m_{\nu _e}sin^4\theta \right)}{\left(cos^2\theta sin^2\theta \right)^2}.$$
$`\left(28\right)`$
The expression for time evolution of $`\nu _1,\nu _2`$ neutrinos (see exp. (25)-(27) with masses $`m_1`$ and $`m_2`$ is
$$\nu _1\left(t\right)=e^{iE_1t}\nu _1\left(0\right),\nu _2\left(t\right)=e^{iE_2t}\nu _2\left(0\right),$$
$`\left(29\right)`$
where
$$E_k^2=\left(p^2+m_k^2\right),k=1,2.$$
If neutrinos are propagating without interactions, then
$$\begin{array}{c}\nu _e\left(t\right)=cos\theta e^{iE_1t}\nu _1\left(0\right)+sin\theta e^{iE_2t}\nu _2\left(0\right),\\ \nu _\mu \left(t\right)=sin\theta e^{iE_1t}\nu _1\left(0\right)+cos\theta e^{iE_2t}\nu _2\left(0\right).\end{array}$$
$`\left(30\right)`$
Using the expression for $`\nu _1`$ and $`\nu _2`$ from (25), and putting it into (20), one can get the following expression:
$$\nu _e\left(t\right)=\left[e^{iE_1t}cos^2\theta +e^{iE_2t}sin^2\theta \right]\nu _e\left(0\right)+$$
$$+\left[e^{iE_1t}e^{iE_2t}\right]sin\theta \mathrm{cos}\theta \nu _\mu \left(0\right),$$
$`\left(31\right)`$
$$\nu _\mu \left(t\right)=\left[e^{iE_1t}sin^2\theta +e^{iE_2t}cos^2\theta \right]\nu _\mu \left(0\right)+$$
$$+\left[e^{iE_1t}e^{iE_2t}\right]sin\theta cos\theta \nu _e\left(0\right).$$
The probability that neutrino $`\nu _e`$ produced at time $`t=0`$, will be transformed into $`\nu _\mu `$ at time $`t`$, is an absolute value of amplitude $`\nu _\mu \left(0\right)`$ in (31) squared, i. e.
$$P\left(\nu _e\nu _\mu \right)=\left(\nu _\mu \left(0\right)\nu _e\left(t\right)\right)^2=$$
$$=\frac{1}{2}\mathrm{sin}^2\left(2\theta \right)\left[1cos\left(\left(m_2^2m_1^2\right)/2p\right)t\right],$$
$`\left(32\right)`$
where it is supposed that $`pm_1,m_2`$ and $`E_kp+m_k^2/2p`$.
The expression (32) presents the probability of neutrino flavor oscillations. The angle $`\theta `$ (mixing angle) characterizes the value of mixing. The probability $`P\left(\nu _e\nu _\mu \right)`$ is a periodical function of distances where the period is determined by the following expression:
$$L_o=2\pi \frac{2p}{m_2^2m_1^2}.$$
$`\left(33\right)`$
And probability $`P\left(\nu _e\nu _e\right)`$ that the neutrino $`\nu _e`$ produced at time $`t=0`$ is preserved as $`\nu _e`$ neutrino at time $`t`$, is given by the absolute value of the amplitude of $`\nu _e\left(0\right)`$ in (31) squared. Since the states in (31) are normalized states, then
$$P\left(\nu _e\nu _e\right)+P\left(\nu _e\nu _\mu \right)=1.$$
$`\left(34\right)`$
So, we see that flavor oscillations caused by nondiagonality of the neutrinos mass matrix violate the law of the $`\mathrm{}_e`$ and $`\mathrm{}_\mu `$ lepton number conservations. However in this case, as one can see from exp. (34), the full lepton numbers $`\mathrm{}=\mathrm{}_e+\mathrm{}_\mu `$ are conserved.
It is necessary to stress that neutrino oscillations in this scheme (mechanism) are virtual if neutrino masses are different and real if neutrino masses are equal and these oscillations are preserved within the uncertainty relations.
### 2.2 The Case of Three Neutrino Mixings (Oscillations) in the Charge Mixings Scheme
In the case of three neutrinos we can choose parameterization of the mixing matrix $`V`$ in the form proposed by Maiani :
$$V=\left(\begin{array}{ccc}1& 0& 0\\ 0& c_\gamma & s_\gamma \\ 0& s_\gamma & c_\gamma \end{array}\right)\left(\begin{array}{ccc}c_\beta & 0& s_\beta \\ 0& 1& 0\\ s_\beta & 0& c_\beta \end{array}\right)\left(\begin{array}{ccc}c_\theta & s_\theta & 0\\ s_\theta & c_\theta & 0\\ 0& 0& 1\end{array}\right),$$
$`\left(35\right)`$
$$c_{e\mu }=\mathrm{cos}\theta s_{e\mu }=\mathrm{sin}\theta ,c_{e\mu }^2+s_{e\mu }^2=1;$$
$$c_{e\tau }=\mathrm{cos}\beta ,s_{e\tau }=\mathrm{sin}\beta ,c_{e\tau }^2+s_{e\tau }^2=1;$$
$`\left(36\right)`$
$$c_{\mu \tau }=\mathrm{cos}\gamma ,s_{\mu \tau }=\mathrm{sin}\gamma ,c_{\mu \tau }^2+s_{\mu \tau }^2=1.$$
It is not difficult to come to consideration of the case of three neutrino types $`\nu _e,\nu _\mu ,\nu _\tau `$.
For the first and second families (at $`\nu _e,\nu _\mu `$ neutrino oscillations) we get
$$\mathrm{cos}\theta =\mathrm{cos}\theta _{\nu _e\nu _\mu }=\frac{g_1}{\sqrt{g_1^2+g_2^2}},$$
$$sin\left(2\theta \right)=\frac{2g_1g_2}{g_1^2+g_2^2}.$$
$`\left(37\right)`$
Then the probability of $`\nu _e\nu _e`$ is given by the following expression:
$$P\left(\nu _e\nu _e\right)=1\mathrm{sin}^2\left(2\theta \right)sin^2\left(\pi t\left(m_2^2m_1^2\right)/2p_{\nu _e}\right),$$
$`\left(38\right)`$
In the case $`g_1g_2`$
$$\mathrm{sin}\theta _{\nu _e\nu _\mu }\mathrm{cos}\theta _{\nu _e\nu _\mu }\frac{1}{\sqrt{2}}.$$
$`\left(39\right)`$
For the first and third families (at $`\nu _e,\nu _\tau `$ neutrino oscillations) we get
$$\mathrm{cos}\beta =\mathrm{cos}\beta _{\nu _e\nu _\tau }=\frac{g_1}{\sqrt{g_1^2+g_3^2}},$$
$$sin\left(2\beta \right)=\frac{2g_1g_3}{g_1^2+g_3^2}.$$
$`\left(40\right)`$
Then the probability of $`\nu _e\nu _e`$ is given by the following expression:
$$P\left(\nu _e\nu _e\right)=1\mathrm{sin}^2\left(2\beta \right)sin^2\left(\pi t\left(m_3^2m_1^2\right)/2p_{\nu _e}\right),$$
$`\left(41\right)`$
In the case $`g_1g_3`$
$$\mathrm{cos}\beta _{\nu _e\nu _\tau }\mathrm{sin}\beta _{\nu _e\nu _\tau }\frac{1}{\sqrt{2}}.$$
$`\left(43\right)`$
For the second and third families (at $`\nu _\nu ,\nu _\tau `$ neutrino oscillations) we get
$$\mathrm{cos}\gamma =\mathrm{cos}\gamma _{\nu _\mu \nu _\tau }=\frac{g_2}{\sqrt{g_2^2+g_3^2}},$$
$$sin\left(2\gamma \right)=\frac{2g_2g_3}{g_2^2+g_3^2}.$$
$`\left(44\right)`$
Then the probability of $`\nu _\mu \nu _\mu `$ is given by the following expression:
$$P\left(\nu _\mu \nu _\mu \right)=1\mathrm{sin}^2\left(2\gamma \right)sin^2\left(\pi t\left(m_3^2m_2^2\right)/2p_{\nu _\mu }\right),$$
$`\left(45\right)`$
In the case $`g_2g_3`$
$$\mathrm{cos}\gamma _{\nu _\mu \nu _\tau }\mathrm{sin}\gamma _{\nu _\mu \nu _\tau }\frac{1}{\sqrt{2}}.$$
$`\left(46\right)`$
So the neutrino mixings (oscillations) appear due to the fact that at neutrino production the eigenstates of the weak interactions (i.e. $`\nu _e,\nu _\mu ,\nu _\tau `$ neutrino states) are generated but not the eigenstates of the weak interaction violating lepton numbers (i.e. $`\nu _1,\nu _2,\nu _3`$ neutrino states). And when neutrinos are passing through vacuum they are converted into superpositions of $`\nu _1,\nu _2,\nu _3`$ neutrinos and through these intermediate states the oscillations (transitions) between $`\nu _e,\nu _\mu ,\nu _\tau `$ neutrinos are realized.
## 3 Conclusion
It is necessary to note that in physics all the processes are realized through dynamics. Unfortunately, in the standard theory of neutrino oscillations based on the masses mixings scheme (mechanism), the dynamics is absent. Probably, this is an indication of the fact that this scheme is incomplete one, i.e., this scheme requires a physical substantiation.
In this work neutrino oscillations generated by the weak interaction couple constant (charge) mixings where considered (as it takes place in the model of vector dominance or in the electroweak interactions model at vector boson mixings). Expressions for angle mixings and lengths of oscillations were obtained. The expressions of probabilities for three neutrino oscillations were given. Neutrino oscillations in this scheme (mechanism) are virtual if neutrino masses are different and real–if neutrino masses are equal.
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4. Beshtoev Kh. M., JINR Communication E2-2004-58, Dubna, 2004;
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# QCD Analytic Perturbation Theory. From integer powers to any power of the running coupling
## I Introduction
A fundamental goal of perturbative QCD is to provide a microscopic description of hadronic short-distance phenomena that yields reliable predictions to be compared with experimental data of increasing precision. While singularities on the timelike axis in the complex $`Q^2`$ plane of hadronic observables are related to physical particles (or resonances), the appearance of singularities on the spacelike axis are unphysical and may violate causality. On the other hand, the expansion of hadronic quantities at large momentum transfer $`Q^2`$ can be safely calculated in terms of a power-series expansion in the running strong coupling $`\alpha _s(Q^2)`$ by virtue of asymptotic freedom. But the one-loop running coupling contains at $`Q^2=\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ ($`\mathrm{\Lambda }_{\mathrm{QCD}}\mathrm{\Lambda }`$ in the following) a ghost singularity—the Landau pole—that spoils its analyticity structure. To restore analyticity and ensure causality in the whole $`Q^2`$ plane, this pole has to be removed. With most available experimental data on several exclusive processes being at rather low $`Q^2`$ values, the Landau-singularity problem is not only of academic interest, but affects significantly perturbative predictions in the low-to-medium $`Q^2`$ domain. The reason is that—lacking all-order perturbative expressions—one has to resort to a renormalization-scheme choice that makes the uncalculated higher-order corrections negligible and adopt a renormalization scale that reflects the typical parton virtualities in the considered process. The latter procedure, however, may result into a scale in the region of only a few $`\mathrm{\Lambda }`$, where the application of perturbation theory for the conventional running coupling without infrared (IR) protection against the Landau pole becomes inapplicable—a prominent example being the Brodsky–Lepage–Mackenzie scale-setting procedure BLM83 . Different strategies have been suggested over the years as how to minimize the dependence on the renormalization scheme and scale setting—unavoidable in any perturbative calculation beyond the leading order—and obtain reliable and stable results in the low-momentum regime (see, for example, Ref. BPSS04 for a recent extensive discussion of these issues in terms of the electromagnetic pion form factor and references cited therein).
In a series of papers during the last few years Shirkov and Solovtsov (SS) SS97 ; Shi98 ; SS99 ; Shi01 ; SS01 have developed an approach which enables the removal of the Landau singularity without introducing extraneous IR regulators, like an effective gluon mass PP79 ; Cor82 ; GHN93 ; MS92 ; MaSt93 . The analyticity of the coupling in the spacelike region is achieved by a nonperturbative, power-behaved term that contains no other scale than $`\mathrm{\Lambda }`$ and leaves the ultraviolet (UV) behavior of the running coupling unchanged. At zero-momentum transfer the Shirkov–Solovtsov coupling assumes a universal value that depends solely on renormalization-group constants. Using dispersion relations, this scheme was both generalized (in approximate form) to higher-loop orders and also extended to the timelike regime DVS00unp ; Shi00b ; Shi04 ; SS98 ; BRS00 ; Shi01 ; SS01 ; Gru97 ; GGK98 ; Mag99 ; KM01 ; KM03 ; Mag03u , encompassing previous incomplete attempts Rad96 ; KP82 in this direction, and amounting to the theoretical framework of Analytic Perturbation Theory (APT). There have been a number of parallel developments by various authors during the past several years to avoid the Landau pole using different “analytization” techniques, prime examples being Refs. KP82 ; Rad96 ; CS93 ; BB95 ; BBB95 ; DMW96 ; CMNT96 ; Gru97 ; GGK98 ; Magn00 ; BRS00 ; Gar01 ; Nes03 ; NP04 ; ANP05 .
Two other major challenges, connected with—first—the implementation of the Shirkov–Solovtsov “analytization” to three-point functions beyond the leading order of perturbation theory and—second—the extension to non-integer powers of the coupling, remained open, or at least partially open. Indeed, in the first case, extensive analyses SSK99 ; SSK00 ; BPSS04 have shown that the “analytization” principle has to be generalized to accommodate a second scale, serving as a factorization scale, or in order to include evolution effects comprising typical logarithms to some fractional power. Technically speaking, this means to extend the assertion of analyticity from the level of the coupling (and its powers) to the level of the whole reaction amplitude. This requirement was formalized by Karanikas and Stefanis (KS) in KS01 ; Ste02 in an attempt to calculate power corrections to the pion form factor and the Drell–Yan process. The systematic development of a perturbative expansion in terms of fractional powers of the coupling—the second major challenge—is the goal of the present investigation, the main focus being placed on the methodology towards improving perturbative higher-order calculations in QCD. This goal has been accomplished and will be described in this paper. A specific application of the KS “analytization” principle to the pion’s electromagnetic form factor at NLO accuracy is given in fully worked out detail in BKS05 . Other applications will follow in future publications in conjunction with the inclusion of heavy-flavor thresholds and the extension to the timelike regime.
The paper is organized as follows. In Sec. II we first review the key features of the original Analytic Perturbation Theory of Shirkov and Solovtsov, highlighting those properties pertaining to the generalization of the approach to fractional powers of the coupling. The actual extension of the approach to fractional—in fact, real—powers of the coupling is performed in Sec. III. This section describes in three subsections the new “analytization” technique, based on the Laplace transform, the verification of the analytic properties of the obtained results, and the way to include products of powers of the coupling with powers of logarithms. Moreover, we provide here approximate expressions for two-loop quantities in terms of one-loop analytic-coupling images and their index derivatives that can be extremely useful in practical calculations. Section IV is devoted to the validation of the developed theoretical framework of the Fractional Analytic Perturbation Theory (FAPT) and includes a table where we collect the algorithmic rules to connect the new analytic framework to the standard QCD perturbative power-series expansion. Our conclusions are drawn in Sec. V, while important technical details are presented in three appendices.
## II Original Analytic Perturbation Theory
In the analytic perturbation-theory approach of Shirkov and Solovtsov, the power-series expansion in the running coupling is given up in favor of a non-power series (functional) expansion. This can be written generically in terms of numerical coefficients $`d_m`$ in the following way MS97 ; Shi98
$`{\displaystyle \underset{m}{}}d_ma_{(l)}^m(Q^2){\displaystyle \underset{m}{}}d_m𝒜_m^{(l)}(Q^2),`$ (1)
where the “normalized” coupling $`a=b_0\alpha _s/(4\pi )`$ ($`b_0`$ is the first coefficient of the QCD $`\beta `$-function—see Appendix B) has been introduced instead of $`\alpha _s`$ in order to simplify intermediate calculations and because then the analytic coupling $`𝒜_1`$ is bounded from above by unity Shi98 . In the above expression, the superscript $`m`$ on $`a_{(l)}^m`$ appears on the left-hand side (LHS) as a power, whereas on the right-hand side (RHS) the subscript $`m`$ on $`𝒜_m^{(l)}`$ denotes the index of the functional expansion;<sup>1</sup><sup>1</sup>1In the following, a calligraphic notation is used to denote analytic images. $`(l)`$ denotes the loop order. For the sake of simplicity, we will avoid to indicate the loop-order index explicitly because we mostly work in the one-loop approximation; deviations, if needed, will be labelled by appropriate superscripts or subscripts in parentheses, like in Eq. (1). The conversion to analytic images of the coupling is achieved in terms of the functions
$`𝒜_m^{(l)}(Q^2)\left[a_{(l)}^m(Q^2)\right]_{\text{an}}`$ (2)
according to the general prescription
$`\left[f(Q^2)\right]_{\text{an}}={\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\text{Im}\left[f(\sigma )\right]}{\sigma +Q^2iϵ}}𝑑\sigma .`$ (3)
For the one-loop running coupling
$`a={\displaystyle \frac{1}{\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)}},`$ (4)
we have
$`𝒜_1(Q^2)\left[a(Q^2)\right]_{\text{an}}={\displaystyle \frac{1}{\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)}}{\displaystyle \frac{1}{Q^2/\mathrm{\Lambda }^21}}`$ (5)
and $`𝒜_1^{(1)}(0)=1`$. Employing the variable $`L=\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)`$, which naturally appears in perturbative QCD (pQCD) calculations, we can recast $`a`$ and $`𝒜`$ in terms of $`L`$ to obtain
$`a^1(L)={\displaystyle \frac{1}{L}},𝒜_1(L)={\displaystyle \frac{1}{L}}{\displaystyle \frac{1}{e^L1}}.`$ (6)
In this context, amplitudes (depending on a single scale $`Q^2`$) perturbatively expanded in terms of the powers of the running coupling map on a non-power series expansion Shi98 ; MS97 :
$`F(a)={\displaystyle \underset{n}{}}f_na^n(L)(L)={\displaystyle \underset{n}{}}f_n𝒜_n(L),`$ (7)
where $`f_n`$ are numbers in minimal subtraction renormalization schemes. By construction, the set $`\{𝒜_n\}`$ constitutes a linear space, which, however, is not equipped with the multiplication operation of its elements. Therefore, the product $`𝒜_n𝒜_m`$ has no rigorous meaning here. The standard algebra is recovered only for the main asymptotic contribution (cf. Eq. (6)) at $`L\mathrm{}`$, when $`\{𝒜_n\}\{a^n\}`$ SS97 ; Shi98 ; SS99 .
Let us now turn to the properties of this map and of the space $`\{𝒜_n\}`$. There are several points to note about them.
1. The map should have the property of isomorphism, i. e., it should conserve the linear structure of the original space:
$`a^0𝒜_01.`$ (8)
2. Renormalization-group summation leads to contributions like
$`f(a)=a^\nu ,\text{where}\nu \text{is real},`$ (9)
necessitating the introduction of the analytic images of $`f(a)`$: $`\left[f(a)\right]_{\mathrm{an}}=[a^\nu ]_{\mathrm{an}}`$. These are exactly those terms needed to supply the original linear space $`\{𝒜_n\}`$ with the completeness property as regards the differential operator with respect to the real index $`\nu `$.
3. Motivated by the typical logarithmic contributions, appearing in loop calculations in standard pQCD, we consider the “analytization” of terms of the sort
$`f(a)=`$ $`\{\begin{array}{cc}& a^\nu \mathrm{ln}(a);\hfill \\ & a^\nu L^m,m=1,2,\hfill \end{array}`$ (12)
giving rise to the corresponding analytic images
$`_\nu (a)[a^\nu \mathrm{ln}(a)]_{\mathrm{an}},_{\nu ,m}(L)[a^\nu L^m]_{\mathrm{an}}.`$ (13)
Expression (9) is universal and allows one to apply any one-loop renormalization-group results to APT. In fact, the corresponding renormalization factor $`Z`$, associated with the renormalizable quantity $`B`$, $`B(Q^2)=B(\mu ^2)Z(Q^2)/Z(\mu ^2)`$, reduces in the one-loop approximation to
$$Za^\nu (L)|_{\nu =\nu _0\gamma _0/\left(2b_0\right)},$$
where $`\gamma _0`$ is the coefficient of the one-loop anomalous dimension. Therefore, we have
$`\left[B(Q^2)\right]_{\text{an}}\left[a^\nu (L)\right]_{\text{an}}=𝒜_\nu (L).`$ (14)
The next two functions $`_\nu (a)`$ and $`_{\nu ,m}(L)`$ appear in NLO of pQCD and also in light-cone sum rules SY99 ; BMS02 and reflect the specific features of these calculations. An example of the first kind in connection with a NLO calculation of the electromagnetic pion form factor is treated in BKS05 , while the investigation of such terms in the context of light-cone sum rules will be considered in a future publication. To be more specific, we will consider below terms of the form $`\left(a_{(2)}\right)^\nu L`$ and $`\left(a_{(2)}\right)^\nu L^2`$—rather than deal with the series containing the constant coefficients $`f_n`$, given in Eq. (7).
One possible way in generalizing the presented original APT formalism to non-integer (fractional) values of the index $`\nu `$ is to construct the spectral density
$`\rho _\nu (\sigma )={\displaystyle \frac{1}{\pi }}\text{Im}\left[a^\nu (\sigma )\right]`$ (15)
for $`\nu `$. Indeed, substituting
$`a(\sigma )={\displaystyle \frac{1}{L(\sigma )i\pi }},L(\sigma )=\mathrm{ln}(\sigma /\mathrm{\Lambda }^2),`$ (16)
for the one-loop running coupling into Eq. (15), we can obtain by a straightforward calculation a closed-form expression for the spectral density in the form
$`\rho _\nu (\sigma )={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{sin}(\nu \phi )}{\left[\pi ^2+L^2(\sigma )\right]^{\nu /2}}},\phi =\mathrm{arccos}\left({\displaystyle \frac{L(\sigma )}{\sqrt{L^2(\sigma )+\pi ^2}}}\right).`$ (17)
A result similar to that has been derived in the context of Electrodynamics in the early article of Ref. Shi60 . It was later re-invented in QCD by Oehme Oeh90 and used by Magradze in Mag99 . To get now the desired analytic coupling for some fractional index, one has to insert this expression back into Eq. (3) and perform the integral numerically, loosing, alas, this way the possibility to reveal the mathematical properties of this function. Let us emphasize at this point that the extension of this procedure to the two-loop order for the first integer values of $`\nu `$ has been done in Refs. Shi98 ; SS99 ; Mag99 , while the inclusion of still higher-loops Mag00 seems feasible.<sup>2</sup><sup>2</sup>2Indeed a partial result for a few fractional $`\nu `$ values has already been obtained—Shirkov, private communication. However, this approach, based on the spectral density (15), is restricted to the specific structure of the Shirkov–Solovtsov APT.
## III Fractional Analytic Perturbation Theory
### III.1 A new generalization technique to include fractional indices
In this subsection, we formulate and outline another procedure to continue the integer index of the analytic coupling to fractional values. First, to generate higher indices at the one-loop level of the analytic images within APT, or, equivalently, higher powers of the standard running coupling within conventional pQCD, we follow Shi98 and write
$`\left({\displaystyle \genfrac{}{}{0pt}{}{𝒜_n(L)}{a^n(L)}}\right)={\displaystyle \frac{1}{(n1)!}}\left({\displaystyle \frac{d}{dL}}\right)^{n1}\left({\displaystyle \genfrac{}{}{0pt}{}{𝒜_1(L)}{a^1(L)}}\right).`$ (18)
Second, to facilitate the transition to fractional index values, it is instrumental to employ the Laplace representation of both types of couplings—the analytic, $`𝒜_1(L)`$, and the conventional one, $`a(L)`$, (both at the one-loop order)—and define at $`L>0`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{𝒜_1(L)}{a^1(L)}}\right)={\displaystyle _0^{\mathrm{}}}e^{Lt}\left({\displaystyle \genfrac{}{}{0pt}{}{\stackrel{~}{𝒜}_1(t)}{\stackrel{~}{a}_1(t)}}\right)𝑑t.`$ (19)
The advantage of this representation is that it transforms the result of a differential operator into an algebraic expression containing monomials. Then, applying Eq. (18) to (19), we get
$`𝒜_n(L)={\displaystyle \frac{1}{(n1)!}}\left({\displaystyle \frac{d}{dL}}\right)^{n1}𝒜_1(L)={\displaystyle _0^{\mathrm{}}}e^{Lt}\left[{\displaystyle \frac{t^{n1}}{(n1)!}}\stackrel{~}{𝒜}_1(t)\right]𝑑t,`$ (20)
so that we establish the correspondence
$`𝒜_n(L)\stackrel{~}{𝒜}_n(t)={\displaystyle \frac{t^{n1}}{(n1)!}}\stackrel{~}{𝒜}_1(t),`$ (21)
whereas for the case of the conventional pQCD coupling, one has the evident Laplace conjugates $`\stackrel{~}{a}_n`$
$`a^1(L){\displaystyle \frac{1}{L}}`$ $``$ $`\stackrel{~}{a}_1=1,`$ (22)
$`a^n(L)`$ $``$ $`\stackrel{~}{a}_n={\displaystyle \frac{t^{n1}}{(n1)!}}\stackrel{~}{a}_1.`$ (23)
Equation (21) enables us to generalize $`𝒜_n(L)`$ to any real index $`\nu `$. To do so, let us introduce the following definition for the Laplace conjugate $`\stackrel{~}{𝒜}_\nu (t)`$:
$`\stackrel{~}{𝒜}_\nu (t)\stackrel{def}{=}{\displaystyle \frac{t^{\nu 1}}{\mathrm{\Gamma }(\nu )}}\stackrel{~}{𝒜}_1(t).`$ (24)
At this stage of the continuation in the index $`\nu `$, we have based our considerations solely on the first relation in Eq. (20). Therefore, the Laplace conjugate (24) remains valid for any non-power perturbative expansion satisfying this relation, reiterating that this holds true at the one-loop level. To complete the generalization process, we should obtain an expression for $`\stackrel{~}{𝒜}_1(t)`$, based on Eq. (6). This gives the result
$`𝒜_1(L)={\displaystyle \frac{1}{L}}{\displaystyle \frac{1}{e^L1}}\stackrel{~}{𝒜}_1(t)=1{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\delta (tm),`$ (25)
which can be verified by a straightforward calculation.
Let us pause for a moment to make some useful remarks concerning the behavior of the two parts of Eq. (25). One should note the strong difference in the behavior of these functions with respect to the logarithmic term of standard perturbation theory, on the one hand,
$$\frac{1}{L}1,$$
and the pole remover appearing in APT, on the other,
$$\frac{1}{e^L1}\underset{m=1}{\overset{\mathrm{}}{}}\delta (tm).$$
Thus, one can define $`𝒜_\nu (L)`$ according to Eq. (20), and, then, using Eqs. (24) and (25), arrive at
$`𝒜_\nu (L)={\displaystyle _0^{\mathrm{}}}e^{Lt}{\displaystyle \frac{t^{\nu 1}}{\mathrm{\Gamma }(\nu )}}\stackrel{~}{𝒜}_1(t)𝑑t={\displaystyle \frac{1}{L^\nu }}{\displaystyle \frac{1}{\mathrm{\Gamma }(\nu )}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}e^{Lm}m^{\nu 1}.`$ (26)
The series on the RHS of the latter equality coincides with the definition of the Lerch transcendental function BE53 $`\mathrm{\Phi }(z,\nu ^{},i)`$ at $`\nu ^{}=\nu 1<0`$ for $`i=1`$, i.e.,
$`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{z^m}{m^{1\nu }}}=z\mathrm{\Phi }(z,1\nu ,1).`$ (27)
The analytic continuation of $`\mathrm{\Phi }(z,s,1)`$ in the variables $`z,s`$, adopting the notation of Batemann and Erdelýi BE53 , determines $`\mathrm{\Phi }`$ as an analytic function of the variable $`z`$ in the plane with a cut along $`(1,\mathrm{})`$ for any fixed $`s`$ (see for more details in Appendix A.<sup>3</sup><sup>3</sup>3The transcendental Lerch function $`\mathrm{\Phi }(z,s,1)`$ is included in the widespread programs “Mathematica 5” and “Maple 7”.). Finally, $`𝒜_\nu `$ in Eq. (26) can be rewritten in the form of an analytic function with respect to both variables $`\nu `$ and $`L`$; viz.,
$`𝒜_\nu (L)={\displaystyle \frac{1}{L^\nu }}{\displaystyle \frac{e^L}{\mathrm{\Gamma }(\nu )}}\mathrm{\Phi }(e^L,1\nu ,1).`$ (28)
We state here and prove in Appendix A that $`𝒜_\nu `$ is an entire function in $`\nu `$.
### III.2 Analytic properties
To assess the analytic properties of Eq. (28), it is useful to recast the Lerch function $`\mathrm{\Phi }(z,\nu ,1)`$ via (see BE53 , Eq. (1.10.14) and also Olver74 , Chapt. 8)
$`z\mathrm{\Phi }(z,\nu ,1)F(z,\nu )`$ (29)
entailing
$`𝒜_\nu (L)={\displaystyle \frac{1}{L^\nu }}{\displaystyle \frac{F(e^L,1\nu )}{\mathrm{\Gamma }(\nu )}},`$ (30)
where the first term in Eq. (30) corresponds to the standard PT, while the second one expresses the pole remover. Note that for a positive integer index, $`\nu =m2`$, one has the relation BE53
$`F(z,1m)=(1)^mF({\displaystyle \frac{1}{z}},1m),`$ (31)
so that substituting Eq. (31) in (30), one arrives at
$`𝒜_m(L)`$ $`=`$ $`(1)^m𝒜_m(L)`$ (32)
that confirms the specific symmetry relations worked out in Shi98 . From relation (32) and Eq. (65) one obtains the explicit asymptotic expression for $`𝒜_m(L)`$ at $`L\mathrm{}`$
$`𝒜_m(L\mathrm{})=(1)^m𝒜_m(|L|\mathrm{})=(1)^m/|L|^m+𝒪\left(1/|L|^m\right).`$ (33)
This estimate can be extended to any real value $`\nu >1`$ of the index $`m`$. To make the content of Eq. (30) more transparent, we display in Fig. 1(a) the graphs of the analytic coupling for indices from $`3`$ to 0 and values of $`L`$ in the range $`3`$ to 3.
Appealing to Eqs. (34) and (35) for negative values of the index $`\nu `$, one makes sure that, for $`L=1`$ (red thick dotted line above the zero line), this function is equal to unity for all integer values of the index $`\nu =m`$. On the other hand, for $`L<1`$, the value of $`𝒜_m(L)`$ depends on whether or not $`m`$ is even or odd. For even values, it is positive, whereas for odd values it is negative, therefore giving rise to oscillations shown in Fig. 1(a). Note also that for values of $`L1`$, the oscillatory behavior of the graphs for $`𝒜_\nu (L)`$ starts to be much less pronounced (red thick broken lines) because $`L^m`$ is positive for all positive values of $`m`$. The opposite behavior is exhibited for $`L<0`$, as one sees from the blue broken lines.
From Fig. 1(b), we observe that, in the region where $`𝒜_\nu (L)`$ is smaller than unity (as explicitly indicated in the figure), this function is monotonic in $`\nu `$ for $`\nu 2`$. On the other hand, in the region where $`𝒜_\nu (L)>1`$—possible only for $`\nu <1`$ and $`L<0`$—this function starts to be non-monotonic in $`\nu `$, so that there are two different points $`\nu _1`$ and $`\nu _2`$, both corresponding to the same value of $`𝒜_\nu (L)`$. Focusing on the values of $`𝒜_\nu (L)`$ for $`L>0`$, we see that all curves are monotonic in $`L`$ and are bounded by an envelope represented by $`𝒜_\nu (0)`$ (the blue thick solid line in Fig. 1(b)). If we consider only the interval of $`\nu (0,1)`$, then the monotonicity property extends also to the negative values of $`L`$.
Contrary to that case, the coupling $`𝒜_m(L)`$ oscillates in $`L`$ Shi98 for higher values of $`m>2`$. These oscillations are not visible in Fig. 1(b) because of the smallness of the corresponding amplitudes. They appear due to rather general reasons:
(i) the asymptotic conditions given by Eq. (33): $`𝒜_m(\mathrm{})=𝒜_m(\mathrm{})=0`$ for $`m2`$;
(ii) the differential relation between $`𝒜_m`$ and $`𝒜_1`$, expressed in Eq. (18).
Therefore, $`𝒜_{m+2}`$ has $`m`$ zeros in the vicinity of the former “Landau pole” ($`L=0`$) Shi98 —see Fig. 2. This property is rather unexpected from the point of view of standard power-series perturbation theory and will be discussed below in connection with Eq. (39).
This oscillation property of the coupling extends to $`𝒜_\nu (L)`$ for all real values of the index $`\nu 2`$.
To reveal the relevance of this representation for physical applications, let us now consider $`𝒜`$ for some particular values of the index $`\nu `$. For the case of a negative index, the $`𝒜_\nu `$ play the role of the “inverse powers” of $`𝒜_1`$ that may be considered as the images of $`a_s^\nu `$. Then, expression (30) can be rewritten in the form
$`𝒜_\nu (L)`$ $`=`$ $`L^\nu {\displaystyle \frac{\text{Li}\text{ν+1}(e^L)}{\mathrm{\Gamma }(\nu )}},`$ (34)
$`𝒜_m(L)=\underset{\epsilon 0}{lim}𝒜_{m+\epsilon }(L)`$ $`=`$ $`L^m,\text{at}m=0,1,2,\mathrm{},`$ (35)
where we have taken into account that for $`\nu 0`$
$`F(z,\nu )=\text{Li}\text{ν}(z),`$ (36)
with Li<sub>ν</sub> being the well-known polylogarithm function. It is worth remarking here that the “inverse powers” $`𝒜_m(L)=L^m`$ coincide with the inverse powers of the original running coupling $`a^m(L)=L^m`$.
To make explicit the properties of Eq. (30), we convert this equation into a series representation, using Eq. (62), to obtain
$`F(z,1\nu )`$ $`=`$ $`\mathrm{\Gamma }(\nu )\left[\mathrm{ln}\left({\displaystyle \frac{1}{z}}\right)\right]^\nu +{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\zeta (1\nu r){\displaystyle \frac{\mathrm{ln}^r(z)}{r!}}`$ (37)
for $`|\mathrm{ln}(z)|<2\pi `$, where $`\zeta (\nu )`$ is the Riemann $`\zeta `$-function. Now we are in the position to express $`𝒜_\nu `$ in the form of a series, i.e.,
$`𝒜_\nu (L)={\displaystyle \frac{1}{\mathrm{\Gamma }(\nu )}}{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\zeta (1\nu r){\displaystyle \frac{(L)^r}{r!}}\text{for}|L|<2\pi `$ (38)
because the “standard logarithms”, contained in both parts of expression (30), mutually cancel, as one verifies by substituting Eq. (37) into Eq. (30). Then, we can state the following important corollaries:
1. $`𝒜_\nu (0)={\displaystyle \frac{\zeta (1\nu )}{\mathrm{\Gamma }(\nu )}}`$ is an entire function of $`\nu `$.
Of particular importance are the following values:
$`𝒜_1(0)={\displaystyle \frac{1}{2}},𝒜_2(0)={\displaystyle \frac{1}{12}},𝒜_3(0)=0,𝒜_4(0)={\displaystyle \frac{1}{720}},𝒜_5(0)=0`$ (39)
that coincide with the results provided in Shi98 . Note that $`𝒜_{2n+1}(0)=0`$ for $`n1`$ is due to the property $`\zeta (2n)=0`$, while the set of $`𝒜_{2n}(0)`$ is alternating in sign BE53 . These properties illustrate the details of the coupling oscillations in the vicinity of $`L=0`$ for index values $`m>2`$. A convenient series representation of $`𝒜_m`$ for an integer index $`m`$ is presented and discussed in Appendix A (item 3).
2. Taking into account the relation $`\underset{ϵ0}{lim}{\displaystyle \frac{\zeta (1ϵr)}{\mathrm{\Gamma }(ϵ)}}=\delta _{0r}`$, (see, e.g., Eq. (66)), one can take the limit
$`𝒜_0(L)=\underset{\nu 0}{lim}𝒜_\nu (L)=1,`$ (40)
dispensing with the constraint $`|L|<2\pi `$ and proving assertion (8).
3. Equation (35) can be re-derived from representation (38) in a way similar to that described in the previous item.
The upshot of these considerations is that the linear space $`\{𝒜_n\}`$ is now completed via the inclusion of the elements $`𝒜_\nu `$ for any real values of the index $`\nu `$, so that one can take derivatives with respect to this continuous variable—dubbed ‘index derivative’.
### III.3 Analytic images of products of coupling powers and logarithms
To this point, we have considered only powers of the running coupling, adopting the viewpoint of the Shirkov–Solovtsov APT. Now we are going to consider more complicated expressions, like $`\left(a_{(l)}\right)^\nu L^m`$, where the power $`\nu `$ is a real number and the power $`m`$ is an integer, following the broader “analytization” principle of KS KS01 ; Ste02 .
To compute this image, we have first to determine the image of $`a^\nu \mathrm{ln}(a)`$, which can be rewritten as the derivative of $`a^\nu `$ with respect to $`\nu `$; viz.,
$`a^\nu \mathrm{ln}(a)={\displaystyle \frac{d}{d\nu }}a^\nu .`$ (41)
Due to the linearity of the differential operator, this derivative can be directly applied to any element of the completed space $`\{𝒜_n\}`$ to generate the corresponding image, $`\left[a^\nu \mathrm{ln}(a)\right]_{\text{an}}`$, and define
$`\left[{\displaystyle \frac{d}{d\nu }}a^\nu \right]_{\text{an}}\stackrel{def}{=}{\displaystyle \frac{d}{d\nu }}𝒜_\nu .`$ (42)
In the following, we shall employ for the sake of simplicity a special notation for the derivatives with respect to the index of the non-power expansion and define
$`𝒟^k𝒜_\nu {\displaystyle \frac{d^k}{d\nu ^k}}𝒜_\nu .`$ (43)
From Eqs. (30) and (42), we obtain
$`\left[a^\nu \mathrm{ln}(a)\right]_{\text{an}}={\displaystyle \frac{\mathrm{ln}L}{L^\nu }}{\displaystyle \frac{d}{d\nu }}\left({\displaystyle \frac{F(e^L,1\nu )}{\mathrm{\Gamma }(\nu )}}\right)`$ (44)
and taking multiple derivatives on both sides of Eq. (42), we compute the image of $`a^\nu L^m`$, like in Eq. (44). This procedure applies to any desired degree $`m`$ of such terms.
The extension to higher loops makes use of the APT expansion of higher-loop quantities in terms of one-loop ones. Before doing that, we consider first the image of $`a_{(2)}`$ on the basis of the perturbation expansion, given in Eq. (73) in conjunction with Eq. (43), to obtain the element $`𝒜_1`$ at the two-loop level:
$`𝒜_1^{(2)}(L)=𝒜_1^{(1)}+c_1𝒟𝒜_{\nu =2}^{(1)}+c_1^2\left(𝒟^2+𝒟^11\right)𝒜_{\nu =3}^{(1)}+𝒪\left(𝒟^3𝒜_{\nu =4}^{(1)}\right).`$ (45)
This formula can be readily generalized to any index $`\nu `$:
$`𝒜_\nu ^{(2)}(L)`$ $`=`$ $`𝒜_\nu ^{(1)}+c_1\nu 𝒟𝒜_{\nu +1}^{(1)}+c_1^2\nu \left[{\displaystyle \frac{\left(1+\nu \right)}{2}}𝒟^2+𝒟^11\right]𝒜_{\nu +2}^{(1)}`$ (46)
$`+𝒪\left(𝒟^3𝒜_{\nu +3}^{(1)}\right),`$
where $`c_1=b_1/b_0^2`$ is an auxiliary expansion parameter. The quality of the two-loop approximation for the lowest index (cf. Eq. (45)) and higher indices (cf. Eq. (46)) will be analyzed numerically in the next section. Here it suffices to mention that the achieved accuracy is of the order of about $`1\%`$ down to $`L=0`$.
To construct the image of $`\left(a_{(2)}\right)^\nu L`$, cf. Eq. (12), we first perform the “analytization” of Eq. (75) and then use Eq. (42) to arrive at the final expression
$`\left[\left(a_{(2)}(L)\right)^\nu L\right]_{\text{an}}_{\nu ,1}^{(2)}(L)`$ $`=`$ $`𝒜_{\nu 1}^{(2)}+c_1𝒟𝒜_\nu ^{(2)}+𝒪\left(𝒜_{\nu +1}^{(2)}\right),`$ (47)
which can be recast, by means of the one-loop analytic coupling and with the aid of Eq. (44), in the form
$`_{\nu ,1}^{(2)}(L)=𝒜_{\nu 1}^{(1)}c_1\nu \left[{\displaystyle \frac{\mathrm{ln}(L)\psi (\nu )}{L^\nu }}+\psi (\nu )𝒜_\nu ^{(1)}+{\displaystyle \frac{𝒟F(e^L,1\nu )}{\mathrm{\Gamma }(\nu )}}\right]+𝒪(𝒟^2𝒜_{\nu +1}^{(1)}).`$ (48)
The “analytization” of $`\left(a_{(2)}(L)\right)^\nu L^2`$, expressed in terms of Eq. (77), can be performed in an analogous way with the result
$`_{\nu ,2}^{(2)}(L)`$ $`=`$ $`𝒜_{\nu 2}^{(2)}+2c_1𝒟𝒜_{\nu 1}^{(2)}+c_1^2𝒟^2𝒜_\nu ^{(2)}2c_1^2𝒜_\nu ^{(2)}+𝒪\left(𝒟𝒜_{\nu +1}^{(2)}\right).`$ (49)
The one-loop approximation of this two-loop expression is given by
$`_{\nu ,2}^{(2)}(L)=𝒜_{\nu 2}^{(1)}+c_1\nu 𝒟𝒜_{\nu 1}^{(1)}+c_1^2{\displaystyle \frac{\nu ^2\nu +4}{2}}𝒟^2𝒜_\nu ^{(1)}+𝒪\left(𝒟𝒜_{\nu +1}^{(1)}\right).`$ (50)
A compilation of the required formulae to achieve the “analytization” of powers of the coupling in conjunction with logarithms at the two-loop order, is provided in Appendix B.
Up to now we have studied expressions appearing in fixed-order perturbation theory. But similar considerations apply also to resummed perturbation theory. Indeed, first attempts to apply the “analytization” procedure of Shirkov–Solovtsov were already presented in SSK99 ; SSK00 ; Ste02 . The crucial point here is how to deal with the requirement of analyticity when performing a Sudakov resummation. Because of the non-power series character of APT, resummation of (soft-gluon) logarithms does not lead to exponentiation. The latter can be retained only in the case of the so-called *naive* “analytization” BPSS04 , proposed in SSK99 ; SSK00 . The exact expression for the Sudakov factor is too complicated and too specific to be discussed in the present analysis. We, therefore, consider in Appendix C a simplified version of a ‘toy Sudakov’ factor that, nevertheless, bears the key characteristics pertaining to resummation under the assertion of analyticity.
For clarity, we compare the basic ingredients of FAPT in Table 1 with their counterparts in conventional perturbation theory and APT. More detailed expressions are shown in Table 2 in the next section.
## IV Validation of the new scheme
### IV.1 Analytic verification of the one-loop spectral density
An alternative way to derive Eq. (17) for the spectral density $`\rho _\nu `$, is to compare two different representations for $`𝒜_\nu `$: one given by the dispersion relation, Eq. (3), and the other provided by the Laplace representation, Eq. (20). Then, we get
$`𝒜_\nu (L)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\rho _\nu (\sigma )}{\sigma +Q^2}}𝑑\sigma ={\displaystyle _0^{\mathrm{}}}e^{Lt}\stackrel{~}{𝒜}_\nu (t)𝑑t.`$ (51)
Next, we make a double Borel transformation of both representations, the Laplace one and that of the dispersion integral, the aim being to extract $`\rho _\nu (\sigma )`$. This is done by applying first $`M^2\widehat{B}_{(M^2Q^2)}`$ on both sides of Eq. (51) and then employing
$`M^2\widehat{B}_{(M^2Q^2)}\left({\displaystyle \frac{1}{\sigma +Q^2}}\right)=\mathrm{exp}(\sigma /M^2),M^2\widehat{B}_{(M^2Q^2)}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{Q^2}}\right)^t={\displaystyle \frac{M^2}{\mathrm{\Gamma }(t)}}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{M^2}}\right)^t.`$ (52)
In the second step, we carry out one more Borel transformation, $`\widehat{B}_{(\mathrm{𝟏}\mathbf{/}𝝈1/M^2)}`$, to obtain
$`\rho _\nu (\sigma )={\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{\sigma }}\right)^t{\displaystyle \frac{\mathrm{sin}(\pi t)}{\pi }}\stackrel{~}{𝒜}_\nu (t)𝑑t.`$ (53)
The final step is to substitute in Eq. (53) the expression for $`\stackrel{~}{A}_\nu (t)`$, given by Eq. (24), to arrive at the final result
$`\rho _\nu (\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{\left(\pi ^2+L^2(\sigma )\right)^{\nu /2}}}\mathrm{sin}\left[\nu \mathrm{arccos}\left({\displaystyle \frac{L(\sigma )}{\sqrt{L^2(\sigma )+\pi ^2}}}\right)\right]`$ (54)
$`=`$ $`{\displaystyle \frac{1}{\left(\pi ^2+L^2(\sigma )\right)^{\nu /2}}}\mathrm{sin}\left[\nu \mathrm{arctan}\left({\displaystyle \frac{\pi }{L(\sigma )}}\right)\right]\text{for}L(\sigma )>0,`$ (55)
where $`L(\sigma )=\mathrm{ln}\left(\sigma /\mathrm{\Lambda }^2\right)`$. To gain a more complete understanding of the role of the Landau pole remover in $`\stackrel{~}{𝒜}_\nu `$, it is important to remark that it does not contribute to the spectral density, the reason being that this part is not altering the nature of the discontinuity. The latter is solely determined by the term $`1/L`$. One appreciates that expressions (55) and (17) coincide, as they should, hence establishing the equivalence between the two alternative extensions of the “analytization” procedure to fractional indices. The two-loop approximate expression for the spectral density is given in Appendix B.
### IV.2 Verification of the two-loop approximations
Now look specifically at the quality of the two-loop expansion in FAPT. In doing so, we define the following quantities with the help of an auxiliary parameter $`c_1`$ and the index derivative $`𝒟`$, (as in Eq. (45)):
* NLO, i.e., retaining terms of order $`c_1`$
$`\mathrm{\Delta }_2^{\text{FAPT}}(L)`$ $`=`$ $`1{\displaystyle \frac{𝒜_1^{(1)}(L)+c_1𝒟𝒜_{\nu =2}^{(1)}(L)}{𝒜_1^{(2)}(L)}}`$ (56)
* NNLO, i.e., retaining terms up to order $`c_1^2`$
$`\mathrm{\Delta }_3^{\text{FAPT}}(L)`$ $`=`$ $`1{\displaystyle \frac{𝒜_1^{(1)}(L)+c_1𝒟𝒜_{\nu =2}^{(1)}(L)+c_1^2\left(𝒟^2+𝒟^11\right)𝒜_{\nu =3}^{(1)}(L)}{𝒜_1^{(2)}(L)}}.`$ (57)
For the corresponding quantities within the standard QCD perturbation theory, we use Eq. (73) to obtain
* NLO, i.e., retaining terms of order $`c_1`$
$`\mathrm{\Delta }_2^{\text{PT}}(L)=1{\displaystyle \frac{a_{(1)}(L)+c_1a_{(1)}^2(L)\mathrm{ln}a_{(1)}(L)}{a_{(2)}(L)}}`$ (58)
* NNLO, i.e., retaining terms up to order $`c_1^2`$
$`\mathrm{\Delta }_3^{\text{PT}}(L)=1{\displaystyle \frac{a_{(1)}(L)+c_1a_{(1)}^2(L)\mathrm{ln}a_{(1)}(L)+c_1^2a_{(1)}^3(L)\left(\mathrm{ln}^2a_{(1)}(L)+\mathrm{ln}a_{(1)}(L)1\right)}{a_{(2)}(L)}}.`$
(59)
First, let us compare the transition from the NLO (cf. Eq. (56)) to the NNLO (cf. Eq. (57)) in FAPT (see Fig. 3(a)). One appreciates that by taking into account the NNLO terms, a significant improvement of the convergence quality of the FAPT series is achieved. Indeed, even at $`Q^2=\mathrm{\Lambda }^2`$, which corresponds to $`L=0`$, the error of truncating the FAPT series at the NLO is about 5%, while by taking into account the NNLO correction this error becomes even smaller than 0.5%. In Fig. 3(b) we show the relative quality of these approximations concerning the loop expansion between the standard perturbation theory and FAPT, as quantified by Eqs. (57) and (59). One appreciates the strong suppression of $`\mathrm{\Delta }_3^{\text{FAPT}}(L)`$ relative to its conventional analogue in the small $`L`$ region, say, below approximately $`L=2`$.
The same comparison can be realized for $`𝒜_2^{(2)}`$, using Eq. (46). Indeed, we demonstrate in Fig. 4(a) the quality of this FAPT expansion in comparison with the results of the numerical integration of the NLO spectral density $`\rho _2`$ (for more details, we refer to Appendix B and SS99 ) in the dispersion-integral representation, provided by Eq. (3). In this graphics, we also display the results obtained numerically by Magradze in Mag03u . The message from Fig. 4(a) is quite clear. Our analytic (solid line) and our numerical calculation (dashed line) are in mutual support, while the results of Mag03u differ considerably with respect to both the magnitude and the trend of the negative values of $`L`$.
The good convergence of the proposed series for $`𝒜_{1,2}^{(2)}`$ (Eq. (46)), that had been demonstrated above, can be traced to the basis of APT. Indeed, this non-power expansion of the quantities $`𝒜_{1,2}^{(2)}`$ in terms of $`𝒜_m^{(1)}`$ has a finite radius of convergence, the reasons being discussed in Ref. Mik04 .
To give the reader an impression of the dependence on $`L`$ of $`𝒜_\nu ^{(2)}(L)`$, we show in Fig. 4(b) a comparison of this quantity with its counterpart in standard QCD, namely, $`\left[a_{(2)}(L)\right]^\nu `$. For the purpose of illustration, we select the value $`\nu =\gamma _2/(2b_0)0.62`$, which corresponds to the 1-loop evolution exponent of the non-singlet quark operator of index 2, entering a number of applications in DIS and also various exclusive reactions BMS02 ; BMS03 ; BMS01 .
In support of our two-loop approximation (within FAPT), we display in Fig. 5 results for the analytic images $`𝒜_\nu ^{(2)}`$ with $`\nu =k/4`$ and $`k=0,1,\mathrm{},8`$. We observe the same monotonic pattern, i.e., no crossing, of curves, found already for the one-loop case, shown in Fig. 1(b). We investigated numerically the range of negative values of $`L`$ and found that crossing appears only for indices $`\nu <1`$, again in close analogy to the one-loop case.
The pivotal results of this paper are collected in Table 2, where we provide the reader with explicit calculational rules to connect the standard QCD perturbation theory with FAPT. We stress that the presented algorithm has broad applications in phenomenology and can play a major role in the perturbative analysis of observables both in inclusive and exclusive QCD reactions.
## V Conclusions
With hindsight we can say that the requirement of analyticity at the amplitude level of hadronic quantities in QCD, expressed by Karanikas and Stefanis KS01 , is instrumental in improving perturbation-theory calculations. First, as shown in an accompanying paper by two of us (A.P.B. and N.G.S.) together with Karanikas BKS05 , it enables one to minimize both the sensitivity on the renormalization scheme and scale setting and also the dependence on the factorization scale. The reason for this latter advantage is that it includes into the “analytization” procedure not only the powers of the running coupling, but also logarithms (or exponentials) that may contain the momentum scale with respect to which analyticity is required. Second, starting at this point, we have shown in the present work that invoking this analyticity principle gives rise to a generalization of the original APT to fractional powers of the coupling. As a bonus, this approach improves the convergence of the perturbative expansion significantly—see Fig. 3.
Our main goal in this analysis was to work out in mathematical detail the procedure for determining analytic expressions for any real power of the running coupling and delineate the main results. To keep our presentation as general as possible, we have purposefully refrained from considering specific examples and concentrated instead on generic features and expressions. In this vein, we have discussed products of the running coupling (or its powers) with (powers of) logarithms, which are typical for contributions encountered in higher-order corrections of QCD perturbation theory, or when taking into account evolution effects via the renormalization-group equation. Similar logarithmic terms also appear in calculations employing light cone sum rules SY99 ; BMS02 . In an analogous way, we have discussed the resummation of non-power series in the analytic images in order to capture the key features of Sudakov resummation of soft-gluon effects (see Appendix C). All these elements of FATP, required for further applications of this formalism to improve the calculation of any hadronic amplitude at the two-loop level, are collected in Table 2. In this context, we mention that we have developed approximate expressions for the two-loop analytic images in terms of one-loop quantities that can facilitate practical computations significantly.
In conclusion, this report has emphasized rigorous methods rather than specific applications. A first example of the present framework is discussed in BKS05 , focusing on the topic of the renormalization-scheme and factorization-scale independence of the electromagnetic pion form factor, relative to its treatment within the standard QCD perturbation theory or original APT. The methods presented here are intended to be used in the low-to-medium momentum range where the standard perturbative approach faces the problem of the Landau pole and in processes or under circumstances where the original APT is insufficient because it is tied to integer powers of the coupling. We believe that our assortment of analytic expressions for a variety of expressions ranging from any real powers of the coupling to more complicated products containing logarithms, provides sufficient evidence for the usefulness of the approach for higher-order perturbative calculations.
###### Acknowledgements.
We wish to thank D. V. Shirkov for valuable discussions and comments. Two of us (A.P.B. and S.V.M.) are indebted to Prof. Klaus Goeke for the warm hospitality at Bochum University, where the major part of this investigation was carried out. This work was supported in part by the Deutsche Forschungsgemeinschaft, the Heisenberg–Landau Programme (grant 2005), and the Russian Foundation for Fundamental Research (grants No. 03-02-16816, 03-02-04022 and 05-01-00992).
## Appendix A Analytic properties of $`𝚽\mathbf{(}𝒛\mathbf{,}\mathbf{}𝝂\mathbf{,}\mathrm{𝟏}\mathbf{)}`$ and $`𝓐_𝝂`$
1. The function $`\mathrm{\Phi }(z,\nu ,i)`$ can be determined by means of the analytic continuation of the series BE53
$`\mathrm{\Phi }(z,s,i)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{z^m}{(m+i)^s}}\text{for }|z|<1,s1,2,3,\mathrm{}`$ (60)
in both variables $`z`$ and $`s`$. This analytic continuation for every fixed $`s`$, which is not a positive integer, determines $`\mathrm{\Phi }`$ as an analytic function of $`z`$, regular in the plane with a cut along the axis $`(1,\mathrm{})`$, and for every fixed $`z`$ in the cut plane as an analytic function of $`s`$ being regular, except possibly at the points $`s=1,2,3,\mathrm{},`$ as was mentioned in BE53 . We can improve this statement for $`s1`$ using Eq. (36) to obtain
$`\mathrm{\Phi }(z,s,i)`$ $`=`$ $`{\displaystyle \frac{1}{z^i}}\left[z\mathrm{\Phi }(z,s,1){\displaystyle \underset{m1}{\overset{i1}{}}}{\displaystyle \frac{z^m}{m^s}}\right]={\displaystyle \frac{1}{z^i}}\left[\text{Li}\text{s}(z){\displaystyle \underset{m1}{\overset{i1}{}}}{\displaystyle \frac{z^m}{m^s}}\right].`$ (61)
One appreciates that there are no singularities for any positive integer values of $`s`$. Hence, we can conclude that $`\mathrm{\Phi }(z,s,1)`$ is an analytic function in $`s`$ at any fixed $`z`$ on the cut plane. Moreover, the function $`𝒜_\nu (L)`$, see expression (28),
$$𝒜_\nu (L)=\frac{1}{L^\nu }\frac{e^L\mathrm{\Phi }(e^L,1\nu ,1)}{\mathrm{\Gamma }(\nu )},$$
has no poles in $`\nu `$ and is, therefore, an entire function in $`\nu `$.
2. There is a useful series representation for $`\mathrm{\Phi }(z,s,1)`$ BE53 (cf. (37)); viz.,
$`\mathrm{\Phi }(z,s,1)`$ $`=`$ $`{\displaystyle \frac{1}{z}}\left[\mathrm{\Gamma }(1s)\left[\mathrm{ln}\left({\displaystyle \frac{1}{z}}\right)\right]^{s1}+{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}\zeta (sr){\displaystyle \frac{\mathrm{ln}^r(z)}{r!}}\right]`$ (62)
for $`|\mathrm{ln}(z)|<2\pi `$, $`s1,2,3,\mathrm{}`$ that allows one to continue $`\mathrm{\Phi }(z,s,1)`$ for integer positive $`s=m`$ values by means of the limit $`s=m+\epsilon ,\epsilon 0`$ in Eq. (62). To take this limit, one should expand in $`\epsilon `$ the first term in the square brackets in expression (62), which is proportional to $`\mathrm{\Gamma }(1m\epsilon )`$. The other singular term appears in the sum and is proportional to $`\zeta (1+\epsilon )=1/\epsilon \psi (1)+𝒪(\epsilon )`$. The singularities, contained in both these parts, mutually cancel.
3. The expansion of $`𝒜_\nu (L)`$, Eq. (38), is simplified for an integer index $`\nu =m1`$ to read
$`𝒜_m(L)={\displaystyle \frac{1}{(m1)!}}{\displaystyle \underset{r=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{B_{m+r}}{(m+r)r!}}(L)^r\text{for}|L|<2\pi ,`$ (63)
where $`B_m`$ are the Bernoulli numbers. From the property $`B_{2n+1}=0`$ it follows that $`𝒜_{2n}(L)`$ is an even function of its argument, while $`𝒜_{2n+1}(L)`$ is an odd one.
Note here that the pole remover $`F(e^L,1m)`$ in expression (30) reduces to elementary functions for the case of an integer index. Indeed, according to Eq. (18), the operator $`d/dL`$ shifts the second argument of the function $`F`$ by unity, i.e., $`mm+1`$, to get
$$\frac{d}{dL}F(e^L,m)=F(e^L,(m+1)).$$
Taking into account that $`F(e^L,0)=(e^L1)^1`$ and applying the previous relation $`m`$ times we arrive at
$`F(e^L,m)=\left({\displaystyle \frac{d}{dL}}\right)^mF(e^L,0)=\left({\displaystyle \frac{d}{dL}}\right)^m{\displaystyle \frac{1}{e^L1}}.`$ (64)
This representation leads to an exponentially suppressed asymptotic limit for the function $`F(e^L,m)`$; viz.,
$`F(e^L,m)|_{L\pm \mathrm{}}e^{|L|}.`$ (65)
4. We supply here the Lindelöf formula WW27
$`\zeta (\nu )={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{\nu 1}}+{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\mathrm{sin}\left[\nu \mathrm{arctan}(t)\right]dt}{(1+t^2)^{\nu /2}(e^{2\pi t}1)}},`$ (66)
which fixes $`\zeta (\nu )`$ as an analytic function with a simple pole at $`\nu =1`$. This representation for $`\zeta (\nu )`$ has been used in Sec. III.2.
5. Now we are in the position to supply also the analytic images of the coupling in the timelike regime for $`L(s)\mathrm{log}(s/\mathrm{\Lambda }^2)0`$, employing the notation of DVS00 ; Shi01 :
$`𝔄_\nu (s)`$ $`=`$ $`{\displaystyle _s^{\mathrm{}}}{\displaystyle \frac{d\sigma }{\sigma }}\rho _\nu (\sigma )={\displaystyle \frac{1}{\pi }}{\displaystyle _{L(s)}^{\mathrm{}}}𝑑L{\displaystyle \frac{\mathrm{sin}\left[\nu \mathrm{arctan}\left(\pi /L\right)\right]}{\left(\pi ^2+L^2\right)^{\nu /2}}}.`$ (67)
This integral can be evaluated to provide a result analogous to $`𝒜_\nu (L)`$ for the spacelike regime; namely,
$`𝔄_\nu (s)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\left[(\nu 1)\mathrm{arctan}\left(\pi /L(s)\right)\right]}{\pi (\nu 1)\left(\pi ^2+L(s)^2\right)^{(\nu 1)/2}}}.`$ (68)
A similar expression for the timelike coupling has been obtained before in BKM01 using the “contour-improved resummation technique”.
## Appendix B “Analytization” of powers of the coupling multiplied by logarithms
1. The expansion of the $`\beta `$-function in the NLO approximation is given by
$`{\displaystyle \frac{d}{dL}}\left({\displaystyle \frac{\alpha _s(L)}{4\pi }}\right)=b_0\left({\displaystyle \frac{\alpha _s(L)}{4\pi }}\right)^2b_1\left({\displaystyle \frac{\alpha _s(L)}{4\pi }}\right)^3,`$ (69)
where $`L=\mathrm{ln}(\mu ^2/\mathrm{\Lambda }^2)`$ and
$`b_0={\displaystyle \frac{11}{3}}C_\text{A}{\displaystyle \frac{4}{3}}T_\text{R}N_f,b_1={\displaystyle \frac{34}{3}}C_\text{A}^2\left(4C_\text{F}+{\displaystyle \frac{20}{3}}C_\text{A}\right)T_\text{R}N_f`$ (70)
with $`C_\mathrm{F}=\left(N_\mathrm{c}^21\right)/2N_\mathrm{c}=4/3`$, $`C_\mathrm{A}=N_\mathrm{c}=3`$, $`T_\text{R}=1/2`$, and $`N_f`$ denoting the number of flavors. Then, the corresponding two-loop equation for our coupling $`a=b_0\alpha /(4\pi )`$ looks like
$`{\displaystyle \frac{da_{(2)}}{dL}}=a_{(2)}^2(L)\left[1+c_1a_{(2)}(L)\right]\text{with}c_1{\displaystyle \frac{b_1}{b_0^2}}.`$ (71)
The renormalization-group solution of this equation assumes the form
$`{\displaystyle \frac{1}{a_{(2)}}}+c_1\mathrm{ln}\left[{\displaystyle \frac{a_{(2)}}{1+c_1a_{(2)}}}\right]=L.`$ (72)
Then, for the expansion of $`a_{(2)}(L)`$ in terms of $`a_{(1)}(L)=1/L`$ we have, retaining terms of the order $`a_{(1)}^3`$,
$`a_{(2)}`$ $`=`$ $`a_{(1)}+c_1a_{(1)}^2\mathrm{ln}a_{(1)}+c_1^2a_{(1)}^3\left(\mathrm{ln}^2a_{(1)}+\mathrm{ln}a_{(1)}1\right)+𝒪(a_{(1)}^4\mathrm{ln}^3(a_1)).`$ (73)
2. Now, for the product $`\left[a_{(2)}\right]^\nu L`$, we obtain from (72)
$`\left(a_{(2)}\right)^\nu L=\left(a_{(2)}\right)^{\nu 1}+\left(a_{(2)}\right)^\nu c_1\mathrm{ln}\left[{\displaystyle \frac{a_{(2)}}{1+c_1a_{(2)}}}\right].`$ (74)
Expanding the logarithmic term $`\mathrm{ln}[1+c_1a_{(2)}]`$, while retaining terms of order $`a_{(2)}^{\nu 1}`$, $`a_{(2)}^\nu \mathrm{ln}(a_{(2)})`$; viz.,
$`\left(a_{(2)}\right)^\nu L`$ $`=`$ $`\left(a_{(2)}\right)^{\nu 1}+c_1\left(a_{(2)}\right)^\nu \mathrm{ln}(a_{(2)})𝒪(a_{(2)}^{\nu +1})`$ (75)
and, finally, expanding the coupling $`a_{(2)}`$ in terms $`a=a_{(1)}`$, we find
$`\left(a_{(2)}\right)^\nu L`$ $`=`$ $`a^{\nu 1}+\nu a^\nu c_1\mathrm{ln}(a)+𝒪(a^{\nu +1}\mathrm{ln}^2(a)).`$ (76)
Calculating $`\left(a_{(2)}\right)^\nu L^2`$ in an analogous way, we derive
$`\left(a_{(2)}\right)^\nu L^2`$ $`=`$ $`\left(a_{(2)}\right)^{\nu 2}\left[1+c_1a_{(2)}\mathrm{ln}\left(a_{(2)}\right)\right]^22c_1^2a_{(2)}^\nu 𝒪(a_{(2)}^{\nu +1}\mathrm{ln}(a_{(2)}))`$ (77)
$`=`$ $`a^{\nu 2}+\nu a^{\nu 1}c_1\mathrm{ln}(a)+\left({\displaystyle \frac{\nu ^2\nu +4}{2}}\right)a^\nu c_1^2\mathrm{ln}^2(a)+𝒪(a^\nu \mathrm{ln}(a)).`$
3. We consider here the spectral density $`\rho _\nu (\sigma )`$ beyond the leading-order approximation. At the $`l`$-loop level, $`\rho _\nu ^{(l)}(\sigma )`$ can always be presented in the same form as for the leading order, given in Eq. (17), i.e.
$`\rho _\nu ^{(l)}(\sigma )={\displaystyle \frac{1}{\pi }}\text{Im}\left[a_{(l)}^\nu (\sigma )\right]={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{sin}[\nu \phi _{(l)}(\sigma )]}{\left(R_{(l)}(\sigma )\right)^\nu }},`$ (78)
but keeping in mind that the phase $`\phi _{(l)}`$ and the radial part $`R_{(l)}`$ have a multi-loop content. At the two-loop level, one should, strictly speaking, deal with the imaginary part of the Lambert function $`W_1`$ (see Mag99 ) because the exact solution of Eq. (72) can be realized in terms of the Lambert function. Instead of following this procedure, we can alternatively take the well-known first-iteration solution of Eq. (72) that provides us with sufficient accuracy the following result:
$`{\displaystyle \frac{1}{a_{(2)}(L)}}{\displaystyle \frac{1}{a_{(2)}^{\text{iter}}(L)}}=L+c_1\mathrm{ln}\left(L+c_1\right).`$ (79)
For this approximate solution $`a_{(2)}^{\text{iter}}`$, we have
$`\left(R_{(2)}(\sigma )\right)^2`$ $`=`$ $`\left(L(\sigma )+c_1\mathrm{ln}\left(\sqrt{(L(\sigma )+c_1)^2+\pi ^2}\right)\right)^2+\left(\pi +c_1\varphi (L(\sigma ))\right)^2,`$ (80)
$`\phi _{(2)}(\sigma )`$ $`=`$ $`\mathrm{arccos}\left[{\displaystyle \frac{L(\sigma )+c_1\mathrm{ln}\left(\sqrt{(L(\sigma )+c_1)^2+\pi ^2}\right)}{R_{(2)}(\sigma )}}\right],`$ (81)
$`\varphi (L(\sigma ))`$ $`=`$ $`\mathrm{arccos}\left[{\displaystyle \frac{L(\sigma )+c_1}{\sqrt{(L(\sigma )+c_1)^2+\pi ^2}}}\right]`$ (82)
with $`L(\sigma )=\mathrm{ln}\left(\sigma /\mathrm{\Lambda }^2\right)`$. The spectral density $`\rho _{\nu =1}^{\text{(2)-iter}}(\sigma )`$ with the phase and the radial part from Eqs. (80)–(82) appears to be very close to the numerical, but exact result for $`\rho _1^{\text{(2)}}(\sigma )`$, based on $`W_1`$—see, e.g., SS99 .
## Appendix C “Analytization” of the toy model for Sudakov resummation
Here we discuss the analytic image of expression $`F_\text{S}(x,L)\mathrm{exp}\left[xa_{(l)}(L)\right]`$, which originates as a part of the procedure of the Sudakov resummation, where $`x`$ is a free parameter. We consider the following example, serving as a “toy model” for this resummation:
$`\left\{F_\text{S}(x,L)\right\}_{\mathrm{an}}\left\{1+{\displaystyle \underset{m=1}{}}{\displaystyle \frac{\left[xa_{(l)}(L)\right]^m}{m!}}\right\}_{\mathrm{an}}=1+{\displaystyle \underset{m=1}{}}(x)^m{\displaystyle \frac{𝒜_m^{(l)}(L)}{m!}}.`$ (83)
One can verify that for the asymptotic limits of $`L`$, Eq. (83) reduces to the evident forms:
$`1+{\displaystyle \underset{m=1}{}}(x)^m{\displaystyle \frac{𝒜_m^{(l)}(L)}{m!}}=`$ $`\{\begin{array}{ccc}1x& \text{for}& L\mathrm{};\hfill \\ 1& \text{for}& L+\mathrm{}.\hfill \end{array}`$ (86)
The first asymptote on the RHS of Eq. (86) appears due to the equality $`𝒜_m^{(l)}(\mathrm{})=\delta _{1m}`$, see, for instance, Shi98 . The second asymptote is due to the property that in the UV regime APT reduces to the standard perturbation theory. Both properties are illustrated in Fig. 6(a). In Fig. 6(b) we compare different versions of the toy Sudakov function, obtained within the present framework (solid line), using the *naive* “analytization” SSK99 ; SSK00 ; BPSS04 (dashed line), and conventional QCD perturbation theory (dotted line).
On the other hand, restricting the loop order to $`l=1`$, one can derive an explicit expression for Eq. (83) in the region of $`L>0`$, which is based on the Laplace representation given by Eq. (20) and expression (25) for $`\stackrel{~}{𝒜}_1`$, namely,
$`F_\text{S}^{\text{(an)}}(x,L)=e^{x/L}+\sqrt{x}{\displaystyle \underset{m=1}{}}e^{Lm}{\displaystyle \frac{J_1(2\sqrt{xm})}{\sqrt{m}}},L>0.`$ (87)
The perturbative part of $`\stackrel{~}{𝒜}_m`$ in (83) reproduces exactly the asymptotic expression $`\mathrm{exp}\left[xa_{(1)}(L)\right]`$ on the RHS of Eq. (87), while the pole remover generates the sum of the exponents in $`L`$, weighted by the Bessel functions, $`J_1`$, exhibiting how the large $`L`$ behavior is violated.
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# 2005 International Linear Collider Workshop - Stanford, U.S.A. Frequency Scanned Interferometry for ILC Tracker Alignment
## I Introduction
The motivation for this project is to design a novel optical system for quasi-real time alignment of tracker detector elements used in High Energy Physics (HEP) experiments. A.F. Fox-Murphy et.al. from Oxford University reported their design of a frequency scanned interferometer (FSI) for precise alignment of the ATLAS Inner Detector ox98 ; ox04 . Given the demonstrated need for improvements in detector performance, we plan to design an enhanced FSI system to be used for the alignment of tracker elements in the next generation of electron positron Linear Collider detectors. Current plans for future detectors require a spatial resolution for signals from a tracker detector, such as a silicon microstrip or silicon drift detector, to be approximately 7-10 $`\mu m`$orangebook . To achieve this required spatial resolution, the measurement precision of absolute distance changes of tracker elements in one dimension should be on the order of 1 $`\mu m`$. Simultaneous measurements from hundreds of interferometers will be used to determine the 3-dimensional positions of the tracker elements.
The University of Michigan group constructed two demonstration Frequency Scanned Interferometer (FSI) systems with laser beams transported by air or single-mode optical fiber in the laboratory for initial feasibility studies. Absolute distance was determined by counting the interference fringes produced while scanning the laser frequencystone99 . The main goal of the demonstration systems was to determine the potential accuracy of absolute distance measurements that could be achieved under controlled conditions. Secondary goals included estimating the effects of vibrations and studying error sources crucial to the absolute distance accuracy. The main contents of this proceedings article come from our published paperfsi05 . However, new material in this paper includes a description of a dual-laser system and a possible optical alignment for a silicon tracker detector.
## II Principles
The intensity $`I`$ of any two-beam interferometer can be expressed as
$$\begin{array}{c}I=I_1+I_2+2\sqrt{I_1I_2}\mathrm{cos}(\varphi _1\varphi _2)\end{array}$$
$`(1)`$
where $`I_1`$ and $`I_2`$ are the intensities of the two combined beams, and $`\varphi _1`$ and $`\varphi _2`$ are the phases. Assuming the optical path lengths of the two beams are $`L_1`$ and $`L_2`$, the phase difference in Eq. (1) is $`\mathrm{\Phi }=\varphi _1\varphi _2=2\pi |L_1L_2|(\nu /c)`$, where $`\nu `$ is the optical frequency of the laser beam, and c is the speed of light.
For a fixed path interferometer, as the frequency of the laser is continuously scanned, the optical beams will constructively and destructively interfere, causing “fringes”. The number of fringes $`\mathrm{\Delta }N`$ is
$$\begin{array}{c}\mathrm{\Delta }N=|L_1L_2|(\mathrm{\Delta }\nu /c)=L\mathrm{\Delta }\nu /c\end{array}$$
$`(2)`$
where $`L`$ is the optical path difference between the two beams, and $`\mathrm{\Delta }\nu `$ is the scanned frequency range. The optical path difference (OPD for absolute distance between beamsplitter and retroreflector) can be determined by counting interference fringes while scanning the laser frequency.
## III Demonstration System of FSI
A schematic of the FSI system with a pair of optical fibers is shown in Fig.1. The light source is a New Focus Velocity 6308 tunable laser (665.1 nm $`<\lambda <`$ 675.2 nm). A high-finesse ($`>200`$) Thorlabs SA200 F-P is used to measure the frequency range scanned by the laser. The free spectral range (FSR) of two adjacent F-P peaks is 1.5 GHz, which corresponds to 0.002 nm. A Faraday Isolator was used to reject light reflected back into the lasing cavity. The laser beam was coupled into a single-mode optical fiber with a fiber coupler. Data acquisition is based on a National Instruments DAQ card capable of simultaneously sampling 4 channels at a rate of 5 MS/s/ch with a precision of 12-bits. Omega thermistors with a tolerance of 0.02 K and a precision of 0.01 $`mK`$ are used to monitor temperature. The apparatus is supported on a damped Newport optical table.
In order to reduce air flow and temperature fluctuations, a transparent plastic box was constructed on top of the optical table. PVC pipes were installed to shield the volume of air surrounding the laser beam. Inside the PVC pipes, the typical standard deviation of 20 temperature measurements was about $`0.5mK`$. Temperature fluctuations were suppressed by a factor of approximately 100 by employing the plastic box and PVC pipes.
Detectors for HEP experiments must usually be operated remotely for safety reasons because of intensive radiation, high voltage or strong magnetic fields. In addition, precise tracking elements are typically surrounded by other detector components, making access difficult. For practical HEP application of FSI, optical fibers for light delivery and return are therefore necessary.
The beam intensity coupled into the return optical fiber is very weak, requiring ultra-sensitive photodetectors for detection. Considering the limited laser beam intensity and the need to split into many beams to serve a set of interferometers, it is vital to increase the geometrical efficiency. To this end, a collimator is built by placing an optical fiber in a ferrule (1mm diameter) and gluing one end of the optical fiber to a GRIN lens. The GRIN lens is a 0.25 pitch lens with 0.46 numerical aperture, 1 mm diameter and 2.58 mm length which is optimized for a wavelength of 630nm. The density of the outgoing beam from the optical fiber is increased by a factor of approximately 1000 by using a GRIN lens. The return beams are received by another optical fiber and amplified by a Si femtowatt photoreceiver with a gain of $`2\times 10^{10}V/A`$.
## IV Multiple-Distance-Measurement Techniques
For a FSI system, drifts and vibrations occurring along the optical path during the scan will be magnified by a factor of $`\mathrm{\Omega }=\nu /\mathrm{\Delta }\nu `$, where $`\nu `$ is the average optical frequency of the laser beam and $`\mathrm{\Delta }\nu `$ is the scanned frequency range. For the full scan of our laser, $`\mathrm{\Omega }67`$. Small vibrations and drift errors that have negligible effects for many optical applications may have a significant impact on a FSI system. A single-frequency vibration may be expressed as $`x_{vib}(t)=a_{vib}\mathrm{cos}(2\pi f_{vib}t+\varphi _{vib})`$, where $`a_{vib}`$, $`f_{vib}`$ and $`\varphi _{vib}`$ are the amplitude, frequency and phase of the vibration, respectively. If $`t_0`$ is the start time of the scan, Eq. (2) can be re-written as
$$\begin{array}{c}\mathrm{\Delta }N=L\mathrm{\Delta }\nu /c+2[x_{vib}(t)\nu (t)x_{vib}(t_0)\nu (t_0)]/c\end{array}$$
$`(3)`$
If we approximate $`\nu (t)\nu (t_0)=\nu `$, the measured optical path difference $`L_{meas}`$ may be expressed as
$$\begin{array}{c}L_{meas}=L_{true}4a_{vib}\mathrm{\Omega }\mathrm{sin}[\pi f_{vib}(tt_0)]\times \\ \mathrm{sin}[\pi f_{vib}(t+t_0)+\varphi _{vib}]\end{array}$$
$`(4)`$
where $`L_{true}`$ is the true optical path difference in the absence of vibrations. If the path-averaged refractive index of ambient air $`\overline{n}_g`$ is known, the measured distance is $`R_{meas}=L_{meas}/(2\overline{n}_g)`$.
If the measurement window size $`(tt_0)`$ is fixed and the window used to measure a set of $`R_{meas}`$ is sequentially shifted, the effects of the vibration will be evident. We use a set of distance measurements in one scan by successively shifting the fixed-length measurement window one F-P peak forward each time. The arithmetic average of all measured $`R_{meas}`$ values in one scan is taken to be the measured distance of the scan (although more sophisticated fitting methods can be used to extract the central value). For a large number of distance measurements $`N_{meas}`$, the vibration effects can be greatly suppressed. Of course, statistical uncertainties from fringe and frequency determination, dominant in our current system, can also be reduced with multiple scans. Averaging multiple measurements in one scan, however, provides similar precision improvement to averaging distance measurements from independent scans, and is faster, more efficient, and less susceptible to systematic errors from drift. In this way, we can improve the distance accuracy dramatically if there are no significant drift errors during one scan, caused, for example, by temperature variation. This multiple-distance-measurement technique is called ’slip measurement window with fixed size’, shown in Fig.2. However, there is a trade off in that the thermal drift error is increased with the increase of $`N_{meas}`$ because of the larger magnification factor $`\mathrm{\Omega }`$ for a smaller measurement window size.
In order to extract the amplitude and frequency of the vibration, another multiple-distance-measurement technique called ’slip measurement window with fixed start point’ is used, as shown in Fig.2. In Eq. (3), if $`t_0`$ is fixed, the measurement window size is enlarged one F-P peak for each shift, an oscillation of a set of measured $`R_{meas}`$ values indicates the amplitude and frequency of vibration. This technique is not suitable for distance measurement because there always exists an initial bias term, from $`t_0`$, which cannot be determined accurately in our current system.
## V Absolute Distance and Vibration Measurement
The typical measurement residual versus the distance measurement number in one scan using the above technique is shown in Fig.3(a), where the scanning rate was 0.5 nm/s and the sampling rate was 125 kS/s. Measured distances minus their average value for 10 sequential scans are plotted versus number of measurements ($`N_{meas}`$) per scan in Fig.3(b). The standard deviations (RMS) of distance measurements for 10 sequential scans are plotted versus number of measurements ($`N_{meas}`$) per scan in Fig.3(c). It can be seen that the distance errors decrease with an increase of $`N_{meas}`$. The RMS of measured distances for 10 sequential scans is 1.6 $`\mu m`$ if there is only one distance measurement per scan ($`N_{meas}=1`$). If $`N_{meas}=1200`$ and the average value of 1200 distance measurements in each scan is considered as the final measured distance of the scan, the RMS of the final measured distances for 10 scans is 41 nm for the distance of 449828.965 $`\mu m`$, the relative distance measurement precision is 91 ppb.
The standard deviation (RMS) of measured distances for 10 sequential scans is approximately 1.5 $`\mu m`$ if there is only one distance measurement per scan for closed box data. By using the multiple-distance-measurement technique, the distance measurement precisions for various closed box data with distances ranging from 10 cm to 70 cm collected in the past year are improved significantly; precisions of approximately 50 nanometers are demonstrated under laboratory conditions, as shown in Table 1. All measured precisions listed in Table 1. are the RMS’s of measured distances for 10 sequential scans. Two FSI demonstration systems, ’air FSI’ and ’optical fiber FSI’, are constructed for extensive tests of multiple-distance-measurement technique, ’air FSI’ means FSI with the laser beam transported entirely in the ambient atmosphere, ’optical fiber FSI’ represents FSI with the laser beam delivered to the interferometer and received back by single-mode optical fibers.
Based on our studies, the slow fluctuations are reduced to a negligible level by using the plastic box and PVC pipes to suppress temperature fluctuations. The dominant error comes from the uncertainties of the interference fringes number determination; the fringes uncertainties are uncorrelated for multiple distance measurements. In this case, averaging multiple distance measurements in one scan provides a similar precision improvement to averaging distance measurements from multiple independent scans. But, for open box data, the slow fluctuations are dominant, on the order of few microns in our laboratory. The measurement precisions for single and multiple distance open-box measurements are comparable, which indicates that the slow fluctuations cannot be adequately suppressed by using the multiple-distance-measurement technique. A dual-laser FSI systemcoe2001 intended to cancel the drift error is currently under study in our laboratory .
In order to test the vibration measurement technique, a piezoelectric transducer (PZT) was employed to produce vibrations of the retroreflector. For instance, the frequency of the controlled vibration source was set to $`1.01\pm 0.01`$ Hz with amplitude $`9.5\pm 1.5`$ nanometers. The magnification factors, distance measurement residuals and corrected measurement residuals for 2000 measurements in one scan are shown in Fig.3(d), Fig.3(e) and Fig.3(f), respectively. The extracted vibration frequencies and amplitudes using this technique, $`f_{vib}=1.025\pm 0.002`$ Hz, $`A_{vib}=9.3\pm 0.3`$ nanometers, agree well with the expectation values.
In addition, vibration frequencies at 0.1, 0.5, 1.0, 5, 10, 20, 50, 100 Hz with controlled vibration amplitudes ranging from 9.5 nanometers to 400 nanometers were studied extensively using our current FSI system. The measured vibrations and expected vibrations all agree well within the 10-15% level for amplitudes, 1-2% for frequencies, where we are limited by uncertainties in the expectations. Vibration frequencies far below 0.1 Hz can be regarded as slow fluctuations, which cannot be suppressed by the above analysis techniques.
Detailed information about estimation of major error sources for the absolute distance measurement and limitation of our current FSI system is provided elsewherefsi05 .
## VI Dual-Laser FSI System
A dual-laser FSI system has been built in order to reduce drift error and slow fluctuations occuring during the laser scan. Two lasers are operating simultaneously, the two laser beams are coupled into one optical fiber but isolated by using two choppers. The principle of the dual-laser technique is shown in the following. For the first laser, the measured distance $`D_1=D_{true}+\mathrm{\Omega }_1\times ϵ_1`$, and $`ϵ`$ is drift error during the laser scanning. For the second laser, the measured distance $`D_2=D_{true}+\mathrm{\Omega }_2\times ϵ_2`$. Since the two laser beams travel the same optical path during the same period, the drift errors $`ϵ_1`$ and $`ϵ_2`$ should be very comparable. Under this assumption, the true distance can be extracted using the formula $`D_{true}=(D_2\rho \times D_1)/(1\rho )`$, where, $`\rho =\mathrm{\Omega }_2/\mathrm{\Omega }_1`$, the ratio of magnification factors from two lasers.
The laser beams are isolated by choppers periodically, so only half the fringes are recorded for each laser, degrading the distance measurement precision. Missing fringes during chopped intervals for each laser must be recovered through robust interpolation algorithms. The chopper edge transitions make this interpolation difficult. Several techniques are under study.
## VII A Possible Silicon Tracker Alignment System
One possible silicon tracker alignment system is shown in Fig.4. The left plot shows lines of sight for alignment in R-Z plane of the tracker barrel, the middle plot for alignment in X-Y plane of the tracker barrel, the right plot for alignment in the tracker forward region. Red lines/dots show the point-to-point distances need to be measured using FSIs. There are 752 point-to-point distance measurements in total for the alignment system. More studies are needed to optimize the distance measurments grid.
###### Acknowledgements.
This work is supported by the National Science Foundation and the Department of Energy of the United States.
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# Combinatorial triangulations of homology spheres
## 1 Introduction and results
All the simplicial complexes considered in this paper are finite. We say that a simplicial complex $`K`$ triangulates a topological space $`X`$ (or $`K`$ is a triangulation of $`X`$) if $`X`$ is homeomorphic to the geometric carrier $`|K|`$ of $`K`$.
The vertex-set of a simplicial complex $`K`$ is denoted by $`V(K)`$. If $`K`$, $`L`$ are two simplicial complexes, then a simplicial isomorphism from $`K`$ to $`L`$ is a bijection $`\pi :V(K)V(L)`$ such that for $`\sigma V(K)`$, $`\sigma `$ is a face of $`K`$ if and only if $`\pi (\sigma )`$ is a face of $`L`$. The complexes $`K`$, $`L`$ are called (simplicially) isomorphic when such an isomorphism exists. We identify two simplicial complexes if they are isomorphic.
A simplicial complex $`K`$ is called pure if all the maximal faces of $`K`$ have the same dimension. A maximal face in a pure simplicial complex is also called a facet.
If $`\sigma `$ is a face of a simplicial complex $`K`$ then the link of $`\sigma `$ in $`K`$, denoted by $`\mathrm{Lk}_K(\sigma )`$ (or simply by $`\mathrm{Lk}(\sigma )`$), is by definition the simplicial complex whose faces are the faces $`\tau `$ of $`K`$ such that $`\tau `$ is disjoint from $`\sigma `$ and $`\sigma \tau `$ is a face of $`K`$.
A subcomplex $`L`$ of a simplicial complex $`K`$ is called an induced (or full ) subcomplex of $`K`$ if $`\sigma K`$ and $`\sigma V(L)`$ imply $`\sigma L`$. The induced subcomplex of $`K`$ on the vertex set $`U`$ is denoted by $`K[U]`$.
For a commutative ring $`R`$, a simplicial complex $`K`$ is called $`R`$-acyclic if $`|K|`$ is $`R`$-acyclic, i.e., $`\stackrel{~}{H}_q(|K|,R)=0`$ for all $`q0`$ (where $`\stackrel{~}{H}^q(|K|,R)`$ denotes the reduced homology).
By a subdivision of a simplicial complex $`K`$ we mean a simplicial complex $`K^{}`$ together with a homeomorphism from $`|K^{}|`$ onto $`|K|`$ which is facewise linear. Two simplicial complexes $`K`$ and $`L`$ are called combinatorially equivalent (denoted by $`KL`$) if they have isomorphic subdivisions. So, $`KL`$ if and only if $`|K|`$ and $`|L|`$ are piecewise-linear (pl) homeomorphic (see ).
For a set $`U`$ with $`d+1`$ elements, let $`K`$ be the simplicial complex whose faces are all the non-empty subsets of $`U`$. Then $`K`$ triangulates the $`d`$-dimensional closed unit ball. This complex is called the standard $`d`$-ball and is denoted by $`\mathrm{\Delta }_{d+1}^d(U)`$ or simply by $`\mathrm{\Delta }_{d+1}^d`$. A polyhedron is called a pl $`d`$-ball if it is pl homeomorphic to $`|\mathrm{\Delta }_{d+1}^d|`$. A simplicial complex $`X`$ is called a combinatorial $`d`$-ball if it is combinatorially equivalent to $`\mathrm{\Delta }_{d+1}^d`$. So, $`X`$ is a combinatorial $`d`$-ball if and only if $`|X|`$ is a pl $`d`$-ball.
For a set $`V`$ with $`d+2`$ elements, let $`S`$ be the simplicial complex whose faces are all the non-empty proper subsets of $`V`$. Then $`S`$ triangulates the $`d`$-sphere. This complex is called the standard $`d`$-sphere and is denoted by $`S_{d+2}^d(V)`$ or simply by $`S_{d+2}^d`$. A polyhedron is called a pl $`d`$-sphere if it is pl homeomorphic to $`|S_{d+2}^d|`$. A simplicial complex $`X`$ is called a combinatorial $`d`$-sphere if it is combinatorially equivalent to $`S_{d+2}^d`$. So, $`X`$ is a combinatorial $`d`$-sphere if and only if $`|X|`$ is a pl $`d`$-sphere.
A simplicial complex $`K`$ is called a combinatorial $`d`$-manifold if the link of each vertex is a combinatorial $`(d1)`$-sphere. A simplicial complex $`K`$ is a combinatorial $`d`$-manifold if and only if $`|K|`$ is a closed pl $`d`$-manifold (see ).
If a triangulation $`K`$ of a space $`X`$ is a combinatorial manifold then $`K`$ is called a combinatorial triangulation of $`X`$. If $`K`$ is a triangulation of a 3-manifold then the link of a vertex is a triangulation of the 2-sphere and all triangulations of the 2-sphere are combinatorial 2-spheres. So, any triangulation of a 3-manifold is a combinatorial triangulation.
Let $`\tau \sigma `$ be two faces of a simplicial complex $`K`$. We say that $`\tau `$ is a free face of $`\sigma `$ if $`\sigma `$ is the only face of $`K`$ which properly contains $`\tau `$. (It follows that $`dim(\sigma )dim(\tau )=1`$ and $`\sigma `$ is a maximal simplex in $`K`$.) If $`\tau `$ is a free face of $`\sigma `$ then $`K^{}:=K\{\tau ,\sigma \}`$ is a simplicial complex. We say that there is an elementary collapse of $`K`$ to $`K^{}`$. We say $`K`$ collapses to $`L`$ and write $`K\stackrel{^\mathrm{s}}{}L`$ if there exists a sequence $`K=K_0`$, $`K_1,\mathrm{}`$, $`K_n=L`$ of simplicial complexes such that there is an elementary collapse of $`K_{i1}`$ to $`K_i`$ for $`1in`$ (see ). If $`L`$ consists of a 0-simplex (a point) we say that $`K`$ is collapsible and write $`K\stackrel{^\mathrm{s}}{}0`$. Clearly, if $`K\stackrel{^\mathrm{s}}{}L`$ then $`|K||L|`$ as polyhedra and hence $`|K|`$ and $`|L|`$ have the same homotopy type (see ). So, if a simplicial complex $`K`$ is collapsible then $`|K|`$ is contractible and hence, in particular, $`K`$ is $`\text{}_2`$-acyclic. Here we prove :
###### Theorem 1
. If a $`\text{}_2`$-acyclic simplicial complex has $`7`$ vertices then it is collapsible.
As an application of Theorem 1, we prove our main result - a recognition theorem for combinatorial spheres :
###### Theorem 2
. Let $`M`$ be an $`n`$-vertex combinatorial triangulation of a $`\text{}_2`$-homology $`d`$-sphere. Suppose $`M`$ has an $`m`$-vertex combinatorial $`d`$-ball as an induced subcomplex, where $`nm+7`$. Then $`M`$ is a combinatorial sphere.
In consequence we get the following.
###### Corollary 3
. Let $`M`$ be an $`n`$-vertex combinatorial $`d`$-manifold. If $`|M|`$ is a $`\text{}_2`$-homology sphere and $`nd+8`$ then $`M`$ is a combinatorial sphere.
###### Corollary 4
. Let $`M`$ be a $`(d+9)`$-vertex combinatorial triangulation of a $`\text{}_2`$-homology $`d`$\- sphere. If $`M`$ is not a combinatorial sphere then $`M`$ can not admit any bistellar $`i`$-move for $`i<d`$.
Since by the universal coefficient theorem any integral homology sphere is a $`\text{}_2`$-homology sphere, Theorem 2, Corollary 3 and Corollary 4 remain true if we replace $`\text{}_2`$-homology by integral homology in the hypothesis. In particular, we have :
###### Corollary 5
. Let $`M`$ be an $`n`$-vertex combinatorial triangulation of an integral homology $`d`$-sphere.
1. If $`nd+8`$ then $`M`$ is a combinatorial sphere.
2. If $`n=d+9`$ and $`M`$ is not a combinatorial sphere then $`M`$ can not admit any bistellar $`i`$-move for $`i<d`$.
###### Remark 1
. Corollary 3 is clearly trivial for $`d2`$. In , Brehm and Kühnel proved that any $`n`$-vertex combinatorial $`d`$-manifold is a combinatorial $`d`$-sphere if $`n<3d/2+3`$ and it is either a combinatorial $`d`$-sphere or a cohomology projective plane if $`n=3d/2+3`$. So, Corollary 3 has new content only for $`3d8`$.
###### Remark 2
. Another result in says that any $`n`$-vertex combinatorial $`d`$-manifold is simply connected for $`n2d+2`$. Since a simply connected integral homology sphere is a sphere for $`d3`$, and since for $`d4`$ all combinatorial triangulations of $`d`$-spheres are combinatorial spheres, this result implies that all combinatorial triangulations of integral homology $`d`$-spheres ($`d3,4`$) with $`2d+2`$ vertices are combinatorial spheres. This is stronger than Corollary 5 $`(a)`$ for $`d6`$. Thus Corollary 5 $`(a)`$ has new content only for $`d=3,4,5`$.
###### Remark 3
. In \[8, p. 35\], Lutz presented a 12-vertex combinatorial triangulation of the lens space $`L(3,1)`$. (It is mentioned in \[7, p. 79\] that Brehm obtained a 12-vertex combinatorial triangulation of $`L(3,1)`$ earlier.) Since $`L(3,1)`$ is a $`\text{}_2`$-homology $`3`$-sphere ($`H_1(L(3,1),\text{})=\text{}_3`$, $`H_2(L(3,1),\text{})=0`$), Corollary 3 is sharp for $`d=3`$.
It follows from Corollary 3 that 12 is the least number of vertices required to triangulate $`L(3,1)`$. It follows from Corollary 4 that a 12-vertex combinatorial triangulation of $`L(3,1)`$ can not admit any bistellar $`i`$-move for $`0i2`$.
###### Remark 4
. Recall that the Dunce Hat is the topological space obtained from the solid triangle $`abc`$ by identifying the oriented edges $`\stackrel{}{ab}`$, $`\stackrel{}{bc}`$ and $`\stackrel{}{ac}`$. The following is a triangulation of the Dunce Hat using 8 vertices.
Since this example is contractible but not collapsible, it follows that the bound 7 in Theorem 1 is best possible.
###### Remark 5
. Let $`H^{\mathrm{\hspace{0.17em}3}}`$ be the non-orientable 3-manifold obtained from $`S^{\mathrm{\hspace{0.17em}2}}\times [0,1]`$ by identifying $`(x,0)`$ with $`(x,1)`$. It follows from works of Walkup \[14, Theorems 3, 4\] that if $`K`$ is a combinatorial $`3`$-manifold and $`|K|`$ is not homeomorphic to $`S^{\mathrm{\hspace{0.17em}3}}`$, $`S^{\mathrm{\hspace{0.17em}2}}\times S^{\mathrm{\hspace{0.17em}1}}`$ or $`H^{\mathrm{\hspace{0.17em}3}}`$ then $`f_1(K)4f_0(K)+8`$ and hence $`f_0(K)11`$. Thus if $`M`$ ($`S^{\mathrm{\hspace{0.17em}3}}`$) is a $`\text{}_2`$-homology 3-sphere then at least 11 vertices are needed for any combinatorial triangulation of $`M`$. Now, Corollary 3 implies that at least 12 vertices are needed. In , Björner and Lutz have presented a 16-vertex combinatorial triangulation of the Poincaré homology 3-sphere.
In , we have shown that all combinatorial triangulations of $`S^{\mathrm{\hspace{0.17em}4}}`$ with at most 10 vertices are combinatorial $`4`$-spheres. Now, Corollary 3 implies that all combinatorial triangulations of $`S^{\mathrm{\hspace{0.17em}4}}`$ with at most 12 vertices are combinatorial spheres. So, any combinatorial triangulation (if it exists) of $`S^{\mathrm{\hspace{0.17em}4}}`$ which is not a combinatorial sphere requires at least 13 vertices.
###### Remark 6
. The conclusion in Corollary 4 (namely, that certain combinatorial manifolds do not admit any proper bistellar move) appears to be a strong structural restriction. We owe to F. H. Lutz the information that the smallest known combinatorial sphere (other than a standard sphere) not admitting any proper bistellar move is a 16-vertex 3-sphere.
## 2 Preliminaries and Definitions.
For a simplicial complex $`K`$, the maximum $`k`$ such that $`K`$ has a $`k`$-face is called the dimension of $`K`$. An 1-dimensional simplicial complex is called a graph. A simplicial complex $`K`$ is called connected if $`|K|`$ is connected.
For $`i=1,2,3`$, the $`i`$-faces of a simplicial complex are also called the edges, triangles and tetrahedra of the complex, respectively. For a face $`\sigma `$ in a simplicial complex $`K`$, the number of vertices in $`\mathrm{Lk}_K(\sigma )`$ is called the degree of $`\sigma `$ in $`K`$ and is denoted by $`\mathrm{deg}_K(\sigma )`$.
If the number of $`i`$-simplices of a $`d`$-dimensional simplicial complex $`K`$ is $`f_i(K)`$, then the vector $`f=(f_0,\mathrm{},f_d)`$ is called the $`f`$-vector of $`K`$ and the number $`\chi (K):=_{i=0}^d(1)^if_i(K)`$ is called the Euler characteristic of $`K`$. If $`f_{k1}=\left(\genfrac{}{}{0pt}{}{f_0}{k}\right)`$ then $`K`$ is called $`k`$-neighbourly.
For two simplicial complexes $`K`$, $`L`$ with disjoint vertex sets, the join $`KL`$ is the simplicial complex $`KL\{\sigma \tau :\sigma K,\tau L\}`$.
If $`K`$ is a $`d`$-dimensional simplicial complex then define the pure part of $`K`$ as the simplicial complex whose simplices are the sub-simplices of the $`d`$-simplices of $`K`$.
A $`d`$-dimensional pure simplicial complex $`K`$ is called a weak pseudomanifold if each $`(d1)`$-face is contained in exactly two facets of $`K`$. A $`d`$-dimensional weak pseudomanifold $`K`$ is called a pseudomanifold if for any pair $`\tau `$, $`\sigma `$ of facets, there exists a sequence $`\tau =\tau _0,\mathrm{},\tau _n=\sigma `$ of facets of $`K`$, such that $`\tau _{i1}\tau _i`$ is a $`(d1)`$-simplex of $`K`$ for $`1in`$. In other words, a weak pseudomanifold is a pseudomanifold if and only if it does not have any weak pseudomanifold of the same dimension as a proper subcomplex. Clearly, any connected combinatorial manifold is a pseudomanifold.
For $`n3`$, the $`n`$-vertex combinatorial 1-sphere ($`n`$-cycle) is the unique $`n`$-vertex 1-dimensional pseudomanifold and is denoted by $`S_n^{\mathrm{\hspace{0.17em}1}}`$.
A $`d`$-dimensional pure simplicial complex $`K`$ is called a weak pseudomanifold with boundary if each $`(d1)`$-face is contained in 1 or 2 facets of $`K`$ and there exists a $`(d1)`$-face of degree 1. The boundary $`K`$ of $`K`$ is by definition the pure simplicial complex whose facets are the degree one $`(d1)`$-faces of $`K`$.
A simplicial complex $`K`$ is called a combinatorial $`d`$-manifold with boundary if the link of each vertex is either a combinatorial $`(d1)`$-sphere or a combinatorial $`(d1)`$-ball and there exists a vertex whose link is a combinatorial $`(d1)`$-ball. A simplicial complex $`K`$ is a combinatorial $`d`$-manifold with boundary if and only if $`|K|`$ is a compact pl $`d`$-manifold with non-empty boundary. Clearly, if $`K`$ is a combinatorial $`d`$-manifold with boundary then $`K\mathrm{}`$ and $`\mathrm{Lk}_K(v)=(\mathrm{Lk}_K(v))`$, for $`vV(K)`$. Therefore, $`K`$ is a combinatorial $`(d1)`$-manifold. Clearly, if $`K`$ is a combinatorial $`d`$-ball ($`d>0`$) then $`K`$ is a combinatorial $`d`$-manifold with boundary and $`K`$ is a combinatorial $`(d1)`$-sphere.
###### Example 1
. Some weak pseudomanifolds on 6 or 7 vertices.
$`\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}1}},\mathrm{},\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}5}}`$ are combinatorial spheres. $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$ triangulates the real projective plane. $`\mathrm{{\rm Y}}_1`$, $`\mathrm{{\rm Y}}_2`$ are the smallest examples of weak pseudomanifolds which are not pseudomanifolds.
The following results (which we need later) follow from the classification of all 2-dimensional weak pseudomanifolds on $`7`$ vertices (e.g., see ).
###### Proposition 2.1
. Let $`K`$ be an $`n`$-vertex $`2`$-dimensional weak pseudomanifold. If $`n6`$ then $`K`$ is isomorphic to $`S_4^{\mathrm{\hspace{0.17em}2}}`$, $`S_3^{\mathrm{\hspace{0.17em}1}}S_2^{\mathrm{\hspace{0.17em}0}}`$, $`S_2^{\mathrm{\hspace{0.17em}0}}S_2^{\mathrm{\hspace{0.17em}0}}S_2^{\mathrm{\hspace{0.17em}0}}`$, $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$ or $`\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}1}}`$ above.
###### Proposition 2.2
. Let $`K`$ be a $`7`$-vertex $`2`$-dimensional weak pseudomanifold. If the number of facets of $`K`$ is $`10`$ then $`K`$ is isomorphic to $`S_5^{\mathrm{\hspace{0.17em}1}}S_2^{\mathrm{\hspace{0.17em}0}}`$, $`\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}2}},\mathrm{},\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}5}}`$, $`\mathrm{{\rm Y}}_1`$ or $`\mathrm{{\rm Y}}_2`$ above.
Let $`X`$ be a pure simplicial complex of dimension $`d1`$. Let $`A`$ be a set of size $`d+2`$ such that $`A`$ contains at least one and at most $`d+1`$ facets of $`X`$. (It follows that all except at most one element of $`A`$ are vertices of $`X`$.) Define the pure $`d`$-dimensional simplicial complex $`\kappa _A(X)`$ as follows. The facets of $`\kappa _A(X)`$ are (i) the facets of $`X`$ not contained in $`A`$ and (ii) the $`(d+1)`$-subsets of $`A`$ which are not facets of $`X`$. $`\kappa _A`$ is said to be a generalized bistellar move. Clearly $`\kappa _A(\kappa _A(X))=X`$. Let $`\beta =\{xA:A\{x\}X\}`$ and $`\alpha =A\beta `$. Then $`\alpha X`$ and $`\beta \kappa _A(X)`$. The set $`\beta `$ is called the core of $`A`$. If $`\alpha `$ is an $`i`$-simplex of $`X`$ then $`\kappa _A`$ is also called a generalized bistellar $`i`$-move. Observe that if $`d`$ is even and $`\kappa _A`$ is a generalized bistellar $`(d/2)`$-move then $`f_d(\kappa _A(X))=f_d(X)`$.
Now suppose $`X`$ is a weak pseudomanifold, and $`A`$, $`\alpha `$ and $`\beta `$ are as above. Notice that (a) either $`\alpha `$ is a $`d`$-simplex in $`X`$ or $`V(\mathrm{Lk}_X(\alpha ))\beta `$ and (b) if $`\beta X`$ then $`\mathrm{Lk}_{\kappa _A(X)}(\beta )=\mathrm{Lk}_X(\beta )S_{i+1}^{i1}(\alpha )S_{i+1}^{i1}(\alpha )`$ (and therefore $`\kappa _A(X)`$ is not a combinatorial manifold even if $`X`$ is so). We shall say that $`\kappa _A`$ is a bistellar move if (bs1) $`\beta X`$ and (bs2) either $`\alpha `$ is a $`d`$-simplex in $`X`$ or $`V(\mathrm{Lk}_X(\alpha ))=\beta `$ (and hence $`\mathrm{Lk}_X(\alpha )`$ is the standard sphere on the vertex set $`\beta `$). If $`1id1`$ then a bistellar $`i`$-move is called a proper bistellar move. Observe that if $`X`$ is a combinatorial $`d`$-manifold then (bs2) holds for any $`(d+2)`$-subset $`A`$. If a generalized bistellar move is not a bistellar move then it is called singular.
Two weak pseudomanifolds are called bistellar equivalent if there exists a finite sequence of bistellar moves leading from one to the other. Let $`\kappa _A`$ be a bistellar move on $`X`$. If $`X_1`$ is obtained from $`X`$ by starring () a new vertex in $`\alpha `$ and $`X_2`$ is obtained from $`\kappa _A(X)`$ by starring a new vertex in $`\beta `$ then $`X_1`$ and $`X_2`$ are isomorphic. Thus if $`X`$ and $`Y`$ are bistellar equivalent then $`XY`$. In , Pachner proved the following : Two combinatorial manifolds are bistellar equivalent if and only if they are combinatorially equivalent.
###### Example 2
. Let the notations be as in Example 1.
1. Let $`A=\{1,2,5,6\}V(\text{}P_6^{\mathrm{\hspace{0.17em}2}})`$. Put $`R=\kappa _A(\text{}P_6^{\mathrm{\hspace{0.17em}2}})`$. Then $`R`$ is not a weak pseudomanifold. Observe that (bs1) is not satisfied here and hence $`\kappa _A`$ is a singular bistellar move. Note that the automorphism group $`A_5`$ of $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$ is transitive on the $`4`$-subsets of its vertex set. In consequence, all singular bistellar 1-moves on $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$ yield isomorphic simplicial complexes.
2. Let $`B=\{2,3,6,7\}V(\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}2}})`$. Then $`\kappa _B(\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}2}})`$ is the union of two spheres with one common edge $`67`$. Here (bs1) is not satisfied.
3. Let $`C=\{1,2,3,6\}V(\mathrm{{\rm Y}}_1)`$. Then $`\kappa _C(\mathrm{{\rm Y}}_1)=\mathrm{{\rm Y}}_2`$. Here also (bs1) is not satisfied and $`\kappa _C(\mathrm{{\rm Y}}_1)\mathrm{{\rm Y}}_1`$ but $`\kappa _C(\mathrm{{\rm Y}}_1)`$ is a weak pseudomanifold.
4. Let $`D=\{1,2,3,6\}V(\mathrm{{\rm Y}}_2)`$. Then $`\kappa _D(\mathrm{{\rm Y}}_2)=\mathrm{{\rm Y}}_1`$. Here (bs2) is not satisfied.
5. If $`E=\{2,3,4,6\}V(\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}4}})`$ then $`\kappa _E(\mathrm{\Sigma }_4)`$ is a $`7`$-vertex pseudomanifold with 12 facets. In this case, (bs1) is not satisfied.
6. Let $`F=\{2,3,4,6\}V(\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}2}})`$. Then $`\kappa _F`$ is a bistellar move and $`\kappa _F(\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}2}})=\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}3}}`$.
Let $`LK`$ be simplicial complexes. The simplicial neighbourhood of $`L`$ in $`K`$ is the subcomplex $`N(L,K)`$ of $`K`$ whose maximal simplices are those maximal simplices of $`K`$ which intersect $`V(L)`$. Clearly, $`N(L,K)`$ is the smallest subcomplex of $`K`$ whose geometric carrier is a topological neighbourhood of $`|L|`$ in $`|K|`$. The induced subcomplex $`C(L,K)`$ on the vertex-set $`V(K)V(L)`$ is called the simplicial complement of $`L`$ in $`K`$.
Suppose $`P^{}P`$ are polyhedra and $`P=P^{}B`$, where $`B`$ is a pl $`k`$-ball (for some $`k1`$). If $`P^{}B`$ is a pl $`(k1)`$-ball then we say that there is an elementary collapse of $`P`$ to $`P^{}`$. We say that $`P`$ collapses to $`Q`$ and write $`PQ`$ if there exists a sequence $`P=P_0,P_1,\mathrm{},P_n=Q`$ of polyhedra such that there is an elementary collapse of $`P_{i1}`$ to $`P_i`$ for $`1in`$. If $`Q`$ is a point we say that $`P`$ is collapsible and write $`P0`$. For two simplicial complexes $`K`$ and $`L`$, if $`K\stackrel{^\mathrm{s}}{}L`$ then clearly $`|K||L|`$. A regular neighbourhood of a polyhedron $`P`$ in a pl $`d`$-manifold $`M`$ is a $`d`$-dimensional submanifold $`W`$ with boundary such that $`WP`$ and $`W`$ is a neighbourhood of $`P`$ in $`M`$. The following is a direct consequence of the Simplicial Neighbourhood Theorem (\[11, Theorem 3.11\]).
###### Proposition 2.3
. Let $`K`$ be a combinatorial $`d`$-manifold with boundary. Suppose $`K`$ is an induced subcomplex of $`K`$. Let $`L`$ be the simplicial complement of $`K`$ in $`K`$. Then $`|K||L|`$.
Proof. Let $`M`$ be a pl d-manifold such that $`|K|`$ is in the interior of $`M`$ (we can always find such $`M`$, e.g., one such $`M`$ can be obtained from $`|K|(|K|\times [0,1])`$ by identifying $`(x,0)`$ with $`x|K|`$).
Since $`L=C(K,K)`$, $`|L||K||K|`$ and hence $`|K|`$ is a neighbourhood of $`|L|`$ in $`\mathrm{int}(M)`$. Again, since $`L`$ is the simplicial complement of $`K`$ in $`K`$ and $`K`$ is an induced subcomplex of $`K`$, $`C(L,K)=K`$. Finally, since $`K`$ is an induced subcomplex of dimension $`d1`$, each $`d`$-simplex of $`K`$ intersects $`V(L)`$. This implies that $`N(L,K)=K`$.
Let $`P=|L|`$, $`A=|K|`$ and $`J=K`$. Then $`A=|K|`$ and $`N^{^{^{}}}(L,K):=N(L,K)C(L,K)=J`$. Thus (i) $`P`$ is a compact polyhedron in the interior of the pl manifold $`M`$, (ii) $`A`$ is a neighbourhood of $`P`$ in $`\mathrm{int}(M)`$, (iii) $`A`$ is a compact pl manifold with boundary and (iv) $`(K,L,J)`$ are triangulations of $`(A,P,A)`$ where $`L`$ is an induced subcomplex of $`K`$, $`K=N(L,K)`$ and $`J=N^{^{^{}}}(L,K)`$. Then, by the Simplicial Neighbourhood Theorem, $`A`$ is a regular neighbourhood of $`P`$. Hence $`AP`$. $`\mathrm{}`$
We need the following well-known results (see \[11, Lemma 1.10, Corollaries 3.13, 3.28\]) later.
###### Proposition 2.4
. Let $`B`$, $`D`$ be pl $`d`$-balls and $`h:BD`$ a pl homeomorphism. Then $`h`$ extends to a pl homeomorphism $`h_1:BD`$.
###### Proposition 2.5
. Let $`S`$ be a pl $`d`$-sphere. If $`BS`$ is a pl $`d`$-ball then the closure of $`SB`$ is a pl $`d`$-ball.
###### Proposition 2.6
. A collapsible pl manifold with boundary is a pl ball.
Question . Is it true that under the hypothesis of Proposition 2.3, we have $`K\stackrel{^\mathrm{s}}{}L`$ ?
## 3 $`\text{}_2`$-acyclic simplicial complexes.
In this section we prove Theorem 1.
###### Lemma 3.1
. Let $`X`$ be a $`7`$-vertex simplicial complex. Suppose $`(a)`$ $`X`$ is $`\text{}_2`$-acyclic, $`(b)`$ $`X`$ is not collapsible, and $`(c)`$ $`X`$ is minimal subject to $`(a)`$ and $`(b)`$ $`(`$i.e., $`X`$ has no proper subcomplex satisfying $`(a)`$ and $`(b))`$. Then $`X`$ is pure of dimension $`d=2`$ or $`3`$ and each $`(d1)`$-face of $`X`$ occurs in at least two facets.
Proof. Notice that, because of the minimality assumption, $`X`$ has no free face. Clearly, $`dim(X)5`$, since otherwise $`X`$ is a combinatorial ball. Suppose $`dim(X)=5`$. By minimality, each 4-face of $`X`$ is in 0 or $`2`$ facets. Since $`X`$ has 7 vertices, it follows that each 4-face is in 0 or 2 facets. Therefore the pure part $`Y`$ of $`X`$ is a 7-vertex 5-dimensional weak pseudomanifold and hence $`Y=S_7^{\mathrm{\hspace{0.17em}5}}X`$. Then $`H_5(X,\text{}_2)0`$, a contradiction. Thus $`dim(X)4`$.
Suppose, if possible, $`dim(X)=4`$. Let $`Y`$ be the pure part of $`X`$. Then, each 3-face of $`Y`$ occurs in at least two facets. If $`\mathrm{\#}(V(Y))6`$, then $`Y=S_6^{\mathrm{\hspace{0.17em}4}}`$ and hence $`H_4(X,\text{}_2)0`$, a contradiction. Thus $`V(Y)=V(X)`$ has size 7. Define a binary relation $``$ on $`V(Y)`$ by $`y_1y_2`$ if $`V(Y)\{y_1,y_2\}`$ is not a facet of $`Y`$. Since each 3-face of $`Y`$ is in at least two facets, it follows that $``$ is an equivalence relation with at least two equivalence classes. Therefore either there is an equivalence class $`W`$ of size 6 or else we can write $`V(Y)=V_1V_2`$, where $`V_1`$, $`V_2`$ are unions of $``$-classes and $`\mathrm{\#}(V_1)2`$, $`\mathrm{\#}(V_2)2`$. In consequence $`Y`$ (and hence $`X`$) contains a 4-sphere as a subcomplex : the standard sphere on $`W`$ or the join of the standard spheres on $`V_1`$ and $`V_2`$. Therefore $`H_4(X,\text{}_2)0`$, a contradiction. Thus $`dim(X)3`$.
If $`dim(X)=1`$ then $`X`$ is a $`\text{}_2`$-acyclic connected graph and hence is a tree. But any tree has end vertices and hence is collapsible, a contradiction. So, $`dim(X)=2`$ or 3.
Since $`\stackrel{~}{H}_0(X,\text{}_2)=0`$, $`X`$ is connected. Since $`X`$ has no free vertex, it follows that each vertex of $`X`$ is in at least two edges.
Next we show that $`X`$ has no maximal edge. Suppose, on the contrary, $`X`$ has a maximal edge $`e`$. Then $`Y:=X\{e\}`$ is a subcomplex of $`X`$. We claim that $`Y`$ is disconnected. If not, then there is a subcomplex $`K=S_n^{\mathrm{\hspace{0.17em}1}}`$ of $`X`$ containing the edge $`e`$. The formal sum of the edges in $`K`$ is an 1-cycle over $`\text{}_2`$ which is not a boundary since it involves the maximal edge $`e`$. Hence $`H_1(X,\text{}_2)0`$, a contradiction. So, $`Y`$ is disconnected. Since each vertex of $`X`$ is in at least two edges, it follows that each component of $`Y`$ has $`3`$ vertices. Since $`X`$ has seven vertices, it follows that some component of $`Y`$ has exactly three vertices and contains an $`S_3^{\mathrm{\hspace{0.17em}1}}`$. If these three vertices span a 2-face then its edges are free in $`X`$, contradicting minimality. In the remaining case $`X`$ has an induced $`S_3^{\mathrm{\hspace{0.17em}1}}`$ whose edges are maximal, contradicting $`\text{}_2`$-acyclicity of $`X`$.
In case $`dim(X)=2`$, this shows that $`X`$ is pure. In case $`dim(X)=3`$, we proceed to show that $`X`$ has no maximal 2-face, proving that it is pure in that case too.
Suppose, on the contrary, that $`dim(X)=3`$ and $`X`$ has a maximal 2-face $`\mathrm{\Delta }=abc`$. Let’s say that an edge of $`X`$ is good if it is in a tetrahedron of $`X`$, and call it bad otherwise. First suppose that all three edges in $`\mathrm{\Delta }`$ are good. Since $`X`$ has no free triangle, each vertex in the link of an edge has degree 0 or $`2`$ and hence there are at least three vertices of degree $`2`$ in the link of a good edge. Since $`\mathrm{\Delta }`$ is maximal, it follows that the link of each of the three edges in $`\mathrm{\Delta }`$ has $`3`$ vertices outside $`\mathrm{\Delta }`$. Since, there are only four vertices outside $`\mathrm{\Delta }`$, it follows from the pigeonhole principle that there is a common vertex $`x`$ outside $`\mathrm{\Delta }`$ which occurs in the link of all three edges in $`\mathrm{\Delta }`$. Hence $`S_4^{\mathrm{\hspace{0.17em}2}}(\mathrm{\Delta }\{x\})`$ is a subcomplex of $`X`$. The sum of the four triangles in this $`S_4^{\mathrm{\hspace{0.17em}2}}`$ is a 2-cycle (with $`\text{}_2`$ coefficients) which can not be the boundary of a 3-chain since one of these triangles is maximal. Therefore $`H_2(X,\text{}_2)0`$, a contradiction. Thus $`\mathrm{\Delta }`$ contains at least one bad edge.
We claim that $`\mathrm{\Delta }`$ can’t have more than one bad edges. Suppose, on the contrary, that $`ab`$ and $`ac`$ are bad edges in $`X`$. Notice that (arguing as in the proof of the case $`dim(X)=4`$), if a 3-dimensional simplicial complex on $`6`$ vertices has $`2`$ tetrahedra through each triangle then it contains a combinatorial $`S^{\mathrm{\hspace{0.17em}3}}`$. Therefore the pure part $`Y`$ of $`X`$ must have seven vertices. In particular $`aY`$. Since $`ab`$ and $`ac`$ are bad edges, $`b,c\mathrm{Lk}_Y(a)`$ and hence $`\mathrm{deg}_Y(a)4`$. Therefore $`\mathrm{Lk}_Y(a)=S_4^{\mathrm{\hspace{0.17em}2}}`$. Hence we can apply an improper bistellar move to $`Y`$ to remove the vertex $`a`$, yielding a 6-vertex 3-dimensional simplicial complex $`\stackrel{~}{Y}`$ with $`2`$ tetrahedra through each triangle. Hence $`\stackrel{~}{Y}`$ has an $`S^{\mathrm{\hspace{0.17em}3}}`$ as a subcomplex, so that $`H_3(Y,\text{}_2)=H_3(\stackrel{~}{Y},\text{}_2)0`$. Therefore $`H_3(X,\text{}_2)0`$, a contradiction. Thus $`\mathrm{\Delta }`$ contains exactly one bad edge, say $`ab`$. Hence $`ac`$ and $`bc`$ are good edges.
Since $`X`$ has no free edge, there is a second triangle, say $`abd`$, through $`ab`$. Since $`ab`$ is a bad edge, $`abd`$ is maximal. By the above argument, $`ad`$ and $`bd`$ are good edges. If both $`acd`$ and $`bcd`$ are triangles of $`X`$ then $`X`$ has $`S_4^{\mathrm{\hspace{0.17em}2}}(a,b,c,d)`$ as a subcomplex, and at least one of the triangles of this $`S_4^{\mathrm{\hspace{0.17em}2}}`$ is maximal in $`X`$, yielding the contradiction $`H_2(X,\text{}_2)0`$ as before. Therefore, without loss of generality, we may assume $`bcdX`$. Note that $`a`$ is an isolated vertex in $`\mathrm{Lk}_X(bc)`$ and $`d`$ does not occur in $`\mathrm{Lk}_X(bc)`$. Since $`bc`$ is a good edge, it follows that all three vertices outside $`\{a,b,c,d\}`$ (say $`x`$, $`y`$ and $`z`$) occur in $`\mathrm{Lk}_X(bc)`$. Similarly, $`x,y,z\mathrm{Lk}_X(bd)`$. Again, the good edges $`ac`$ and $`ad`$ have at most one non-isolated vertex from $`\{a,b,c,d\}`$ in their links, hence each of them has at least two of $`x`$, $`y`$, $`z`$ in their links. Therefore, there is one vertex, say $`x`$, which occurs in the link of all the four edges $`ac`$, $`bc`$, $`ad`$, $`bd`$. Hence $`S_2^{\mathrm{\hspace{0.17em}0}}(c,d)S_3^{\mathrm{\hspace{0.17em}1}}(a,b,x)`$ is a subcomplex of $`X`$. Since one of the triangles in this 2-sphere is maximal, it follows that $`H_2(X,\text{}_2)0`$, a contradiction. Thus $`X`$ has no maximal triangles nor maximal edges, so $`X`$ is pure.
Finally, the last assertion follows from purity and minimality of $`X`$. $`\mathrm{}`$
###### Lemma 3.2
. Let $`X`$ be a $`7`$-vertex $`2`$-dimensional $`\text{}_2`$-acyclic simplicial complex. Then $`X`$ is collapsible.
Proof. Let $`X`$ be a minimal counter example. Let $`f_i`$, $`0i2`$, be the number of $`i`$-faces in $`X`$. Since $`X`$ is $`\text{}_2`$-acyclic, $`\chi (X)=1`$. Thus, $`f_0=7`$ and $`f_1=f_2+6`$.
For $`i0`$, let $`e_i`$ be the number of edges of degree $`i`$ in $`X`$. By Lemma 3.1, $`e_i=0`$ for $`i1`$. Two-way counting yields
$$\underset{i=2}{\overset{5}{}}e_i=f_1=f_2+6,\underset{i=2}{\overset{5}{}}ie_i=3f_2.$$
Hence
$$e_3+3e_5e_3+2e_4+3e_5=f_212.$$
(1)
Let’s say that an edge of $`X`$ is odd (respectively even) if it lies in an odd (respectively even) number of triangles. Note that each graph has an even number of vertices of odd degree. Applying this trivial observation to the vertex links of $`X`$, we conclude that each vertex of $`X`$ is in an even number of odd edges. Thus the total number $`e_3+e_5`$ of odd edges is $`=0`$ or $`3`$. If there is no odd edge then the sum of all the triangles gives a non-zero element of $`H_2(X,\text{}_2)`$, a contradiction. So, $`e_3+e_53`$. Combining this with (1), we get $`f_215`$ and hence $`f_121=\left(\genfrac{}{}{0pt}{}{7}{2}\right)`$. Hence $`f_1=21`$, $`f_2=15`$, $`e_3=3`$, $`e_4=e_5=0`$.
Since each vertex is in an even number of odd edges, it follows that the three odd edges form a triangle $`\mathrm{\Delta }`$, which may or may not be in $`X`$.
If $`\mathrm{\Delta }`$ is in $`X`$, then the sum of the remaining triangles gives a non-zero element of $`H_2(X,\text{}_2)`$, a contradiction. If $`\mathrm{\Delta }`$ is not in $`X`$ then (as each of the three edges in $`\mathrm{\Delta }`$ has three vertices in its link and there are four vertices outside $`\mathrm{\Delta }`$) by the pigeonhole principle there is a vertex $`x\mathrm{\Delta }`$ such that $`x`$ occurs in the link of each of the three edges in $`\mathrm{\Delta }`$. Then the sum of all the triangles excepting the three triangles in $`\mathrm{\Delta }\{x\}`$ gives a non-zero element of $`H_2(X,\text{}_2)`$, a contradiction. $`\mathrm{}`$
###### Lemma 3.3
. Let $`U`$ be a $`2`$-dimensional pure simplicial complex on $`7`$ vertices. Suppose the number of triangles in $`U`$ is $`10`$ and each edge of $`U`$ is in an even number of triangles. Then either $`U`$ is the union of two combinatorial spheres $`(`$on $`4`$ or $`5`$ vertices$`)`$ with no common triangle, or $`U`$ is isomorphic to one of $`S_4^{\mathrm{\hspace{0.17em}2}}`$, $`S_3^{\mathrm{\hspace{0.17em}1}}S_2^{\mathrm{\hspace{0.17em}0}}`$, $`S_2^{\mathrm{\hspace{0.17em}0}}S_2^{\mathrm{\hspace{0.17em}0}}S_2^{\mathrm{\hspace{0.17em}0}}`$, $`S_5^{\mathrm{\hspace{0.17em}1}}S_2^{\mathrm{\hspace{0.17em}0}}`$, $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$, $`\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}1}},\mathrm{},\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}5}}`$ or $`R`$ $`(`$of Example 1 and Example 2 $`(a))`$.
Proof. Let $`𝒮`$ be the list of simplicial complexes in the statement of this lemma. We find by inspection that $`𝒮`$ is closed under generalized bistellar 1-moves.
If $`f_0(U)5`$ then $`U`$ is a weak pseudomanifold and hence, by Proposition 2.1, $`U𝒮`$. So assume $`f_0(U)=6`$ or 7. The proof is by induction on the number $`n(U)`$ of degree 4 edges in $`U`$. If $`n(U)=0`$ then $`U`$ is a weak pseudomanifold and hence, by Propositions 2.1 and 2.2, $`U𝒮`$. So let $`n(U)>0`$ and suppose that we have the result for all smaller values of $`n(U)`$.
By the assumption, all the edges of $`U`$ are of degree 2 or 4. Therefore, a two-way counting yields $`4n(U)+2(f_1(U)n(U))=3f_2(U)30`$. Thus, $`n(U)+f_1(U)15`$. Therefore,
$$f_1(U)<15,$$
(2)
showing that $`U`$ has at least one non-edge. Fix an edge $`ab`$ of degree 4 in $`U`$. Let $`W`$ be the link of $`ab`$. If each pair of vertices in $`W`$ formed an edge in $`U`$ then $`f_1(U)`$ would be $`15`$, contradicting (2). So, there exist $`c,dW`$ such that $`cd`$ is a non-edge in $`U`$.
Let $`A=\{a,b,c,d\}`$. Then $`\kappa _A`$ is a generalized bistellar 1-move and hence $`\kappa _A(U)`$ also satisfies the hypothesis of the lemma, and $`n(\kappa _A(U))=n(U)1`$. Therefore, by the induction hypothesis, $`\kappa _A(U)𝒮`$. Since $`𝒮`$ is closed under generalized bistellar 1-moves, $`U=\kappa _A(\kappa _A(U))𝒮`$. $`\mathrm{}`$
###### Lemma 3.4
. Let $`X`$ be a $`7`$-vertex $`3`$-dimensional simplicial complex. Suppose $`(a)`$ $`X`$ is $`\text{}_2`$-acyclic, $`(b)`$ $`X`$ is not collapsible, and $`(c)`$ $`X`$ is minimal subject to $`(a)`$ and $`(b)`$. Then the $`f`$-vector of $`X`$ is $`(7,20,30,16)`$, $`(7,21,32,17)`$, $`(7,21,33,18)`$, $`(7,21,34,19)`$ or $`(7,21,35,20)`$.
Proof. For $`0i3`$, let $`f_i`$ be the number of $`i`$-faces of $`X`$. For $`i0`$, let $`t_i`$ be the number of triangles of degree $`i`$ in $`X`$. By Lemma 3.1, we have $`t_i=0`$ for $`i1`$. Two way counting yields
$$\underset{i=2}{\overset{4}{}}t_i=f_2,\underset{i=2}{\overset{4}{}}it_i=4f_3$$
and hence
$$t_3t_3+2t_4=4f_32f_2.$$
(3)
Say that a triangle of $`X`$ is odd (respectively even) if it is in an odd (respectively even) number of tetrahedra of $`X`$. By the same argument as in Lemma 3.2, each edge is in an even number of odd triangles, so that the number $`t_3`$ of odd triangles is 0 or $`4`$.
If there is no odd triangle then the sum of all the tetrahedra gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction. So, $`t_34`$. Combining this with (3) we get
$$2f_3f_22.$$
(4)
Since $`X`$ is $`\text{}_2`$-acyclic, by a result of Stanley (), $`X`$ has a 2-dimensional subcomplex $`Y`$ such that the $`f`$-vector of $`X`$ equals the $`f`$-vector of a cone over $`Y`$. (In , the author uses the vanishing of the reduced cohomology groups as his definition of acyclicity, while we have used the homology definition. However, since the coefficient ring used is a field, these two definitions coincide.) Let $`(g_0,g_1,g_2)`$ be the $`f`$-vector of $`Y`$. Thus, $`g_0=6`$ and
$$f_1=g_1+6,f_2=g_1+g_2,f_3=g_2.$$
(5)
Hence (4) yields
$$g_2g_1+2.$$
(6)
Let $`m=\left(\genfrac{}{}{0pt}{}{6}{2}\right)g_1`$, $`n=\left(\genfrac{}{}{0pt}{}{6}{3}\right)g_2`$ be the number of non-edges and non-triangles of $`Y`$, respectively. Since each non-edge is in exactly four non-triangles and any two non-edges are shared by at most one non-triangle, we have $`n4m\left(\genfrac{}{}{0pt}{}{m}{2}\right)`$. Also, from (6) we get $`nm+3`$. Hence $`m+34m\left(\genfrac{}{}{0pt}{}{m}{2}\right)`$ or $`(m1)(m6)0`$. So, either $`m1`$ or $`m6`$.
First suppose $`m6`$, i.e., $`g_19`$. If each edge of $`Y`$ was in $`3`$ triangles then we would have $`g_2g_1`$, contradicting (6). So, there is an edge of $`Y`$ contained in four triangles, together covering all the nine edges of $`Y`$. But, apart from the four triangles already seen, no three of these nine edges form a triangle of $`Y`$. Thus $`g_2=4`$, $`g_1=9`$ – contradicting (6). So, $`m1`$, i.e., $`g_1=14`$ or 15.
If $`g_1=14`$ then the four triangles through the missing edge are missing from $`Y`$, so that $`g_216`$. Thus, by (6), $`(g_1,g_2)=(14,16)`$, $`(15,17)`$, $`(15,18)`$, $`(15,19)`$ or $`(15,20)`$. The lemma now follows from (5). $`\mathrm{}`$
###### Lemma 3.5
. Let $`X`$ be a $`7`$-vertex $`3`$-dimensional $`\text{}_2`$-acyclic simplicial complex. Then $`X`$ is collapsible.
Proof. Let $`X`$ be a minimal counter example. As before, each edge is in an even number of odd triangles. Let $`f_i`$’s and $`t_j`$’s be as in the proof of Lemma 3.4. Then, by Lemma 3.4, $`t_3+2t_4=4f_32f_210`$ and hence the number $`t_3`$ of odd triangles is $`10`$.
Let $`U`$ denote the pure 2-dimensional simplicial complex whose facets are the odd triangles of $`X`$. Then each edge of $`U`$ is in an even number of triangles of $`U`$. Therefore, by Lemma 3.3, we get the following cases :
Case 1 : $`U`$ is the union of two combinatorial spheres with no common triangle (on 4 or 5 vertices), say on vertex sets $`A`$ and $`B`$.
First suppose $`\mathrm{\#}(A)=\mathrm{\#}(B)=4`$. If both $`A`$ and $`B`$ are 3-faces in $`X`$ then the pure simplicial complex $`\stackrel{~}{X}`$ whose facets are those of $`X`$ other than $`A`$, $`B`$ is a 3-dimensional weak pseudomanifold. This implies that the sum of all the tetrahedra, excepting $`A`$ and $`B`$, gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction. So, without loss of generality $`AX`$.
Since each of the four triangles inside $`A`$ is of degree 3 in $`X`$, the three vertices (say $`x`$, $`y`$, $`z`$) outside $`A`$ occur in the link of all the four triangles. Then the 3-sphere $`S_4^{\mathrm{\hspace{0.17em}2}}(A)S_2^{\mathrm{\hspace{0.17em}0}}(x,y)`$ occurs as a subcomplex of $`X`$, forcing $`H_3(X,\text{}_2)0`$, a contradiction.
In the remaining case $`\mathrm{\#}(A)=4`$, $`\mathrm{\#}(B)=5`$ (since $`U`$ has at most 10 triangles, the case $`\mathrm{\#}(A)=\mathrm{\#}(B)=5`$ does not arise). Write $`B=\{b_1,b_2,b_3,x,y\}`$ and $`U=S_4^{\mathrm{\hspace{0.17em}2}}(A)(S_3^{\mathrm{\hspace{0.17em}1}}(b_1,b_2,b_3)S_2^{\mathrm{\hspace{0.17em}0}}(x,y))`$. As above, we must have $`AX`$.
If both $`b_1b_2b_3x`$ and $`b_1b_2b_3y`$ are in $`X`$, then the sum of the 3-faces other than $`A`$, $`b_1b_2b_3x`$ and $`b_1b_2b_3y`$ gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction. So, without loss of generality, $`b_1b_2b_3xX`$. Since the triangles of $`S_3^{\mathrm{\hspace{0.17em}1}}(b_1,b_2,b_3)S_2^{\mathrm{\hspace{0.17em}0}}(x,y))`$ are degree 3 triangles in $`X`$, it follows that $`b_1b_2xy`$, $`b_1b_3xy`$, $`b_2b_3xyX`$. Then the sum of the tetrahedra other than $`A`$ and these three tetrahedra gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction.
Case 2 : $`U=S_4^{\mathrm{\hspace{0.17em}2}}`$. We get a contradiction as in Case 1.
Case 3 : $`U=S_3^{\mathrm{\hspace{0.17em}1}}S_2^{\mathrm{\hspace{0.17em}0}}`$. We get a contradiction as in Case 1.
Observation 1 : As $`t_38`$ in the remaining cases, we have $`2f_3f_24`$ and hence only the following two possibilities survive for the $`f`$-vector of $`X`$ : $`(7,21,34,19)`$ and $`(7,21,35,20)`$. Therefore $`X`$ has at most one missing triangle and at most one triangle of degree 4, and these two cases are exclusive. It follows that, if $`x`$ is a vertex not covered by the odd triangles, then $`\mathrm{Lk}_X(x)`$ is a 6-vertex 2-dimensional neighbourly weak pseudomanifold. But, from Proposition 2.1, we see that $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$ is the only possibility. Thus, $`\mathrm{Lk}_X(x)=\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$. This implies that if $`V_1V(U)`$ is a 3-set then exactly one of $`V_1`$ and $`V(U)V_1`$ is a simplex in $`\mathrm{Lk}_X(x)`$. In particular, any two triangles in $`\mathrm{Lk}_X(x)`$ intersect.
Case 4 : $`U=S_2^{\mathrm{\hspace{0.17em}0}}(a_1,a_2)S_2^{\mathrm{\hspace{0.17em}0}}(b_1,b_2)S_2^{\mathrm{\hspace{0.17em}0}}(c_1,c_2)`$. Then the odd triangles of $`X`$ are $`a_ib_jc_k`$, $`1i,j,k2`$. If $`\{a_1a_2b_jc_k:\mathrm{\hspace{0.17em}1}j,k2\}X`$, then the sum of the remaining tetrahedra gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction. So, without loss of generality, $`a_1a_2b_1c_1X`$. As $`a_1b_1c_1`$, $`a_2b_1c_1`$ are degree 3 triangles, it follows that $`a_1b_1b_2c_1`$, $`a_2b_1b_2c_1X`$. If both $`a_1b_1b_2c_2`$ and $`a_2b_1b_2c_2`$ are in $`X`$ then $`X\{a_ib_1b_2c_k:1i,k2\}`$, hence we get a contradiction as before. So, without loss of generality, $`a_2b_1b_2c_2X`$.
Since $`a_1a_2b_1c_1`$, $`a_2b_1b_2c_2X`$ and $`a_1b_1c_1`$, $`a_2b_2c_2`$ are degree 3 triangles, it follows that these two disjoint triangles occur in the link of $`x`$. But this contradicts Observation 1.
Case 5 : $`U=\mathrm{\Sigma }_1`$ of Example 1. Thus, the odd triangles are $`125`$, $`126`$, $`156`$, $`235`$, $`236`$, $`345`$, $`346`$ and $`456`$. If $`1256`$, $`3456X`$ then, since $`125`$ and $`346`$ are degree 3 triangles, they are disjoint triangles in $`\mathrm{Lk}_X(x)`$, contradicting Observation 1. So, without loss of generality, $`1256X`$.
If $`3456X`$ then, since $`345`$, $`346`$, $`456`$ are degree 3 triangles, $`2345`$, $`2346`$, $`2456X`$. Then the sum of all the tetrahedra, excepting $`1256`$, $`2345`$, $`2346`$, $`2456`$, gives a non-zero element of $`H_3(X,\text{}_2)`$. So, $`3456X`$.
If $`2356X`$, then the sum of all the tetrahedra, excepting $`1256`$, $`2356`$, $`3456`$, gives a non-zero element of $`H_3(X,\text{}_2)`$. Therefore $`2356X`$.
Since $`235`$ and $`236`$ are degree 3 triangles, $`2345`$, $`2346X`$. First suppose that at least one of $`1356`$, $`2456`$ is in $`X`$. Without loss, say $`2456X`$. Then the sum of all the tetrahedra, excepting $`1256`$, $`2456`$, $`2345`$, $`2346`$, gives a non-zero element of $`H_3(X,\text{}_2)`$. Thus $`1356`$, $`2456X`$. Then, since $`156`$, $`456`$ are degree 3 triangles, $`156x`$, $`456xX`$.
Since $`2356,2456X`$, $`x\mathrm{Lk}_X(256)`$, i.e., $`256xX`$. Similarly, looking at $`356`$, we conclude that $`356xX`$. Thus, $`56x`$ is a degree 4 triangle in $`X`$. But this is not possible since, by Observation 1, $`\mathrm{Lk}_X(x)`$ is $`\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$.
Observation 2 : In the remaining cases, $`t_3=10`$ and hence the $`f`$-vector of $`X`$ is $`(7,21,35,20)`$. In consequence, $`t_4=0`$. Thus all triangles are of degree 2 or 3. Since $`f_3=\left(\genfrac{}{}{0pt}{}{7}{3}\right)`$, each edge in $`X`$ has degree 5. Thus if $`e`$ is an edge outside $`U`$ then the link of $`e`$ is a pentagon ($`S_5^{\mathrm{\hspace{0.17em}1}}`$).
Case 6 : $`U=\text{}P_6^{\mathrm{\hspace{0.17em}2}}`$. In this case, all the 4-sets of vertices not containing $`x`$ contain exactly two odd triangles each. In particular, all the tetrahedra of $`X`$ not containing $`x`$ contain exactly two odd triangles each. Trivially, each tetrahedron through $`x`$ contains at most one odd triangle. Thus, letting $`\alpha _i`$, $`i0`$, denote the number of tetrahedra of $`X`$ containing exactly $`i`$ odd triangles, we have $`\alpha _2=2010=10`$ and $`\alpha _0+\alpha _1=10`$. But two way counting yields $`\alpha _1+2\alpha _2=10\times 3=30`$. Hence $`\alpha _1=10`$, $`\alpha _0=0`$. Thus $`x`$ occurs in the link of each odd triangle and hence $`\mathrm{Lk}_X(x)=U`$. Therefore the 10 tetrahedra of $`X`$ not passing through $`x`$ add up to a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction.
Case 7 : $`U=R`$ of Example 2 $`(a)`$. Thus, the odd triangles are $`123`$, $`124,125,126,135,146`$, $`236`$, $`245`$, $`345`$ and $`346`$. We claim that $`\mathrm{Lk}_X(12)`$ . If, for instance, $`1236X`$ then, as $`123`$, $`126`$, $`236`$ are degree 3 triangles, $`x`$ belongs to the link of each of these triangles. Then $`\mathrm{Lk}_X(2x)`$ , contradicting Observation 2. This proves the claim.
Since $`3,4,5,6`$ are of degree 3 and $`x`$ is of degree 2 in $`\mathrm{Lk}_X(12)`$, it follows that $`\mathrm{Lk}_X(12)=`$ or $`=`$ .
In the first case, $`125,126\mathrm{Lk}_X(x)`$. Hence, by Observation 1, $`345`$, $`346\mathrm{Lk}_X(x)`$. Since these two are degree 3 triangles, it follows that $`\mathrm{Lk}_X(345)=\{1,2,6\}`$ and $`\mathrm{Lk}_X(346)=\{1,2,5\}`$. Since $`1`$, $`2`$ are of degree 2 in $`\mathrm{Lk}_X(34)`$, this forces $`\mathrm{Lk}_X(34)=`$ and hence $`x\mathrm{Lk}_X(34)`$. This is a contradiction since $`X`$ is 3-neighbourly.
In the second case, $`125,126\mathrm{Lk}_X(x)`$ and hence, by Observation 1, $`345,346\mathrm{Lk}_X(x)`$. That is, $`5x`$, $`6x\mathrm{Lk}_X(34)`$. Also, as $`34\mathrm{Lk}_X(12)`$, we have $`12\mathrm{Lk}_X(34)`$. Since $`5,6`$ are of degree 3 and $`1`$, $`2`$, $`x`$ are of degree 2 in $`\mathrm{Lk}_X(34)`$, it follows that $`\mathrm{Lk}_X(34)=`$ . Hence $`1345,2345,345xX`$. Also, as $`123`$ is a degree 3 triangle and $`1234X`$, we have $`1235X`$. Thus $`\mathrm{Lk}_X(35)`$. Since $`1`$, $`4`$ are of degree 3 while $`2`$, $`6`$, $`x`$ are of degree 2 in this link, it follows that $`\mathrm{Lk}_X(35)=`$ . Hence $`356xX`$. Then $`\mathrm{Lk}_X(3x)`$, contradicting Observation 2.
Claim : In the remaining cases, if $`F`$ is a set of four vertices of $`U`$ containing at least two odd triangles, then either $`FX`$ or $`FV(\mathrm{Lk}_U(x))`$ for some vertex $`x`$.
In these cases, $`V(U)=V(X)`$. If $`FX`$ contains two odd triangles, then on the average, a vertex outside $`F`$ occurs in the links (in $`X`$) of $`\frac{3\times 2+2\times 2}{3}>3`$ of the four triangles inside $`F`$. Thus there is a vertex $`x`$ in the link of all these triangles. If $`FV(\mathrm{Lk}_U(x))`$ for this $`x`$, then choose a vertex $`yF`$ such that $`xyU`$. Then $`\mathrm{Lk}_X(xy)S_3^{\mathrm{\hspace{0.17em}1}}(F\{y\})`$, contradicting Observation 2. This proves the claim.
Case 8 : $`U=S_5^{\mathrm{\hspace{0.17em}1}}(\text{}_5)S_2^{\mathrm{\hspace{0.17em}0}}(u,v)`$. In this case, the above claim implies that $`X`$ contains the five tetrahedra $`\{u,v,i,i+1\}`$, $`i\text{}_5`$. Then the sum of the remaining fifteen tetrahedra gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction.
Case 9 : $`U=\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}2}}`$ of Example 1. Thus, the odd triangles are $`126`$, $`127`$, $`167`$, $`236`$, $`237`$, $`346`$, $`347`$, $`456`$, $`457`$ and $`567`$. By the above claim, $`1267`$, $`2367`$, $`3467`$, $`4567X`$. Then the sum of the remaining sixteen tetrahedra gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction.
Case 10 : $`U=\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}3}}`$ of Example 1. Thus, the odd triangles are $`126`$, $`127`$, $`167`$, $`234`$, $`237`$, $`246`$, $`347`$, $`456`$, $`457`$ and $`567`$. By the claim, $`1267`$, $`2347`$, $`4567X`$.
If $`2467X`$ then the sum of all the tetrahedra, excepting $`1267`$, $`2347`$, $`4567`$, $`2467`$, gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction. So, $`2467X`$. Then, $`\mathrm{Lk}_X(246)=\{1,3,5\}`$.
Since $`\mathrm{deg}(247)=2`$ and $`2347X`$, assume without loss of generality, that $`2457X`$ and $`1247X`$. Then $`\mathrm{Lk}_X(127)=\{3,5,6\}`$.
So, $`2456`$, $`2457X`$ and $`\mathrm{deg}(245)=2`$. Hence $`2345X`$. Then $`\mathrm{Lk}_X(234)=\{1,6,7\}`$.
Now, $`1234`$, $`1237X`$ and $`\mathrm{deg}(123)=2`$. Therefore, $`1236X`$. Then $`\mathrm{Lk}_X(126)=\{4,5,7\}`$. This implies that $`\mathrm{Lk}_X(25)`$, a contradiction to Observation 2.
Case 11 : $`U=\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}4}}`$ of Example 1. Thus, the odd triangles are $`124`$, $`127`$, $`145`$, $`156`$, $`167`$, $`234`$, $`237`$, $`347`$, $`457`$ and $`567`$. By the claim, $`1247`$, $`1457`$, $`1567`$, $`2347X`$. Then the sum of the remaining sixteen tetrahedra gives a non-zero element of $`H_3(X,\text{}_2)`$, a contradiction.
Case 12 : $`U=\mathrm{\Sigma }_{\mathrm{\hspace{0.17em}5}}`$ of Example 1. Thus, the odd triangles are $`123`$, $`126`$, $`135`$, $`156`$, $`234`$, $`246`$, $`345`$, $`457`$, $`467`$, $`567`$. By the claim, $`1234`$, $`1235`$, $`1246`$, $`1256`$, $`1345`$, $`2345`$, $`3457`$, $`4567`$ $`X`$. Thus $`\mathrm{Lk}_X(14)`$ and $`\mathrm{Lk}_X(25)`$ . Since $`14`$ and $`25`$ are not in $`U`$, Observation 2 implies that $`\mathrm{Lk}_X(14)=`$ and $`\mathrm{Lk}_X(25)=`$ . Thus $`1457`$, $`2457X`$. Then the triangle $`457`$ is of degree 4 in $`X`$, a contradiction. This completes the proof. $`\mathrm{}`$
Proof of Theorem 1 . Let $`Y`$ be a minimal counter example. So, $`Y`$ is an $`n`$-vertex (for some $`n7`$) $`\text{}_2`$-acyclic simplicial complex which is not collapsible to any proper subcomplex.
If $`n<7`$ then choose a facet $`\alpha `$ of $`Y`$ and an element $`vV(Y)`$. Let $`\stackrel{~}{Y}`$ be obtained from $`Y`$ by the bistellar $`d`$-move $`\kappa _{\alpha \{v\}}`$, where $`d`$ is the dimension of $`Y`$. Then $`\stackrel{~}{Y}`$ is an $`(n+1)`$-vertex $`\text{}_2`$-acyclic simplicial complex. Since $`Y`$ has no free face, $`\stackrel{~}{Y}`$ has no free face and hence $`\stackrel{~}{Y}`$ is not collapsible to any proper subcomplex. Repeating this construction (if necessary) we get a 7-vertex $`\text{}_2`$-acyclic simplicial complex $`X`$ which is not collapsible to any proper subcomplex. Then, by Lemma 3.1, $`X`$ is of dimension $`2`$ or $`3`$. But, this is not possible by Lemmas 3.2 and 3.5. This completes the proof. $`\mathrm{}`$
## 4 Homology spheres.
A connected $`d`$-dimensional weak pseudomanifold is called a normal pseudomanifold if the links of all the simplices of dimension up to $`d2`$ are connected. Observe that if $`X`$ is a normal pseudomanifold then $`X`$ is a pseudomanifold. (If not then, since $`X`$ is connected, there exist two intersecting facets $`\tau `$, $`\sigma `$ for which there is no sequence of facets $`\tau =\tau _0,\mathrm{},\tau _n=\sigma `$ such that $`\tau _{i1}\tau _i`$ is a face of codimension 1 for $`1in`$. Choose $`\tau `$, $`\sigma `$ among all such pairs such that $`dim(\tau \sigma )`$ is maximum. Then $`dim(\tau \sigma )d2`$ and $`\mathrm{lk}_X(\tau \sigma )`$ is not connected, a contradiction.) Notice that all the links of positive dimensions (i.e., the links of simplices of dimension $`d2`$) in a normal $`d`$-pseudomanifold are normal pseudomanifolds (and hence are pseudomanifolds). Clearly, any triangulation of a connected closed manifold is a normal pseudomanifold.
###### Lemma 4.1
. Let $`Y`$ be a $`d`$-dimensional normal pseudomanifold. Let $`Y_1`$ be a proper induced subcomplex of $`Y`$ which is pure of dimension $`d`$. Put $`L=C(Y_1,Y)`$ and $`Y_2=N(L,Y)`$. Then $`(a)`$ $`Y_1`$, $`Y_2`$ are weak pseudomanifolds with boundary, $`(b)`$ $`Y_2`$ is an induced subcomplex of $`Y_2`$ and $`(c)`$ $`Y_2=Y_1=Y_1Y_2`$.
Proof . Since $`Y`$ is a pseudomanifold and $`Y_1Y`$ is pure of maximum dimension, $`Y_1`$ is a weak pseudomanifold with boundary. Since the maximal simplices of $`Y_2`$ are those maximal simplices of $`Y`$ which intersect $`V(L)`$, $`Y_2`$ is pure of dimension $`d`$ and each $`d`$-simplex of $`Y`$ is either in $`Y_1`$ or in $`Y_2`$ but not in both. This implies that $`Y_2`$ is a weak pseudomanifold with boundary. This proves $`(a)`$.
Let $`V_1=V(Y_1)`$, $`V_2=V(L)`$. Then $`V(Y)=V_1V_2`$. Now, $`\tau `$ is a facet of $`Y_2`$ $``$ there exists a unique $`d`$-face $`\sigma _2Y_2`$ containing $`\tau `$ $``$ there exists a unique $`d`$-face $`\sigma _1Y_1`$ containing $`\tau `$ $``$ $`\tau `$ is a facet of $`Y_1`$. Therefore, $`Y_2=Y_1Y_1Y_2`$.
Conversely, let $`\tau `$ be an $`i`$-simplex in $`Y_1Y_2`$. If possible, let $`\tau Y_1=Y_2`$. Then $`\tau Y_1Y_1`$. Therefore, $`\mathrm{lk}_{Y_1}(\tau )`$ is a $`(di1)`$-dimensional weak pseudomanifold (without boundary) and $`\mathrm{lk}_{Y_1}(\tau )\mathrm{lk}_Y(\tau )`$. If $`i=d1`$ then each of $`\mathrm{lk}_{Y_1}(\tau )`$ and $`\mathrm{lk}_Y(\tau )`$ consists of two vertices and hence $`\mathrm{lk}_{Y_1}(\tau )=\mathrm{lk}_{Y_1}(\tau )`$. Now, assume that $`id2`$. Since $`Y`$ is a normal pseudomanifold, $`\mathrm{lk}_Y(\tau )`$ is a pseudomanifold (of dimension $`1`$) and hence $`\mathrm{lk}_{Y_1}(\tau )=\mathrm{lk}_Y(\tau )`$. This implies that $`\mathrm{star}_Y(\tau )Y_1`$ and hence $`\mathrm{star}_Y(\tau )L=\mathrm{}`$. Now, $`\tau Y_2`$ and $`Y_2`$ is pure. Therefore, there exists a $`d`$-simplex $`\sigma Y_2`$ such that $`\tau \sigma `$. Since $`\sigma Y_2`$, $`\sigma V(L)=\mathrm{}`$. This implies that $`\mathrm{star}_Y(\tau )L\mathrm{}`$, a contradiction. Therefore, $`\tau Y_1=Y_2`$. So, $`Y_1Y_2=Y_1=Y_2`$. This proves (c).
Since $`Y_2=Y_1`$, $`Y_2Y_2[V_1]=Y_2[V_1]Y[V_1]=Y_2[V_1]Y_1Y_2Y_1=Y_2`$. Thus, $`Y_2=Y_2[V_1]=Y_2[V_1V(Y_2)]`$. This proves $`(b)`$. $`\mathrm{}`$
###### Lemma 4.2
. Let $`X`$ be a connected combinatorial $`d`$-manifold. Let $`X_1`$ be an induced subcomplex of $`X`$ which is a combinatorial $`d`$-ball. Put $`L=C(X_1,X)`$ and $`X_2=N(L,X)`$. Then
1. $`X_2`$ is a connected combinatorial $`d`$-manifold with boundary.
2. $`|X_2||L|`$.
3. If, further, $`L`$ is collapsible then $`X`$ is a combinatorial sphere.
Proof . Let $`V_1=V(X_1)`$, $`V_2=V(L)`$. Then $`V(X)=V_1V_2`$. As in the proof of Lemma 4.1, $`X_2`$ is pure of dimension $`d`$ and each $`d`$-simplex of $`X`$ is either in $`X_1`$ or in $`X_2`$ but not in both.
Let $`v`$ be a vertex of $`X_2`$. Notice that $`vX_1X_1`$ $``$ $`\mathrm{Lk}_{X_1}(v)\mathrm{Lk}_X(v)`$ are $`(d1)`$-spheres $``$ $`\mathrm{Lk}_{X_1}(v)=\mathrm{Lk}_X(v)`$ $``$ $`vX_2`$, a contradiction. So, either $`vV_2`$ or $`vX_1`$.
If $`vV_2`$ then each $`d`$-simplex of $`X`$ containing $`v`$ is in $`X_2`$ and hence $`\mathrm{Lk}_{X_2}(v)=\mathrm{Lk}_X(v)`$ is a combinatorial $`(d1)`$-sphere.
If $`vX_1`$ then $`(Y,Y_1,Y_2):=(\mathrm{Lk}_X(v),\mathrm{Lk}_{X_1}(v),\mathrm{Lk}_{X_2}(v))`$ satisfies the hypothesis of Lemma 4.1. Therefore, by Lemma 4.1, $`\mathrm{Lk}_{X_1}(v)\mathrm{Lk}_{X_2}(v)=(\mathrm{Lk}_{X_2}(v))`$. This implies that the closure of $`|\mathrm{Lk}_X(v)||\mathrm{Lk}_{X_1}(v)|`$ in $`|\mathrm{Lk}_X(v)|`$ is $`|\mathrm{Lk}_{X_2}(v)|`$. Since $`|\mathrm{Lk}_X(v)|`$ is a pl $`(d1)`$-sphere and $`|\mathrm{Lk}_{X_1}(v)|`$ is a pl $`(d1)`$-ball, by Proposition 2.5, $`|\mathrm{Lk}_{X_2}(v)|`$ is a pl $`(d1)`$-ball. Thus, $`\mathrm{Lk}_{X_2}(v)`$ is a combinatorial $`(d1)`$-ball.
Thus $`X_2`$ is a combinatorial $`d`$-manifold with boundary such that $`X_2`$ ($`=X_1`$, by Lemma 4.1) is connected. Therefore, if $`X_2`$ were disconnected, it would have a $`d`$-dimensional weak pseudomanifold as a component. This is not possible since $`X`$ is a $`d`$-dimensional pseudomanifold. Therefore $`X_2`$ is connected. This proves $`(a)`$.
As $`L=X[V_2]`$, we have $`LX_2`$ and hence $`L=X_2[V_2]`$. Since, by Lemma 4.1, $`X_2`$ is the induced subcomplex of $`X_2`$ on $`V_1V(X_2)`$, this implies that $`L`$ is the simplicial complement of $`X_2`$ in $`X_2`$. Then, by Proposition 2.3, $`|X_2||L|`$. This proves $`(b)`$.
Now, if $`L\stackrel{^\mathrm{s}}{}0`$ then $`|L|0`$ and hence $`|X_2|0`$. So, by Proposition 2.6, $`|X_2|`$ is a pl ball.
Let $`\sigma `$ be a $`d`$-simplex in $`S_{d+2}^d`$. Let $`B_1=|\sigma |`$ and $`B_2=|S_{d+2}^d\{\sigma \}|`$. Then $`B_1`$ and $`B_2`$ are pl $`d`$-balls. Let $`f_2:B_2|X_2|`$ be a pl homeomorphism. Let $`f=f_2|_{B_2}`$. Since $`B_1=B_2`$ and $`(|X_1|)=|X_1|=|X_2|`$, $`f:B_1(|X_1|)`$ is a pl homeomorphism. By Proposition 2.4, there exists a pl homeomorphism $`f_1:B_1|X_1|`$ such that $`f_1|_{B_1}=f=f_2|_{B_2}`$. Then $`f_1f_2`$ is a pl homeomorphism from $`|S_{d+2}^d|`$ to $`|X|`$. This proves $`(c)`$. $`\mathrm{}`$
###### Lemma 4.3
. Let $`X`$ be a combinatorial triangulation of a $`\text{}_2`$-homology $`d`$-sphere. Let $`X_1`$ be an induced subcomplex of $`X`$ which is a combinatorial $`d`$-ball. Let $`L=C(X_1,X)`$ and $`X_2=N(L,X)`$. Then $`X_2`$ is $`\text{}_2`$-acyclic.
Proof . Let $`J=X_1X_2`$. Then, by Lemma 4.1, $`J=X_1`$. So, $`J`$ is a combinatorial $`(d1)`$-sphere. Therefore, $`H_{d1}(J,\text{}_2)=\text{}_2`$ and $`\stackrel{~}{H}_q(J,\text{}_2)=0`$ for all $`qd1`$. Also $`\stackrel{~}{H}_q(X_1,\text{}_2)=0`$ for all $`q0`$. For $`q1`$, we have the following exact Mayer-Vietoris sequence of homology groups with coefficients in $`\text{}_2`$ (see ) :
$$\mathrm{}H_{q+1}(X)H_q(J)H_q(X_1)H_q(X_2)H_q(X)\stackrel{~}{H}_{q1}(J)\mathrm{}$$
(7)
Now, $`H_d(X,\text{}_2)=\text{}_2`$ and $`\stackrel{~}{H}_q(X,\text{}_2)=0`$ for $`qd`$. By Lemma 4.2, $`|X_2|`$ is a connected $`d`$-manifold with non-trivial boundary. Therefore, $`H_d(X_2,\text{}_2)=0`$ and $`H_0(X_2,\text{}_2)=\text{}_2`$. Then, by (7), $`H_q(X_2,\text{}_2)=0`$ for $`0<q<d1`$ and for $`q=d1`$ we get the following short exact sequence of abelian groups :
$$0\text{}_2\text{}_2H_{d1}(X_2,\text{}_2)0.$$
Clearly, this implies $`H_{d1}(X_2,\text{}_2)=0`$. Thus, $`\stackrel{~}{H}_q(X_2,\text{}_2)=0`$ for all $`q0`$. $`\mathrm{}`$
Proof of Theorem 2. Let $`X_1`$ be an $`m`$-vertex induced subcomplex of $`M`$ which is a combinatorial $`d`$-ball. Let $`L=C(X_1,M)`$ and $`X_2=N(L,M)`$. Then, by Part $`(b)`$ of Lemma 4.2, $`|X_2||L|`$.
Again, by Lemma 4.3, $`X_2`$ is $`\text{}_2`$-acyclic and hence $`L`$ is $`\text{}_2`$-acyclic. Since $`nm+7`$, the number of vertices in $`L`$ is $`7`$. Therefore, by Theorem 1, $`L`$ is collapsible. Then, by Part $`(c)`$ of Lemma 4.2, $`M`$ is a combinatorial sphere. $`\mathrm{}`$
Proof of Corollary 3. If $`\sigma `$ is a $`d`$-simplex of $`M`$ then the induced subcomplex $`\mathrm{\Delta }_{d+1}^d(\sigma )`$ is a $`(d+1)`$-vertex combinatorial $`d`$-ball. Therefore, by Theorem 2, $`M`$ is a combinatorial sphere. $`\mathrm{}`$
Proof of Corollary 4. Assume, if possible, that $`M`$ admits a bistellar $`i`$-move $`\kappa _A`$ for some $`i<d`$. Let $`\beta `$ be the core of $`A`$ and $`\alpha =A\beta `$. Then $`M[A]=\mathrm{\Delta }_{i+1}^i(\alpha )S_{di+1}^{di1}(\beta )`$ is a $`(d+2)`$-vertex combinatorial $`d`$-ball. Therefore, by Theorem 2, $`M`$ is a combinatorial sphere, a contradiction. This proves the corollary. $`\mathrm{}`$
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# Reference potential approach to the quantum-mechanical inverse problem: II. Solution of Krein equation
## 1. Introduction
In this paper, a recently proposed reference potential approach to the one-dimensional quantum mechanical inverse problem is further developed. Let us briefly recall that the starting idea is to choose a suitable reference potential for the system. For this fixed potential, it is always possible to calculate all its spectral characteristics, which means that in this artificial way one can obtain the complete set of information (otherwise inaccessible) needed to uniquely solve the inverse problem: 1) full energy spectrum of the bound states; 2) full energy dependence of the phase shift for the scattering states; 3) norming constants of the regular energy eigenfunctions for all bound states. Of course, there is no need to regain the potential which is already known by definition. However, the quantities related to the reference potential can be used as zeroth approximations to the real spectral characteristics of the system. A further step might be, for example, to calculate another potential whose discrete energy levels would exactly fit with their actually observed values, so that at least in this sense the new potential would be more realistic than the initial reference potential.
To illustrate the method, as previously, Xe<sub>2</sub> molecule is under study, and the same three-component exactly solvable reference potential is used. We have already ascertained the full energy dependence of the phase shift for the scattering states and demonstrated its excellent agreement with the celebrated Levinson theorem . In addition, we are provided with full knowledge of the Jost function, which is the most important spectral characteristic of the system. Thus, we are prepared to attack the most serious computational-technical problem, solution of an integral equation, which would enable to uniquely ascertain the potential. For this purpose a combined approach is used. In Section 2 we recall some useful properties of the Gelfand-Levitan method , which enable to separate the main problem into two independent parts: 1) calculating an auxiliary potential $`V_0(r)`$ with exactly the same spectral density for positive energies as the reference potential $`V(r)`$, but with no bound states; 2) calculating a series of auxiliary potentials $`V_k(r),`$ $`k=1,2,\mathrm{},n`$ (i.e., adding, one by one, all bound states, and keeping their norming constants), also having the same spectral density for positive energies. The second step is, in fact, less complicated and is only briefly discussed in this paper. Our main goal is to ascertain the auxiliary potential $`V_0(r)`$ with no bound states, starting from the known Jost function for positive energies. This is the subject of Section 3 where the Krein method will be used. Some concluding remarks and a discussion of the further perspectives of the method form the content of Section 4.
## 2. Step-by-step use of Gelfand-Levitan method
First, let us recall general solution scheme in the frame of Gelfand-Levitan method (see, e.g., for more details). One’s aim is to solve the integral equation
(1)
$$K(r,r^{})+G(r,r^{})+\underset{0}{\overset{r}{}}K(r,s)G(s,r^{})𝑑s=0,$$
whose kernel
(2)
$$G(r,r^{})=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{sin}\left(kr\right)\mathrm{sin}\left(kr^{}\right)}{k^2}𝑑\sigma ,$$
is determined by the quantity $`d\sigma \left[{\displaystyle \frac{d\rho _2(E)}{dE}}{\displaystyle \frac{d\rho _1(E)}{dE}}\right]dE,`$ which contains the the difference of two spectral densities
(3)
$$\frac{d\rho (E)}{dE}=\{\genfrac{}{}{0.0pt}{}{\pi ^1\sqrt{E}|F(E)|^2,E0,}{\underset{n}{}C_n\delta (EE_n),\text{ }E<0,}$$
one of them ($`{\displaystyle \frac{d\rho _1(E)}{dE}}`$) being related to a known potential ($`V_1(r)`$). In this sense Gelfand-Levitan method also represents a reference potential approach. The simplest possibility is to take $`V_1(r)0`$, and correspondingly, $`{\displaystyle \frac{d\rho _1(E)}{dE}}={\displaystyle \frac{\sqrt{E}}{\pi }}`$. Then
(4)
$$G(r,r^{})=\frac{2}{\pi }\underset{0}{\overset{\mathrm{}}{}}\mathrm{sin}\left(kr\right)\mathrm{sin}\left(kr^{}\right)g(k)𝑑k+\underset{n}{}\frac{C_n}{4\gamma _n^2}\mathrm{sinh}\left(\gamma _nr\right)\mathrm{sinh}\left(\gamma _nr^{}\right),$$
where $`\gamma _n^2={\displaystyle \frac{2mE_n}{\mathrm{}^2}}`$, $`E_n`$ being the bound levels and $`C_n`$, their related norming constants. The characteristic function
(5)
$$g(k)\frac{1}{\left|F(k)\right|^2}1$$
is determined by the modulus of the Jost function $`F(k)`$. Thus, in principle, one can solve Eq. (1) and then calculate the potential
(6)
$$V(r)=2C\frac{d}{dr}K(r,r),\text{ }C\frac{\mathrm{}^2}{2m}.\text{ }$$
This scheme might seem simple, but its actual realization stumbles upon serious computational-technical difficulties. Now, let us assume that we somehow managed to ascertain an auxiliary potential $`V_0(r)`$ with exactly the same spectral density for positive energies as the desired potential $`V(r)`$, but having no bound states. In this case one easily finds that
(7)
$$G(r,r^{})=\underset{j}{}C_j\phi _0(i\gamma _j,r)\phi _0(i\gamma _j,r^{}),$$
where $`\phi _0(i\gamma _j,r)`$ are the regular solutions (NB! not the real eigenfunctions!) related to the auxiliary potential $`V_0(r).`$ Next, let us introduce another auxiliary potential $`V_1(r)`$ with just one bound eigenvalue $`E_1`$ (note that numeration of the levels here starts from the highest-energy one, and therefore, $`E_n`$ corresponds to the zeroth level), and again, with exactly the same spectral density for positive energies ($`\left|F(k)\right|`$ remains the same). Thereafter, one introduces an auxiliary potential $`V_2(r)`$ with two levels ($`E_1`$ and $`E_2`$), etc., until he comes to the desired potential $`V_n(r)=V(r)`$.
It can be proved that
(8)
$$V_j(r)V_0(r)=2C\left\{\mathrm{ln}\left[detC_j(r)\right]\right\}^{\prime \prime },\text{ }j=1,2,\mathrm{},n$$
and the corresponding regular solution
(9)
$$\phi _j(k,r)=det\left|\genfrac{}{}{0.0pt}{}{C_j(r)\text{ }\mathrm{\Psi }_j(r)}{\beta _j(k,r)\text{ }\phi _0(k,r)}\right|\left[detC_j(r)\right]^1.$$
Here, a $`j\times j`$ matrix $`C_j(r)I+\underset{0}{\overset{r}{}}R^{(j)}(s)𝑑s`$ ($`I`$ is the unit matrix) with the elements of $`R^{(j)}`$ being
(10)
$$R_{lm}^{(j)}=C_l\phi _0(i\gamma _l,r)\phi _0(i\gamma _m,r),\text{ (}l,m=1,2,\mathrm{},j\text{)}$$
a column vector (do not confuse the norming constants $`C_l`$ with the matrices $`C_j(r)`$)
(11)
$$\mathrm{\Psi }_j(r)=\left(\genfrac{}{}{0.0pt}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{C_1\phi _0(i\gamma _1,r)}{C_2\phi _0(i\gamma _2,r)}}}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{C_j\phi _0(i\gamma _j,r)}}}}}}}\right),$$
and a row vector
(12)
$$\beta _j(k,r)=\left(\underset{0}{\overset{r}{}}\phi _0(i\gamma _1,s)\phi _0(k,s)𝑑s\text{ }\underset{0}{\overset{r}{}}\phi _0(i\gamma _2,s)\phi _0(k,s)𝑑s\text{ …}\underset{0}{\overset{r}{}}\phi _0(i\gamma _j,s)\phi _0(k,s)𝑑s\right)$$
have been introduced.
To be more specific, let us examine the simplest case of just one bound state ($`j=1`$). Then $`C_j(r)`$ reduces to scalar and one gets
(13)
$$V_1(r)=V_0(r)2C\left\{\mathrm{ln}\left[1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s\right]\right\}^{\prime \prime }.$$
As can be shown , the regular solution $`\phi _1(i\gamma _1,r)`$ related to $`V_1(r)`$ (i.e., the real confined eigenfunction) reads
(14)
$$\phi _1(i\gamma _1,r)=\frac{\phi _0(i\gamma _1,r)}{1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s}.$$
Let us prove that $`C_1`$, as needed, is the norming constant of $`\phi _1(i\gamma _1,r)`$. Indeed, slightly rearranging Eq. (14): $`1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s={\displaystyle \frac{\phi _0(i\gamma _1,r)}{\phi _1(i\gamma _1,r)}}`$, differentiating both sides: $`C_1\phi _0^2(i\gamma _1,r)={\displaystyle \frac{\phi _0^{}\phi _1\phi _0\phi _1^{}}{\phi _1^2}},`$ i.e., $`C_1\phi _1^2(i\gamma _1,r)=\left(\frac{1}{\phi _0}\right)^{}\phi _1{\displaystyle \frac{1}{\phi _0}}\phi _1^{}=\left({\displaystyle \frac{\phi _1}{\phi _0}}\right)^{}`$, and integrating both sides of the latter equation, one gets $`\underset{0}{\overset{\mathrm{}}{}}C_1\phi _1^2(i\gamma _1,r)𝑑r=1`$, because any regular solution $`\phi (r)r`$ as $`r0`$, and therefore, $`lim_{r0}\left({\displaystyle \frac{\phi _1}{\phi _0}}\right)=1`$.
Now, let us analyze a more complicated case of two bound states ($`j=1`$). According to Eqs. (8) and (13),
(15)
$$V_2(r)V_1(r)=V_2(r)V_0(r)\left[V_1(r)V_0(r)\right]=2C\left\{\mathrm{ln}\left[\frac{detC_2(r)}{1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s}\right]\right\}^{\prime \prime }.$$
Here, the $`2\times 2`$ matrix $`C_2(r)=\left({\displaystyle \genfrac{}{}{0.0pt}{}{C_{11}(r)\text{ }C_{12}(r)}{C_{21}(r)\text{ }C_{22}(r)}}\right)`$ has the elements
(16)
$$C_{11}(r)=1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s,\text{ }C_{12}(r)=C_1\underset{0}{\overset{r}{}}\phi _0(i\gamma _1,s)\phi _0(i\gamma _2,s)𝑑s$$
$$C_{21}(r)=\frac{C_2}{C_1}C_{12}(r),\text{ }C_{22}(r)=1+C_2\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _2,s)𝑑s,$$
and therefore, the argument of the logarithm in Eq. (15) reads
(17)
$$f(r)C_{22}\frac{C_1C_2\left[\underset{0}{\overset{r}{}}\phi _0(i\gamma _1,s)\phi _0(i\gamma _2,s)𝑑s\right]^2}{1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s}=$$
$$=1+C_2\left\{\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _2,s)𝑑s\frac{C_1\left[\beta _{12}(r)\right]^2}{1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s}\right\},$$
where $`\beta _{12}(r)\underset{0}{\overset{r}{}}\phi _0(i\gamma _1,s)\phi _0(i\gamma _2,s)𝑑s`$.
As can be easily proved,
(18)
$$f(r)=1+C_2\underset{0}{\overset{r}{}}\phi _1^2(i\gamma _2,s)𝑑s.$$
Indeed, using Eq. (14) and a more general formula
$$\phi _1(k,r)=\phi _0(k,r)C_1\phi _0(i\gamma _1,r)\times \frac{\underset{0}{\overset{r}{}}\phi _0(i\gamma _1,s)\phi _0(k,s)𝑑s}{1+C_1\underset{0}{\overset{r}{}}\phi _0^2(i\gamma _1,s)𝑑s},$$
and taking $`k=i\gamma _2`$, one gets
$$\phi _1(i\gamma _2,r)=\phi _0(i\gamma _2,r)C_1\phi _1(i\gamma _1,r)\beta _{12}(r),$$
i.e.,
$`\phi _1^2(i\gamma _2,r)`$ $`=\phi _0^2(i\gamma _2,r)+\left[C_1\phi _1(i\gamma _1,r)\beta _{12}(r)\right]^22C_1\phi _0(i\gamma _2,r)\phi _1(i\gamma _1,r)\beta _{12}(r)=`$
$`=\phi _0^2(i\gamma _2,r)C_1\left[{\displaystyle \frac{\phi _1(i\gamma _1,r)}{\phi _0(i\gamma _1,r)}}\beta _{12}^2(r)\right]^{},`$
which, after integrating from 0 to $`r`$, and taking account of Eq. (14), proves the claim.
Thus, in full analogy with Eq. (13), one can write
(19)
$$V_2(r)=V_1(r)2C\left\{\mathrm{ln}\left[1+C_2\underset{0}{\overset{r}{}}\phi _1^2(i\gamma _2,s)𝑑s\right]\right\}^{\prime \prime },$$
and more generally,
(20)
$$V_k(r)=V_{k1}(r)2C\left\{\mathrm{ln}\left[1+C_k\underset{0}{\overset{r}{}}\phi _{k1}^2(i\gamma _k,s)𝑑s\right]\right\}^{\prime \prime },\text{ }k=1,2,\mathrm{},n.$$
This way one gets rid of inconvenient matrix equations and can move towards the goal step-by-step, introducing one bound state at a time.
The described procedure can be reverted , i.e., starting from the potential $`V_k(r)`$, one can construct the auxiliary potential $`V_{k1}(r)`$, removing the zeroth level ($`E_k`$ according to the numeration used). The new potential reads
(21)
$$V_{k1}(r)=V_k(r)+2C\left\{\frac{2\mathrm{\Psi }_0^{(k)}(r)\left(\mathrm{\Psi }_0^{(k)}(r)\right)^{}}{\underset{r}{\overset{\mathrm{}}{}}\left(\mathrm{\Psi }_0^{(k)}(s)\right)^2𝑑s}+\left[\frac{\left(\mathrm{\Psi }_0^{(k)}(r)\right)^2}{\underset{r}{\overset{\mathrm{}}{}}\left(\mathrm{\Psi }_0^{(k)}(s)\right)^2𝑑s}\right]^2\right\},\text{ (}k=n,n1,\mathrm{},1\text{)}$$
where $`\mathrm{\Psi }_0^{(k)}(r)`$ is the eigenfunction of the zeroth level, not necessarily normalized (note that the norming constant is absent here). One can remove, one by one, all bound states until he comes to the potential $`V_0(r)`$ with no bound states. An illustration to Eq. (21) can be seen in Fig. 1.
Thus, we have demonstrated that there is no need to solve Eq. (1) all at once. It is probably much easier to first solve this equation for the auxiliary potential $`V_0(r)`$ and then, step-by-step, introduce the bound states as described above. The Jost functions $`F_n(k)`$ and $`F_0(k)`$ related to the potentials $`V_n(r)`$ and $`V_0(r)`$, respectively, are connected by a simple formula
(22)
$$F_n(k)=F_0(k)\underset{j=1}{\overset{n}{}}\left(\frac{ki\gamma _j}{k+i\gamma _j}\right),$$
which demonstrates that these potentials, indeed, have exactly the same spectral densities for positive energies . On the other hand, Eq. (22) illustrates the general rule that any zero of the Jost function corresponds to a bound state.
To end this section, let us recall important asymptotic formulas (see for details)
(23)
$$V_n(r)V_0(r)4C\left(\underset{j=1}{\overset{n}{}}C_j\right)r,\text{ }r0,$$
(24)
$$V_n(r)V_0(r)\frac{2C}{C_1}\left(2\gamma _1\right)^5\mathrm{exp}(2\gamma _1r),\text{ }r\mathrm{}.$$
From Eq. (23) one can infer that in the immediate vicinity of the zero point $`r=0`$ the potentials $`V_n(r)`$ and $`V_0(r)`$ practically coincide. Indeed, all norming constants $`C_j=\left[\underset{0}{\overset{\mathrm{}}{}}\phi _j^2(i\gamma _j,r)𝑑r\right]^1`$are extremely small quantities, since they are related to the regular solutions proportional to $`r`$ as $`r0`$, and vanishing as $`r\mathrm{}`$, but achieving very large absolute values between these asymptotic regions.
## 3. Krein method
Thus, in view of the results of previous section, our main goal is to accurately ascertain the auxiliary potential $`V_0(r)`$ with no bound states, provided that its spectral density $`{\displaystyle \frac{d\rho _0(E)}{dE}}=\pi ^1\sqrt{E}\left|F(E)\right|^2`$ is known. This is indeed the case, since we have carefully calculated the modulus of the Jost function for the reference potential $`V(r)`$ , and the whole idea is to use the same quantity ($`\left|F(E)\right|`$) for the auxiliary potential $`V_0(r)`$ as well. Since there are no bound states, the kernel of the Gelfand-Levitan equation, according to Eq. (4), becomes analytically very simple:
(25)
$$G(r,r^{})=H(rr^{})H(r+r^{}),$$
where a new function
(26)
$$H(r)\pi ^1\underset{0}{\overset{\mathrm{}}{}}g(k)\mathrm{cos}(kr)𝑑k,$$
with $`g(k)`$ given by Eq. (5), has been introduced. The function $`H(r)=H(r)`$ is very important for the further treatment, wherefore a special designation, Krein $`H`$-function, will be used for this item henceforward. A motivation for this name stems from a very useful method for solving the inverse problem, which has been elaborated by Krein . This method is based on a Fredholm-type integral equation
(27)
$$\mathrm{\Gamma }_{2r}(r^{})+H(r^{})+\underset{0}{\overset{2r}{}}\mathrm{\Gamma }_{2r}(s)H(sr^{})𝑑s=0.$$
It can be shown that the solutions of Gelfand-Levitan and Krein equations are linked by a simple formula (cf. with Eq. (25))
(28)
$$K(r,r^{})=\mathrm{\Gamma }_{2r}(rr^{})\mathrm{\Gamma }_{2r}(r+r^{}),$$
and the desired potential reads
(29)
$$V_0(r)=4C\left\{\left[G(x)\right]^2\frac{dG(x)}{dx}\right\},\text{ }x2r,$$
where $`G(x)\mathrm{\Gamma }_{2r}(2r).`$ A special notation, Krein $`G`$-function, will used for this quantity henceforward.
Without any doubt, Krein method is well suited for our purposes, but there is still a lot of analytical and computational-technical work to do. First, we have to accurately ascertain the Krein $`H`$-function, which, according to Eq. (26), is simply the Fourier cosine transform of the characteristic function $`g(k)`$ whose full energy dependence has already been ascertained . One might think that the $`H`$-function can be determined using the well-known fast Fourier transform technique. This, however, is an erroneous view, because (as we demonstrate below) the $`H`$-function has to be calculated in a wide distance range with very small step, to ensure the correct asymptotic behavior of the resulting potential. Fortunately, the problem can be solved accurately and quite easily with the help of solely analytic means.
### 3.1. Calculation of $`H`$-function
The function $`g(k)`$ is shown in Fig. 2. In addition to the overall curve, some characteristic slices can be seen in the insets. There exists a range $`k(0,k_1)`$ where $`g(k)=1`$ (with high accuracy), and for this range one immediately gets the relevant component of the $`H`$-function
$$H_0(r)=\frac{\mathrm{sin}(k_1r)}{\pi r}.$$
Now, let us see how the $`g(k)`$ curve passes through the ”critical” region around $`k_0=\sqrt{V(0)/C}`$. As it happens (see the upper inset in Fig. 2), in a narrow range $`k(k_1,k_2)`$ ($`k_0`$ is also located within this range) $`g(k)+1`$ can be nicely approximated by a sum of several Gaussians:
(30)
$$g(k)=1+\underset{j}{}a_j\mathrm{exp}\left[\frac{1}{2}\left(\frac{k\stackrel{~}{k}_j}{b_j}\right)^2\right],$$
where the parameters $`a_j,b_j`$ and $`\stackrel{~}{k}_j`$ can be determined from a least-squares fit. For the reference potential examined here, four such components have been introduced, and the value $`k_2=`$ 19230 Å<sup>-1</sup> has been chosen. Thus, in the above formula one can replace $`k_1`$ with $`k_2`$ (due to -1 in Eq. (30)), i.e.,
(31)
$$H_1(r)=\frac{\mathrm{sin}(k_2r)}{\pi r},$$
while the sum of Gaussians gives another component of the Krein $`H`$-function, $`H_2(r)=`$ $`\underset{j}{}H_2^{(j)}(r),`$ where all constituents $`H_2^{(j)}(r)`$ can be ascertained analytically. Indeed, one can introduce a new independent variable $`x{\displaystyle \frac{k\stackrel{~}{k}_j}{\sqrt{2}b_j}},`$ and calculate
(32)
$$H_2^{(j)}(r)=\frac{\sqrt{2}a_jb_j}{\pi }\underset{x_1}{\overset{x_2}{}}\mathrm{exp}(x^2)\mathrm{cos}(\sqrt{2}b_jxr+\stackrel{~}{k}_jr)dx=\frac{a_jb_j}{\sqrt{2\pi }}\mathrm{exp}(\frac{1}{2}b_j^2r^2)\times $$
$$\times \left\{\mathrm{cos}(\stackrel{~}{k}_jr)\left[\mathrm{Re}\mathrm{erf}(y_2)\mathrm{Re}\mathrm{erf}(y_1)\right]\mathrm{sin}(\stackrel{~}{k}_jr)\left[\mathrm{Im}\mathrm{erf}(y_2)\mathrm{Im}\mathrm{erf}(y_1)\right]\right\},$$
where $`x_l{\displaystyle \frac{k_l\stackrel{~}{k}_j}{\sqrt{2}b_j}}`$ and $`y_l=x_l{\displaystyle \frac{ib_jr}{\sqrt{2}}}`$($`l=1,2`$ and $`i`$ is the imaginary unit). The error function, for any complex argument $`z,`$ can be evaluated in terms of the confluent hypergeometric functions
(33)
$$\mathrm{erf}(z)=\frac{\sqrt{2}}{\pi }z\mathrm{\Phi }(\frac{1}{2},\frac{3}{2};z^2),$$
where $`\mathrm{\Phi }(a,c;x)=1+{\displaystyle \frac{ax}{1!c}}+{\displaystyle \frac{a(a+1)x^2}{2!c(c+1)}}+\mathrm{}`$ For large arguments another expression is more convenient:
(34)
$$\mathrm{erf}(z)=1\frac{\mathrm{exp}(z^2)}{\sqrt{\pi }z}\left[1\frac{1}{2z^2}+\frac{13}{\left(2z^2\right)^2}\frac{135}{\left(2z^2\right)^3}+\mathrm{}\right].$$
Consequently, any constituent of the Krein $`H`$-function expressed by Eq. (32) can be easily ascertained with any desired accuracy.
Next one can introduce an arbitrary (but still reasonable) boundary point $`k_3`$, and approximate $`g(k)`$ in the range $`k(k_2,k_3)`$ as follows:
(35)
$$g(k)=\underset{j}{}a_j\mathrm{exp}\left[b_j(k\stackrel{~}{k}_j)\right].$$
This brings along another component of the $`H`$-function $`H_3(r)=`$ $`\underset{j}{}H_3^{(j)}(r)`$ with constituents
(36)
$$H_3^{(j)}(r)=\frac{a_j}{\pi }\underset{k_2}{\overset{k_3}{}}\mathrm{exp}[b_j(k\stackrel{~}{k}_j)]\mathrm{cos}\left(kr\right)dk=\frac{a_j\mathrm{exp}(b_j\stackrel{~}{k}_j)}{\pi \left(b_j^2+r^2\right)}\times $$
$$\times \left\{\mathrm{exp}(b_jk_3)\left[r\mathrm{sin}\left(k_3r\right)b_j\mathrm{cos}\left(k_3r\right)\right]\mathrm{exp}(b_jk_2)\left[r\mathrm{sin}\left(k_2r\right)b_j\mathrm{cos}\left(k_2r\right)\right]\right\}.$$
Introducing new suitable boundary points $`k_4`$, $`k_5`$, etc., the approximation of $`g(k)`$ in the form of Eq. (35) can be continued until the conventional starting point $`k_a`$ of the asymptotic region (see the treatment below). The number of these boundary points, as well as the number of exponents in any particular interval, is, of course, a subject for probes and trials. We have introduced four such intervals and used a different three-exponent approximation in any of them. As has been carefully checked, this ensures the accuracy of at least 6 significant digits for the calculated Krein $`H`$-function in the whole physical domain. The parameters of all components are given in Table 1.
For the remaining part of the $`k`$-space the asymptotic formula
(37)
$$\mathrm{ln}\left|F(k)\right|=\frac{a_2}{k^2}+\frac{a_4}{k^4}+\frac{a_6}{k^6}+\mathrm{},\text{ }kk_a,$$
can be used, where
(38)
$$a_2=\frac{V(0)}{4C},\text{ }a_4=\frac{2\left[V(0)\right]^2CV^{\prime \prime }(0)}{16C^2}.$$
The coefficient $`a_6,`$ as well as the coefficients for higher-order terms, can also be calculated in terms of the reference potential and its derivatives, but the resulting expressions are rather complicated and inconvenient for practical use. Instead, we only introduced just one additional term, $`{\displaystyle \frac{a_6}{k^6}}`$, and determined the coefficient $`a_6`$ from the general demand (see Eq. (23)) that the potentials $`V_0(r)`$ and $`V(r)`$ should coincide as $`r0`$. Such physically well motivated choice of $`a_6`$ is indeed possible, as will be explained below. Thus, within this approximation, the asymptotic part of $`g(k)`$ reads
(39)
$$g(k)=\frac{b_1}{k^2}+\frac{b_2}{k^4}+\frac{b_3}{k^6},\text{ }kk_a,$$
where
(40)
$$b_1=2a_2\text{}b_2=2(a_4a_2^2)\text{}b_3=2(a_62a_2a_4+\frac{2}{3}a_2^3).$$
The relevant asymptotic component of the $`H`$-function becomes
$$H_a(r)\pi ^1\underset{k_a}{\overset{\mathrm{}}{}}g(k)\mathrm{cos}(kr)𝑑k.$$
Since the expression for $`H_a(r)`$ will contain the sine integral $`\mathrm{Si}(x_a)\underset{0}{\overset{x_a}{}}{\displaystyle \frac{\mathrm{sin}t}{t}}𝑑t,`$ let us recall a useful formula
(41)
$$\mathrm{Si}(x)=\frac{\pi }{2}\frac{i}{2}\mathrm{exp}(ix)\mathrm{\Psi }(1,1;ix)+\frac{i}{2}\mathrm{exp}(ix)\mathrm{\Psi }(1,1;ix),$$
where $`i`$ is the imaginary unit, and the function $`\mathrm{\Psi }(a,c;z)`$ is a particular solution of the confluent hypergeometric equation introduced by Tricomi. For a large argument it can be evaluated from the asymptotic series
(42)
$$\mathrm{\Psi }(a,c;z)=z^a\underset{n=0}{\overset{N}{}}\frac{(a)_n(ac+1)_n}{n!(z)^n},$$
where $`(a)_n\mathrm{\Gamma }(a+n)/\mathrm{\Gamma }(a)=a(a+1)(a+2)\mathrm{}(a+n1)`$ is the Pochhammer symbol, and $`N`$ must not be too large. If Eq. (41) is usable, i.e., in the case of sufficiently large $`r`$, one gets the following expression:
(43)
$$\pi H_a(r)=\left(\frac{b_1A_1}{1!k_a}\frac{b_2A_2}{3!k_a^3}+\frac{b_3A_3}{5!k_a^5}\right)\frac{\mathrm{sin}(x_a)}{x_a}$$
$$\left(\frac{b_1B_1}{1!k_a}\frac{b_2B_2}{3!k_a^3}+\frac{b_3B_3}{5!k_a^5}\right)\frac{\mathrm{cos}(x_a)}{x_a^2},\text{ }x_ak_ar>>1,$$
where
$$A_i=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(1)^{i+j}\left[2(i+j)1\right]!}{\left(x_a\right)^{2j}},\text{ }B_i=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(1)^{i+j}\left[2(i+j)\right]!}{\left(x_a\right)^{2j}},\text{ }i=1,2,3.$$
If Eq. (41) cannot be used, one can apply the universal expansion
(44)
$$\mathrm{Si}(x_a)=x_a\frac{\left(x_a\right)^3}{33!}+\frac{\left(x_a\right)^5}{55!}\mathrm{},$$
to get another formula
(45)
$$\pi H_a(r)=\left(\frac{b_1r}{1!}\frac{b_2r^3}{3!}+\frac{b_3r^5}{5!}\right)\left[\mathrm{Si}(x_a)\frac{\pi }{2}\right]+$$
$$\left(\frac{b_1X_0}{1!k_a}\frac{b_2X_1}{3!k_a^3}+\frac{b_3X_2}{5!k_a^5}\right)\mathrm{cos}(x_a)r\left(\frac{b_2Y_0}{3!k_a^2}\frac{b_3Y_1}{5!k_a^4}\right)\mathrm{sin}(x_a),$$
(46)
$$X_i=\underset{j=0}{\overset{i}{}}(1)^{ij}\left[2(ij)\right]!\left(x_a\right)^{2j},\text{ }Y_i=\underset{j=0}{\overset{i}{}}(1)^{ij}\left[2(ij)+1\right]!\left(x_a\right)^{2j}.$$
From Eq. (45) one can infer that near the zero point $`r=0`$
$$H_a(r)\pi ^1\left(\frac{b_1}{k_a}+\frac{b_2}{3k_a^3}+\frac{b_3}{5k_a^5}\right)\frac{b_1r}{2},$$
and consequently, according to Eqs. (38) and (40), the derivative at zero point $`\left[H_a(0)\right]^{}={\displaystyle \frac{b_1}{2}}={\displaystyle \frac{V(0)}{4C}}.`$ Let us prove that the same relation holds for the overall Krein $`H`$-function. Indeed, as $`r0`$, one can always choose a value for $`k_a`$ which is large enough, so that Eq. (39) can be used, but on the other hand, small enough, so that $`\mathrm{sin}(kr)kr`$, if $`k(0,k_a].`$ Integrating by parts (note that $`g(\mathrm{})=0`$), one gets from Eq. (26)
$$\pi H(r0)=\underset{0}{\overset{k_a}{}}g^{}(k)k𝑑k+\frac{2}{r}\underset{k_a}{\overset{\mathrm{}}{}}\left(\frac{b_1}{k^3}+\frac{2b_2}{k^5}+\frac{3b_3}{k^7}+\mathrm{}\right)\mathrm{sin}(kr)𝑑k.$$
Only the second term gives contribution to $`H^{}(0)`$, and after few elementary transformations one comes to the desired result
(47)
$$H(r0)=\frac{b_1r}{2}+\pi ^1\underset{0}{\overset{\mathrm{}}{}}g(k)𝑑k,$$
which proves that
(48)
$$H^{}(0)=\frac{b_1}{2}=\frac{V(0)}{4C}.$$
Thus, in the case examined here, the Krein $`H`$-function can be, indeed, ascertained analytically. To get the overall $`H`$-function, one just sums the components described above: $`H(r)=\underset{j}{}H_j(r).`$ The result for different distance regions is shown in Figs. 3 and 4. As can be seen, $`H(r)`$ is a rapidly oscillating function with decaying amplitude. The period of oscillations stabilizes quite rapidly and remains very close to the characteristic value $`L={\displaystyle \frac{2\pi }{k_2}}`$ (see Eq. (31)), while the amplitude slowly approaches zero as $`r\mathrm{}.`$
### 3.2. Solution of Krein equation
Before setting about the main task, we have to fix the coefficient $`b_3`$ in Eq. (39). To this end, let us take into consideration that according to Eq. (27), $`G(0)=H(0).`$ Also, as repeatedly mentioned (see Eq. (23)), the potentials $`V_0(r)`$ and $`V(r)`$ should coincide as $`r0`$. One therefore can replace $`V_0(r)`$ with $`V(r)`$ in Eq. (29), when studying the region very close to $`r=0`$. In addition, the Krein $`G`$-function there can be approximated by a quadratic function: $`G(x)=a+bx+cx^2`$ ($`x=2r`$), and consequently, $`G^{}(x)=b+2cx.`$ In this region a pseudo-Morse approximation for the reference potential is used: $`V(r)=V_0+A_0\mathrm{exp}(\alpha _0x)\sqrt{A_0\epsilon _0}\mathrm{exp}(\alpha _0x/2).`$ Here, $`A_0D_0\mathrm{exp}(2\alpha _0r_0)`$ and $`\epsilon _0D_0/4`$ (see for more details). Therefore, the parameters $`a,`$ $`b,`$ $`c`$ can be ascertained directly from Eq. (29). Indeed, one uses the relations $`N_1a^2b={\displaystyle \frac{V(0)}{4C}},`$ $`N_2cab={\displaystyle \frac{\alpha _0(A_0\sqrt{A_0\epsilon _0}/2)}{8C}},`$ $`N_3b^2+2ac={\displaystyle \frac{\alpha _0^2(4A_0\sqrt{A_0\epsilon _0})}{32C}},`$ then solves the equation $`\left(a^2N_1\right)^2+2a\left[a(a^2N_1)+N_2\right]N_3=0`$ to ascertain the parameter $`a=G(0)=H(0)`$, and thereafter finds $`b`$ and $`c`$. Having fixed these parameters, and also the value of $`k_a`$, the coefficient $`b_3`$ can be quite easily determined on a trial-by-trial basis. This way the value $`b_3=`$ -5.883044$``$10<sup>24</sup><sup>-6</sup>) has been fixed, which corresponds to $`k_a=75000`$ Å<sup>-1</sup>.
Now, let us proceed with solution of Eq. (27). This equation can be discreticized using, for example, a four-point quadrature rule
(49)
$$\underset{kh}{\overset{(k+3)h}{}}f(x)𝑑x=\frac{3h}{8}f(kh)+\frac{9h}{8}f(\left[k+1\right]h)+\frac{9h}{8}f(\left[k+2\right]h)+\frac{3h}{8}f(\left[k+3\right]h),$$
which is exact for $`f(x)`$ a cubic polynomial. Applying this formula to Eq. (27), one gets the following system of linear equations:
(50)
$$\left(1+\mathrm{\Delta }H_0\right)\mathrm{\Gamma }_{3n,0}+3\left(\mathrm{\Delta }H_1\mathrm{\Gamma }_{3n,1}+\mathrm{\Delta }H_2\mathrm{\Gamma }_{3n,2}\right)+2\mathrm{\Delta }H_3\mathrm{\Gamma }_{3n,3}+$$
$$\text{ }+3\left(\mathrm{\Delta }H_4\mathrm{\Gamma }_{3n,4}+\mathrm{\Delta }H_5\mathrm{\Gamma }_{3n,5}\right)+2\mathrm{\Delta }H_6\mathrm{\Gamma }_{3n,6}+\mathrm{}+\mathrm{\Delta }H_{3n}\mathrm{\Gamma }_{3n,3n}=H_0$$
$$\mathrm{\Delta }H_1\mathrm{\Gamma }_{3n,0}+\left(1+3\mathrm{\Delta }H_0\right)\mathrm{\Gamma }_{3n,1}+3\mathrm{\Delta }H_1\mathrm{\Gamma }_{3n,2}+2\mathrm{\Delta }H_2\mathrm{\Gamma }_{3n,3}+\mathrm{}+\mathrm{\Delta }H_{3n1}\mathrm{\Gamma }_{3n,3n}=H_1$$
$$\text{ }\mathrm{\Delta }H_2\mathrm{\Gamma }_{3n,0}+3\mathrm{\Delta }H_1\mathrm{\Gamma }_{3n,1}+\left(1+3\mathrm{\Delta }H_0\right)\mathrm{\Gamma }_{3n,2}+\mathrm{}=H_2$$
$$\text{ }\mathrm{\Delta }H_3\mathrm{\Gamma }_{3n,0}+3(\mathrm{\Delta }H_2\mathrm{\Gamma }_{3n,1}+\mathrm{\Delta }H_1\mathrm{\Gamma }_{3n,2})+\left(1+2\mathrm{\Delta }H_0\right)\mathrm{\Gamma }_{3n,3}+\mathrm{}=H_3$$
$``$
$``$
$``$
$$\mathrm{\Delta }H_{3n}\mathrm{\Gamma }_{3n,0}+3(\mathrm{\Delta }H_{13n}\mathrm{\Gamma }_{3n,1}+\mathrm{\Delta }H_{23n}\mathrm{\Gamma }_{3n,2})+\left(1+\mathrm{\Delta }H_0\right)\mathrm{\Gamma }_{3n,3n}+\mathrm{}=H_{3n}$$
or, in a more compact form,
(51)
$$\mathrm{\Gamma }_{3n,k}+\mathrm{\Delta }(H_k\mathrm{\Gamma }_{3n,0}+3H_{1k}\mathrm{\Gamma }_{3n,1}+3H_{2k}\mathrm{\Gamma }_{3n,2}+2H_{3k}\mathrm{\Gamma }_{3n,3}+\mathrm{}$$
$$+H_{3nk}\mathrm{\Gamma }_{3n,3n})=H_k,\text{ }k=0,1,2,\mathrm{},3n.$$
Here $`\mathrm{\Delta }{\displaystyle \frac{3h}{8}},H_k=H_k=H(kh),`$ and $`\mathrm{\Gamma }_{3n,k}\mathrm{\Gamma }_{2r}(kh)`$ (see Eq. (27)).
For any argument $`x=2r=3nh,`$ Eq. (50) can be solved with the help of Gaussian elimination procedure, which is appropriate here, since only the last element of the solution vector, $`\mathrm{\Gamma }_{3n,3n}=\mathrm{\Gamma }_{2r}(2r)=G(x)`$ is actually needed to calculate the potential according to Eq. (29). Alternatively, one may rewrite Eq. (51) in a matrix form
(52)
$$\left(I+\mathrm{\Delta }U\right)G=H,$$
where $`I`$ denotes (3$`n+1`$)$`\times `$(3$`n+1`$) unit matrix,
$$G\left(\genfrac{}{}{0.0pt}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{\mathrm{\Gamma }_{3n,0}}{\mathrm{\Gamma }_{3n,1}}}}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{\mathrm{\Gamma }_{3n,3n}}}}}}}}\right),\text{ }H\left(\genfrac{}{}{0.0pt}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{H_0}{H_1}}}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{H_{3n}}}}}}}}\right),$$
and
$$U\left(\genfrac{}{}{0.0pt}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{H_0\text{ 3}H_1\text{ 3}H_2\text{ 2}H_3\text{ … 3}H_{3n1}\text{ }H_{3n}}{H_1\text{ 3}H_0\text{ .3}H_1\text{ 2}H_2\text{ …3}H_{3n2}\text{ }H_{3n1}}}}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{{\displaystyle \genfrac{}{}{0.0pt}{}{}{H_{3n}\text{ 3}H_{3n1}\text{ 3}H_{3n2}\text{ 2}H_{3n3}\text{ … 3}H_1\text{ }H_0}}}}}}}\right).$$
The solution of Eq. (52) reads
(53)
$$G=\left(I+\mathrm{\Delta }U\right)^1H=H+\mathrm{\Delta }UH\mathrm{\Delta }^2U^2H+\mathrm{},$$
and therefore,
(54) $`\mathrm{\Gamma }_{3n,3n}`$ $`=H_{3n}+\mathrm{\Delta }\left(H_{3n}H_0+\text{3}H_{3n1}H_1+\text{3}H_{3n2}H_2+\text{2}H_{3n3}H_3+\mathrm{}+H_0H_{3n}\right)`$
$`\mathrm{\Delta }^2{\displaystyle \underset{i,j}{}}g_ig_jH_{3ni}H_{\left|ij\right|}H_j+\mathrm{}`$
The coefficients $`g_i`$ (here and henceforward) are defined as follows: $`g_i=1,`$ if $`i=0`$ or $`i=m3n;`$ $`g_i=2,`$ if $`i=3k`$ and $`k=1,2,\mathrm{},n1;`$ $`g_i=3`$ in any other case.
According to Eqs. (6) and (28), $`V(r)=4C\left\{{\displaystyle \frac{d\mathrm{\Gamma }_x(0)}{dx}}{\displaystyle \frac{d\mathrm{\Gamma }_x(x)}{dx}}\right\}.`$ Comparing this with Eq. (29), one gets an important relation
(55)
$$\frac{d\mathrm{\Gamma }_x(0)}{dx}=\left[G(x)\right]^2,\text{ (}x=2r\text{)}$$
Thus, instead of searching for the last element of the solution vector $`G`$, one may calculate its first element $`\mathrm{\Gamma }_{m0}`$ ($`m=3n`$), which in some sense is more convenient. Indeed, using such an approach, we can make use of the results of previous calculations. Namely, as can be proved
(56)
$$\mathrm{\Gamma }_{m+3,0}=\mathrm{\Gamma }_{m0}+\mathrm{\Delta }\left(H_m^2+\text{3}H_{m+1}^2+\text{3}H_{m+2}^2+H_{m+3}^2\right)\mathrm{\Delta }^2S_m+\mathrm{}$$
where
(57)
$$S_m=2\underset{j=0}{\overset{m1}{}}g_jH_j(H_mH_{mj}+3H_{m+1}H_{m+1j}+3H_{m+2}H_{m+2j}+H_{m+3}H_{m+3j})+$$
$$H_0(3H_m^2+9H_{m+1}^2+9H_{m+2}^2+H_{m+3}^2)+6H_1(H_{m+3}H_{m+2}+3H_{m+2}H_{m+1}+2H_{m+1}H_m)+$$
$$+6H_2(H_{m+3}H_{m+1}+2H_{m+2}H_m)+4H_3H_{m+3}H_m,\text{ }m=0,1,2,\mathrm{}$$
From Eqs. (55) and (56) one gets another formula
(58)
$$8\left[G(x)\right]^2=H_m^2+\text{3}H_{m+1}^2+\text{3}H_{m+2}^2+H_{m+3}^2\mathrm{\Delta }S_m+\mathrm{},\text{ }x=mh,$$
which may prove very useful, if the higher order terms can be ignored.
In Fig. 5 one can see the calculated Krein $`G`$-function in the range from 0 to 10<sup>-6</sup> Å. Eq. (50) has been solved by Gaussian elimination, and an extremely small step $`h=`$ 10<sup>-9</sup> was used to ensure high accuracy. For comparison, another curve (actually, almost straight line) is depicted, which exactly corresponds to the reference potential $`V(r)`$, i.e., it is the solution of Eq. (29) interpreted as Riccati equation:
(59)
$$\frac{dG(x)}{dx}=\left[G(x)\right]^2\frac{V(x/2)}{4C},\text{ }G(0)=H(0).$$
The curves practically coincide, which is a clear evidence of the validity of the approach. Indeed, this seemingly trivial calculation is, in fact, very sensitive to even minor inaccuracies in calculating the Krein $`H`$-function, which would result in undesired and unphysical discrepancies, e.g., oscillations, of the $`G`$-function. Since no such discrepancies are seen at small distances, one may expect that both the Jost function and the $`H`$-function have been ascertained quite correctly. This in turn gives ground to hope that the potential can be correctly ascertained at longer distances as well. This, however, is not at all an easy task, and was not attempted here.
## 4. Conclusion
In this paper, as well as in the previous one , the treatment was more concentrated on principles rather than the methods and techniques of computation. We demonstrated that the proposed reference potential approach enables one to accurately ascertain the important spectral characteristics, needed to uniquely solve the quantum-mechanical inverse problem. This way, one may get reasonable initial guesses to the real spectral characteristics of the system, which can be used, for example, to calculate another potential (Bargmann potential) whose Jost function differs from the initial one only by a rational factor. Let us briefly analyze the simplest case when this factor reads
(60)
$$\frac{kia}{k+ib}=\frac{kia}{k+ia}\frac{k+ia}{k+ib},$$
where both, $`a`$ and $`b,`$ are real and positive. In fact, it means that a discrete eigenvalue $`b^2`$ is replaced by $`a^2.`$ As we see, Eq. (60) involves two operations. One of them, introducing a new bound state $`E=a^2`$, can be performed as described in Section 2, while the second factor, $`{\displaystyle \frac{k+ia}{k+ib}}`$, causes an additional deformation of the initial potential
(61)
$$\mathrm{\Delta }V(r)V_2(r)V_1(r)=2C\left\{\mathrm{ln}\frac{W[f_1(ia,r),\phi _1(ib,r)]}{b^2a^2}\right\}^{\prime \prime },$$
where the symbol $`W`$ denotes Wronskian determinant and $`f_1(ia,r)`$ is the Jost solution of the Schrödinger equation ($`f_1(ia,r)\mathrm{exp}(ar)`$ as $`r\mathrm{}`$). Both $`f_1(ia,r)`$ and the regular solution $`\phi _1(ib,r)`$ are related to the potential $`V_1(r).`$ Thus, if $`a^2`$ is expected to be more realistic eigenvalue than $`b^2`$, one may hope that $`V_2(r)`$ is more realistic potential than $`V_1(r)`$ as well.
The proposed approach is based on various analytic procedures. For example, we used piecewise analytic approximation for the characteristic function $`g(k)`$ (see Eq. (5) and Fig. 2), and derived simple formulas for relevant constituents of the Krein $`H`$-function. In addition, we proved some general relations, Eqs. (47) and (48), regarding $`H`$-function near the zero point $`r=0,`$ and fixed an appropriate value for $`G(0)=H(0)=\pi ^1\underset{0}{\overset{\mathrm{}}{}}g(k)𝑑k`$ (see Section 3.2), which ensures the correct behavior of the resulting potential in this region.
One of our goals was to promote Krein method to solve the inverse problem. This method is especially suitable to ascertain the auxiliary potential $`V_0(r)`$ with no bound states, starting from the known Jost function. Having ascertained $`V_0(r),`$ one can build up a series of auxiliary potentials: $`V_1(r),`$ $`V_2(r),\mathrm{},`$ $`V_n(r)`$, introducing, one by one, all bound states with known eigenvalues. Several possibilities of solving the main integral equation, Eq. (27), have been proposed and discussed, but not yet fully exploited. From the computational-technical point of view the problem is much more complicated than it might seem at first sight. Indeed, to ascertain, for example, the correct position of the most right-side point in Fig. 5, a system of 999 equations has been solved. One can imagine that it is not so easy to extend calculations to much larger distances than shown in Fig. 5. Perhaps, to bridge over these technical difficulties, one can take some advantage of Eqs. (54) and (56), which seem to be straightforward and useful solution schemes. On the other hand, rapid development of parallel computing and Grid technology, as well as prospects of quantum computing, also suggest some optimism for further research in this field.
## Acknowledgement
The research described in this paper has been supported by Grants No 5863 and 5549 from the Estonian Science Foundation.
## Figure captions
1. Abraham-Moses trick applied to three-component model potential for Xe<sub>2</sub> (in ground electronic state). Each step consists of removing the zeroth level along with calculating the new potential according to Eq. (21). The depth of the potential well is about 24.3 meV and the original reference potential (curve 1) has 24 levels. Although only 5 lowest partner potentials (all having the same spectral density for positive energies) are shown, they roughly demonstrate that all partner potentials should coincide as $`r0.`$ In addition, the presented curves might help to imagine how the auxiliary potential with no bound states would look like.
2. Demonstration of the characteristic function $`g(k)`$ given by Eq. (5). The overall $`g(k)`$ curve shown in the main figure seems to decay rapidly, but this impression is deceptive, because the $`k`$ scale is logarithmic. As can be seen, $`g(k)=1`$ (with high accuracy), if $`kk_0=\sqrt{V(0)/C}.`$ The insets demonstrate how nicely the different intervals can be described by a sum of Gaussians (upper inset) or exponents (lower inset) according to Eqs. (30) and (35), respectively. Parameters of the different components are given in Table 1, and Eq. (39) has been used for the region $`kk_a=75000`$ Å.
3. Calculated Krein $`H`$-function in the immediate vicinity of the zero point $`r=0`$.
4. Another demonstration of the Krein $`H`$-function. The upper graph starts where Fig. 3 ends, while the lower graph starts where the upper one ends. Note that the period of oscillations is nearly constant and very close to the characteristic value $`L={\displaystyle \frac{2\pi }{k_2}}`$ (see Eq. (31)).
5. Calculated $`G`$-function as a solution of Eq. (50) (solid curve), i.e., corresponding to the auxiliary potential $`V_0(r)`$ with no bound states. Note that the $`G`$-curve falls much slower than the $`H`$-curve ascends (cf. with Fig. 3). Another curve (open circles) has been calculated according to Eq. (59), and is therefore directly related to the reference potential $`V(r)`$. Since these differently calculated curves practically coincide, the two potentials in question also coincide in the range depicted.
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# Longitudinal versus transverse spheroidal vibrational modes of an elastic sphere
## I Introduction
With the explosion of interest in the optical properties of nanoparticles, the classic elastic mechanical problem of the vibrational modes of a free continuum sphere has found a new context for application. The problem was formally and numerically solved back in 1882.Lamb (1882) Nanoparticles, i.e. spherical clusters of atoms ranging in diameter from 1 nm to 100 nm, have sufficiently few atoms that the continuum approximation can be questioned.Cheng et al. (2005a, b) Even so, it is acceptable to ignore the effects of the discrete crystal lattice for the few vibrational modes with lowest frequency as long as the nanoparticle diameter exceeds several nanometers.
Inelastic light scattering of a continuous laser beam shining on the nanoparticle permits detection of the mechanical vibrations since the changing size and shape of the nanoparticle modulates the polarizability of the nanoparticle, so that the monochromatic incident light turns into scattered light with sidebands shifted up and down by the frequency of the vibrations. This can be seen using experimental setups of the Raman and Brillouin type.
For theoretical convenience the material is assumed to be homogeneous, isotropic and linear. The outer surface of the sphere is free from externally imposed stresses and this situation will be referred to as the “free sphere model” (FSM). The original paper by LambLamb (1882) classified the FSM modes of vibration into two classes, now called “torsional” (TOR) and “spheroidal” (SPH). The distinctive feature of torsional modes is that the material density does not vary. In other words, the divergence of the displacement field is zero. Furthermore, the spherical symmetry permits classification of the modes by angular momentum number $`\mathrm{}0`$. (However, later on we will show that there is value in considering $`\mathrm{}`$ to be a continuous variable.) There is no dependence of the frequency on the $`z`$ angular momentum $`m`$. Beyond this, the modes are indexed by $`n0`$. It is convenient to let $`p`$ denote either SPH or TOR, to indicate individual modes by $`(p,\mathrm{},n)`$ and their frequencies by $`\omega _{p\mathrm{}n}`$.
What we explore in this paper is an additional classification of the SPH modes beyond that employed since Lamb. In particular, SPH modes can be classified (albeit approximately and subjectively) as either being primarily longitudinal (SPH<sub>L</sub>) or primarily transverse (SPH<sub>T</sub>) in nature. The specific meaning of this will be explained further on. This is not a sharp division, and actual modes fall somewhere in between the two ideals. However, the contrast is sufficiently sharp that this new distinction among SPH FSM modes as SPH<sub>L</sub> or SPH<sub>T</sub> is a very important tool.
In a recent theoretical paper, G. Bachelier and A. MlayahBachelier and Mlayah (2004) predicted that (SPH,$`\mathrm{}=2`$,$`n`$) modes with differing values of $`n`$ contribute to the Raman spectrum in a highly non-uniform way. In this paper we will show that this can be explained using the previously mentionned distinction between SPH modes. They pointed out that there are two separate mechanisms that couple (SPH,$`\mathrm{}=2`$) acoustic vibrations to the surface plasmon resonance and in turn lead to Raman scattering. First, change of the particle shape and second, modulation of the density leading to change of optical properties through the deformation potential.
Section II reprises the formalism necessary for the FSM solution. In Section III, we show explicitly what we mean by SPH<sub>L</sub> and SPH<sub>T</sub>. In Section IV, we illustrate the natural appearance of SPH<sub>L</sub> and SPH<sub>T</sub> modes in the high frequency limit. Section V discusses these results and their connection with inelastic light scattering experiments.
## II The Free Sphere Model
Vibrational modes of a free linear elastic continuum homogeneous isotropic sphere were found by Lamb in 1882.Lamb (1882)
For a mode with angular frequency $`\omega `$, the displacement of material point $`\stackrel{}{r}`$ from its equilibrium position is $`\stackrel{}{u}(\stackrel{}{r})cos(\omega t)`$. For a $`m=0`$ TOR mode, $`\stackrel{}{u}`$ = A $`\times (\stackrel{}{r}j_{\mathrm{}}(k_Tr)P_{\mathrm{}}(cos\theta ))`$, where $`j_{\mathrm{}}`$ are spherical Bessel functions of the first kind and $`P_{\mathrm{}}`$ are Legendre polynomials. For a $`m=0`$ SPH mode, $`\stackrel{}{u}`$ = $`\stackrel{}{u}_L`$ \+ $`\stackrel{}{u}_T`$ where
$$\stackrel{}{u}_L(r,\theta )=Bj_{\mathrm{}}(k_Lr)P_{\mathrm{}}(cos\theta )$$
(1)
and
$$\stackrel{}{u}_T(r,\theta )=C\times \times (\stackrel{}{r}j_{\mathrm{}}(k_Tr)P_{\mathrm{}}(cos\theta ))$$
(2)
where $`A`$, $`B`$ and $`C`$ are real coefficients, $`v_Lk_L`$ = $`v_Tk_T`$ = $`\omega `$, and $`v_T`$ and $`v_L`$ are the transverse and longitudinal speeds of sound.
Modes with $`z`$ angular momentum $`m0`$ have a different functional form.
$`R`$ is the nanoparticle radius. If $`\sigma _{ij}`$ is the stress tensor, the boundary conditions at $`r`$ = $`R`$ are $`\sigma _{rr}`$ = $`\sigma _{r\theta }`$ = 0. It is convenient to introduce dimensionless frequencies $`\eta `$ = $`k_TR`$ and $`\xi `$ = $`k_LR`$. Following Eringen,Eringen and Suhubi (1975) application of these boundary conditions determines the allowed SPH vibrational frequencies as zeroes of a 2 $`\times `$ 2 determinant for $`\mathrm{}>0`$.
$$\mathrm{\Delta }_{\mathrm{}}=\begin{array}{cc}T_{11}& T_{13}\\ T_{41}& T_{43}\end{array}$$
(3)
where
$`T_{11}`$ $`=`$ $`\left(\mathrm{}^2\mathrm{}{\displaystyle \frac{\eta ^2}{2}}\right)j_{\mathrm{}}(\xi )+2\xi j_{\mathrm{}+1}(\xi )`$
$`T_{13}`$ $`=`$ $`\mathrm{}(\mathrm{}+1)\left\{(\mathrm{}1)j_{\mathrm{}}(\eta )\eta j_{\mathrm{}+1}(\eta )\right\}`$
$`T_{41}`$ $`=`$ $`(\mathrm{}1)j_{\mathrm{}}(\xi )\xi j_{\mathrm{}+1}(\xi )`$
$`T_{43}`$ $`=`$ $`\left(\mathrm{}^21{\displaystyle \frac{\eta ^2}{2}}\right)j_{\mathrm{}}(\eta )+\eta j_{\mathrm{}+1}(\eta )`$
For $`\mathrm{}=0`$, the allowed vibrational frequencies are the zeroes of $`T_{11}`$.
Noting that the displacement fields are real-valued, it is appropriate to use the following inner product between two displacement fields $`u_A`$ and $`u_B`$:Murray and Saviot (2004)
$$(u_A|u_B)=\frac{_{r<R}\stackrel{}{u_A}(\stackrel{}{r})\stackrel{}{u}_B(\stackrel{}{r})\rho d^3\stackrel{}{r}}{_{r<R}\rho d^3\stackrel{}{r}}$$
(4)
A normalization condition (such as ($`u`$$`|`$$`u`$) = 1) would typically determine the final values of $`B`$ and $`C`$. But the details of the condition do not affect the results reported here. The displacement field $`\stackrel{}{u}(\stackrel{}{r})`$ for some selected modes are depicted in Fig. 1.
## III Spheroidal Mode Longitudinality
Isotropic elastic materials differ in their Poisson ratio, $`\nu `$, which is related to $`x=v_T/v_L`$ through $`x=\sqrt{(12\nu )/(22\nu )}`$. Figure 2 shows how the dimensionless frequency, $`\eta `$, of the SPH $`\mathrm{}=2`$ FSM modes varies with $`v_T/v_L`$. It is quite apparent that some modes keep the same $`\eta `$ as $`v_T/v_L`$ is varied. However, other modes change frequency as $`v_T/v_L`$ changes. There are transition points where a given mode changes from being constant to varying with $`v_T/v_L`$.
This pattern visible in Fig. 2 motivates the search for a numerical criterion to permit this contrast among modes to be quantified. We adopt the starting point that in some sense some modes are more transverse in nature (SPH<sub>T</sub>) while others are more longitudinal (SPH<sub>L</sub>). We then coin the term “longitudinality”, denoted by $`L`$, for a quantity that varies on a scale from 0 to 1 with 0 being purely transverse and 1 being purely longitudinal. There is no single obvious way of doing this. Rather, we have evaluated a number of quantities as candidates for the best measure of longitudinality, of which we present four which work well. These will be denoted $`L1`$, $`L2`$, $`L3`$, and $`L4`$.
Consider a particular SPH mode with indices $`\mathrm{}`$ and $`n`$. Its frequency is $`\omega (v_L,v_T)`$. Define $`L1`$ by
$$L1=\frac{v_L}{\omega }\frac{\omega }{v_L}=\frac{x}{\eta }\frac{d\eta }{dx}=1\frac{v_T}{\omega }\frac{\omega }{v_T}$$
(5)
Noting that $`\stackrel{}{u}(\stackrel{}{r})=\stackrel{}{u}_T(\stackrel{}{r})+\stackrel{}{u}_L(\stackrel{}{r})`$, we define $`L2`$ as $`(u_L|u_L)/(u|u)`$, and $`L3`$ as $`1\left((u_T|u_T)/(u|u)\right)`$. But note also that $`(u_L|u_L)+(u_T|u_T)(u|u)`$ since $`(u_L|u_T)0`$.
Given a fixed value of $`v_T/v_L`$ and $`n`$, $`\eta `$ may be considered to be a continuous function of $`\mathrm{}`$, as in Fig. 3. In terms of this $`\eta (\mathrm{})`$, define
$$L4=\frac{v_T}{v_Lv_T}\left(\frac{2}{\pi }\frac{d\eta }{d\mathrm{}}1\right)$$
(6)
Let $`<..>_V`$ and $`<..>_S`$ denote averages over the nanoparticle volume and surface, respectively. In particular, $`(u|u)=<u_r^2+u_\theta ^2+u_\varphi ^2>_V`$. Some other measures of interest are as follows: URV = $`<u_r^2>_V/(u|u)`$, URS = $`<u_r^2>_S/(u|u)`$, UTS = $`<u_\theta ^2+u_\varphi ^2>_S/(u|u)`$, and U2S = URS + UTS.
Note that all of these quantities are defined in such a way as to be independent of $`m`$.
Except at low $`\eta `$, Fig. 4 shows that L1 and L2 are in close agreement. L3 and L4 are not plotted, but also agree closely except at low $`\eta `$ (see Table 1). Generally, a given mode either has all of L1, L2, L3, and L4 low, or else all high. It is thus possible to classify modes as SPH<sub>L</sub> or SPH<sub>T</sub>. Rarely, there are cases where the values of L1, L2, L3, and L4 are in the intermediate range, such as in Fig. 4 for $`\mathrm{}=3`$ for the two modes near $`\eta =13`$. Such modes which are neither clearly SPH<sub>L</sub> nor SPH<sub>T</sub> always occur in pairs. The reason for this is explained in Section IV.
Table 2 provides additional information about the modes. The dimensionless frequency $`\eta `$ is provided for convenience, as is the ratio of coefficients $`B`$ and $`C`$.
In principle, $`C/B`$ could be expected to provide useful information about whether a mode is SPH<sub>L</sub> or SPH<sub>T</sub>. In the extreme case that $`C=0`$, the mode is evidently SPH<sub>L</sub>, and likewise when $`B=0`$ the mode is SPH<sub>T</sub>. But the $`C/B`$ values do not exhibit an informative pattern.
There is a strong contrast in the values of URS for different modes. However, it does not correlate to whether the mode is SPH<sub>L</sub> or SPH<sub>T</sub> except at high enough $`\eta `$. URS is an interesting quantity because it is the one we have to monitor for the surface deformation mechanism except for $`\mathrm{}=0`$ modes.
Group theoretical argumentsDuval (1992) show that only SPH modes with $`\mathrm{}`$=0 and $`\mathrm{}`$=2 are Raman active. This assumes that the nanoparticle is perfectly spherical in shape and spherically symmetric in all of its properties. The basic nature of the $`\mathrm{}=0`$ modes is much more clear because of their simplicity and symmetry. Consequently, the modes (SPH,$`\mathrm{}=2`$,$`n`$) are of primary interest when trying to understand Raman intensities.
From the value of L2 $``$ 0.14 in Tab. 2, the displacement of (SPH,2,0) is mostly due to the $`u_T`$ term and not the $`u_L`$ term. Its squared displacement due to the $`u_L`$ term alone is just 14% of the total. The $`u_T`$ term has zero divergence. Therefore, (SPH,2,0) doesn’t have much divergence. So the effect of changing density on the dielectric constant through the deformation potential may not be significant to the overall Raman intensity.
On the other hand, based on its URS of $``$ 0.66 and UTS of $``$ 0.20, the surface displacement of (SPH,2,0) is strongly along $`r`$ and only weakly along $`\theta `$ as Fig. 1 illustrates. Note that, $`r`$ surface displacement changes the nanoparticle shape, while $`\theta `$ displacement does not.
The (SPH,2,1) mode differs from (SPH,2,0) in several ways. From the $`L1`$ value of 0.2281 in Fig. 4, we can see that the frequency of (SPH,2,1) depends more on $`v_L`$. Also, $`L20.416`$ in Fig. 4 shows that (SPH,2,1) has more of a $`u_L`$ component, even if it is still weaker than the $`u_T`$ part. But this means that (SPH,2,1) can have much more divergence than (SPH,2,0). So the deformation potential mechanism can modulate the dielectric constant. But it is very interesting to notice from the URS value of $``$ 0.00 in Tab. 2 (more precisely, 0.0003) that (SPH,2,1) causes negligible radial movement of the surface. So (SPH,2,1) barely changes the shape of the nanoparticle, as Fig. 1 shows.
(SPH,2,3) has strong $`v_L`$ dependence ($`L10.8475`$) in Fig. 4 as well as a strong $`u_L`$ component ($`L20.766`$). So it is clear that it is SPH<sub>L</sub>. It’s surface displacement is mostly along $`r`$ and not $`\theta `$ from its URS value of 0.71 and UTS $``$ 0.04. So (SPH,2,3) will strongly affect the shape of the nanoparticle surface, as shown in Fig. 1.
$`u_L`$ and $`u_T`$ take on simpler forms as $`\eta `$ becomes larger. For large $`\eta `$, the $`u_L`$ term has primarily radial displacement, while the $`u_T`$ term corresponds to displacement in the $`\theta `$ direction.
For the lowest modes, the situation is qualitatively different. Consider (SPH,2,0) with $`v_T/v_L=0.5`$. Suppose to simplify this discussion we normalize the displacement field so that $`(u|u)=1`$. Then $`(u_L|u_L)0.14`$. However, $`(u_T|u_T)1.85`$. So L3 for (SPH,2,0) is actually $``$ -0.85, making it “ultra transverse” by that measure. It seems quite odd that the $`u_T`$ term alone has a magnitude much greater than that of the overall motion. The resolution of this puzzle is that $`u_L`$ and $`u_T`$ are not mutually orthogonal with respect to the inner product of Eq. 4. In fact, $`(u_T|u_L)`$ $``$ -0.50. According to the usual vector relation, $`\stackrel{}{a}\stackrel{}{b}`$ = $`\stackrel{}{a}\stackrel{}{b}cos\theta _{ab}`$, the “angle” between $`u_L`$ and $`u_T`$ is $``$ 165 degrees for the (SPH,2,0) mode. This angle is nearly unchanged as $`v_T/v_L`$ varies. Thus, $`u_L`$ and $`u_T`$ are nearly antiparallel vectors in the function space of vector fields within the nanoparticle interior. It can be said, then, that the functional forms of $`u_L`$ and $`u_T`$ are actually relatively similar. This is a bit of a surprise since one is curl-free while the other is divergence free. This angle between $`u_L`$ and $`u_T`$ rapidly approaches 90 degrees as $`\eta `$ increases (i.e. for modes with higher $`n`$).
As Fig. 2 shows, the starkness of the contrast between SPH<sub>L</sub> and SPH<sub>T</sub> modes is at its best when $`v_T/v_L`$ is lower. For materials with high $`v_T/v_L`$ such as Si and Ge, FSM modes tend more to be mixtures of SPH<sub>L</sub> and SPH<sub>T</sub>, especially at low $`\eta `$. But the concept of longitudinality is quite applicable to materials such as Au and Ag.
## IV High frequency mode classification
The reason for the dichotomy of SPH modes as SPH<sub>T</sub> and SPH<sub>L</sub> can be simply explained in the high frequency limit. Consider $`\mathrm{\Delta }_{\mathrm{}}`$, the 2 $`\times `$ 2 determinant in Eq. 3, and its four matrix elements. Note that $`\xi /\eta `$ = $`v_T/v_L`$. So at high frequency, both $`\eta `$ and $`\xi `$ are large. In that case, $`T_{11}`$ and $`T_{43}`$ will be much larger than $`T_{13}`$ and $`T_{41}`$ because of their terms including factors of $`\eta ^2`$. Consequently, $`\mathrm{\Delta }_{\mathrm{}}`$ is very well approximated by $`T_{11}T_{43}`$. Since normal modes correspond to zeroes of $`\mathrm{\Delta }_{\mathrm{}}`$, it is clear that there will be two sets of modes: those which are approximately zeroes of $`T_{11}`$ and $`T_{43}`$ respectively. The first group are SPH<sub>L</sub> and the second group are SPH<sub>T</sub>.
The roots of $`T_{11}`$ and $`T_{43}`$ are plotted versus $`\mathrm{}`$ in Fig. 3 with lines with crosses for SPH<sub>L</sub> modes and lines with empty squares for SPH<sub>T</sub> modes. The lines with full circles are the exact FSM mode frequencies.
There are three kinds of situations where we don’t expect this approximation to be valid: (1) for low $`\eta `$ (2) when it is not true that $`\eta \mathrm{}`$ and (3) where longitudinal ($`T_{11}`$) and transverse ($`T_{43}`$) modes for a given $`\mathrm{}`$ are close – i.e. when the associated curves cross. Except in the previously mentioned places, the agreement between FSM and our approximation is quite good. The low $`\eta `$ situation corresponds specifically to the similar prefactors of $`j_{\mathrm{}}`$ for T<sub>11</sub> and T<sub>43</sub> not being large. It is apparent that T<sub>11</sub> and T<sub>43</sub> can only be useful as estimators of SPH<sub>L</sub> and SPH<sub>T</sub> mode frequencies when $`\eta \mathrm{}`$. This is confirmed from inspection of the lower right portion of Fig. 3.
For large $`x`$, $`j_{\mathrm{}}(x)\mathrm{sin}(x\mathrm{}\frac{\pi }{2})/x`$. Therefore, for large $`\xi `$ and $`\eta `$, the roots of $`T_{11}`$ can be approximated by $`\xi \mathrm{}\frac{\pi }{2}+(1+n_L)\pi `$ and the roots of $`T_{43}`$ by $`\eta \mathrm{}\frac{\pi }{2}+n_T\pi `$ where $`n_L0`$ and $`n_T0`$ are integers. These lead to remarkably compact approximate expressions for SPH<sub>L</sub> and SPH<sub>T</sub> FSM frequencies in Hertz, respectively:
$$f\frac{v_L}{d}(\frac{\mathrm{}}{2}+n_L+1)$$
(7)
$$f\frac{v_T}{d}(\frac{\mathrm{}}{2}+n_T)$$
(8)
where $`d`$ = $`2R`$. These expressions are very suggestive of the formula for acoustic standing waves in a one dimensional system of length $`d`$. Table 2 shows the value of either $`n_L`$ or $`n_T`$ for each mode.
The behaviour observed in Fig. 2 becomes simple to explain. To a good approximation, SPH FSM modes are either SPH<sub>L</sub> or SPH<sub>T</sub>. This approximation is considered here to be good because it gives the right number of vibrational modes and it predicts their frequency with a reasonable accuracy.
“Anti-crossing” is observed in Fig. 2 each time the variation of the frequency of a SPH<sub>L</sub> mode crosses the one of a SPH<sub>T</sub> mode. In Fig. 2 there are two kinds of curves: horizontal lines for SPH<sub>T</sub> modes and descending curves for SPH<sub>L</sub> ones. Then, each time these curves come together, an anti-crossing pattern appears for the FSM solutions. In Fig. 3, FSM frequencies $`\eta `$ are plotted versus $`\mathrm{}`$ for a sphere made of a material which has $`v_T/v_L`$ = 0.5. Because the SPH<sub>L</sub> and SPH<sub>T</sub> approximation curves are plotted, the anti-crossing patterns are clearly revealed. The continuation of Bessel functions to non-integer $`\mathrm{}`$ permits relationships among modes for different integer $`\mathrm{}`$ to be clearly seen. This is preferable to the common practice of joining modes on such a graph with hand-drawn straight lines.
## V Discussion:
Normal elastic waves in a solid have a longitudinal acoustic (LA) branch and two transverse acoustic (TA) branches. However, for FSM it seemed there are just two kinds: SPH and TOR. By classifying SPH modes into two kinds (i.e. SPH<sub>L</sub> and SPH<sub>T</sub>), there are now three categories of modes, as we would expect.
We plot in Fig. 5 the mean squared radial surface displacement (URS) at the surface of a 5 nm diameter silver nanoparticle for all SPH $`\mathrm{}=2`$ modes. The magnitude of URS is in good agreement with the calculated Raman intensitiesBachelier and Mlayah (2004). (These calculations took into account the non-linear dispersion of acoustic phonons in silver. As a result, the calculated vibration wavenumbers do not match.) As discussed before, (the SPH,2,0) mode is quite special even if we class it as a SPH<sub>T</sub> mode. It changes the surface shape and therefore contributes significantly to inelastic light scattering. Other harmonics contribute significantly only when their URS is large and this in turn is very well correlated to their SPH<sub>L</sub> nature as can be seen in Fig. 2.
Many experiments have observed peaks in Raman spectra attributed to acoustic phonon vibrations of silver Portalès et al. (2001a, b); Duval et al. (2001); Fujii et al. (1991); Portalès (2001); Courty et al. (2002); Nelet et al. (2004) silicon Fujii et al. (1996); Saviot et al. (2004a) and CdS<sub>x</sub>Se<sub>1-x</sub>. Verma et al. (1999); Saviot et al. (1996); Ivanda et al. (2003); Irmer et al. (2000); Saviot et al. (1998) These studies have regularly succeeded in observing the (SPH,2,0) mode and the (SPH,0,0) mode. A number of studies have seen (SPH,0,$`n`$) with $`n`$ up to 4Nelet et al. (2004). However, there has never been a clear indication of Raman scattering from (SPH,2,1) even though there have been determined efforts to see it.
At the same time, (SPH<sub>L</sub>,2,$`n_L=0`$) seems like a strong candidate to have noticeable Raman scattering, since it has strong radial surface motion as well as a strong $`u_L`$ component that will give it stronger divergence in its interior.
It should be noted that $`\mathrm{}=0`$ modes are always SPH<sub>L</sub>. That is why no full circles are plotted in Fig. 3 on the $`T_{43}`$ root curves at $`\mathrm{}=0`$. This has been the source of many erroneous calculations in the pastSaviot et al. (2004b).
It is often claimedTanaka et al. (1993); Tamura et al. (1982); Tamura and Ichinokawa (1983); Ovsyuk and Novikov (1996) that modes with $`n`$ = 0 are “surface modes” while modes with $`n>0`$ are “inner modes”. The values of U2S in Tab. 2 show that this is a misconception. While this is true for $`\mathrm{}`$ = 0 and 1, for $`\mathrm{}`$ = 2, 3 and 4 it can be seen that (SPH,$`\mathrm{}`$,1) has the strongest surface motion relative to all (SPH,$`\mathrm{}`$,$`n`$).
Although Tab. 2 shows URS to be zero for (SPH,2,1), the more precise value of $`v_T/v_L`$ where URS is zero is 0.488. URS for (SPH,2,1) is only near zero for $`v_T/v_L`$ close to 0.488. However, URS remains small for (SPH,2,1) for materials whose Poisson ratio is close to $`\frac{1}{3}`$ which is true of many common materials. This contradicts a widespread misconceptionOhno et al. (2000); Lomnitz and Nilsen-Hofseth (2005); Heyliger and Jilani (1992) that SPH FSM modes always have a radial displacement component at the surface.
###### Acknowledgements.
D. B. M. acknowledges support from the Natural Sciences and Engineering Research Council of Canada, the Okanagan University College Grant-in-Aid Fund and National Sun Yat-Sen University and thanks L. M. L. Murray and A. S. Laarakker for valuable scientific suggestions.
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# Polarization observables in the semiexclusive photoinduced three-body breakup of 3He.
## I Introduction
The study of polarization phenomena is a natural extension of investigation of unpolarized processes. It provides additional information on details of the underlying nuclear Hamiltonian not available in unpolarized reactions. In nucleon-nucleon (NN) systems the polarized processes provide a necessary data set to construct the NN potentials . The investigation of nucleon-deuteron (Nd) elastic scattering and the deuteron breakup reaction with polarized incoming nuclei or polarization of the outgoing particles measured is indispensable to learn about properties of the three-nucleon (3N) forces. Nowadays spin observables in Nd elastic scattering where the initial deuteron and/or nucleon is polarized and also the polarization of the final particles is measured, are available and can be compared with rigorous theoretical predictions. Also for deuteron breakup such studies were performed, both experimentally as well as theoretically . Altogether, this allowed to test the current models of the nuclear Hamiltonian.
In addition to the strong forces, the electromagnetic processes contain new dynamical ingredients due to the interaction between real or virtual photons with the currents of nuclei. It was found that in such processes contributions to the nuclear current due to meson exchanges play an important role. Studies of polarization observables in photo- and electro-induced processes on the deuteron , as well as in the Nd radiative capture can be used to determine the structure of nuclear currents. The combination of strong and electromagnetic interactions is a demanding test for theoretical models. The results up to now show an overall good agreement of theoretical predictions with the data, however there is still room for improvement . Recently an important progress is observed in experimental investigations of processes with polarized photons. The high-intensity sources of highly polarized photon beams obtained by the Compton backscattering give hope for future precise data . The analysis of the first measurement of the <sup>3</sup>He breakup using polarized photons at low energies is in progress and was reported recently in .
In this paper we would like to present the results of theoretical investigations of spin observables in kinematically incomplete $`\stackrel{}{\gamma }(^3\stackrel{}{He},N)NN`$ processes in which the incoming photon and/or the <sup>3</sup>He nucleus are polarized. This study is done for three photon laboratory energies $`E_\gamma `$=12, 40 and 120 MeV. For each photon energy, the energy spectrum of the detected outgoing nucleon at different angles has been calculated. We restrict ourselves to photon energies below the pion production threshold and have chosen the above energies as examples of low, intermediate and high photon energies. It was shown in that for those energies one can expect different manifestations of the action of the 3N force in two-body photodisintegration of <sup>3</sup>He. While at low energies the inclusion of the three-nucleon forces decreases the cross section, at higher energies 3N forces act in the opposite direction. At intermediate energies the influence of 3N forces on the two and three-body breakup cross sections is negligible. As will be shown in section III for several spin observables the influence of 3N forces is visible in the semiexclusive spectrum of the outgoing nucleon also at intermediate energies of the incoming photon. The presented results should be a useful guide for future experiments. Up to now, to the best of our knowledge, no such predictions have been published.
In section II we shortly describe the theoretical formalism underlying our calculations and give definitions for the studied spin observables. In section III we present our predictions for three-body breakup. In addition we turn into two-body <sup>3</sup>He breakup and compare our results to a few existing data. We summarize in section IV.
## II Theoretical Framework
The theoretical framework we use is described in detail in Refs. . For the convenience of the reader we briefly summarize the most important steps. The basic nuclear matrix element $`N_{m_i,\tau ,m}^{3\mathrm{N}}`$ is expressed through the state $`\stackrel{~}{U}_\tau ^m`$ which fulfills the Faddeev-like equation
$`\stackrel{~}{U}_\tau ^m=(1+P)j_\tau (\stackrel{}{Q})\mathrm{\Psi }_{{}_{}{}^{3}\mathrm{He}}^m+\left(tG_0P+{\displaystyle \frac{1}{2}}(1+P)V_4^{(1)}G_0(tG_0+1)P\right)\stackrel{~}{U}_\tau ^m.`$ (1)
Here $`j_\tau (\stackrel{}{Q})`$ is a spherical $`\tau `$component of the <sup>3</sup>He electromagnetic current operator, $`t`$ the NN t-matrix, $`G_0`$ the free 3N propagator and $`P`$ the sum of the cyclical and anticyclical permutations of 3 particles. Further $`V_4^{(1)}`$ is that part of the 3NF, which is symmetrical (like the NN $`t`$-matrix) under the exchange of nucleons 2 and 3, and $`\mathrm{\Psi }_{{}_{}{}^{3}\mathrm{He}}^m`$ is the <sup>3</sup>He bound state with spin projection $`m`$. The nuclear matrix element for three-body breakup of <sup>3</sup>He is given via
$`N_{m_i,\tau ,m}^{3\mathrm{N}}={\displaystyle \frac{1}{2}}\mathrm{\Phi }_0^{m_i}(tG_0+1)P\stackrel{~}{U}_\tau ^m,`$ (2)
where $`\mathrm{\Phi }_0^{m_i}`$ is the properly anti-symmetrized (in the two-body subsystem) state of three free nucleons with their spin projections $`m_i`$.
Given the $`N_{m_i,\tau ,m}^{3\mathrm{N}}`$ amplitudes, one can calculate any polarization observables. They are expressed through the nuclear matrix elements with different spin projections carried by the initial photon, the <sup>3</sup>He nucleus, and by the outgoing nucleons.
Chossing the z-axis to be the direction of the incoming photon and allowing for a linear photon polarization $`P_0^\gamma `$ along the x-axis, with the polarization component $`P_0^\gamma =1`$, and for the <sup>3</sup>He target nucleus polarization $`P_0^{{}_{}{}^{3}He}`$ along the y-axis, the cross section in a kinematically incomplete reaction $`\stackrel{}{\gamma }(^3\stackrel{}{He},N)NN`$, when the outgoing nucleon is detected at angles ($`\theta ,\varphi `$) is given by
$`\sigma _{\gamma ,^3He}^{pol}(\theta ,\varphi )`$ $`=`$ $`\sigma _{\gamma ,^3He}^{unpol}(\theta )[1+P_0^\gamma cos(2\varphi )A_x^\gamma (\theta )+P_0^{{}_{}{}^{3}He}cos(\varphi )A_y^{{}_{}{}^{3}He}(\theta )+`$ (3)
$`P_0^\gamma cos(2\varphi )P_0^{{}_{}{}^{3}He}cos(\varphi )C_{x,y}^{\gamma ,^3He}(\theta )+P_0^\gamma sin(2\varphi )P_0^{{}_{}{}^{3}He}sin(\varphi )C_{y,x}^{\gamma ,^3He}(\theta )].`$ (4)
Here the nonvanishing spin observables are the photon ($`A_x^\gamma (\theta )`$) and the <sup>3</sup>He ($`A_y^{{}_{}{}^{3}He}(\theta )`$) analyzing powers, and the spin correlation coefficients $`C_{x,y}^{\gamma ,^3He}(\theta )`$ and $`C_{y,x}^{\gamma ,^3He}(\theta )`$. They can be obtained by measuring the spectra of the outgoing nucleon using a proper combination of $`\varphi `$ angles and are expressed through the nuclear matrix element $`N_{m_i,\tau ,m}^{3\mathrm{N}}`$ by:
$`A_x^\gamma (\theta )`$ $``$ $`{\displaystyle \frac{\underset{m_im}{}(2\mathrm{}\{N_{m_i,1m}N_{m_i,+1m}^{}\})}{_{m_im}(|N_{m_i,+1m}|^2+|N_{m_i,1m}|^2)}}`$ (5)
$`A_y^{{}_{}{}^{3}He}(\theta )`$ $``$ $`{\displaystyle \frac{\underset{m_i}{}(2\mathrm{}\{N_{m_i,1\frac{1}{2}}N_{m_i,1\frac{1}{2}}^{}\}2\mathrm{}\{N_{m_i,+1\frac{1}{2}}N_{m_i,+1\frac{1}{2}}^{}\})}{_{m_im}(|N_{m_i,+1m}|^2+|N_{m_i,1m}|^2)}}`$ (6)
$`C_{x,y}^{\gamma ,^3He}(\theta )`$ $``$ $`{\displaystyle \frac{\underset{m_i}{}(2\mathrm{}\{N_{m_i,1\frac{1}{2}}N_{m_i,+1\frac{1}{2}}^{}\}+2\mathrm{}\{N_{m_i,1\frac{1}{2}}N_{m_i,+1\frac{1}{2}}^{}\})}{_{m_im}(|N_{m_i,+1m}|^2+|N_{m_i,1m}|^2)}}`$ (7)
$`C_{y,x}^{\gamma ,^3He}(\theta )`$ $``$ $`{\displaystyle \frac{\underset{m_i}{}(2\mathrm{}\{N_{m_i,1\frac{1}{2}}N_{m_i,+1\frac{1}{2}}^{}\}+2\mathrm{}\{N_{m_i,1\frac{1}{2}}N_{m_i,+1\frac{1}{2}}^{}\})}{_{m_im}(|N_{m_i,+1m}|^2+|N_{m_i,1m}|^2)}}.`$ (8)
## III Results
We solved Eq.(1) using a momentum space partial wave decomposition and the AV18 nucleon-nucleon potential alone or supplemented with the Urbana IX 3NF . For both parities and the total angular momentum of the 3N system $`J\frac{15}{2}`$ all partial waves with angular momenta in the two-body subsystem $`j3`$ have been used. We refer to for more details on our basis, partial wave decomposition and numerics. The electromagnetic nuclear current operator was taken as the single nucleon current supplemented by the exchange currents of the $`\pi `$\- and $`\rho `$-like nature .
Before presenting the polarization observables, for the sake of completeness, we would like to show the unpolarized cross section for the $`\gamma (^3\mathrm{He},N)NN`$ reaction with the detected outgoing nucleon to be a proton (Fig. 1) or a neutron (Fig. 2). We choose the detection polar angle $`\theta `$ to be $`\theta =30^{}`$, 60, 90, 120 or 150. The spectra at $`\theta =90^{}`$ were already presented in . The structures seen in these spectra originate from an interplay between strong final state interactions, meson exchange currents, phase space factors and the properties of the 3N bound state wave function. For example, for the neutron spectrum at $`E_\gamma =120`$ MeV and $`\theta =90^{}`$ two peaks around $`E_n20`$ and 70 MeV come from the final state interactions between two nucleons. The maximum around 50 MeV comes from the interplay between the two-body currents, the phase space factors and the properties of the <sup>3</sup>He bound state wave function. As is seen in Figs. 1 and 2 that structure depends smoothly on the angle of the outgoing nucleon with the largest cross sections around $`\theta =90^{}`$. The Urbana IX 3NF effects are visible at the lower and the upper energies of the incoming photon, and are nearly negligible at the intermediate energy.
The photon analyzing power $`A_x^\gamma (\theta )`$ are shown in Figs. 3-4. For photon energies $`E_\gamma =12`$ and 40 MeV and detecting protons this observable decreases with increasing proton energy and reaches values -1 and -0.8 at the highest proton energies, respectively. $`A_x^\gamma (\theta )`$ is rather insensitive to the 3NF at these photon energies. However, at $`E_\gamma =120`$ MeV the 3NF effects become sizable, and they change the photon analyzing power by up to $``$10%. The strongest effects are visible at protons energies around 25-50 MeV and at lower detection angles. For the detected neutron $`A_x^\gamma (\theta )`$ reaches values -1 for $`E_\gamma =12`$ and 40 MeV at the upper ends of the spectra. At $`E_\gamma =120`$ MeV the value of the photon analyzing power is small (up to $``$ -0.2) except in of the region of maximal energies of the detected neutrons. At all investigated energies the 3NF effects are negligible when the neutron is detected.
Contrary to a rather small 3NF effects in the photon analyzing power, the <sup>3</sup>He analyzing power $`A_y^{{}_{}{}^{3}He}(\theta )`$ is sensitive to the action of 3N forces (see Figs. 5-6). This is the case especially for the two lowest photon energies and the detected neutron and at $`E_\gamma =12`$ MeV and $`E_\gamma =120`$ MeV when the proton is measured. For the detected neutron the largest 3NF effects of up to 15% are at $`E_\gamma =12`$ MeV and they are seen in the whole neutron spectrum. In the proton case the most interesting situation is the highest photon energy $`E_\gamma =120`$ MeV, where 3NF effects of a magnitude above $``$ 20% are seen nearly for all energies of the detected proton. The action of the 3NF shifts the predictions in the opposite directions for the lowest and the highest photon energy. Unfortunately in the case of the detected proton the 3NF effects occur at relatively small (below 0.1) absolute values of $`A_y^{{}_{}{}^{3}He}(\theta )`$. For the detected neutron 3NF effects occur also for $`A_y^{{}_{}{}^{3}He}(\theta )0.12`$. However, in this case 3NF effects are seen even at intermediate photon energy, at all neutron angles and in the whole energy range. The structure of the spectrum is again due to an interplay of all dynamical components in the nuclear matrix elements. The dependence of the nuclear analyzing power on the direction of the outgoing nucleon is rather smooth, but the shape of the spectra changes significantly for different photon energies.
The spin correlation coefficients $`C_{x,y}^{\gamma ,^3He}(\theta )`$ are shown in Figs. 7-8. In that case the largest 3NF effects ($``$ 15%) occur in the whole spectrum at $`E_\gamma =12`$ MeV when the neutron is measured. Smaller 3NF effects are also visible at other photon energies, however, their magnitude depends on the detection angle. For the measured proton, 3NF effects are negligible at the two higher photon energies. The absolute values of $`C_{x,y}^{\gamma ,^3He}(\theta )`$ for the detected proton (neutron) stays below $``$0.25 ($``$0.1) at $`E_\gamma =12`$ MeV and $``$0.4 ($``$0.25) at the two higher photon energies.
A similar picture arises for the spin correlation coefficients $`C_{y,x}^{\gamma ,^3He}(\theta )`$ (see Figs. 9-10). Here 3NF effects are also visible at higher photon energies. For $`E_\gamma =40`$ MeV and the measured neutron, 3NF effects are largest around the outgoing neutron energy $``$16 MeV and $`\theta =60^{}120^{}`$. The absolute values of $`C_{y,x}^{\gamma ,^3He}(\theta )`$ for neutron detection are below $``$0.1 for $`E_\gamma =12`$ and 40 MeV, and approach up to $``$0.4 for $`E_\gamma =120`$ MeV. For the measured proton $`C_{y,x}^{\gamma ,^3He}(\theta )`$ reaches 0.25, 0.5 and 0.25 for $`E_\gamma =12`$, 40 and 120 MeV, respectively. In the case of the detected proton the small 3NF effects (below 10%) occur at all photon and nucleon energies and at all detection angles.
Now we would like to address the sensitivity of the spin observables to the nuclear current used. To study this, we compare the predictions for the above spin observables at the detection angle $`\theta =90^{}`$ for three different choices of the current operator: the single nucleon current (SNC) only, when the explicit two-body meson exchange currents are added to the SNC, and finally when the current operator is constructed using the Siegert theorem . The Siegert approach will also include 3N currents in the electric multipoles. We should mention, however, that in our realization of the Siegert theorem we keep only single nucleon operators and do not (yet) supplement the magnetic multipoles by the explicit $`\pi `$\- and $`\rho `$\- exchange currents. Also the explicit $`\pi `$\- and $`\rho `$\- currents are not fully consistent with the underlying AV18 NN force, but only with its dominant parts . For a recent investigation filling that gap see . Despite these defects we think that the comparison of our Siegert approach with the explicit use of the $`\pi `$\- and $`\rho `$\- currents will enable us to identify those observables, which are especially sensitive to the choice of two and possibly three-body currents.
The photon analyzing power $`A_x^\gamma (\theta )`$ is insensitive to such a change of the nuclear current at the lowest energy (see Fig. 11). At $`E_\gamma =40`$ MeV only a slight shift of predictions is observed under inclusion of the meson exchange currents. The effects coming from the two models of exchange currents are insignificant for the neutron knockout but lead to a small spread of theoretical predictions for the proton detection. At $`E_\gamma =120`$ MeV one finds a clear difference when the two models including exchange currents are used, and when only the single nucleon current is taken into account. The difference between SNC predictions and explicit $`\pi `$\- and $`\rho `$\- currents (Siegert) results ammounts up to 140% (180%) at E$`{}_{n}{}^{}`$20 MeV, and up to 650% (880%) at E$`{}_{p}{}^{}`$17 MeV, respectively.
For $`A_y^{{}_{}{}^{3}He}(\theta )`$, shown in Fig. 12 the single nucleon current predictions differ from others at all photon energies. While for the detected neutron meson exchange currents play an important role at all studied energies, in the proton case they are important only at $`E_\gamma =120`$ MeV. The differences between Siegert and MEC are visible at all photon energies. At $`E_\gamma =40`$ MeV they reach up to $``$50% for neutron energies around 5-10 MeV. The case of the measured proton around $`E_p15`$ MeV and for $`E_\gamma =120`$ MeV seems to be very interesting, since the different nuclear currents lead to a different sign of $`A_y^{{}_{}{}^{3}He}(\theta )`$ (see Fig. 12). The differences are also seen for the spin correlation coefficients $`C_{x,y}^{\gamma ,^3He}(\theta )`$ and $`C_{y,x}^{\gamma ,^3He}(\theta )`$, presented in Figs. 13-14. For $`C_{x,y}^{\gamma ,^3He}(\theta )`$ and the measured neutron there are clear differences, when using Siegert approach or direct $`\pi `$ and $`\rho `$ currents. They amount up to $``$50% at $`E_\gamma =40`$ MeV. For both cases, the neutron or proton detection, and $`E_\gamma =12`$ MeV the predictions without 3NF differ significantly, while the inclusion of the Urbana IX force leads to an agreement between both predictions. Both spin correlation coefficients are strongly influenced by the meson exchange currents. Even at $`E_\gamma =12`$ MeV single nucleon current predictions differ significantly from results based on the nuclear current supplemented by exchange currents. The role of exchange currents grows with the photon energy. In the case of $`C_{y,x}^{\gamma ,^3He}(\theta )`$, $`E_\gamma =40`$ MeV and proton energies below $`E_p9`$ MeV, we observe different action of the exchange currents when they are included via Siegert or by the explicit $`\pi `$\- and $`\rho `$ exchanges. It shows that this observable is very interesting to study details of the nuclear current operator and deserves experimental efforts. Thus we can state that the spin observables in the $`\stackrel{}{\gamma }(\stackrel{}{{}_{}{}^{3}\mathrm{He}},N)NN`$ reaction can provide valuable information about the nuclear current operator.
Finally, we address ourselves to the $`\stackrel{}{\gamma }(\stackrel{}{{}_{}{}^{3}\mathrm{He}},p)d`$ process and compare our results with the data of Ref. . There the cross section asymmetry
$$\mathrm{\Sigma }\frac{d\sigma _{}d\sigma _{}}{d\sigma _{}+d\sigma _{}},$$
(9)
where $`d\sigma _{}`$ ($`d\sigma _{}`$) is the cross section measured parallely (perpendicularly) to the photon polarization direction, was investigated for linearly polarized photons with energies above 90 MeV. In Fig. 15 we compare our predictions to the data of at the photon energy $`E_\gamma =120`$ MeV. We see that two of the three data points are in good agreement with our theory. Our prediction at the third data point is too low in comparison to data. The 3NF shifts the theory in the right direction into the two data points. Unfortunately, most of the data points taken in are at photon energies above the pion production threshold where our formalism is not adequate. Nevertheless, in Fig. 16 we compare our predictions with data at $`E_\gamma =200`$ MeV in order to check if our predictions at higher energies give at least a qualitative description of the data. We see that while the shape of the theoretical predictions is similar to the shape of the data, the absolute values of the predicted analyzing power are too small by a factor of 2. This probably can be traced back to the missing dynamical ingredients in our theoretical framework, which may become important at such high energies. As was the case at $`E_\gamma =120`$ MeV, also at $`E_\gamma =200`$ MeV the 3NF improves the description of the data. Since our calculations are much more advanced than the one used in Ref. , we would like to point out, that the very good description of the data presented in might be to some extent accidental.
## IV Summary
We investigated all the nonvanishing spin observables in the three-body, semiexclusive <sup>3</sup>He photodisintegration process when the incoming photon and/or the <sup>3</sup>He target nucleus are polarized. We found that the dependence of those spin observables on the angle of the outgoing nucleon is rather smooth and in most cases the shape of the energy spectra slightly changes with the incoming photon energy. In the case of the $`A_y^{{}_{}{}^{3}He}(\theta )`$ analyzing power and the spin correlation coefficients $`C_{x,y}^{\gamma ,^3He}(\theta )`$ and $`C_{y,x}^{\gamma ,^3He}(\theta )`$ clear effects of the 3NF are seen. Some of the observables (e.g. $`C_{y,x}^{\gamma ,^3He}(\theta )`$) are sensitive to the details of the many-body contributions to the nuclear current operator, which we examplified by using the single nucleon current alone and by supplementing it either with explicit inclusion of $`\pi `$\- and $`\rho `$\- meson exchange currents or by applying the Siegert theorem.
The presented results show that the polarization observables for <sup>3</sup>He photodisintegration, even in the relatively simple semiexclusive experiments, could provide valuable data to test the nuclear forces and/or the reaction mechanism. These observables should be studied experimentally. On the other hand, there are observables (e.g. $`A_x^\gamma (\theta )`$ at E<sub>γ</sub>=12 MeV) which are insensitive to the choosen current operator model and to the inclusion of the 3N force. Such observables are natural candidates to test the most simple dynamical ingredients.
###### Acknowledgements.
This work was supported by the Polish Committee for Scientific Research under Grant No. 2P03B0825 and NATO Grant No. PST.CLG.978943. W.G. would like to thank the Polish-German Academical Society. The numerical calculations have been performed on the Cray SV1 and IBM Regatta p690+ of the NIC in Jülich, Germany.
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# MKPH-T-05-06 Incoherent pion photoproduction on the deuteron with polarization observables II: Influence of final state rescattering
## I Introduction
Photoproduction of pions on the deuteron has two main but complementary points of interest. The first one is to obtain information on the elementary reaction on the neutron by using the deuteron as an effective neutron target. A prerequisite for this is that one has reliable control on off-shell and medium effects. In order to minimize such effects, quasi-free kinematics is preferred. The second but not secondary aspect is just the influence of a nuclear environment on the production process, for the study of which off-quasi-free kinematics is better suited.
Pion photoproduction on the deuteron has been studied quite extensively over the past 50 years, starting with early work in ChL51 ; LaF52 ; BlL77 . The role of final state interaction (FSI) has been investigated by Laget Lag78 ; Lag81 applying a diagrammatic approach. The influence of FSI effects were found to be quite small for charged pion photoproduction compared to the neutral channel. A satisfactory agreement with experimental data was achieved for $`\pi ^{}`$ production Be+73 . Subsequently, these results were confirmed by Levchuk et al. LeP96 for the $`d(\gamma ,\pi ^0)np`$ reaction for which the elementary photoproduction operator of Blomqvist and Laget BlL77 was used. This work was improved and extended to charged pion production channels in LeS00 , where a more realistic elementary production operator from the SAID Said and MAID Maid multipole analyses was taken and $`NN`$-rescattering included, based on the Bonn r-space potential MaH87 . The influence of $`NN`$-FSI was confirmed and good agreement with experimental data was achieved. In the threshold region a sizeable effect from $`\pi N`$-rescattering was noted in Le+00 which arises from intermediate charged pion production with subsequent charge exchange rescattering on the spectator nucleon. The influence of $`NN`$\- and $`\pi N`$-rescattering on polarization observables has been investigated in LeP96 for the GDH sum rule in the $`\pi ^0np`$-channel as well as in LoS00 for target asymmetries in the $`\pi ^{}pp`$-reaction.
FSI effects in incoherent pion photoproduction were also studied by Darwish et al. DaA03a . The same approach was then applied to the spin asymmetry with respect to circularly polarized photons and vector polarized deuterons DaA03b , which determines the much discussed Gerasimov-Drell-Hearn sum rule GDH . However, the approach was limited to the $`\mathrm{\Delta }(1232)`$-resonance region in view of a relatively simple elementary production operator, based on an effective Lagrangian approach from Schmidt et al. ScA96 . A puzzling result of this work was that the influence of FSI on the total cross sections for charged pion production resulted in a slight decrease in the $`\mathrm{\Delta }`$-resonance region in contrast to previous work Lag78 ; Lag81 ; LeS00 where a slight increase was found. Recently, this work was extended in a series of papers Dar04a ; Dar05a ; Dar05b ; Dar05c ; DaS05 to a study of various polarization asymmetries of the semi-exclusive differential cross section for $`\stackrel{}{d}(\stackrel{}{\gamma },\pi )NN`$. The semi-exclusive beam asymmetry $`\mathrm{\Sigma }`$ for linearly polarized photons and the target asymmetries $`T_{IM}`$ with respect to polarized deuterons were considered in Dar04a and beam-target asymmetries in Dar05a , in both cases only in impulse approximation. Final state interaction effects were subsequently discussed in Dar05b ; Dar05c ; DaS05 . Unfortunately, many of these results are based on incorrect expressions for polarization observables as pointed out recently in ArF05 , where a thorough derivation of these observables is given.
Thus the present work was motivated firstly to use a better elementary production operator from the MAID model Maid , allowing one to go to higher photon energies and also to give a more reliable description of the threshold region. Secondly, we would like to clarify the role of the final state interaction (FSI) in view of the above mentioned differences in the role of FSI effects. Thirdly, the increasing importance of polarization observables requires a more thorough and reliable treatment as done in Dar04a ; Dar05a ; Dar05b ; Dar05c ; DaS05 . In the present work we consider again besides the impulse approximation (IA) complete rescattering in the final two-body subsystems, i.e. in the $`NN`$\- and $`\pi N`$-subsystems. Results on the spin asymmetry of the total cross section, based on the present approach, and its contribution to the Gerasimov-Drell-Hearn sum rule have already been reported in ArF04 .
In the next section we briefly review the basic formalism for the general differential cross section with inclusion of polarization observables as derived in ArF05 . Furthermore, the essential ingredients for the calculation of the $`T`$-matrix in the impulse approximation (IA) and the rescattering contributions are described here. The results on the unpolarized differential cross section for the semi-exclusive process $`\stackrel{}{d}(\stackrel{}{\gamma },\pi )NN`$ as well as all beam, target, and beam-target asymmetries will be presented and discussed in Sect. III together with a comparison to existing data. Finally, we will conclude in Sect. IV with a summary and an outlook. The separation of the the various asymmetries of the semi-exclusive differential cross section are discussed in Appendix A, and a modified impulse approximation is given in Appendix B.
## II The formalism
To begin with, we will briefly outline the kinematic framework of the reaction under study, namely
$$\gamma (k,\stackrel{}{\epsilon }_\mu )+d(p_d)\pi (q)+N_1(p_1)+N_2(p_2),$$
(1)
defining the notation of the four-momenta of the participating particles. The circular polarization of the photon is denoted by $`\stackrel{}{\epsilon }_\mu `$ ($`\mu =\pm 1`$). For the description of cross sections and polarization observables we take as reference frame the laboratory frame and as independent variables for the characterization of the final state the outgoing pion momentum $`\stackrel{}{q}=(q,\theta _q,\varphi _q)`$ and the spherical angles $`\mathrm{\Omega }_p=(\theta _p,\varphi _p)`$ of the relative momentum $`\stackrel{}{p}=(\stackrel{}{p}_1\stackrel{}{p}_2)/2=(p,\mathrm{\Omega }_p)`$ of the two outgoing nucleons. The coordinate system is chosen as right-handed with $`z`$-axis along the photon momentum $`\stackrel{}{k}`$. According to the convention of ArF05 , we distinguish in general three planes: (i) the photon plane spanned by the photon momentum and the direction of maximal linear photon polarization, which defines the direction of the $`x`$-axis, (ii) the pion plane, spanned by photon and pion momenta, which intersects the photon plane along the $`z`$-axis with an angle $`\varphi _q`$, and (iii) the nucleon plane spanned by total and relative momenta of the two nucleons. It intersects the pion plane along the total momentum of the two nucleons (see Fig. 1). In case the linear photon polarization vanishes, one can choose $`\varphi _q=0`$ and then photon and pion planes coincide.
### II.1 The $`T`$-matrix
All observables are determined by the $`T`$-matrix elements of the electromagnetic pion production current $`J_{\gamma \pi }`$ between the initial deuteron and the final $`\pi NN`$ states
$$T_{sm_s,\mu m_d}(q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)=^{()}\stackrel{}{q},\stackrel{}{p};sm_s|\stackrel{}{\epsilon }_\mu \stackrel{}{J}_{\gamma \pi }(0)|\stackrel{}{d};1m_d,$$
(2)
where $`s`$ and $`m_s`$ denote the total spin and its projection on the relative momentum of the outgoing two nucleons and $`m_d`$ correspondingly the deuteron spin projection on the chosen $`z`$-axis. As is shown in ArF05 , the dependence on $`\varphi _q`$ can be split of, i.e.
$$T_{sm_s\mu m_d}(q,\mathrm{\Omega }_q,\mathrm{\Omega }_p)=e^{i(\mu +m_dm_s)\varphi _q}t_{sm_s\mu m_d}(q,\theta _q,\theta _p,\varphi _{pq}).$$
(3)
Thus the small $`t`$-matrix depends besides on $`q`$, $`\theta _q`$, and $`\theta _p`$ only on the difference of the azimuthal angles of $`\stackrel{}{q}`$ and $`\stackrel{}{p}`$, i.e. on $`\varphi _{pq}=\varphi _p\varphi _q`$. The small $`t`$-matrix elements are the basic quantities which determine the differential cross section and asymmetries. The latter are listed explicitly in ArF05 .
For the calculation of the $`T`$-matrix we start from the impulse approximation (IA) to which the contributions from $`NN`$\- and $`\pi N`$-rescattering are added. Possible two-body contributions to the electromagnetic interaction are neglected. Thus the treatment is completely analogous to previous work on incoherent $`\pi `$\- and $`\eta `$-photoproduction on the deuteron DaA03a ; FiA97 to which the reader is referred for formal details. Then the $`T`$-matrix is given by the sum
$$T_{sm_s\mu m_d}=T_{sm_s\mu m_d}^{IA}+T_{sm_s\mu m_d}^{NN}+T_{sm_s\mu m_d}^{\pi N}.$$
(4)
For the IA contribution, which describes the production on one nucleon while the other acts as a spectator, one has
$`T_{sm_s\mu m_d}^{IA}`$ $`=`$ $`\stackrel{}{q},\stackrel{}{p},sm_s|\left[t_{\gamma \pi }(1)+t_{\gamma \pi }(2)\right]|\mathrm{\hspace{0.17em}1}m_d`$ (5)
$`=`$ $`{\displaystyle \underset{m_s^{}}{}}(sm_s|\stackrel{}{p}_1|t_{\gamma \pi }(W_{\gamma N_1})|\stackrel{}{p}_2\varphi _{m_s^{}m_d}(\stackrel{}{p}_2)|\mathrm{\hspace{0.17em}1}m_s^{}(12)),`$
where $`t_{\gamma \pi }`$ denotes the elementary pion photoproduction operator, which we take from the MAID model, $`W_{\gamma N_1}`$ the invariant energy of the $`\gamma N`$ system, and $`\stackrel{}{p}_{1/2}=(\stackrel{}{k}\stackrel{}{q})/2\pm \stackrel{}{p}`$, and $`\varphi _{m_sm_d}(\stackrel{}{p})`$ is related to the internal deuteron wave function in momentum space by
$$\stackrel{}{p},1m_s|1m_d^{(d)}=\varphi _{m_sm_d}(\stackrel{}{p})=\underset{L=0,2}{}\underset{m_L}{}i^L(Lm_L\mathrm{\hspace{0.17em}1}m_s|1m_d)u_L(p)Y_{Lm_L}(\widehat{p}).$$
(6)
In view of the fact that for neutral pion production the major influence of $`NN`$-rescattering arises from the non-orthogonality of the final $`NN`$-plane wave to the deuteron ground state, we have considered in addition a modified IA-amplitude, for which the deuteron component is projected out from the final plane wave (see Appendix B for details).
The two rescattering contributions have a similar structure
$`T_{sm_s\mu m_d}^{NN}`$ $`=`$ $`\stackrel{}{q},\stackrel{}{p},sm_s|t_{NN}G_{NN}[t_{\gamma \pi }(W_{\gamma N_1})+t_{\gamma \pi }(W_{\gamma N_2})]|\mathrm{\hspace{0.17em}1}m_d,`$ (7)
$`T_{sm_s\mu m_d}^{\pi N}`$ $`=`$ $`\stackrel{}{q},\stackrel{}{p},sm_s|t_{\pi N}G_{\pi N}[t_{\gamma \pi }(W_{\gamma N_1})+t_{\gamma \pi }(W_{\gamma N_2})]|\mathrm{\hspace{0.17em}1}m_d,`$ (8)
where $`t_{NN}`$ and $`t_{\pi N}`$ denote respectively the $`NN`$ and $`\pi N`$ scattering matrices and $`G_{NN}`$ and $`G_{\pi N}`$ the corresponding free two-body propagators. For the actual evaluation, the scattering matrices are expanded into partial waves, and the expansion is then truncated at a certain angular momentum such that convergence is achieved.
### II.2 The differential cross section including polarization asymmetries
The general five-fold differential cross section $`d^5\sigma /dqd\mathrm{\Omega }_qd\mathrm{\Omega }_p`$ including beam and target polarization has been derived in ArF05 , and we refer to this work for details. In the present work we are interested in the semi-exclusive reaction $`\stackrel{}{d}(\stackrel{}{\gamma },\pi )NN`$ where only the produced pion is detected. This means integration of the five-fold differential cross section over $`\mathrm{\Omega }_p`$, yielding as semi-exclusive differential cross section ArF05
$`{\displaystyle \frac{d^3\sigma }{dqd\mathrm{\Omega }_q}}`$ $`=`$ $`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}[1+P_l^\gamma \{\stackrel{~}{\mathrm{\Sigma }}^l\mathrm{cos}2\varphi _q+{\displaystyle \underset{I=1}{\overset{2}{}}}P_I^d{\displaystyle \underset{M=I}{\overset{I}{}}}\stackrel{~}{T}_{IM}^l\mathrm{cos}[M\varphi _{qd}2\varphi _q\delta _{I1}\pi /2]d_{M0}^I(\theta _d)\}`$ (9)
$`+{\displaystyle \underset{I=1}{\overset{2}{}}}P_I^d{\displaystyle \underset{M=0}{\overset{I}{}}}(\stackrel{~}{T}_{IM}^0\mathrm{cos}[M\varphi _{qd}\delta _{I1}\pi /2]+P_c^\gamma \stackrel{~}{T}_{IM}^c\mathrm{sin}[M\varphi _{qd}+\delta _{I1}\pi /2])d_{M0}^I(\theta _d)],`$
where $`\varphi _{qd}=\varphi _q\varphi _d`$. Explicit expressions for the asymmetries $`\stackrel{~}{\mathrm{\Sigma }}^l`$, $`\stackrel{~}{T}_{IM}`$, and $`\stackrel{~}{T}_{IM}^{c/l}`$ are listed in the appendix of ArF05 . Furthermore, the photon polarization is characterized by the degree of circular polarization $`P_c^\gamma `$ and the degree of linear polarization $`P_l^\gamma `$, where the $`x`$-axis has been chosen in the direction of maximum linear polarization. The deuteron target is characterized by four parameters, namely the vector and tensor polarization parameters $`P_1^d`$ and $`P_2^d`$, respectively, and by the orientation angles $`\theta _d`$ and $`\varphi _d`$ of the deuteron orientation axis with respect to which the deuteron density matrix has been assumed to be diagonal.
We would like to point out that in forward and backward pion emission the following asymmetries vanish at $`\theta _q=0`$ or $`\pi `$
$$\stackrel{~}{\mathrm{\Sigma }}^l=0,\stackrel{~}{T}_{IM}^{0,c}=0\text{for}M0,\text{and}T_{IM}^l=0\text{for}M2,$$
(10)
because of helicity conservation, i.e. in this case the cross section should not depend on $`\varphi _q`$.
In the next section we will present results for the case that only the direction of the outgoing pion is measured and not its momentum. Then the corresponding differential cross section $`d^2\sigma /d\mathrm{\Omega }_q`$ is given by an expression formally analogous to (9), where unpolarized cross section and asymmetries are replaced by
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}`$ $``$ $`{\displaystyle \frac{d^2\sigma _0}{d\mathrm{\Omega }_q}}={\displaystyle _{q_{min}(\theta _q)}^{q_{max}(\theta _q)}}𝑑q{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}},`$ (11)
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{\mathrm{\Sigma }}^l(q,\theta _q)`$ $``$ $`{\displaystyle \frac{d^2\sigma _0}{d\mathrm{\Omega }_q}}\mathrm{\Sigma }^l(\theta _q)={\displaystyle _{q_{min}(\theta _q)}^{q_{max}(\theta _q)}}𝑑q{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{\mathrm{\Sigma }}^l(q,\theta _q),`$ (12)
$`{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{T}_{IM}^\alpha (q,\theta _q)`$ $``$ $`{\displaystyle \frac{d^2\sigma _0}{d\mathrm{\Omega }_q}}T_{IM}^\alpha (\theta _q)={\displaystyle _{q_{min}(\theta _q)}^{q_{max}(\theta _q)}}𝑑q{\displaystyle \frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}}\stackrel{~}{T}_{IM}^\alpha (q,\theta _q),\alpha \{0,l,c\}.`$ (13)
The upper and lower integration limits are given by
$`q_{max}(\theta _q)`$ $`=`$ $`{\displaystyle \frac{1}{2b}}\left(a\omega \mathrm{cos}\theta _q+E_{\gamma d}\sqrt{a^24bm_\pi ^2}\right),`$ (14)
$`q_{min}(\theta _q)`$ $`=`$ $`\mathrm{max}\{0,{\displaystyle \frac{1}{2b}}\left(a\omega \mathrm{cos}\theta _qE_{\gamma d}\sqrt{a^24bm_\pi ^2}\right)\},`$ (15)
where
$`a`$ $`=`$ $`W_{\gamma d}^2+m_\pi ^24m_N^2,`$ (16)
$`b`$ $`=`$ $`W_{\gamma d}^2+\omega ^2\mathrm{sin}^2\theta _q,`$ (17)
$`W_{\gamma d}^2`$ $`=`$ $`m_d(m_d+2\omega ),`$ (18)
$`E_{\gamma d}`$ $`=`$ $`m_d+\omega .`$ (19)
Finally, in the total cross section only a few polarization observables survive, namely one has ArF05
$$\sigma (P_l^\gamma ,P_c^\gamma ,P_1^d,P_2^d)=\sigma _0\left[1+P_2^d\overline{T}_{20}^{\mathrm{\hspace{0.17em}0}}\frac{1}{2}(3\mathrm{cos}^2\theta _d1)+P_c^\gamma P_1^d\overline{T}_{10}^c\mathrm{cos}\theta _d+P_l^\gamma P_2^d\overline{T}_{22}^l\mathrm{cos}(2\varphi _d)\frac{\sqrt{6}}{4}\mathrm{sin}^2\theta _d\right],$$
(20)
where unpolarized total cross section and asymmetries are given by
$`\sigma _0`$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }_q_{q_{min}(\theta _q)}^{q_{max}(\theta _q)}𝑑q\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}},`$ (21)
$`\sigma _0\overline{T}_{IM}^\alpha `$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }_q_{q_{min}(\theta _q)}^{q_{max}(\theta _q)}𝑑q\frac{d^3\sigma _0}{dqd\mathrm{\Omega }_q}\stackrel{~}{T}_{IM}^\alpha },`$ (22)
with $`\alpha \{0,l,c\}`$. This concludes the formal part.
## III Results and discussion
For the calculation of the $`NN`$-rescattering contribution we have taken the separable representation of the realistic Paris potential from HaP85 and included all partial waves up to $`{}_{}{}^{3}D_{3}^{}`$. Also the deuteron wave function was calculated using this potential. In principle, any other realistic potential, e.g. the Bonn r-space potential MaH87 , could be used as well, because the results do not depend sensitively on the potential model as was found in DaA03a . However, one comment is in order with respect to the question, whether the use of such a nonrelativistic $`NN`$-potential can be justified in view of the high energies involved, because the potential is fit to $`NN`$-scattering data up to nucleon lab kinetic energies of $`T_{NN}=330`$ MeV only. But even up to $`T_{NN}=500`$ MeV, it reproduces the phase shifts reasonably well because the inelasticity parameters are still small. Since the calculation of $`NN`$-rescattering requires $`NN`$-scattering amplitudes also at considerably higher energies, where we still use the same separable representation, one might expect some serious error. On the other hand, the size of this error depends on the relative size of that part of the phase space, where the energy of the $`NN`$-subsystem exceeds the region of validity. In order to estimate the corresponding error, we present in Fig. 2 the cross section $`d\sigma /dT_N`$ as a function of the equivalent nucleon kinetic lab energy $`T_N`$ of $`NN`$-scattering. One readily notes that even for a photon energy of 800 MeV the dominant part of the cross section corresponds to nucleon laboratory kinetic energies of less than 500 MeV. Thus, the use of such a realistic potential is justified.
Similarly, $`\pi N`$-rescattering is evaluated using a realistic separable representation of the $`\pi N`$-interaction from NoB90 and taking into account all partial waves up to $`l=2`$. We have evaluated the semi-exclusive differential cross section including the various polarization asymmetries in IA alone and with inclusion of $`NN`$\- and $`\pi N`$-rescattering. As already mentioned, the elementary pion photoproduction amplitude is taken from the MAID model. Since it is parametrized in terms of the Chew-Goldberger-Low-Nambu amplitudes (CGLN) ChG57 , defined in the $`\gamma N`$ c.m. frame, we first had to transform this amplitude into a general frame of reference. This is achieved by introducing invariant amplitudes and establishing relations to the CGLN amplitudes in the c.m. frame. This is described in detail in SaA04 and needs not to be repeated here. Furthermore, for the evaluation of the MAID amplitudes the invariant $`\pi N`$-energy and the pion angle in the $`\pi N`$ c.m. system have to be specified. For this purpose we assume that the four-momenta $`q`$ and $`p_f`$ of pion and active nucleon in the final state obey the on-shell condition. Then the corresponding c.m. variables are obtained by a Lorentz transformation with $`\stackrel{}{\beta }=(\stackrel{}{q}+\stackrel{}{p}_f)/(w_\pi +E_f)`$. The initial four-momentum $`p_i`$ of the active nucleon is determined by assuming four-momentum conservation at the elementary vertex, i.e. $`p_i=q+p_fk`$.
### III.1 Total cross sections
We begin the discussion of the results with the total cross sections for the three charge states displayed in Fig. 3, where we have plotted in addition the corresponding elementary cross section for comparison. The threshold region is separately plotted in Fig. 4.
One readily notes that for charged pion production the rescattering effects are in general quite small. Only close to threshold they lead to a sizeable enhancement (see left panel of Fig. 4), mainly by $`NN`$-rescattering while $`\pi N`$-rescattering is almost negligible.
The significant role of $`NN`$-FSI in photoproduction of mesons at very low energies has been noted previously in Nob67 ; Lag81 ; FiA97 . We would only like to mention that this effect has a kinematical rather than a dynamical origin. As pointed out in FiA97 , in IA the energy needed for pion production below the free nucleon threshold is provided exclusively through the high momentum of a nucleon moving towards the incoming photon. As a result, the IA predicts an anomalous suppression of the cross section, because of a small probability for finding a nucleon with high momentum in the deuteron wave function. Thus, $`NN`$-rescattering provides a mechanism to balance the strong mismatch between the momentum needed to produce the pion and the characteristic internal nucleon momentum in the deuteron, so that the strong suppression appearing in IA can be avoided. The same reasoning is true for $`\pi ^0`$ channel. But in this case the below mentioned role of orthogonality turns out to be very important so that the resulting FSI effect becomes destructive already at 10 MeV and higher above threshold.
At higher energies, near and above the maximum the cross sections of charged pion production are reduced by FSI by a few percent. Therefore, the main difference to the elementary cross section comes from the Fermi motion leading to a slight reduction and a shift of the maximum, and a broadening of the whole distribution. As one can see in Fig. 3, the result of these features is that the energy integrated cross section
$$I(E)=\underset{E_{th}}{\overset{E}{}}\sigma (E_\gamma )𝑑E_\gamma $$
(23)
as function of the upper integration limit $`E`$, is preserved over a wide energy region. In other words, the integral at $`E=1`$ GeV has approximately the same value for the reaction on the free nucleon and on the deuteron. Concerning the role of FSI in the $`\pi ^\pm `$ channels at higher energies, one can assume that the interaction between the emitted particles leads basically to a redistribution of events in phase space, so that in the absence of absorption the overall yield of particles does not change and the FSI effect in the total cross section remains insignificant. It is, however, not the case for the near threshold region where, as was discussed above, the IA amplitude turns out to be anomalously suppressed in the available phase space.
In contrast to charged pion production, one notes quite large FSI effects from $`NN`$-rescattering in $`\pi ^0`$-production as displayed in the right panel of Fig. 3. However, these effects are mainly due to the elimination of the spurious coherent contribution in the IA. That is nicely demonstrated by the results for the modified IA, also shown in Fig. 3, where the deuteron wave function component in the final plane wave is projected out (see Appendix B for details). The additional rescattering contributions beyond the orthogonality effect are comparable to what was found in charged pion production, except near threshold where one finds a much smaller influence of FSI. With respect to the elementary cross section, one notes a sizeable reduction for incoherent $`\pi ^0`$-production on the deuteron, because part of the strength goes to the coherent channel $`d(\gamma ,\pi ^0)d`$. The spurious admixture of the coherent channel in the IA cross section was also discussed in Sio01 .
The polarization observables of the total cross section are shown in Fig. 5 only for $`\pi ^{}`$\- and $`\pi ^0`$-production since the ones for $`\pi ^+`$\- and $`\pi ^{}`$-production are quite similar. The vector asymmetry for circularly polarized photons $`\sigma _0\overline{T}_{10}^c`$ (middle panels), which governs the spin asymmetry of the GDH sum rule, reproduces the results reported in ArF04 . The sensitivity to FSI is weak, similar to the total cross sections. The tensor asymmetries for $`\pi ^{}`$-production are sizeable in the near threshold region, exhibiting a sharp peak in absolute values. $`\sigma _0\overline{T}_{20}^{\mathrm{\hspace{0.17em}0}}`$ reaches a second quite broad maximum above the $`\mathrm{\Delta }`$-resonance, while $`\sigma _0\overline{T}_{22}^l`$ remains small. FSI is more notable than in the total cross section. For $`\pi ^0`$ production, the tensor asymmetries show quite a different behavior compared to $`\pi ^{}`$-production. They are quite small in general, almost vanishing in the near threshold region. Furthermore, they show huge FSI effects beyond the orthogonality effect.
### III.2 Differential cross section
Angular distributions are shown in Fig. 6 for $`\pi ^{}`$-production and in Fig. 7 for $`\pi ^0`$-production. Like for the total cross section, one notes for $`\pi ^{}`$-production a small size of rescattering effects. At the lowest energy (250 MeV) they decrease the cross section at forward angles and increase it in backward direction. At 350 MeV one finds a slight enhancement near the maximum, while at the two higher energies a sizeable decrease is found only in the extreme forward direction. For comparison, we show also the elementary cross sections, and again the smoothing effect of the Fermi motion is apparent.
Compared our calculation with the predictions of Lag81 and LeS00 we note a significant disagreement for the IA at forward angles. At $`E_\gamma `$=500 MeV our IA cross section exhibits a visible rise which is governed mainly by the pion photoelectric term in the elementary amplitude. At the same time the calculations of LeS00 (see dotted curve in the lower left panel of Fig. 6 at $`E_\gamma `$=500 MeV) predict a strong reduction in the same angular region, so that the corresponding IA result at $`\theta _\pi =0`$ is about 10 times smaller than ours.
A possible reason of this discrepancy has been discussed in DaA03a . Namely, it was claimed in LeS00 that the drastic reduction of the cross section compared to the elementary reaction is a manifestation of the Pauli principle, leading to a strong suppression of the $`pp`$ states in the region of small relative momenta in this subsystem. On the contrary, we would like to note that this suppression takes place only in the triplet state ($`s=1`$) of the emitted protons, whereas the singlet part ($`s=0`$), where the processes on the individual nucleons can coherently enhance each other, peaks at $`\theta _\pi =0`$ similar to the elementary cross section. The resulting angular distribution shows some reduction at forward angles, compared to the elementary reaction, which, however, is not as large as that exhibited by the triplet part only and is, therefore, much smaller compared to the prediction of Lag81 ; LeS00 .
On the other hand, inclusion of FSI leads in Lag81 ; LeS00 to a significant increase at small angles compared to the IA cross section, whereas in our case the effect of FSI is small and destructive. It is, therefore, interesting to note that the difference between the full results of the present work (see Fig. 24 below) and of LeS00 is rather small, so that both models describe the data equally well.
As was mentioned above, in the $`\pi ^0`$ channel the orthogonality effect appears so that the role of FSI is approximately determined by the fraction of the coherent reaction $`\gamma d\pi ^0d`$ in the inclusive $`\pi ^0`$ photoproduction. Fig. 7 demonstrates huge spurious contribution of the coherent process in IA which is eliminated by applying the modified IA. This result could be expected because the effect of orthogonality should be especially visible in the region of small pion angles, where the momentum transferred to the nucleons is minimal and the overlap between initial and final wave functions is most important. The additional FSI effects are quite small. Only at 350 MeV one notes a slight increase of the cross section near the maximum around 60 and at 780 MeV a more sizeable increase.
### III.3 Beam asymmetry for linearly polarized photons
As next we turn to the beam asymmetry $`\mathrm{\Sigma }^l`$ displayed in Fig. 8 for the three charge states at various energies in the $`\mathrm{\Delta }`$-resonance region and above. In all of these asymmetries one notes a relatively small influence from FSI. We will consider first the results for $`\pi ^{}`$-production in the top panels of Fig. 8. For the lowest three energies the photon asymmetry $`\mathrm{\Sigma }^l`$ is negative. It is quite small below the $`\mathrm{\Delta }`$-region and increases considerably in absolute size with increasing energy, becoming quite broad around 350 Mev but more forward peaked at 500 Mev. At the highest energy of 780 MeV one notes a different behavior. The deep minimum in forward direction turns into a broad positive distribution above 40. The three lowest energies are also representative for $`\pi ^+`$-production (middle panels of Fig. 8). However at the highest energy the asymmetry remains negative, but exhibits a secondary minimum around 90.
The $`\mathrm{\Sigma }^l`$-asymmetry for $`\pi ^0`$-production, shown in the lowest panels of Fig. 8, exhibits quite a different behavior compared to charged pion production. One finds a broad, structureless and sinus-shaped negative distribution with a maximum around 80 to 90 with increasing amplitude by about a factor of two going from 250 to 500 MeV. At 780 MeV the width of the distribution becomes smaller and the minimum moves towards forward angles. In general also here the influence of rescattering is small. It is mainly due to the removal of the spurious coherent contribution. The relatively largest influence appears at the lowest energy of $`E_\gamma =250`$ MeV.
Compared to the results in Dar04a ; Dar05a ; Dar05c one readily notes quite substantial differences for both, charged and neutral pion production, both in angular behavior and also in absolute size, in particular at higher energies. For $`\pi ^{}`$-production a larger influence of FSI is found in DaS05 which we cannot confirm. Furthermore, $`\mathrm{\Sigma }^l`$ does not vanish in DaS05 at $`\theta _q=0`$ and $`\pi `$ as it should, although it is small. Strangely enough, the authors mention this feature as a notable effect. The origin of these differences is not clear. In any case, the comparison suffers from the questionable formal expressions used for the calculation of $`\mathrm{\Sigma }^l`$ which, moreover, differ in the various publications Dar04a ; Dar05c ; DaS05 .
### III.4 Target asymmetries for oriented deuteron
The target asymmetries $`T_{IM}^0`$ are shown in Figs. 9 through 11 for the three charge states, respectively. In view of the similarity of the asymmetries for the two charge states except for $`T_{11}^0`$ at 780 MeV, we will restrict the discussion to $`\pi ^{}`$\- and $`\pi ^0`$-production. In general FSI effects are again small, the largest appear in $`T_{21}^0`$. The vector asymmetry $`T_{11}^0`$ is positive up to 500 MeV and shows a broad distribution over the whole angular range. But at 780 MeV an oscillatory behavior develops with sizeable amplitude. Only at the lowest energy of $`E_\gamma =250`$ MeV one notes some notable influence at forward angles from FSI. With respect to Dar04a ; Dar05a ; DaS05 we find an opposite sign, a smaller amplitude and also a somewhat different angular behavior.
The tensor target asymmetry $`T_{20}^0`$ exhibits a pronounced sharp peak at 0 and a rapid fall-off with increasing angles, remaining quite small above 30. The results are similar to Dar04a ; Dar05a ; DaS05 except that in contrast to the small negative values of Dar04a ; Dar05a at backward angles, we find small positive values. Moreover, we find smaller FSI influences at backward angles than in DaS05 . $`T_{21}^0`$ peaks at small angles around 20 for $`E_\gamma =250`$ MeV which disappears at the three higher energies becoming a broader distribution with a considerably smaller size. FSI shows some notable influence. Again we find significant differences to Dar04a ; Dar05a with respect to shape, size, and FSI effects, in particular at the lowest energy. Finally, $`T_{22}^0`$ exhibits a prominent peak in forward direction which becomes sharper and moves towards smaller angles with increasing photon energy. This is in qualitative agreement with Dar04a ; Dar05a although the size is different and our FSI effects are much smaller.
The target asymmetries for $`\pi ^0`$-production in Fig. 11 show quite a different behavior compared to charged pion production. The structure of $`T_{11}^0`$ changes significantly with energy. While at $`E_\gamma =250`$ MeV one finds a maximum around 130, one notes a forward negative minimum and a backward positive maximum at $`E_\gamma =350`$ MeV, at $`E_\gamma =500`$ MeV a broad and quite flat negative distribution, and finally at 780 MeV a forward maximum and a negative minimum around 80. The tensor asymmetries are much more sensitive to FSI. This is particularly apparent in $`T_{20}^0`$ exhibiting a forward negative minimum in IA which turns into a positive forward peak when FSI is switched on. Also $`T_{21}^0`$ shows such a drastic influence from FSI. $`T_{22}^0`$ is much smaller than for charged pion production and shows an oscillatory behavior. FSI effects are noticeable again. For these asymmetries the differences to the results of Dar04a ; Dar05a are again quite significant.
### III.5 Beam-target asymmetries for circularly polarized photons
As next we will discuss the double polarization asymmetries $`T_{IM}^c`$ for circularly polarized photons and oriented deuterons. They are in a certain sense complementary to the target asymmetries $`T_{IM}^0`$ because, while one is the real part, the other is the imaginary part of the basic quantities $`V_{IM}^1`$ as defined in ArF05 . Since the results for the two charged pions are again in general quite similar, we display the results only for $`\pi ^{}`$ in Fig. 12 and for $`\pi ^0`$ in Fig. 13. In contrast to what has been claimed in Dar05a , all of them are nonvanishing. As a sideremark, although $`T_{10}^c`$ should vanish according to Dar05a , a non-vanishing spin asymmetry is reported in Dar05b . For both, $`\pi ^0`$ and $`\pi ^{}`$, the vector asymmetry $`T_{10}^c`$ is quite sizeable in forward and backward direction and also around 90 for $`\pi ^0`$. It is this asymmetry which determines the Gerasimov-Drell-Hearn sum rule. The influence of FSI is quite marginal, in particular for $`\pi ^{}`$. For both charge states the energy dependence is weak. The other vector asymmetry $`T_{11}^c`$ shows a rather different behavior. For $`\pi ^{}`$ it is predominantly negative with a slight preference of the backward direction, especially at higher energies, while for $`\pi ^0`$ it has a positive maximum around 30 and a negative minimum around 150, almost independent of the energy. Considerably smaller are the tensor asymmetries $`T_{21}^c`$ and $`T_{22}^c`$, but they are a little more sensitive to FSI. This is particularly apparent for $`T_{21}^c`$ in $`\pi ^0`$ production above the $`\mathrm{\Delta }`$-resonance region.
### III.6 Beam-target asymmetries for linearly polarized photons
The beam-target asymmetries $`T_{IM}^l`$ for linearly polarized photons and polarized deuterons for $`\pi ^{}`$ and $`\pi ^0`$ production are shown in Figs. 14 through 17. Also these asymmetries are very similar for the two charge states $`\pi ^{}`$ and $`\pi ^+`$. The vector asymmetries $`T_{1M}^l`$ do not vanish at all as claimed in Dar04a . They are small for $`\pi ^{}`$ at the lowest energy of $`E_\gamma =250`$ MeV, but become sizeable for the higher energies (see Fig. 14). They are considerably smaller for $`\pi ^0`$ production as shown in Fig. 15 but more sensitive to FSI. Of the corresponding tensor asymmetries in Figs. 16 and 17, $`T_{22}^l`$ is by far the largest for both charge states $`\pi ^{}`$ and $`\pi ^0`$, exhibiting a pronounced forward peak which becomes slightly sharper with increasing energy. The size decreases significantly when going from $`T_{22}^l`$ to $`T_{22}^l`$. By the way, they are certainly not equal as claimed in Dar04a . While $`T_{22}^l`$ is restricted to vanish at $`\theta _q=0`$ and $`\pi `$, this is not the case for $`T_{22}^l`$. For $`\pi ^{}`$ production one notes some FSI effects in $`T_{2\pm 1}^l`$, whereas for $`\pi ^0`$ quite drastic influences from FSI can be seen, leading even to a sign change in some of them. A comparison to the results in Dar04a makes no sense because of the above mentioned wrong formal expressions in Dar04a .
### III.7 Comparison with experiment
We now will turn to a comparison with experimental data, where available. Fig. 18 shows the total cross sections for $`\pi ^{}`$\- and $`\pi ^0`$-production. The agreement with experimental data for $`\pi ^{}`$-production is satisfactory although not perfect. The theory is a little high in the maximum but low in the dip region between the $`\mathrm{\Delta }`$ and the second resonance region. For $`\pi ^0`$, one notes some slight overshooting in the maximum and a sizeable overestimation above the $`\mathrm{\Delta }`$ in the second resonance region. This discrepancy was already discussed in Kru99 where it was noted that the smearing and damping of the second resonance peak can hardly be explained by the Fermi motion effect alone. According to our results, the broadening of the resonance structure in pion production on the deuteron is quite significant. This effect is readily seen in Fig. 19, where we present differential cross sections for $`\pi ^0`$-production on the deuteron as functions of the photon energy for fixed pion angle. The cross section refers to the equivalent $`\gamma N`$ c.m. system, where the nucleon is at rest in the deuteron. The solid curves are our full calculation multiplied by energy independent factors, whose values are listed in the various panels. One readily sees that the Fermi motion and to some extent FSI leads to a disappearing of the resonance structure. The ratio of the theoretical cross section to the experimental result is about 1.75, except for $`\mathrm{cos}\theta _{\pi ^0}^{}=0.9`$ where it is about 3. However the origin of this factor is unclear.
One possible source could be the neglected interaction of nucleon resonances with the spectator nucleon which might result in a broadening of the resonances due to additional inelastic processes. In fact, within a simple model calculation for the $`\mathrm{\Delta }`$-resonance a significant lowering of the total cross section from such an interaction was found in ReA05 . A further source could be the neglected inelasticity of the final $`NN`$-interaction at higher energies. On the other hand the question arises, why the same lowering is not observed in the $`\pi ^{}`$-photoproduction where the agreement with the data is quite satisfactory.
The threshold region is shown in Fig. 20 for $`\pi ^+`$-production. Good agreement with the data is achieved, comparable in quality to a recent precision calculation in chiral perturbation theory LeB05 . The crucial role of $`NN`$-FSI at low energies, leading to a strong enhancement over the IA, is furthermore demonstrated by the differential cross section with respect to the relative energy of the two final neutrons in $`\pi ^+`$-production in Fig. 21 for two different kinematical situations, pion emission in more forward (left panel) and in more backward direction (right panel). The $`{}_{}{}^{1}S_{0}^{}`$-state of $`nn`$-scattering near threshold, a manifestation of the so-called antibound state as companion to the deuteron, is nicely resolved and also quite well reproduced by the theory if $`NN`$-rescattering is included. A similar result has been reported in LeC04 .
With respect to differential cross sections, a comparison for $`\pi ^{}`$-production is exhibited in Fig. 22 and for $`\pi ^0`$-production in Fig. 23. For $`\pi ^{}`$-production one notes a satisfactory agreement with experimental data from Benz et al. Be+73 , whereas for $`\pi ^0`$-production one finds for the three highest energies a slight overestimation of the theory in the maximum and at forward angles an underestimation.
A comparison to recent data from the LEGS-collaboration (LEGS-exp.L3b) Sandorfi on the unpolarized semi-exclusive differential cross section and $`\mathrm{\Sigma }^l`$ for $`\pi ^{}`$ production is shown in Fig. 24 for two energies near the $`\mathrm{\Delta }`$-region. One notes very little influence from FSI and quite a good agreement of the theoretical description with the data. Our results on $`\mathrm{\Sigma }^l`$ are similar to the IA calculation of Lee and Sato Lee , but at variance to Dar05c where also larger FSI effect were quoted. But the latter results are questionable for reasons mentioned above. Another comparison for $`\mathrm{\Sigma }^l`$ with data from the LEGS-collaboration Sandorfi are exhibited in Fig. 25 for constant pion emission angle as function of the photon energy. In this case the size is underestimated by the theory, although FSI shifts the results for IA in the right direction but not enough. A much better almost perfect agreement with respect to $`\mathrm{\Sigma }^l`$ is shown in Fig. 26 for $`\pi ^0`$-production.
Finally, we show in Fig. 27 for $`\pi ^{}`$-production a comparison between theory and experiment for the semi-exclusive differential spin asymmetry $`d^2(\sigma ^P\sigma ^A)/d\mathrm{\Omega }_q`$ with respect to circularly polarized photons and the deuteron spin oriented parallel (P) or antiparallel (A) to the photon spin. This spin asymmetry is related to the beam-target asymmetry $`T_{10}^c`$ according to
$$\frac{d^2(\sigma ^P\sigma ^A)}{d\mathrm{\Omega }_q}=\sqrt{6}\frac{d^2\sigma _0}{d\mathrm{\Omega }_q}T_{10}^c.$$
(24)
Compared to the predictions in Dar05b we find the depth of the minimum at 0 almost independent of the energy and a small positive asymmetry between 30 and 60. This spin asymmetry has been measured by the A2-collaboration Pedroni . However, the analysis of the data is not yet completed. Preliminary data were shown in Dar05b without authorization and thus are not shown here. Compared to these, the agreement is quite satisfactory, although in the angular region, where the data are available, the spin asymmetry is very small, almost compatible with zero. Thus it would be very desirable to have additional data at more forward angles, where the theoretical asymmetry exhibits a pronounced minimum. This concludes the discussion of results.
## IV Conclusion and outlook
In the present work we have exploited the role of polarization observables in incoherent pion photoproduction on the deuteron with particular emphasis on the influence of final state interaction in the $`NN`$\- and $`\pi N`$-subsystems of the final state. In the unpolarized total and semi-exclusive differential cross section $`d^2\sigma /d\mathrm{\Omega }_q`$, where only the direction of the produced pion is measured, the influence of final state rescattering is quite small for charged pion production for photon energies up to 1 GeV. For $`\pi ^0`$-production the influence is much larger. However, the dominant part of FSI-effect arises from the removal of a spurious coherent contribution in the impulse approximation when $`NN`$-rescattering is switched on. This is demonstrated by a modified IA, where the deuteron wave function component in the final $`NN`$-plane wave is projected out. The remaining FSI-effect is comparable to charged pion production.
As polarization observables we have considered all beam, target and beam-target asymmetries of the semi-exclusive differential cross section. Many of them are quite sizeable, in particular the photon asymmetry $`\mathrm{\Sigma }^l`$ and the various vector asymmetries. The tensor asymmetries are in general considerably smaller. They are often quite insensitive to final state rescattering. Only a few, $`T_{21}`$, and $`T_{21}^c`$, show a larger influence in charged pion production.
A very interesting and still open question concerns the disagreement between theoretical and experimental results for $`\gamma d\pi ^0np`$ in the second resonance region. Although our calculation explains the strong smearing of the resonance structure, the data are overestimated by about a factor of 1.5 (see Figs. 18 and 19). Hopefully, new measurements of the ratio
$$R=\frac{d\sigma (\gamma ,\pi ^0n)}{d\sigma (\gamma ,\pi ^0p)}$$
(25)
for quasifree photoproduction on neutrons and protons can clarify the situation. The old data from Bacci ; Hemmi , pointing to $`R=1`$ in the second resonance region, seem to be in disagreement with the results of Kru99 .
Future theoretical improvements should be devoted to the inclusion of two-body effects in the photoproduction amplitude, e.g. the interaction between a resonance and the spectator nucleon, inclusion of inelasticities in the $`NN`$-interaction and the role of relativistic effects at higher energies. The problem of off-shell effects for the elementary amplitude is another unsolved task. Furthermore, there is urgent need for a unified description of single and double pion production.
###### Acknowledgements.
We would like to thank Michael Schwamb for interesting discussions and a careful reading of the manuscript. This work was supported by the Deutsche Forschungsgemeinschaft (SFB 443).
## Appendix A Separation of polarization asymmetries
In this appendix we will discuss how the various polarization asymmetries of the semi-incusive differential cross section can be separated by a proper variation of the photon polarization parameters ($`P_l^\gamma `$ and $`P_c^\gamma `$), the deuteron polarization parameters ($`P_1^d`$ and $`P_2^d`$), the polarization angles ($`\theta _d`$ and $`\varphi _d`$), and the dependence on the pion azimuthal angle $`\varphi _q`$, similar to what has been described in ArL05 . To this end we write the differential cross section of (9) as follows
$$S(P_l^\gamma ,P_c^\gamma ,P_1^d,P_2^d)=S_0\left[1+P_1^dA_d^V+P_2^dA_d^T+P_l^\gamma (A_\gamma ^l+P_1^dA_{\gamma d}^{lV}+P_2^dA_{\gamma d}^{lT})+P_c^\gamma (P_1^dA_{\gamma d}^{cV}+P_2^dA_{\gamma d}^{cT})\right],$$
(A1)
with $`S_0=S(0,0,0,0)`$ as unpolarized differential cross section. Furthermore, we have introduced as generalized single polarization $`A`$-asymmetries
$`A_d^V(\theta _q,\varphi _{qd},\theta _d)`$ $`=`$ $`\stackrel{~}{T}_{11}^0(\theta _q)\mathrm{sin}\varphi _{qd}d_{10}^1(\theta _d),`$ (A2)
$`A_d^T(\theta _q,\varphi _q,\varphi _{qd},\theta _d)`$ $`=`$ $`{\displaystyle \underset{M=0}{\overset{2}{}}}\stackrel{~}{T}_{2M}^0(\theta _q)\mathrm{cos}(M\varphi _{qd})d_{M0}^2(\theta _d),`$ (A3)
$`A_\gamma ^l(\theta _q,\varphi _q)`$ $`=`$ $`\stackrel{~}{\mathrm{\Sigma }}^l(\theta _q)\mathrm{cos}2\varphi _q,`$ (A4)
and double polarization $`A`$-asymmetries
$`A_{\gamma d}^{lV}(\theta _q,\varphi _q,\varphi _{qd},\theta _d)`$ $`=`$ $`{\displaystyle \underset{M=1}{\overset{1}{}}}\stackrel{~}{T}_{1M}^l(\theta _q)\mathrm{sin}(M\varphi _{qd}2\varphi _q)d_{M0}^1(\theta _d)`$ (A5)
$`=`$ $`{\displaystyle \underset{M=1}{\overset{1}{}}}\stackrel{~}{T}_{1M}^l(\theta _q)[\mathrm{sin}(M\varphi _{qd})\mathrm{cos}(2\varphi _q)\mathrm{cos}(M\varphi _{qd})\mathrm{sin}(2\varphi _q)]d_{M0}^1(\theta _d),`$
$`A_{\gamma d}^{lT}(\theta _q,\varphi _q,\varphi _{qd},\theta _d)`$ $`=`$ $`{\displaystyle \underset{M=2}{\overset{2}{}}}\stackrel{~}{T}_{2M}^l(\theta _q)\mathrm{cos}(M\varphi _{qd}2\varphi _q)d_{M0}^2(\theta _d)`$ (A6)
$`=`$ $`{\displaystyle \underset{M=2}{\overset{2}{}}}\stackrel{~}{T}_{2M}^l(\theta _q)[\mathrm{cos}(M\varphi _{qd})\mathrm{cos}(2\varphi _q)+\mathrm{sin}(M\varphi _{qd})\mathrm{sin}(2\varphi _q)]d_{M0}^2(\theta _d),`$
$`A_{\gamma d}^{cV}(\theta _q,\varphi _{qd},\theta _d)`$ $`=`$ $`{\displaystyle \underset{M=0}{\overset{1}{}}}\stackrel{~}{T}_{1M}^c(\theta _q)\mathrm{cos}(M\varphi _{qd})d_{M0}^1(\theta _d),`$ (A7)
$`A_{\gamma d}^{cT}(\theta _q,\varphi _{qd},\theta _d)`$ $`=`$ $`{\displaystyle \underset{M=1}{\overset{2}{}}}\stackrel{~}{T}_{2M}^c(\theta _q)\mathrm{sin}(M\varphi _{qd})d_{M0}^2(\theta _d).`$ (A8)
Each of these generalized $`A`$-asymmetries are functions of $`\theta _q`$ and of some of the variables $`\varphi _q`$, $`\varphi _{qd}`$, and $`\theta _d`$. The first step is to isolate them by evaluating linear combinations of cross sections for different values of the appropriate polarization parameters ($`P_{l,c}^\gamma `$, $`P_{1,2}^d`$). For example, $`A_\gamma ^l`$ is obtained from the cross section difference for unpolarized deuterons and linearly polarized photons ($`P_l^\gamma >0`$) minus the unpolarized one ($`P_l^\gamma =0`$), i.e.
$$A_\gamma ^l=\frac{1}{P_l^\gamma S_0}\left[S(P_l^\gamma ,0,0,0)S_0\right].$$
(A9)
The others can be obtained by the following linear combinations of cross sections
$`A_d^V`$ $`=`$ $`{\displaystyle \frac{1}{2P_1^dS_0}}\left[S(0,0,P_1^d,P_2^d)S(0,0,P_1^d,P_2^d)\right],`$ (A10)
$`A_d^T`$ $`=`$ $`{\displaystyle \frac{1}{2P_2^dS_0}}\left[S(0,0,P_1^d,P_2^d)+S(0,0,P_1^d,P_2^d)2S_0\right],`$ (A11)
$`A_{\gamma d}^{cV}`$ $`=`$ $`{\displaystyle \frac{1}{4P_c^\gamma P_1^dS_0}}\left[S(0,P_c^\gamma ,P_1^d,P_2^d)S(0,P_c^\gamma ,P_1^d,P_2^d)S(0,P_c^\gamma ,P_1^d,P_2^d)+S(0,P_c^\gamma ,P_1^d,P_2^d)\right],`$ (A12)
$`A_{\gamma d}^{cT}`$ $`=`$ $`{\displaystyle \frac{1}{4P_c^\gamma P_2^dS_0}}\left[S(0,P_c^\gamma ,P_1^d,P_2^d)S(0,P_c^\gamma ,P_1^d,P_2^d)+S(0,P_c^\gamma ,P_1^d,P_2^d)S(0,P_c^\gamma ,P_1^d,P_2^d)\right],`$ (A13)
$`A_{\gamma d}^{lV}`$ $`=`$ $`{\displaystyle \frac{1}{2P_l^\gamma P_1^dS_0}}\left[S(P_l^\gamma ,0,P_1^d,P_2^d)S(0,0,P_1^d,P_2^d)S(P_l^\gamma ,0,P_1^d,P_2^d)+S(0,0,P_1^d,P_2^d)\right],`$ (A14)
$`A_{\gamma d}^{lT}`$ $`=`$ $`{\displaystyle \frac{1}{2P_l^\gamma P_2^dS_0}}[S(P_l^\gamma ,0,P_1^d,P_2^d)S(0,0,P_1^d,P_2^d)+S(P_l^\gamma ,0,P_1^d,P_2^d)S(0,0,P_1^d,P_2^d)`$ (A15)
$`2P_l^\gamma A_\gamma ^lS_0].`$
Two of the generalized $`A`$-asymmetries in (A2) through (A8) yield directly one asymmetry each, namely $`\stackrel{~}{\mathrm{\Sigma }}^l`$ from $`A_\gamma ^l`$ and $`\stackrel{~}{T}_{11}^0`$ from $`A_d^V`$, i.e.
$`\stackrel{~}{\mathrm{\Sigma }}^l`$ $`=`$ $`A_\gamma ^l(\varphi _q=0),`$ (A16)
$`\stackrel{~}{T}_{11}^0`$ $`=`$ $`\sqrt{2}A_d^V(\varphi _{qd}=\pi /2,\theta _d=\pi /2).`$ (A17)
One could also obtain $`\stackrel{~}{\mathrm{\Sigma }}^l`$ from taking for linearly polarized photons the cross section difference of pions in the photon plane to pions perpendicular to this plane, i.e.
$$\stackrel{~}{\mathrm{\Sigma }}^l=\frac{1}{P_l^\gamma }\frac{S(P_l^\gamma >0,0,0,0)|_{\varphi _q=0}S(P_l^\gamma >0,0,0,0)|_{\varphi _q=\pi /2}}{S(P_l^\gamma >0,0,0,0)|_{\varphi _q=0}+S(P_l^\gamma >0,0,0,0)|_{\varphi _q=\pi /2}}.$$
(A18)
The remaining general asymmetries $`A_d^T`$ and $`A_{\gamma d}^{(l/c)(V/T)}`$ contain linear combinations of the asymmetries $`\stackrel{~}{T}_{2M}^0`$ and $`\stackrel{~}{T}_{IM}^{l/c}`$. In order to separate the latter, one can exploit the dependence of the $`A`$-asymmetries on the angular variables $`\varphi _q`$, $`\varphi _{qd}`$, and $`\theta _d`$. This is achieved, following the analogous problem in deuteron electrodisintegration ArL05 , by observing that the general functional form of an $`A`$-asymmetry is
$$A^I(\varphi _q,\varphi _{qd},\theta _d)=\underset{M=I}{\overset{I}{}}\alpha _{IM}(\varphi _q,\varphi _{qd})d_{M0}^I(\theta _d),(I=1,2),$$
(A19)
where
$$\alpha _{IM}(\varphi _q,\varphi _{qd})=c_{IM}(\varphi _q)\mathrm{cos}M\varphi _{qd}+s_{IM}(\varphi _q)\mathrm{sin}M\varphi _{qd},$$
(A20)
and the $`\varphi _q`$-dependent functions $`c_{IM}(\varphi _q)`$ and $`s_{IM}(\varphi _q)`$ have either the form
$$a_0+a_1\mathrm{cos}\varphi _q+a_2\mathrm{cos}2\varphi _q$$
(A21a)
or
$$b_1\mathrm{sin}\varphi _q+b_2\mathrm{sin}2\varphi _q.$$
(A21b)
One should note, that for $`A_{\gamma d}^{c(V/T)}`$ the $`\varphi _q`$-dependence is absent, i.e. $`c_{IM}0`$ or $`s_{IM}0`$, and the sum over $`M`$ in (A19) runs from 0 through $`I`$. For a given $`I`$ the $`M`$-components $`\alpha _{IM}(\varphi _q,\varphi _{qd})`$ of the asymmetry $`A^I(\varphi _q,\varphi _{qd},\theta _d)`$ can be separated by a proper choice of $`\theta _d`$ exploiting the properties of the small $`d_{M0}^I`$-functions. For $`I=1`$ (vector asymmetries), taking $`\theta _d=0`$ or $`\pi /2`$, i.e. $`d_{M0}^1(0)=\delta _{M0}`$ or $`d_{M0}^1(\pi /2)=M/\sqrt{2}`$, yields $`\alpha _{10}`$ or $`\alpha _{11}`$, respectively, and for the tensor asymmetries ($`I=2`$) one may first choose $`\theta _d=0`$ yielding with $`d_{M0}^2(0)=\delta _{M0}`$ directly $`\alpha _{20}`$. The latter being determined, then setting $`\theta _d=\pi /4`$ and $`\pi /2`$, one can obtain the remaining two terms $`\alpha _{21}`$ and $`\alpha _{22}`$. For the separation of $`\alpha _{21}`$ and $`\alpha _{22}`$ one can also choose $`\theta _d=\theta _d^0=\text{arcos}(1/\sqrt{3})`$ together with $`\varphi _{qd}`$ and $`\varphi _{qd}+\pi `$. Then the sum and difference of the corresponding asymmetries result in $`\alpha _{21}`$ and $`\alpha _{22}`$, respectively.
In the next step, in order to separate the two contributions $`c_{IM}`$ and $`s_{IM}`$ in (A20), one can take first $`\varphi _{qd}=0`$ giving $`c_{IM}`$ and then $`\varphi _{qd}=\pi /(2M)`$ for $`M0`$ which yields directly $`s_{IM}`$. The remaining separation of the coefficients $`a_n`$ or $`b_n`$ in (A21) is then achieved by appropriate choices of $`\varphi _q`$. This completes the separation.
## Appendix B A modified impulse approximation
For incoherent $`\pi ^0`$-photoproduction in IA, the $`NN`$-final state is described by a plane wave $`|\stackrel{}{p},sm_s`$ which is not orthogonal to the deuteron wave function. Thus the IA matrix element contains a spurious contribution from coherent $`\pi ^0`$-photoproduction, the size of which is governed by the overlap between the plane wave and the deuteron wave function which is just the deuteron wave function in momentum space (see Eq. (6))
$$\varphi _{m_sm_d}(\stackrel{}{p})=\stackrel{}{p},1m_s|1m_d^{(d)}0,$$
(B1)
whereas, if the interaction between the nucleons is properly taken into account, the overlap with the final $`NN`$-scattering wave vanishes. Therefore, one can expect that a large fraction of $`NN`$-rescattering effects in incoherent $`\pi ^0`$-photoproduction arises from the elimination of the spurious coherent contribution.
This spurious contribution can be avoided by applying a modified IA, in which one uses a modified $`NN`$-final state wave function, where the deuteron wave function component is projected out by the replacement
$$|\stackrel{}{p},1m_s|\stackrel{}{p},1m_s\underset{m_d}{}|1m_d^{(d)}1m_d|\stackrel{}{p},1m_s=|\stackrel{}{p},1m_s\underset{m_d}{}|1m_d^{(d)}\varphi _{m_sm_d}^{}(\stackrel{}{p}).$$
(B2)
This means the following replacement for the IA matrix element
$$\stackrel{}{p},1m_s|T|1m_d^{(d)}\stackrel{}{p},1m_s|T|1m_d^{(d)}\underset{m_d^{}}{}\varphi _{m_sm_d^{}}(\stackrel{}{p})^{(d)}1m_d^{}|T|1m_d^{(d)}.$$
(B3)
The matrix element $`{}_{}{}^{(d)}1m_d^{}|T|1m_d_{}^{(d)}`$ corresponds to the one for coherent pion photoproduction $`\gamma d\pi ^0d`$ in the off-shell region. The comparison of the original IA with the modified IA reveals then what fraction of the whole FSI-effect arises from the non-orthogonality.
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# A nested embedding theorem for Hardy-Lorentz spaces with applications to coefficient multiplier problems
## 1 Introduction
In this paper we characterize coefficient multipliers between certain types of analytic function spaces on the open unit disk. We are primarily concerned with multipliers having one of the non-locally convex Hardy-Lorentz spaces $`H^{p,q},`$ $`0<p<1,`$ $`0<q\mathrm{},`$ for the domain space. For such multipliers we will consider a variety of target spaces including Lebesgue sequence spaces, other Hardy spaces, and various analytic function spaces of mixed norm type. Our method depends upon a nested embedding theorem for Hardy-Lorentz spaces (Theorem 4.1) obtained through interpolation from embedding theorems of Hardy and Littlewood and of Flett. Thus, the strategy is to trap $`H^{p,q}`$ between a pair of mixed norm spaces of Bergman-type and then deduce multiplier results for $`H^{p,q}`$ from corresponding known multiplier results for the endpoint spaces. The paper is organized as follows. In Section 2 we define the Hardy-Lorentz spaces and the analytic mixed norm spaces. Also included in this section are some results from interpolation, fractional calculus, and $`H^p`$\- theory needed for the sequel. Our primary references for Lorentz spaces and Hardy spaces are and , respectively. Section 3 covers preliminary material on coefficient multipliers. In Section 4 we state and prove the embedding theorem for $`H^{p,q}`$. We then indicate how this theorem may be used to obtain the duality results of . In addition, we determine the Abel dual of $`H^{p,\mathrm{}}`$. In Sections 5, 6, and 7, respectively, we find multipliers of $`H^{p,q}`$ into the Lebesgue sequence spaces $`\mathrm{}^s,`$ into mixed norm spaces of Bergman-type, and into certain Hardy spaces. In Section 8 we discuss the case when the target space is an analytic Lipschitz or Zygmund space, a Bloch space, or BMOA.
Throughout the paper $`𝔻`$ will denote the open unit disk in the complex plane and $`𝕋`$ will denote its boundary. The symbol $`H(𝔻)`$ is used to denote the space of analytic functions on $`𝔻`$. If $`X`$ and $`Y`$ are toplogical spaces with $`XY`$ we write $`XY`$ to indicate continuous inclusion. All vector spaces are assumed to be complex. By a Frechet space we mean a locally convex F-space. If $`E`$ is a topological vector space then $`E^{}`$ denotes the topological dual space of $`E`$ consisting of all continuous linear functionals on $`E`$. The symbol $`AB`$ is used to indicate the existence of absolute positive constants $`C_j,j=1,2,`$ such that $`C_1A/BC_2`$.
## 2 Hardy-Lorentz Spaces and Mixed Norm <br>Spaces
Let $`m`$ denote normalized Lebesgue measure on $`𝕋`$ and let $`L^0(m)`$ be the space of complex-valued Lebesgue measurable functions on $`𝕋`$. For $`fL^0(m)`$ and $`s0`$, we write $`\lambda _f(s)=m(\{z𝕋:|f(z)|>s\}`$ for the distribution function and $`f^{}(s)=inf(\{t0:\lambda _f(t)s\}`$ for the decreasing rearrangement of $`|f|`$, each taken with respect to $`m`$. Let $`0<p,q\mathrm{}`$. For the reader’s convenience, we recall the definition of the Lorentz spaces $`L^{p,q}(m)`$. The Lorentz functional $`||||_{p,q}`$ is defined at $`fL^0(m)`$ by $`f_{p,q}=(_0^1[f^{}(s)s^{\frac{1}{p}}]^q\frac{ds}{s})^{1/q}`$ for $`0<q<\mathrm{}`$ and $`f_{p,\mathrm{}}=sup_{s0}[f^{}(s)s^{\frac{1}{p}}]`$. The corresponding Lorentz space is $`L^{p,q}(m)=\{fL^0(m):f_{p,q}<\mathrm{}\}`$. Since $`f_{p,p}=f_p,`$ where $`f_p`$ denotes the standard $`L^p`$-functional on $`L^0(m)`$, the Lorentz spaces form a 2-parameter array $`\{(L^{p,q}(m),||||_{p,q})\}_{0<p,q\mathrm{}}`$ of quasi-Banach spaces containing the Lebesgue space scale $`\{(L^p(m),f_p)\}_{0<p\mathrm{}}`$ as the main diagonal. Inclusions among the Lorentz spaces are given by
$$L^{p,q}(m)L^{p,r}(m),0<p\mathrm{},0<qr\mathrm{},$$
( 2.1)
and, since $`m(𝕋)<\mathrm{}`$,
$$L^{r,s}(m)L^{p,q}(m),0<p<r\mathrm{},0<q,s\mathrm{}.$$
( 2.2)
The space $`L^{p,q}(m)`$ is separable if and only if $`q\mathrm{}`$. The class of functions $`fL^0(m)`$ satisfying $`lim_{s0}[f^{}(s)s^{\frac{1}{p}}]=0`$ is a separable closed subspace of $`L^{p,\mathrm{}}(m)`$ which is denoted by $`L_0^{p,q}(m)`$. We observe here that for $`q\mathrm{},`$ the space $`L^{\mathrm{},q}(m)=0`$. In the sequel we will follow the convention that in all discussions concerning the space $`L^{p,q}(m)`$ it is assumed that $`q=\mathrm{}`$ whenever $`p=\mathrm{}`$.
For $`w𝔻,`$ and $`fH(𝔻),`$ the function $`f_w`$ is defined on $`|z|<1/|w|`$ by$`f_w(z)=f(wz)`$. The space of continuous complex-valued functions on $`𝕋`$ will be denoted by $`C(𝕋)`$. The function $`f_w`$ is considered as both an analytic function on the disk $`|z|<1/|w|,`$ and as a function in $`C(𝕋)`$. For $`0<r<1,`$ the functions $`f_r`$ are called the dilations of $`f`$. Recall that the means $`M_p(r,f)`$ are defined in the usual way by $`M_p(r,f)=(_𝕋|f(r(z)|^pdm(z))^{1/p}`$, $`0<p<\mathrm{}`$ and $`M_{\mathrm{}}(r,f)=sup_{z𝕋}|f(rz)|`$. The Hardy space $`H^p`$ is defined as $`H^p=\{fH(𝔻):f_{H^p}<\mathrm{}\}`$, where $`f_{H^p}=sup_{0<r<1}M_p(r,f).`$ The Nevanlinna class $`N`$ is the subclass of functions $`fH(𝔻)`$ for which $`sup_{0<r<1}_𝕋\mathrm{log}^+|f_r(z)|dm(z)<\mathrm{}.`$ Functions in $`N`$ are known to have non-tangential limits $`m\text{-a.e. on}𝕋`$. Consequently every $`fH(𝔻)`$ determines a boundary value function which we also denote by $`f`$. Thus $`f(z)=lim_{r1^{}}f_r(z)`$, $`m`$-a.a. $`z𝕋.`$ The Smirnov class $`N^+`$ is the subclass of $`N`$ consisting of those functions $`f`$ for which $`lim_{r1^{}}_𝕋\mathrm{log}^+|f_r(z)|dm(z)=_𝕋\mathrm{log}^+|f(z)|dm(z).`$
It follows from standard $`H^p`$-theory that a function $`fH(𝔻)`$ belongs to $`H^p`$ if and only if $`fN^+`$ with boundary value function in $`L^p(m),`$ in which case $`f_{H^p}=f_p,`$ . Motivated by this characterization of $`H^p`$ we define the Hardy-Lorentz space $`H^{p,q}`$, $`0<p,q\mathrm{}`$ to be the space of functions $`fN^+`$ with boundary value function in $`L^{p,q}(m)`$ and we put $`f_{H^{p,q}}=f_{p,q}`$. Then $`\{(H^{p,q},||||_{H^{p,q}})\}_{0<p,q\mathrm{}}`$ is an array of quasi-Banach spaces of analytic functions on $`𝔻`$ with the standard Hardy space scale as the main diagonal. As with $`L^{p,q}(m),`$ $`H^{p,q}`$ is separable if and only if $`q\mathrm{}`$. The functions in $`H^{p,\mathrm{}}`$ with boundary value function in $`L_0^{p,\mathrm{}}(m)`$ form a closed separable subspace of $`H^{p,\mathrm{}}`$ which is denoted by $`H_0^{p,\mathrm{}}`$. Analogs of the inclusion relations (2.1) and (2.2) hold for the Hardy-Lorentz spaces and $`H^{p,q}H_0^{p,\mathrm{}}`$ for all $`q\mathrm{}`$. The polynomials are dense in $`H_0^{p,\mathrm{}}`$ and in $`H^{p,q},`$ $`q\mathrm{}.`$ For $`fH^{p,q},`$ $`q\mathrm{},`$ the dilations $`f_rf`$ in $`H^{p,q}`$ as $`r1^{}`$. If $`fH^{p,\mathrm{}}`$ then the dilations $`f_rf`$ in $`H^{p,\mathrm{}}`$ as $`r1^{}`$ if and only if $`fH_0^{p,\mathrm{}}`$. We note that a similar statement can be made for the disk algebra $`A(𝔻)`$ which is defined as the subspace of $`H^{\mathrm{}}`$ with boundary function in $`C(𝕋)`$. That is, the polynomials are dense in $`A(𝔻)`$ and the dilations $`f_rf`$ in $`H^{\mathrm{}}`$ as $`r1^{}`$ if and only if $`fA(𝔻)`$, .
An important result in the theory of Lorentz spaces is the identification of these spaces with the intermediate spaces arising in the real interpolation theory of the Lebesgue spaces. An analytic analog of this result is given in Theorem 2.1 below. Theorem 2.1 is one of two interpolation theorems needed for the sequel. It was proved in but omitted the endpoint case corresponding to $`H^{\mathrm{}}.`$ The complete version was proved in , see also .
###### Theorem 2.1.
Let $`0<\theta <1`$ and for $`j=0,1`$ let $`0<p_j,q_j\mathrm{}`$ with $`p_0p_1`$.
(i) Set $`\frac{1}{p}=\frac{1\theta }{p_0}+\frac{\theta }{p_1}`$. Then for every $`0<q\mathrm{}`$ we have, with equivalent quasinorms,
$$(H^{p_0,q_0},H^{p_1,q_1})_{\theta ,q}=H^{p,q}.$$
(ii) Set $`\frac{1}{q}=\frac{1\theta }{q_0}+\frac{\theta }{q_1}`$. Then for every $`0<p<\mathrm{}`$ we have, with equivalent quasinorms,
$$(H^{p,q_0},H^{p,q_1})_{\theta ,q}=H^{p,q}.$$
The second collection of domain spaces $`EH(𝔻)`$ that we define are spaces of mixed “ norm ” type. Before introducing these spaces, we describe the fractional calculus that we will be using. Let $`0<\beta <\mathrm{}`$ and suppose $`fH(𝔻)`$ with Taylor series representation
$$f(z)=\underset{n=0}{\overset{\mathrm{}}{}}a_nz^n,z𝔻$$
( 2.3)
The fractional derivative and fractional integral of $`f`$ of order $`\beta `$ are the functions respectively defined at $`z𝔻`$ by $`f^{[\beta ]}(z)=_{n=0}^{\mathrm{}}\frac{\mathrm{\Gamma }(n+\beta +1)}{n!}z^n`$ and $`f_{[\beta ]}(z)=_{n=0}^{\mathrm{}}\frac{n!}{\mathrm{\Gamma }(n+\beta +1)}z^n`$, where $`\mathrm{\Gamma }`$ is the gamma function. The symbols $`D^\beta `$ and $`D_\beta `$ stand for the associated operators defined on $`H(𝔻)`$ by $`D^\beta (f)=f^{[\beta ]}`$ and $`D_\beta (f)=f_{[\beta ]}`$, $`fH(𝔻)`$. We adopt the convention that for $`\mathrm{}<\beta <0`$, $`f^{[\beta ]}=f_{[\beta ]}`$, $`f_{[\beta ]}=f^{[\beta ]}`$ and similarly for $`D^\beta `$ and $`D_\beta `$. The operators $`D^0`$ and $`D_0`$ are understood to be the identity on $`H(𝔻)`$ and $`f^{[0]}=f_{[0]}=f`$.
Suppose then that $`0<p,q\mathrm{},0<\alpha <\mathrm{}`$ and let $`fH(𝔻)`$. We set $`f_{H(p,q,\alpha )}=(_0^1M_p(r,f)^q(1r)^{q\alpha 1}𝑑r)^{1/q},q\mathrm{}`$ and $`||f||_{H(p,\mathrm{},\alpha )}=sup_{0<r<1}[M_p(r,f)(1r)^\alpha `$. Then the weighted mixed Bergman space$`H(p,q,\alpha )`$ is defined as $`H(p,q,\alpha )=\{fH(𝔻):f_{H(p,q,\alpha )}<\mathrm{}\}`$. We also define $`H_0(p,\mathrm{},\alpha )`$ to be the subspace of functions $`fH(p,\mathrm{},\alpha )`$ satisfying $`M_p(r,f)(1r)^\alpha 0`$ as $`r1^{}`$. In this notation, $`H(p,p,1/p)`$ is the standard Bergman space $`A^p=\{fH(𝔻):_𝔻|f(z)|^p𝑑\nu (z)<\mathrm{}\}`$, where $`\nu `$ is Lebesgue measure on $`𝔻`$. For $`\mathrm{}<\beta <\mathrm{}`$, we set $`f_{H(p,q,\alpha ,\beta )}=f^{[\beta ]}_{H(p,q,\alpha )}`$. The weighted mixed Bergman-Sobolev space $`H(p,q,\alpha ,\beta )`$ is then defined as $`H(p,q,\alpha ,\beta )=\{fH(𝔻):f_{H(p,q,\alpha ,\beta )}<\mathrm{}\}`$. Similarly $`H_0(p,\mathrm{},\alpha ,\beta )`$ is the subspace of $`H(p,\mathrm{},\alpha ,\beta )`$ consisting of those functions $`f`$ satisfying $`f^{[\beta ]}H_0(p,\mathrm{},\alpha )`$. The spaces $`(H(p,q,\alpha ),||||_{H(p,q,\alpha )})`$ and $`(H(p,q,\alpha ,0),||||_{H(p,q,\alpha ,0)})`$ are of course identical and we will continue to use the former notation when $`\beta =0`$. As with $`H^{p,q}`$, $`H(p,q,\alpha ,\beta )`$ is separable if and only if $`q\mathrm{}`$. The space $`H_0(p,\mathrm{},\alpha ,\beta )`$ is a closed separable subspace of $`H(p,\mathrm{},\alpha ,\beta )`$. The polynomials are dense in $`H(p,q,\alpha ,\beta ),q\mathrm{}`$ and in $`H_0(p,\mathrm{},\alpha ,\beta )`$. If $`fH(p,q,\alpha ,\beta ),q\mathrm{}`$, then $`f_rf`$ in $`H(p,q,\alpha ,\beta )`$ as $`r1^{}`$. On the other hand, for $`fH(p,\mathrm{},\alpha ,\beta )`$, $`f_rf`$ in $`H(p,\mathrm{},\alpha ,\beta )`$ as $`r1^{}`$ if and only if $`fH_0(p,\mathrm{},\alpha ,\beta )`$, see , . The spaces $`H(p,q,\alpha ,\beta )`$ are often called mixed norm spaces. In the sequel we will simply say $`H(p,q,\alpha ,\beta )`$ is a Bergman-Sobolev spaces and $`H(p,q,\alpha )`$ is a Bergman space. Many authors use an equivalent definition of these spaces obtained by replacing the fractional calculus operators $`D^\beta `$ and $`D_\beta ,0\beta <\mathrm{}`$, in the definition of $`H(p,q,\alpha ,\beta )`$ with the multiplier operators $`J^\beta `$ and $`J_\beta `$ defined at a function $`fH(𝔻)`$ with Taylor series representation (2.3) by $`J^\beta (f)(z)=_{n=0}^{\mathrm{}}(n+1)^\beta a_nz^n\text{ and }J_\beta (f)(z)=_{n=0}^{\mathrm{}}(n+1)^\beta a_nz^n,z𝔻`$. Then the mixed norm space obtained using $`J^\beta `$ or $`J_\beta `$ is identical to the space $`H(p,q,\alpha ,\beta )`$ as previously defined. For the equivalence of the fractional calculus operators with the multiplier operators in defining the spaces $`H(p,q,\alpha ,\beta )`$ as well as proofs of the following results, the reader is referred to , , and .
###### Lemma 2.1.
Let $`0<p,q\mathrm{}`$, $`0<\alpha ,\beta <\mathrm{}`$. Then the following mappings are continuous surjective isomorphisms.
(i) $`D^\beta :H(p,q,\alpha )H(p,q,\alpha +\beta ),`$
(ii) $`D^\beta :H_0(p,\mathrm{},\alpha )H_0(p,\mathrm{},\alpha +\beta ),`$
(iii) $`D_\beta :H(p,q,\alpha )H(p,q,\alpha \beta )\text{ for }\beta <\alpha `$,
(iv) $`D_\beta :H_0(p,\mathrm{},\alpha )H_0(p,\mathrm{},\alpha \beta )\text{ for }\beta <\alpha `$.
Lemma 2.1 and the definition of $`H(p,q,\alpha ,\beta )`$ imply the following.
###### Lemma 2.2.
Let $`0<p,q\mathrm{},`$ $`0<\alpha <\mathrm{},\mathrm{}<\beta <\mathrm{}`$. Then for $`\mathrm{}<\gamma <\alpha ,`$ the following identifications hold with equivalent quasinorms.
(i) $`H(p,q,\alpha ,\beta )=H(p,q,\alpha \gamma ,\beta \gamma ),`$
(ii) $`H_0(p,\mathrm{},\alpha ,\beta )=H_0(p,\mathrm{},\alpha \gamma ,\beta \gamma ).`$
In particular, Lemma 2.2 implies $`H(p,q,\alpha ,\beta )=H(p,q,\alpha \beta )`$ and $`H_0(p,\mathrm{},\alpha ,\beta )=H_0(p,\mathrm{},\alpha \beta )`$ if $`\mathrm{}<\beta <\alpha `$. Lemma 2.3 which follows represents an extension of Lemma 2.1 to the spaces $`H(p,q,\alpha ,\beta )`$.
###### Lemma 2.3.
Let $`0<p,q\mathrm{}`$, $`0<\alpha <\mathrm{},\mathrm{}<\beta ,\gamma <\mathrm{}`$. Then the following mappings are continuous surjective isomorphisms.
(i) $`D^\gamma :H(p,q,\alpha ,\beta )H(p,q,\alpha +\gamma ,\beta )`$ for $`\gamma >\alpha ,`$
(ii) $`D^\gamma :H_0(p,\mathrm{},\alpha ,\beta )H_0(p,\mathrm{},\alpha +\gamma ,\beta )`$ for $`\gamma >\alpha `$,
(iii) $`D^\gamma :H(p,q,\alpha ,\beta )H(p,q,\alpha ,\beta \gamma ),`$
(iv) $`D^\gamma :H_0(p,\mathrm{},\alpha ,\beta )H_0(p,\mathrm{},\alpha ,\beta \gamma ).`$
The second interpolation result we need is for the Bergman-Sobolev spaces and is due to Fabrega and Ortega, see .
###### Theorem 2.2.
Let $`0<p<\mathrm{},0<q_j\mathrm{},0<\alpha _j,\beta <\mathrm{},j=0,1\text{ and suppose }\alpha _0\alpha _1`$. Let $`0<\theta <1,0<q\mathrm{}`$ and set $`\alpha =(1\theta )\alpha _0+\theta \alpha _1`$. Then we have, with equivalent quasinorms,
$$(H(p,q_0,\alpha _0,\beta ),H(p,q_1,\alpha _1,\beta ))_{\theta ,q}=H(p,q,\alpha ,\beta ).$$
We also need the following embedding theorems. The first of these is due to Flett and indicates how Bergman-Sobolev spaces embed in the standard Hardy spaces. The second result is a well-known embedding theorem of Hardy and Littlewood, see .
###### Theorem 2.3.
Let $`0<p<s<\mathrm{}`$, $`0<qs`$, $`\beta >1/p1/s`$. Then
$$H(p,q,\beta +1/s1/p,\beta )H^s.$$
###### Theorem 2.4.
Let $`0<p<s\mathrm{}`$, $`pt\mathrm{}`$. Then
$$H^pH(s,t,1/p1/s).$$
## 3 Multipliers
Let $`_0`$ denote the set of nonnegative integers and let $`W`$ denote the space of complex sequences indexed by $`_0`$. We always consider $`W`$ as being equipped with the topology of pointwise convergence. With this topology, $`W`$ is a Frechet space. A topological vector space $`X`$ satisfying $`XW`$ is called a K-space. An FK-space is a K-space which is also an F-space. In particular, spaces which are both K-spaces and Frechet spaces are FK-spaces. $`W`$ is also a topological algebra under the natural product of coordinate-wise multiplication. Thus, for $`w=\{w_n\},\lambda =\{\lambda _n\}`$, the product $`\lambda w`$ is defined by $`\lambda w=\{\lambda _nw_n\}`$. It will sometimes be convenient to use the symbol $`B`$ for the product map so that $`B(\lambda ,w)=\lambda w`$. Then $`B:W\times WW`$ is a continuous bilinear operator. For fixed $`\lambda W`$, we will write $`B_\lambda `$ for the continuous linear operator $`B_\lambda :WW`$ defined by $`B_\lambda (w)=\lambda w`$, $`wW`$. Suppose now that $`E`$ and $`X`$ are a pair of vector subspaces of $`W`$. An element $`\lambda W`$ is said to be a multiplier of $`E`$ into $`X`$ if $`\lambda wX`$ for every $`wE`$. The set of multipliers from $`E`$ into $`X`$ is denoted by either of the symbols $`(E,X)`$ or $`E^X`$. Thus $`\lambda (E,X)`$ if and only if the linear operator $`B_\lambda `$ maps $`E`$ into $`X`$. (Consequently, the bilinearity of $`B`$ gives $`(E,X)=_{wE}(B_w^1(X)`$). If $`E`$ and $`X`$ are FK-spaces, an argument based on the Closed Graph Theorem shows that $`(E,X)`$, or more precisely $`\{B_\lambda :\lambda (E,X)\}`$, is a subspace of $``$$`(E,X)`$, the space of continuous $`X`$-valued linear operators on $`E`$. The space $`(E,X)`$ is sometimes called the $`X`$-dual of E. The second $`X`$-dual of $`E`$ is the space $`E^{XX}=(E^X)^X`$. If $`E^{XX}=E`$ then E is said to be $`X`$-reflexive or $`X`$-perfect. We record some of the basic properties of multiplier spaces in the form of a lemma. We omit the obvious proof.
###### Lemma 3.1.
Let $`A`$, $`B`$, $`C`$, $`E`$ be vector subspaces of $`W`$ with $`AB`$. Then
(i) $`B^CA^C,`$
(ii) $`C^AC^B,`$
(iii) $`(A,C)(C^E,A^E).`$
For quasi-Banach spaces $`(E,||||_E)`$, $`(X,||||_X)`$ $`W`$, the operator quasinorm is defined at an operator $`L`$$`(E,X)`$ in the standard way by $`L_{(E,X)}`$$`=sup\{||L(w)||_X:wE,||w||_E1\}`$. Then $`(`$$`(E,X),`$ $`||||_{(E,X)})`$ is a quasi-Banach space containing $`(E,X)`$ as a closed subspace. In particular, $`(E,X)`$ is a quasi-Banach space under the quasinorm $`||||_{(E,X)}`$ defined at $`\lambda (E,X)`$ by $`\lambda _{(E,X)}=B_\lambda _{(E,X)}`$.
We regard $`H(𝔻)`$ as a subspace of $`W`$ by identifying functions in $`H(𝔻)`$ with their Taylor coefficient sequences. Thus, a function $`fH(𝔻)`$ with Taylor series representation (2.3) is identified with the sequence $`a=\{a_n\}`$. $`H(𝔻)`$ is a Frechet space when equipped with the topology of uniform convergence on compact subsets of $`𝔻`$. In addition, we note that $`H(𝔻)`$ is a K-space and hence any F-space $`E`$ satisfying $`EH(𝔻)`$ is a FK-space. In it is shown that the product map $`B`$ on $`W\times W`$ restricts to a bilinear operator $`H(𝔻)\times H(𝔻)H(𝔻)`$. The symbol $`c`$ will denote the Cauchy function defined by $`c(z)=(1z)^1,z𝔻`$. Since $`cH(𝔻),`$ it follows that $`(H(𝔻),H(𝔻))=H(𝔻)`$. If $`f`$, $`g`$ $`H(𝔻)`$, then $`B(f,g)`$ is commonly denoted by $`fg`$ and is called the Hadamard product of $`f`$ and $`g`$. Thus, if $`f`$, $`g`$ $`H(𝔻)`$, with Taylor series representations $`f(z)=_{n=0}^{\mathrm{}}a_nz^n`$ and $`g(z)=_{n=0}^{\mathrm{}}b_nz^n`$, $`z𝔻`$, then $`B(f,g)=fgH(𝔻)`$ has Taylor series representation $`B(f,g)(z)=(fg)(z)=_{n=0}^{\mathrm{}}a_nb_nz^n,z𝔻`$ and, as a sequence, $`B(f,g)=fg=\{a_nb_n\}`$.
We introduce some notation. For $`n_0`$ and $`z𝔻`$ we set $`u_n(z)=z^n`$. Let $`fH(𝔻)`$ with Taylor series representation given by (2.3). For $`N_0`$ we write $`S_N(f)`$ for the partial sum function $`S_N(f)(z)=_{n=0}^Na_nz^n,z𝔻.`$ In the sequel we will be interested in the multiplier spaces $`(E,X)`$ where E is a Hardy-Lorentz space and $`X`$ is a FK-space, $`XH(𝔻)`$. At this point we would like to consider $`(E,X)`$ for some specific target spaces $`X`$ and for a fairly general class of domain spaces $`E`$. One of our choices for $`X`$ is the space $`AS(𝔻)`$ of Abel summable sequences. Recall that the element $`w=\{w_n\}W`$ is said to be Abel summable if $`lim_{r1^{}}_{n=0}^{\mathrm{}}w_nr^n`$ exists. The space $`AS(𝔻)`$ is a Frechet space , with respect to the topology induced by the family $`\{\rho _n:n_0\text{ or }n=\mathrm{}\}`$ of seminorms, where $`\{r_n\}`$ is a fixed sequence in $`(0,1)`$ strictly increasing to $`\mathrm{}`$, and for $`n_0`$ and $`w=\{w_n\}AS(𝔻)`$, $`\rho _{\mathrm{}}(w)=sup_{0<r<1}|_{k=0}^{\mathrm{}}w_kr^k|\text{ and }\rho _n(w)=_{k=0}^{\mathrm{}}|w_k|r_n^k.`$ Furthermore, we have $`AS(𝔻)H(𝔻)`$ and hence $`AS(𝔻)`$ is a FK-space. The $`AS(𝔻)`$-dual of a vector subspace $`E`$ of $`W`$ is known as the Abel dual of $`E`$ and will be denoted by $`E^a`$. The second Abel dual of $`E`$ is the space $`(E^a)^a`$ and will be denoted by $`E^{aa}`$. If $`E^{aa}=E`$, then $`E`$ is said to be Abel reflexive. Note that the functional $`\mathrm{\Psi }`$ on $`AS(𝔻)`$ defined by
$$\mathrm{\Psi }(a)=\underset{r1^{}}{lim}\underset{n=0}{\overset{\mathrm{}}{}}a_nr^n,a=\{a_n\}AS(𝔻)$$
( 3.1)
belongs to $`AS(𝔻)^{}.`$
Propositions 3.2 through 3.4 below describe the relationship between the spaces $`E^{}`$, $`(E,H^{\mathrm{}})`$, and $`E^a`$ for a certain general type of FK-space $`E`$. First we need the following.
###### Proposistion 3.1.
Suppose that $`E`$ is a FK-space satisfying
(i) $`EH(𝔻),`$
(ii) $`u_nE\text{ for all }n_0,`$
(iii) $`c_wE\text{ for all }w𝔻,`$
(iv) For all $`w𝔻,`$ $`\{S_N(c_w)\}\text{ converges to }c_w\text{ in }E\text{ as }N\mathrm{}.`$
Then for any topological vector space $`X`$, operator $`T`$$`(E,X)`$, and $`w𝔻`$, the $`X`$-valued series
$$\underset{n=0}{\overset{\mathrm{}}{}}x_nw^n,x_n=T(u_n),n_0,$$
( 3.2)
converges in $`X`$ to $`T(c_w)`$.
###### Proof.
The partial sums of the series (3.2) satisfy
$$\underset{n=0}{\overset{N}{}}x_nw^n=T(S_N(c_w)).$$
( 3.3)
Since T is continuous the lemma follows from (iv) and (3.3).
Suppose then that $`E`$ is a FK-space satisfying (i) through (iv) of Proposition 3.1, that $`X`$ is a topological vector space and $`T`$$`(E,X)`$. Define a function $`g_T`$ on $`𝔻`$ by $`g_T(w)=T(c_w),w𝔻`$. Then $`g_T`$ is a well-defined $`X`$-valued function on $`𝔻`$ with power series representation given by (3.2). The function $`g_T`$ is called the analytic or Cauchy transform of $`T`$. The linear operator $`Tg_T,`$ taking $``$($`E,X)`$ into the space of $`X`$-valued power series on $`𝔻,`$ will also be referred to as the analytic or Cauchy transform on $``$($`E,X)`$. Consider now the case when $`X`$ is a FK-space satisfying $`XH(𝔻)`$ and $`T=B_\lambda `$ for some multiplier $`\lambda (E,X)`$. For this case, if $`\lambda =\{\lambda _n\}`$ as a sequence of complex numbers then $`T(u_n)=B_\lambda (u_n)=\lambda u_n`$. Since $`\lambda u_n`$ is identified with the sequence having $`\lambda _n`$ in the $`n`$-th entry and 0 elsewhere, we will write $`\lambda _n`$ in place of $`T(u_n),`$ $`n_0,`$ and $`g_\lambda `$ in place of $`g_T`$. The other situation we will be considering is when $`X`$ is the complex field and $`T=\phi E^{}`$. For this case we may form the sequence $`\lambda =\{\lambda _n\},`$ where $`\lambda _n=\phi (u_n),n_0`$ and we again write $`g_\lambda `$ in place of $`g_T`$. Proposition 3.1 ensures that $`g_\lambda H(𝔻)`$. But, as was previously noted, $`H(𝔻)=(H(𝔻),H(𝔻))`$. Since $`(H(𝔻),H(𝔻))(E,H(𝔻))`$ by Lemma 3.1(i), it follows that $`\phi E^{}`$ induces the multiplier $`\lambda (E,H(𝔻))`$ with analytic transform $`g_\lambda `$. In fact, by convention, we have the identification $`\lambda g_\lambda `$. Let us note here that it is possible to have $`\lambda =0`$ for $`\phi 0,`$ so that in general the analytic transform on $``$$`(E,X)`$ is not one-to-one. However we do have the following.
###### Proposistion 3.2.
Suppose that in addition to satisfying conditions (i) through (iv) of Proposition 3.1, the FK-space $`E`$ satisfies
(i) $`f_wE\text{ for each }w𝔻\text{ and }fE,`$
(ii) for each $`fE,\text{ the set }\{f_w:w𝔻\}\text{ is bounded in }E,`$
(iii) $`\{S_N(f_w)\}\text{ converges to }f_w\text{ in }E\text{ for every }w𝔻\text{ and }fE.`$
Let $`\phi E^{}`$ with induced multiplier $`\lambda =\{\lambda _n\},\lambda _n=\phi (u_n)`$ and analytic transform $`g_\lambda `$. Then for each $`w𝔻`$ and $`fE`$,
$$\phi (f_w)=(fg_\lambda )(w).$$
( 3.4)
Consequently, $`\lambda (E,H^{\mathrm{}})`$. Furthermore, the analytic transform on $`E^{}`$ is one-to-one whenever $`E`$ satisfies the additional property
(iv) for every $`fE`$, the dilations $`f_r`$, $`0<r<1`$, converge to $`f`$ in $`E`$ as $`r1^{}.`$
###### Proof.
Let $`\phi E^{}`$ and put $`\lambda =\{\lambda _n\},n_0`$. Then , by the comments made in the paragraph following the proof of Proposition 3.1, we have $`g_\phi =g_\lambda (E,H(𝔻))`$. If $`w𝔻`$ and $`fE`$ has Taylor series representation (2.3) then the continuity of $`\phi `$ and conditions (i) and (iii) of Proposition 3.2 yield $`\phi (f_w)=\phi (lim_n\mathrm{}S_N(f_w))=lim_N\mathrm{}\phi (S_N(f_w)=lim_N\mathrm{}\phi (_{n=0}^Na_nw^nu_n)=lim_N\mathrm{}_{n=0}^Na_n\lambda _nw^n=lim_N\mathrm{}S_N(fg_\lambda )(w)=fg_\lambda (w)`$ which is (3.4). Then (3.4), the continuity of $`\phi `$, and condition (ii) imply $`\lambda (E,H^{\mathrm{}})`$. Finally, suppose $`E`$ satisfies (iv). It then follows from this property and (3.4), that if $`\lambda `$ is the induced multiplier for $`\phi _jE^{},j=1,2`$ we have $`\phi _j(f)=\phi _j(lim_{r1^{}}f_r)=lim_{r1^{}}\phi _j(f_r)=lim_{r1^{}}(fg_\lambda )(r)`$, so that $`\phi _1=\phi _2`$. ∎
In view of the last two propositions, we may regard $`E^{}(E,H^{\mathrm{}})`$ for any FK-space satisfying (i) through (iv) of Propositions 3.1 and 3.2. Note that if $`E`$ is a FK-space with $`A(𝔻)E`$ then $`E`$ satisfies conditions (ii) through (iv) of Proposition 3.1 and conditions (i) through (iii) of Proposition 3.2. Also in it is shown that for the A-spaces, which form a large class of quasi-Banach spaces $`EH(𝔻)`$, condition (iv) of Proposition 3.2 is equivalent to the density of the polynomials in E.
###### Proposistion 3.3.
Suppose $`E`$ is an FK-space. Then $`E^aE^{}`$.
###### Proof.
Let $`\lambda E^a`$ and put $`\phi _\lambda =\mathrm{\Psi }B_\lambda `$, where $`\mathrm{\Psi }`$ is the functional (3.1). The mapping $`\lambda \phi _\lambda `$ is a one-to-one linear operator taking $`E^a`$ into $`E^{}`$ and we may identify $`E^a`$ with its image in $`E^{}`$. ∎
###### Proposistion 3.4.
Suppose $`E`$ is an FK-space satisfying conditions (i) through (iv) of Propositions 3.1 and 3.2. Then $`(E,H^{\mathrm{}})=(E,A(𝔻))E^a`$.
###### Proof.
Let $`\lambda =\{\lambda _n\}(E,H^{\mathrm{}})`$. Let $`fE`$ have Taylor series representation (2.3). By Proposition 3.2(iv) and the continuity of $`B_\lambda `$ we have $`B_\lambda (f_r)B_\lambda (f)`$ in $`H^{\mathrm{}}`$. Hence $`B_\lambda (f)A(𝔻)`$. Therefore, $`(E,H^{\mathrm{}})(E,A(𝔻))`$. By Lemma 3.1(ii) the reverse inclusion holds. Thus $`(E,H^{\mathrm{}})=(E,A(𝔻))`$. Finally, the continuity of $`B_\lambda (f)`$ at 1 gives $`lim_{r1^{}}_{n=0}^{\mathrm{}}a_n\lambda _nr^n=lim_{r1^{}}B_\lambda (f)(r)=B_\lambda (f)(1)`$. Hence $`\lambda E^a`$. ∎
###### Corollary 3.1.
Let $`E`$ be a FK-space satisfying conditions (i) through (iv) of Propositions 3.1 and 3.2. Then $`E^{}=E^a=(E,H^{\mathrm{}})`$.
###### Corollary 3.2.
Let $`0<p,q<\mathrm{}`$. Then
(i) $`(H^{p,q})^{}=(H^{p,q})^a=(H^{p,q},H^{\mathrm{}}),`$
(ii) $`(H_0^{p,\mathrm{}})^{}=(H_0^{p,\mathrm{}})^a=(H_0^{p,\mathrm{}},H^{\mathrm{}}).`$
## 4 A Nested Embedding Theorem For Hardy-Lorentz Spaces
Our main tool for identifying certain multiplier spaces $`(H^{p,q},X)`$ is the following nested embedding theorem for $`H^{p,q}`$. Its proof consists of using Theorems 2.1 and 2.2 to interpolate Theorems 2.3 and 2.4.
###### Theorem 4.1.
Let $`0<p_0<p<s\mathrm{},\mathrm{\hspace{0.17em}0}<qt\mathrm{}`$ and $`\beta >1/p_01/p`$. Then
(i) $`H(p_0,q,\beta +1/p1/p_0,\beta )H^{p,q}H(s,t,1/p1/s).`$
(ii) $`H_0(p_0,\mathrm{},\beta +1/p1/p_0,\beta )H_0^{p,\mathrm{}}H_0(s,\mathrm{},1/p1/s).`$
###### Proof.
(i) Choose $`s_j`$, $`j=0,1`$, such that $`0<p_0<s_0<p<s_1<\mathrm{}`$. Let us further stipulate that for the case $`\beta 1/p_0`$ we require that $`0<s_1<\frac{p_0}{1\beta p_0}`$. This ensures that $`\beta >1/p_01/s_1`$. We can then apply Theorem 2.4 to obtain embeddings
$$H(p_0,s_j,\beta +1/s_j1/p_0,\beta )H^{s_j},j=0,1$$
( 4.1)
Interpolation of (4.1) results in the embeddings
$$(H(p_0,s_0,\beta +1/s_01/p_0,\beta ),H(p_0,s_1,\beta +1/s_11/p_0))_{\theta ,q}(H^{s_0},H^{s_1})_{\theta ,q},$$
( 4.2)
for all $`0<\theta <1,0<q\mathrm{}`$. Then the first embedding in (i) follows by choosing $`\theta `$ to satisfy $`\frac{1}{p}=\frac{1\theta }{s_0}+\frac{\theta }{s_1}`$ and applying Theorems 2.1 and 2.2 to (4.2). The proof of the second embedding in (i) is similar and is in . We omit the details.
(ii) Again, we prove only the first embedding in (ii) since the proof of the second embedding in (ii) is similar and is also in . Thus, let
$`fH_0(p_0,\mathrm{},\beta +1/p1/p_0,\beta )`$. Then $`fH^{p,\mathrm{}}`$ by Theorem 4.1(i). In order to show $`fH_0^{p,\mathrm{}}`$, it is enough to show the dilations $`f_r`$ converge to $`f`$ in $`H^{p,\mathrm{}}`$ as $`r1^{}`$. But the functions $`f_r`$ converge to $`f`$ in $`H(p_0,\mathrm{},\beta +1/p1/p_0,\beta )`$ and this fact combined with (i) implies $`f_rf`$ in $`H^{p,\mathrm{}}`$. Hence $`fH_0^{p,\mathrm{}}`$. ∎
Recall that if $`E`$ is a quasi-Banach space with separating dual $`E^{}`$, then there exists a unique Banach space $`Y`$ in which $`E`$ embeds as a dense subspace and for which $`Y^{}=E^{}`$. The space $`Y`$ is called the Banach envelope of $`E`$ and is denoted by $`[E]_1`$. In we identified the Banach envelopes and dual spaces of the spaces $`H^{p,q}`$ and $`H_0^{p,\mathrm{}}`$ for indices in the range $`0<p<1`$, $`0<q<\mathrm{}`$. The specific result was the following.
###### Theorem 4.2.
Let $`0<p<1`$, $`0<q<\mathrm{}`$. Set $`q_{}=max(1,q)`$ and let $`q^{}`$ be the Hölder conjugate of $`q_{}`$, $`1/q_{}+1/q^{}=1`$. Then
(i)$`[H^{p,q}]_1=H(1,q_{},1/p1)\text{ and }(H^{p,q})^{}=H(\mathrm{},q^{},1,1/p),`$
(ii) $`[H_0^{p,\mathrm{}}]_1=H_0(1,\mathrm{},1,1/p1)`$ and $`(H_0^{p,\mathrm{}})^{}=H(\mathrm{},1,1,1/p)`$.
Note that for $`p=q`$, (i) becomes
$$[H^p]_1=H(1,1,1/p1)\text{ and }(H^p)^{}=H(\mathrm{},\mathrm{},1,1/p)$$
( 4.3)
This is the well-known Duren-Romberg-Shields Theorem . The proof of Theorem 4.2(i) given in essentially consisted of two steps. The first step was the establishment of the second embedding in Theorem 4.1(i). The second step was a constructive proof of the embedding
$`(H^{p,q})^{}H(\mathrm{},q^{},1,1/p)`$. The proof of Theorem 4.2(ii) in was carried out in an analogous fashion. A short proof of Theorem 4.2 can be based on Theorem 4.1 and Theorem 4.3 below. Theorem 4.3 is a general Banach envelope-duality theorem for separable Bergman-Sobolev spaces and is due to the efforts of several authors. For statements and proofs of Theorem 4.3 for Bergman spaces in some specific cases the reader is referred to , , , , , , , , , , , , , , , , , , , , , . Pavlovic’s paper contains a very general and complete version of Theorem 4.3 for the case $`\beta =0`$. To obtain the result for Bergman-Sobolev spaces one uses the validity of Theorem 4.3 for Bergman spaces with Lemmas 2.1 through 2.3.
###### Theorem 4.3.
Let $`0<p\mathrm{},0<q,\alpha <\mathrm{},\mathrm{}<\beta <\mathrm{}`$. Let $`p_0=min(1,p),p_1=max(1,p)\text{ and }q_1=max(1,q)`$. Let $`1/p_1+1/p_1^{}=1/q_1+1/q_1^{}=1`$. Then
(i) $`[H(p,q,\alpha ,\beta )]_1=H(p_1,q_1,\alpha +1/p_01,\beta ),`$
(ii) $`(H(p,q,\alpha ,\beta ))^{}=H(p_1^{},q_1^{},1,\alpha \beta +1/p_0),`$
(iii) $`[H_0(p,\mathrm{},\alpha ,\beta )]_1=H_0(p_1,\mathrm{},\alpha +1/p_01,\beta ),`$
(iv) $`(H_0(p,\mathrm{},\alpha ,\beta )^{}=H(p_1^{},1,1,\alpha \beta +1/p_0).`$
For $`0<p\mathrm{},0<q,\alpha <\mathrm{},\mathrm{}<\beta <\mathrm{},`$ the spaces $`H(p,q,\alpha ,\beta )`$ and $`H_0(p,\mathrm{},\alpha ,\beta )`$ satisfy conditions (i) through (iv) of Propositions 3.1 and 3.2. So by Corollary 3.1, $`H(p,q,\alpha ,\beta )^{}=H(p,q,\alpha ,\beta )^a=(H(p,q,\alpha ,\beta ),H^{\mathrm{}})`$ and similarly for $`H_0(p,\mathrm{},\alpha ,\beta )`$. Thus duality and Abel duality coincide for the separable Bergman-Sobolev spaces. More explicitly, say in the case of Theorem 4.3(ii), if $`\mathrm{\Lambda }H(p,q,\alpha ,\beta )^{}`$ and $`fH(p,q,\alpha ,\beta )`$ has Taylor series representation (2.3), then the proof of Theorem 4.3(ii) shows that
$$\mathrm{\Lambda }(f)=\underset{r1^{}}{lim}\underset{n=0}{\overset{\mathrm{}}{}}a_n\lambda _nr^n,\lambda _n=\mathrm{\Lambda }(u_n),n_0.$$
( 4.4)
Furthermore, the analytic transform $`g_\lambda `$ of the sequence $`\lambda =\{\lambda _n\}`$ in (4.4) satisfies $`g_\lambda H(p_1,1,1,\alpha \beta +1/p_0),g_\lambda _{H(p_1,q_1^{},1,\alpha \beta +1/p_0)}\mathrm{\Lambda }_{H(p,q,\alpha ,\beta )^{}}`$. Conversely, if $`gH(p_1,1,1,\alpha \beta +1/p_0)`$ has Taylor coefficient sequence $`\lambda =\{\lambda _n\}`$, then we may define $`\mathrm{\Lambda }_g`$ as in (4.4). The resulting functional $`\mathrm{\Lambda }_g`$ belongs to $`H(p,q,\alpha ,\beta )^{}`$, has analytic transform $`g`$, and satisfies
$`\mathrm{\Lambda }_g_{H(p,q,\alpha ,\beta )^{}}g_{H(p_1,1,1,\alpha \beta +1/p_0)}`$. Theorem 4.3 also implies that the spaces $`H(p,q,\alpha ,\beta ),1p\mathrm{},1<q<\mathrm{}`$ are reflexive with the properties of reflexivity and Abel reflexivity being the same for these spaces.
To see how Theorem 4.2 follows from Theorems 4.1 and 4.3, let $`0<p<1`$, choose $`0<p_0<p`$ and take $`s=1,t=\text{max}(1,q)=q_{}`$ in Theorem 4.1(i) to obtain the nested embedding
$$H(p_0,q,\beta +1/p1/p_0,\beta )H^{p,q}H(1,q_{},1/p1)$$
( 4.5)
Applying the functor $`[]_1`$ to (4.5) we find
$$[H(p_0,q,\beta +1/p1/p_0,\beta )]_1[H^{p,q}]_1[H(1,q_{},1/p1)]_1.$$
Since $`H(1,q_{},1/p1)`$ is a Banach space,
$$[H(1,q_{},1/p1)]_1=H(1,q_{},1/p1).$$
( 4.6)
Using Theorem 4.3, we also find
$$[H(p_,q,\beta +1/p1/p_0,\beta )]_1=H(1,q_{},\beta +1/p1,\beta ).$$
( 4.7)
But the spaces on the right-hand sides of (4.6) and (4.7) are identical by Lemma 2.2(i). This establishes the first equality in Theorem 4.2(i) and hence the second inequality as well via Theorem 4.3(ii). The proof of Theorem 4.2(ii) is similar.
Combining Theorems 4.2 and 4.3 one sees that the Banach envelopes of the spaces $`H^{p,q}`$ are reflexive and Abel reflexive for $`0<p<1<q<\mathrm{}`$. Similarly, from Corollary 3.2 and Theorem 4.2, we deduce that for $`0<p<1,`$
$$(H_0^{p,\mathrm{}})^a=(H_0^{p,\mathrm{}})^{}=H(\mathrm{},1,1,1/p).$$
( 4.8)
It then follows from (4.8), Theorem 4.3(ii), and Lemma 2.2(i),that
$$(H_0^{p,\mathrm{}})^{aa}=H(1,\mathrm{},1/p1)$$
( 4.9)
That we also have $`(H^{p,\mathrm{}})^a=H(\mathrm{},1,1,1/p)`$ is a consequence of the following result of Shi .
###### Lemma 4.1.
Let $`0<\alpha <\mathrm{}`$. Then
(i) $`H(1,\mathrm{},\alpha )^a=H_0(1,\mathrm{},\alpha )^a,`$
(ii) $`H(1,\mathrm{},\alpha )^{aa}=H(1,\mathrm{},\alpha ).`$
###### Corollary 4.1.
Let $`0<p<1`$. Then
(i) $`(H^{p,\mathrm{}})^a=H(\mathrm{},1,1,1/p),`$
(ii) $`(H^{p,\mathrm{}})^{aa}=H(1,\mathrm{},1/p1).`$
###### Proof.
(i) Using Lemma 3.1(i) twice, (4.8), Theorems 4.2, 4.3 and Lemma 4.1(i) we obtain
$`H(1,\mathrm{},1/p1)^a(H^{p,\mathrm{}})^a(H_0^{p,\mathrm{}})^a=(H_0^{p,\mathrm{}})^{}`$ $`=`$ $`H_0(1,\mathrm{},1/p1)^{}`$
$`=`$ $`H_0(1,\mathrm{},1/p1)^a`$
$`=`$ $`H(1,\mathrm{},1/p1)^a`$
Since the endpoint spaces in (4) are the same, (i) follows from (4.8) and (4).
(ii) This follows from (4.9) and (4). ∎
## 5 Multipliers of $`H^{p,q}\text{ and }H_0^{p,\mathrm{}}\text{ into }\mathrm{}^s,0<p<1,`$$`0<q<\mathrm{},0<s\mathrm{}.`$
In this section we determine the multiplier spaces $`(E,X)`$ where $`E`$ is either $`H^{p,q}`$ or $`H_0^{p,\mathrm{}},`$ $`0<p<1,0<q\mathrm{}`$ and $`X`$ is $`\mathrm{}^s,0<s\mathrm{}`$. Here $`\mathrm{}^s`$ is the usual Lebesgue sequence space consisting of $`s`$-summable sequences in $`W`$ when $`s\mathrm{},`$ and bounded sequences in $`W`$ when $`s=\mathrm{}`$. Actually we do a little more. If $`E`$ is either $`H^{p,q}`$ or $`H_0^{p,\mathrm{}},`$ $`0<p<1,0<q<\mathrm{},`$ we find $`(E,X)`$ whenever $`X`$ is $`\mathrm{}^s`$-reflexive for some $`0<s\mathrm{}`$. We also find the multiplier spaces $`(H^{p,\mathrm{}},X)`$ for solid target spaces $`X`$. Recall that a vector subspace $`X`$ of $`W`$ is said to be solid if for every $`x=\{x_n\},y=\{y_n\}W,`$ we have $`yX`$ whenever $`xX`$ and $`|y_n||x_n|,n_0`$. Equivalently, $`X`$ is solid if either of the conditions
$$\mathrm{}^{\mathrm{}}(X,X)\text{ or }(\mathrm{}^{\mathrm{}},X)=X$$
( 5.1)
are satisfied. The notation $`s(X)=(\mathrm{}^{\mathrm{}},X)`$ is commonly used. In general, $`s(X)`$ is the largest solid subspace of $`X`$. We note here that for arbitrary spaces $`E`$ and $`X,`$ the multiplier space $`(E,X)`$ is solid whenever the target space $`X`$ is solid. Consequently $`\mathrm{}^s`$-reflexive spaces are solid. Thus the result for $`H^{p,\mathrm{}}`$ is more general than the corresponding result for $`H_0^{p,\mathrm{}}`$.
The determination of $`(E,X)`$ in these cases and others frequently requires using the analytic transform to identify $`(E,X)`$ with a weighted sequence space. If $`X`$ is a vector subspace of $`W`$ and the element $`wW,`$ we define the weighted space $`X_w=B_w^1(X)=\{yW:wyX\}`$. If $`X`$ is a quasi-Banach space with quasinorm $`||||_X`$ then $`(X_w,||||_{X_w})`$is a quasi-Banach space where $`y_{X_w}=yw_X,`$ $`yX_w`$. For $`\mathrm{}<\alpha <\mathrm{}`$, let $`w_\alpha =\{w_\alpha (n):n_0\}`$ be the power sequence defined by $`w_\alpha (0)=1`$ and $`w_\alpha =n^\alpha `$ for $`n0`$. In this case we will write $`(X_\alpha ,||||_{X_\alpha })`$ in place of $`(X_w,||||_{X_w})`$. The following lemma gives the relationship between $`(E,X)`$ and $`(E_\alpha ,X_\beta )`$. The proof is purely algebraic, .
###### Lemma 5.1.
Let $`E`$ and $`X`$ be vector subspaces of $`W`$ and let $`\mathrm{}<\alpha ,\beta <\mathrm{}`$. Then
$$(E_\alpha ,X_\beta )=(E_{\alpha \beta },X)=(E,X_{\beta \alpha })=(E,X)_{\beta \alpha }.$$
Of special interest to us are the dyadically blocked sequence spaces $`\mathrm{}(p,q)`$ and their weighted analogs. These spaces are defined as follows. Let
$`0<p,q\mathrm{},`$ set $`I_0=\{0\}`$, and for $`n_0,n>0,`$ set $`I_n=_0[2^{n1},2^n)`$. Then $`\mathrm{}(p,q)`$ is the subspace of $`W`$ consisting of elements $`x=\{x_n\}`$ such that $`x_{\mathrm{}(p,q)}<\mathrm{},`$ where $`x_{\mathrm{}(p,q)}=\{x_k\}_{kI_n}_\mathrm{}^p_\mathrm{}^q.`$ For $`\mathrm{}<\alpha <\mathrm{},`$ we write $`(\mathrm{}(p,q,\alpha ),||||_{\mathrm{}(p,q,\alpha )})`$ for the weighted space $`(\mathrm{}(p,q)_\alpha ,||||_{\mathrm{}(p,q)_\alpha })`$. The spaces $`(\mathrm{}(p,q,\alpha ),||||_{\mathrm{}(p,q,\alpha )})`$ are quasi-Banach spaces and $`(\mathrm{}(p,p,0),||||_{\mathrm{}(p,p,0)})=(\mathrm{}^p,||||_\mathrm{}^p)`$, where $`||||_\mathrm{}^p`$ is the standard quasinorm on $`\mathrm{}^p`$. Furthermore, $`\mathrm{}(p,q,\alpha )\mathrm{}_\alpha ^{\mathrm{}}H(𝔻)`$, and hence $`\mathrm{}(p,q,\alpha )`$ is also a FK-space. The multipliers between these spaces are well-known. The following result is due mainly to Kellogg , see also , , , and . Before stating the theorem we introduce notation. Let $`0<q,s\mathrm{}`$. Then $`qs`$ is the extended real number defined by $`qs=s`$ if $`q=\mathrm{},`$ $`qs=\frac{qs}{qs}`$ if $`0<s<q<\mathrm{},`$ and $`qs=\mathrm{}`$ if $`0<qs\mathrm{}`$.
###### Theorem 5.1.
Let $`0<p,q,r,s\mathrm{}`$, $`\mathrm{}<\alpha ,\beta <\mathrm{}`$. Then
$$(\mathrm{}(p,q,\alpha ),\mathrm{}(r,s,\beta ))=\mathrm{}(pr,qs,\beta \alpha ).$$
We remark that a consequence of Theorem 5.1 is that $`\mathrm{}(p,q,\alpha )`$ is $`\mathrm{}^s`$-reflexive for $`\mathrm{}<\alpha <\mathrm{}`$ and $`0<sp,q\mathrm{}`$. The space $`c_0`$, of null sequences, is an example of a space which fails to be $`\mathrm{}^s`$-reflexive for every $`0<s\mathrm{}`$. In addition to Theorem 5.1 we need a lemma. Lemma 5.2 is due to Aleksandrov , but may also be seen to follow from Theorem 4.1. In the lemma the spaces $`H^{p,\mathrm{}}`$ and $`H_0^{p,\mathrm{}}`$ are considered as sequence spaces of Taylor coefficients.
###### Lemma 5.2.
Let $`0<p<1`$. Then
(i) $`H^{p,\mathrm{}}\mathrm{}_{11/p}^{\mathrm{}},`$
(ii) $`H_0^{p,\mathrm{}}(c_0)_{11/p}.`$
###### Theorem 5.2.
Let $`0<p<1`$ and let $`X`$ be a solid FK-space satisfying $`XH(𝔻)`$.Then $`(H^{p,\mathrm{}},X)=X_{1/p1}.`$
###### Proof.
Since $`X`$ is solid, it follows that $`X_{1/p1}`$ is solid. Therefore, using Lemma 5.2, we find
$$X_{1/p1}=(\mathrm{}^{\mathrm{}},X_{1/p1})=(\mathrm{}_{11/p}^{\mathrm{}},X).$$
( 5.2)
Then the inclusion
$$X_{11/p}(H^{p,\mathrm{}},X)$$
( 5.3)
results from (5.2) and Lemma 5.1.
We obtain the reverse inclusion of (5.3) as follows. Let
$$g(z)=(1z)^{\frac{1}{p}},z𝔻$$
( 5.4)
Then for $`0<t<1,`$
$$g^{}(t)t^{\frac{1}{p}}\text{ and }gH^{p,\mathrm{}}.$$
( 5.5)
The Taylor coefficient sequence of $`g`$ is
$$\left\{\frac{\mathrm{\Gamma }(n+1/p)}{\mathrm{\Gamma }(1/p)n!}\right\}.$$
( 5.6)
A well-known consequence of Stirling’s formula is that
$$\left\{\frac{n^{1/p1}n!}{\mathrm{\Gamma }(n+1/p)}\right\}\mathrm{}^{\mathrm{}}.$$
( 5.7)
Therefore we deduce the reverse inclusion of (5.3) from (5.1) and (5.5) through (5.7). ∎
###### Corollary 5.1.
Let $`0<p<1,0<s\mathrm{}`$. Then $`(H^{p,\mathrm{}},\mathrm{}^s)=\mathrm{}_{1/p1}^s`$.
###### Theorem 5.3.
Let $`0<p<1,0<s\mathrm{}`$. Then $`(H_0^{p,\mathrm{}},\mathrm{}^s)=\mathrm{}_{1/p1}^s`$.
###### Proof.
We prove only the case $`s\mathrm{}`$, the other case being similar. Using Corollary 5.1 and Lemma 3.1(i) we obtain
$$\mathrm{}_{1/p1}^s=(H^{p,\mathrm{}},\mathrm{}^s)(H_0^{p,\mathrm{}},\mathrm{}^s).$$
( 5.8)
Next we show the reverse inclusion of (5.8) holds. Fix $`\lambda =\{\lambda _n\}(H_0^{p,\mathrm{}},\mathrm{}^s)`$. Since $`\mathrm{}^s`$ is solid there is no loss of generality in assuming $`\lambda _n0`$, $`n_0`$. Furthermore, the solidity of $`\mathrm{}_{1/p1}^s`$ and (5.7), show that $`\lambda \mathrm{}_{1/p1}^s`$ if and only if $`\{\frac{\mathrm{\Gamma }(n+1/p)\lambda _n}{\mathrm{\Gamma }(1/p)n!}\}\mathrm{}^s`$. Consider the operator $`B_\lambda `$$`(H_0^{p,\mathrm{}},\mathrm{}^s)`$ corresponding to $`\lambda `$. That is, for $`fH_0^{p,\mathrm{}},`$ with Taylor series representation (2.3) we have
$$B_\lambda (f)=a\lambda =\{a_n\lambda _n\},$$
( 5.9)
and
$$B_\lambda _{(H_0^{p,\mathrm{}},\mathrm{}^s)}=\lambda _{(H_0^{p,\mathrm{}},\mathrm{}^s)}<\mathrm{}.$$
( 5.10)
Let $`g`$ be the Cauchy-type function in (5.4). Then for $`0<r,t<1,`$
$$(g_r)^{}(t)C_1(1r)^{\frac{1}{p}}\chi _{[0,1r)}(t)+C_2t^{\frac{1}{p}}\chi _{(1r,1)}(t)$$
( 5.11)
where the constants $`C_j`$, $`j=1,2`$ are independent of r and t, and $`\chi _A`$ denotes the characteristic function of a set $`A`$. Thus (5.5) and (5.11) imply
$$\underset{0r<1}{sup}g_r_{H^{p,\mathrm{}}}Cg_{H^{p,\mathrm{}}}<\mathrm{}.$$
( 5.12)
For $`0r<1`$, put $`\mathrm{\Phi }(r)=B_\lambda (g_{r^{1/s}})_\mathrm{}^s^s`$. It follows from (5.10) and (5.12) that $`\mathrm{\Phi }`$ is bounded on $`[0,1)`$. From (5.6) and (5.9) we see that $`g_{r^{1/s}}=\{\frac{\mathrm{\Gamma }(n+1/p)\lambda _nr^{n/s}}{\mathrm{\Gamma }(1/p)n!}\}`$ and $`\mathrm{\Phi }(r)=\{\frac{\mathrm{\Gamma }(n+1/p)\lambda _nr^{n/s}}{n!\mathrm{\Gamma }(1/p)}\}_\mathrm{}^s^s`$. It therefore follows that $`\mathrm{\Phi }`$ is an increasing function of $`r`$. From these observations we deduce that $`lim_{r1^{}}\mathrm{\Phi }(r)`$ exists, hence the positive sequence $`\{(\frac{\mathrm{\Gamma }(n+1/p)\lambda _n}{n!\mathrm{\Gamma }(1/p)})^s\}`$is Abel summable and consequently belongs to $`\mathrm{}^1`$. But then $`\{\frac{\mathrm{\Gamma }(n+1/p)\lambda _n}{n!\mathrm{\Gamma }(1/p)}\}\mathrm{}^s`$ which is what we needed to show. ∎
The space $`\mathrm{}^s`$ is a rearrangement invariant quasi-Banach function space with $`_0`$ equipped with counting measure as the underlying measure space. Quasi-Banach function spaces are always solid. A rearrangement invariant quasi-Banach function space $`X`$ is called maximal if every quasinorm bounded increasing sequence in $`X`$ is bounded above in $`X`$, see . It is not hard to see that if $`X`$ is a maximal rearrangement invariant quasi-Banach function space then $`(H_0^{p,\mathrm{}},X)=X_{1/p1}`$. Another situation where we have $`(H_0^{p,\mathrm{}},X)=X_{1/p1}`$ is when $`X`$ is $`\mathrm{}^s`$-reflexive. For $`0<s\mathrm{}`$ and an arbitrary space $`X`$ we use the notation $`X^{K(s)}`$ for $`(X,\mathrm{}^s)`$ and $`X^{K(s)K(s)}`$ for $`(X^{K(s)})^{K(s)}`$. Thus the space $`X`$ is $`\mathrm{}^s`$-reflexive if and only if $`X^{K(s)K(s)}=X`$. We have the following generalization of Theorem 5.3.
###### Theorem 5.4.
Let $`0<p<1`$ and let $`X`$ be a FK-space which is $`\mathrm{}^s`$-reflexive for some $`0<s\mathrm{}`$. Then
$`(H_0^{p,\mathrm{}},X)=X_{1/p1}.`$
###### Proof.
Since $`X`$ is $`\mathrm{}^s`$-reflexive, $`X_{1/p1}`$ is solid. Hence Theorem 5.2 and Lemma 3.1(i) produce the inclusion
$$X_{1/p1}=(H^{p,\mathrm{}},X)(H_0^{p,\mathrm{}},X).$$
( 5.13)
To get the reverse inclusion of (5.13) we use Lemma 3.1(ii), the $`\mathrm{}^s`$-reflexivity of $`X,`$ and the identity
$$(\mathrm{}_{1/p1}^s)^{K(s)}=\mathrm{}_{11/p}^{\mathrm{}}.$$
( 5.14)
Then using (5.2), (5.13), Theorem 5.3, Lemma 3.1(ii), Lemma 5.1, and (5.14) we get
$`(H_0^{p,\mathrm{}},X)`$ $``$ $`(X^{K(s)},(H_0^{p,\mathrm{}})^{K(s)})`$
$`=`$ $`(X^{K(s)},\mathrm{}_{1/p1}^s)`$
$``$ $`(\mathrm{}_{11/p}^{\mathrm{}},X^{K(s)K(s)})`$
$`=`$ $`(\mathrm{}_{11/p}^{\mathrm{}},X)`$
$`=`$ $`X_{1/p1}`$
$``$ $`(H_0^{p,\mathrm{}},X).`$
and the proof is complete. ∎
In Theorem 5.4, the hypothesis that $`X`$ be $`\mathrm{}^s`$-reflexive for some $`0<s\mathrm{}`$ cannot be omitted.
###### Corollary 5.2.
Let $`0<p<1`$. Then $`(H_0^{p,\mathrm{}},c_0)=\mathrm{}_{1/p1}^{\mathrm{}}.`$
Proof: Use Lemma 3.1(i) and (ii), Lemma 5.2(ii), Theorem 5.4 and the identity $`\mathrm{}_{1/p1}^{\mathrm{}}=((c_0)_{11/p},c_0)`$.
We turn now to the study of the multiplier space $`(H^{p,q},X)`$ where $`X`$ is a $`\mathrm{}^s`$-reflexive FK-space. We need two results from the theory of mixed norm spaces. The first of these is part of the folklore. The case $`q=t`$ may be found in .
###### Lemma 5.3.
Let $`0<p2,0<qt\mathrm{},0<\alpha <\mathrm{},\mathrm{}<\beta <\mathrm{}`$. Set $`p_0=min(1,p)`$ and $`p_1=max(1,p)`$. Let $`p_1^{}`$ be the Hölder conjugate of $`p_1`$, $`1/p_1+1/p_1^{}=1`$. Then there is the embedding
$$H(p,q,\alpha ,\beta )\mathrm{}(p_1^{},t,11/p_0+\beta \alpha ).$$
The second result we need is a theorem of Pavlovic characterizing the multiplier spaces $`(H(p,q,\alpha ),\mathrm{}^s)`$ for $`0<p1,0<q,s\mathrm{},0<\alpha <\mathrm{},`$ . See also . For some special cases of the theorem see and .
###### Theorem 5.5.
Let $`0<p1,0<q,s\mathrm{},0<\alpha <\mathrm{}`$. Then
$$(H(p,q,\alpha ),\mathrm{}^s)=\mathrm{}(s,qs,\alpha +1/p1).$$
Lemma 2.3 may be used to obtain the following extension of Theorem 5.5.
###### Corollary 5.3.
Let $`0<p1,0<q,s\mathrm{},0<\alpha <\mathrm{},\mathrm{}<\beta <\mathrm{}`$. Then
$$(H(p,q,\alpha ,\beta ),\mathrm{}^s)=\mathrm{}(s,qs,\alpha \beta +1/p1).$$
Duren and Shields showed that $`(H^p,\mathrm{}^s)=\mathrm{}(s,\mathrm{},1/p1)`$ for $`0<p<1,ps\mathrm{}`$, , , . In Jevtic and Pavlovic showed that
$`(H^p,\mathrm{}^s)=\mathrm{}(s,ps,1/p1)`$ for the case $`0<s<p<1`$. Theorem 5.6 below extends these results to the Hardy-Lorentz space setting.
###### Theorem 5.6.
Let $`0<p<1,0<q<\mathrm{},0<s\mathrm{}`$. Then
$$(H^{p,q},\mathrm{}^s)=\mathrm{}(s,qs,1/p1).$$
( 5.15)
###### Proof.
From Theorem 4.1 we have the embeddings
$$H(p_0,q,\beta +1/p1/p_0,\beta )H^{p,q}H(1,q,1/p1).$$
( 5.16)
where $`\beta >1/p_01/p>0`$. Applying Lemma 3.1(i) to (5.16) and using Corollary 5.3 we have
$`\mathrm{}(s,qs,1/p1)`$ $`=`$ $`(H(1,q,1/p1),\mathrm{}^s)`$
$``$ $`(H^{p,q},\mathrm{}^s)`$
$``$ $`(H(p_0,q,,\beta +1/p1/p_,\beta )`$
$`=`$ $`\mathrm{}(s,qs,1/p1),`$
which establishes (5.15). ∎
Theorem 5.6 generalizes to
###### Theorem 5.7.
Let $`0<p<1`$, $`0<q<\mathrm{}`$. Let $`X`$ be a FK-space which is $`\mathrm{}^s`$-reflexive for some $`0<s\mathrm{}`$ and set $`t=max(q,s)`$. Then
$$(H^{p,q},X)=(\mathrm{}(\mathrm{},t,11/p),X).$$
###### Proof.
Observe that for $`0<q,s\mathrm{}`$,
$$(qs)s=t.$$
( 5.17)
Now apply Lemma 3.1(iii) followed by Theoren 5.6 to obtain
$`(H^{p,q},X)`$ $``$ $`(X^{K(s)},(H^{p,q})^{K(s)})`$
$`=`$ $`(X^{K(s)},\mathrm{}(s,qs,1/p1)).`$
Since $`X`$ is $`\mathrm{}^s`$-reflexive, a second application of Lemma 3.1(iii) to the last space in (5) together with Theorem 5.1 and (5.17) yields
$`(X^{K(s)},\mathrm{}(s,qs,1/p1)`$ $``$ $`(\mathrm{}(s,qs,1/p1)^{K(s)},X^{K(s)K(s)})`$
$`=`$ $`(\mathrm{}(\mathrm{},(qs)s,11/p),X)`$
$`=`$ $`(\mathrm{}(\mathrm{},t,11/p),X).`$
Then, starting with the last space in (5), use Lemma 3.1(iii) three times, first with Lemma 5.3, then with Theorem 4.1, and finally with the Hardy-Lorentz analog of inclusion (2.1). As a result we get
$`\mathrm{}(\mathrm{},t,11/p),X)`$ $``$ $`(H(1,t,1/p1),X)`$
$``$ $`(H^{p,t},X)`$
$``$ $`(H^{p,q},X).`$
Combining (5) through (5) completes the proof. ∎
Theorem 5.4 and Theorem 5.7 may be used to compute the multiplier spaces $`(E,ces(s))`$, where $`E`$ is one of the Hardy-Lorentz spaces $`H_0^{p,\mathrm{}}`$ or $`H^{p,q}`$, $`0<p<1`$, $`0<q\mathrm{}`$ and for $`1<s<\mathrm{}`$, $`ces(s)`$ is the Cesaro sequence space consisting of sequences $`\{x_k\}W`$ satisfying $`_{n=1}^{\mathrm{}}(\frac{1}{n}_{k=1}^n|x_k|)^s<\mathrm{}`$. Since it is known that $`ces(s)=\mathrm{}(1,s,\frac{1}{s}1)`$, , then we have the following result.
###### Corollary 5.4.
Let $`0<p<1`$, $`0<q\mathrm{}`$, $`1<s<\mathrm{}`$. Then
(i) $`(H^{p,q},ces(s))=\mathrm{}(1,qs,1/p+1/s2)`$,
(ii) $`H_0^{p,\mathrm{}},ces(s))=\mathrm{}(1,s,1/p+1/s2)`$.
###### Corollary 5.5.
Let $`0<p<1`$ and suppose $`X`$ is a FK-space which is $`\mathrm{}^s`$-reflexive for some $`0<s\mathrm{}`$. Then for every $`0<qs`$,
$$(H^{p,q},X)=(H^{p,s},X)=(\mathrm{}(\mathrm{},s,11/p),X).$$
Corollary 5.5 asserts that given a FK-space $`X`$ which is $`\mathrm{}^s`$-reflexive for some $`0<s\mathrm{}`$ and a number $`0<p<1,`$ the $`X`$-valued multiplier spaces for the Hardy-Lorentz space scale $`\{H^{p,q}\}_{0<qs}`$ will coincide. This is really due to the fact that the Bergman-Sobolev spaces appearing in the proof of Theorem 5.5 enjoy this property. We conclude this section with a Hardy-Lorentz analog of a well-known result of Hardy and Littlewood for $`H^p`$ spaces.
###### Corollary 5.6.
Let $`0<p<1`$, $`0<q<\mathrm{}`$. Suppose $`fH^{p,q}`$ has Taylor series representation $`f(z)=_{n=0}^{\mathrm{}}a_nz^n,z𝔻`$. Then
$$\left(\underset{n=1}{\overset{\mathrm{}}{}}n^{q(1\frac{1}{p})1}|a_n|^q\right)^{1/q}Cf_{H^{p,q}}.$$
( 5.21)
The diagonal case $`p=q`$ is a result of Hardy and Littlewood and is actually valid for $`0<p2`$. We also mention that for $`fH^p`$, $`0<p1,p<q<\mathrm{},`$ the series (5.21) is known to converge, . Since $`\{n^{1\frac{1}{p}\frac{1}{q}}\}`$ belongs to $`\mathrm{}(q,\mathrm{},1/p1),`$ Corollary 5.5 follows from the identification $`(H^{p,q},\mathrm{}^q)=\mathrm{}(q,\mathrm{},1/p1)`$. Finally we note that (5.21) remains valid if $`H^{p,q}`$ is replaced by $`H(1,q,1/p1)`$.
## 6 Multipliers of $`H^{p,q}`$ and $`H_0^{p,\mathrm{}},`$ $`0<p<1,0<q<\mathrm{}`$ into Bergman-Sobolev spaces
Our main tools in this section are Theorem 4.1 and the following result of Pavlovic , see also .
###### Theorem 6.1.
Let $`0<q,s,t\mathrm{},0<pmin(1,s),0<\alpha ,\beta <\mathrm{}`$. Then
$$(H(p,q,\alpha ),H(s,t,\beta ))=H(s,qt,\beta ,\alpha +1/p1).$$
( 6.1)
Note that by Lemma 2.2, the space on the right-hand side of (6.1) coincides with $`H(s,qt,1,1/p+\alpha \beta )`$.
###### Corollary 6.1.
Let $`0<q,s,t\mathrm{},0<pmin(1,s),0<\alpha ,\beta <\mathrm{},\mathrm{}<\delta ,\gamma <\mathrm{}`$. Then
(i) $`(H(p,q,\alpha ,\delta ),H(s,t,\beta ,\gamma ))=H(s,qt,1,1/p+\alpha \beta +\gamma \delta ),`$
(ii) $`(H_0(p,\mathrm{},\alpha ,\delta ),H(s,t,\beta ,\gamma ))=H(s,t,1,1/p+\alpha \beta +\gamma \delta ),`$
(iii) $`(H_0(p,\mathrm{},\alpha ,\delta ),H_0(s,\mathrm{},\beta ,\gamma ))=H(s,\mathrm{},1,1/p+\alpha \beta +\gamma \delta ),`$
(iv)$`(H(p,\mathrm{},\alpha ,\delta ),H_0(s,\mathrm{},\beta ,\gamma ))=H_0(s,\mathrm{},1,1/p+\alpha \beta +\gamma \delta ),`$
(v) $`(H(p,q,\alpha ,\delta ),H_0(s,\mathrm{},\beta ,\gamma ))=H(s,\mathrm{},1,1/p+\alpha \beta +\gamma \delta )`$ if $`q\mathrm{}.`$
###### Proof.
(i) Let $`g(H(p,q,\alpha ,\delta ),H(s,t,\beta ,\gamma ))`$. By Lemma 2.3 the maps $`D_\delta `$ and $`D^\gamma `$ are continuous surjective isomorphisms
$$D_\delta :H(p,q,\alpha )H(p,q,\alpha ,\delta )$$
$$D^\gamma :H(s,t,\beta ,\gamma )H(s,t,\beta )$$
and hence $`D^\gamma gD_\delta (H(p,q,\alpha ),H(s,t,\beta ))`$. But $`D^\gamma gD_\delta `$ is the same multiplier as $`(g^{[\gamma ]})_{[\delta ]}`$. Therefore, from (6.1), Lemmas 2.2 and 2.3, we deduce
$`(H(p,q,\alpha ,\delta ),H(s,t,\beta ,\gamma ))`$ $``$ $`D_\gamma (D^\delta ((H(p,q,\alpha ),H(s,t,\beta ))))`$
$`=`$ $`D_\gamma (D^\delta (H(s,qt,\beta ,\alpha +1/p1)))`$
$`=`$ $`D_\gamma (H(s,qt,\beta ,\alpha \delta +1/p1))`$
$`=`$ $`H(s,qt,\beta ,\alpha +\gamma \delta +1/p1).`$
Similarly,
$$(H(p,q,\alpha ),H(s,t,\beta ))D_\delta (D^\gamma (H(p,q,\alpha ,\delta ),H(s,t,\beta ,\gamma )).$$
( 6.3)
Therefore (6.3)implies
$`H(s,qt,\alpha +\gamma \delta +1/p1)`$ $`=`$ $`D_\gamma (D^\delta (H(s,qt,\beta ,\alpha +1/p1)))`$
$`=`$ $`D_\gamma (D^\delta ((H(p,q,\alpha ),H(s,t,\beta ))))`$
$``$ $`(H(p,q,\alpha ,\delta ),H(s,t,\beta ,\gamma )).`$
Then (6) and (6) together give us (i).
(ii) This follows from (i) and the monotonicity of the means $`M_k(,f)`$ for $`fH(𝔻),0<k\mathrm{}`$.
(iii) This follows from (ii) and the fact that for every $`0<k\mathrm{},0<\eta <\mathrm{},\mathrm{}<\nu <\mathrm{},`$ the function $`FH(k,\mathrm{},\eta ,\nu )`$ belongs to $`H_0(k,\mathrm{},\eta ,\nu )`$ if and only if $`F_rF`$ in $`H(k,\mathrm{},\eta ,\nu )`$ as $`r1^{}`$.
(iv) First note that the inclusion
$$(H_0(s,\mathrm{},1,1/p+\alpha \beta +\gamma \delta )(H_0(p,\mathrm{},\alpha ,\gamma ),H(s,\mathrm{},\beta ,\gamma ))$$
( 6.5)
follows from (i) and the fact that for every $`0<k\mathrm{},0<\eta <\mathrm{},\mathrm{}<\nu <\mathrm{},`$ the function $`FH(k,\mathrm{},\eta ,\nu )`$ belongs to $`H_0(k,\mathrm{},\eta ,\nu )`$ if and only if $`F_rF`$ in $`H(k,\mathrm{},\eta ,\nu )`$ as $`r1^{}`$.
To prove the reverse inclusion of (6.5) we consider the case $`\delta =\gamma =0`$ first. Let $`g(H(p,\mathrm{},\alpha ),H_0(s,\mathrm{},\beta ))`$. The Cauchy-type function $`F(z)=(1z)^{\alpha \frac{1}{p}},z𝔻`$ belongs to $`H(p,\mathrm{},\alpha )`$. Therefore $`gFH_0(s,\mathrm{},\beta )`$. But for any $`w𝔻`$, $`gF(w)=\mathrm{\Gamma }(\alpha +1/p)^1g^{[\alpha +\frac{1}{p}1]}(w)`$. Hence
$`gH_0(s,\mathrm{},\beta ,\alpha +1/p1)`$ which is equivalent to $`gH_0(s,\mathrm{},1,1/p+\alpha \beta )`$ by Lemma 2.2. For the general case we argue as in the proof of (i), using Lemma 2.3 to write
$$(H(p,\mathrm{},\alpha ,\delta ),H_0(s,\mathrm{},\beta ,\gamma ))=D^\delta (D_\gamma ((H(p,\mathrm{},\alpha ),H_0(s,\mathrm{},\beta ))).$$
Then use the validity of (iv) for the case $`\delta =\gamma =0`$.
(v) The proof is similar to the proof of (iv). ∎
###### Theorem 6.2.
Let $`0<q,s,t\mathrm{},0<\beta <\mathrm{},0<p<min(1,s),\mathrm{}<\gamma <\mathrm{}`$. Then
(i) $`(H^{p,q},H(s,t,\beta ,\gamma ))=H(s,qt,1,1/p+\gamma \beta ),`$
(ii) $`(H^{p,q},H_0(s,\mathrm{},\beta ,\gamma ))=H(s,\mathrm{},1,1/p+\gamma \beta ),q\mathrm{}`$
(iii) $`(H^{p,\mathrm{}},H_0(s,\mathrm{},\beta ,\gamma ))=H_0(s,\mathrm{},1,1/p+\gamma \beta ),`$
(iv) $`(H_0^{p,\mathrm{}},H(s,t,\beta ,\gamma ))=H(s,t,1,1/p+\gamma \beta ),`$
(v) $`(H_0^{p,\mathrm{}},H_0(s,\mathrm{},\beta ,\gamma ))=H(s,\mathrm{},1,1/p+\gamma \beta ).`$
###### Proof.
We prove (i) only, the proofs of (ii) through (v) being similar. Let $`0<p_0<p`$ and $`\delta >1/p_01/p`$. By Theorem 4.1 we have embeddings
$$H(p_0,q,\delta +1/p1/p_0,\delta )H^{p,q}H(s,q,1/p1/s),$$
( 6.6)
By Corollary 6.1, the multiplier spaces $`(H(p_0,q,\delta +1/p1/p_0,\delta ),H(s,t,\beta ,\gamma ))`$ and $`(H(s,q,1/p1/s),H(s,t,\beta ,\gamma ))`$ are both equal to
$`H(s,qt,1,1/p+\gamma \beta )`$. Then (i) follows from this fact and (6.6). ∎
## 7 Multipliers of $`H^{p,q}`$ and $`H_0^{p,\mathrm{}},`$ $`0<p<1,0<q\mathrm{}`$ into Hardy spaces
In this section we consider the multiplier spaces $`(H^{p,q},H^s)`$ and $`(H_0^{p,\mathrm{}},H^s)`$ for $`0<p<1`$, $`0<q\mathrm{}`$, $`0<s<\mathrm{}`$. We are able to determine these spaces for the cases $`0<q\text{ min}(2,s)`$ and $`0<q\mathrm{},s=2`$. Since $`H^2=\mathrm{}^2,`$ the second case was addressed in Section 5. We restate that result in terms of Bergman-Sobolev spaces below. To do this we need the following lemma from , see also , , .
###### Lemma 7.1.
Let $`0<q\mathrm{},0<\alpha <\mathrm{},\mathrm{}<\beta <\mathrm{}.`$ Then
$$H(2,q,\alpha ,\beta )=\mathrm{}(2,q,\beta \alpha ).$$
Using Lemma 7.1, the identification $`H^2=\mathrm{}^2`$, and either Corollary 5.1 or Theorem 5.6, we have the following result.
###### Theorem 7.1.
Let $`0<p<1,0<q\mathrm{}`$. Then
$$(H^{p,q},H^2)=H(2,q2,1,1/p).$$
We turn now to the case $`0<q`$ min$`(2,s)`$. First we record the following.
###### Theorem 7.2.
Let $`0<s\mathrm{},0<p<min(1,s)`$. Then
$$(H^p,H^s)=H(s,\mathrm{},1,1/p).$$
Theorem 7.2 dates back to Hardy and Littlewood who observed that
$`H(s,\mathrm{},1,p)(H^p,H^s)`$ for $`0<p<1s\mathrm{}`$, , . The case $`s=\mathrm{}`$ corresponding to the Duren-Romberg-Shields Theorem of reduces to (4.3). Duren and Shields proved Theorem 7.2 for the case
$`0<p<1s<\mathrm{}`$, . The proof for the case $`0<p<s1`$ is due to Mateljevic and Pavlovic, . Theorem 7.3 below represents an extension of Theorem 7.2 to $`H^{p,q}`$ for $`0<q`$ min$`(2,s)`$. We will need the following lemma. Before stating this result we introduce some notation. For $`0<s<\mathrm{}`$, the Dirichlet-type space $`𝒟^s`$ is defined to be the space $`H(s,s,1,1)`$.
###### Lemma 7.2.
Let $`0<s2t<\mathrm{}`$. Then
(i) $`𝒟^sH^sH(s,2,1,1),`$
(ii) $`H(t,2,1,1)H^t𝒟^t.`$
For statements and proofs of Lemma 7.2 the reader may consult ,,,,, and . Recently, A. Baernstein, D. Girela, and J. A. Pelaez have shown that for all $`0<s<\mathrm{}`$, $`H^s𝒰=𝒟^s𝒰`$, where $`𝒰`$ is the class of univalent functions on $`𝔻`$, .
###### Theorem 7.3.
Let $`0<s<\mathrm{}`$, $`0<p<min(1,s),\text{and }0<qmin(2,s)`$. Then
$$(H^{p,q},H^s)=H(s,\mathrm{},1,1/p).$$
###### Proof.
Proof: Assume first that $`0<s2`$. Then using Lemma 7.2 and Lemma 3.1(ii) we have
$$(H^{p,q},𝒟^s)(H^{p,q},H^s)(H^{p,q},H(s,2,1,1)).$$
( 7.1)
By Theorem 6.2(i), both of the endpoint spaces in (7.1) are equal to $`H(s,\mathrm{},1,1/p)`$. For the case $`2s<\mathrm{}`$, the reverse inclusion of (7.1) holds and the rest of the proof is exactly the same as in the first case. ∎
Theorem 7.3 has the following corollary.
###### Corollary 7.1.
Let $`0<qs2,0<p<min(1,s)`$. Then
$$(H^{p,q},H^{s,q})=\underset{qts}{}(H^{p,t},H^{s,t})$$
( 7.2)
###### Proof.
Proof: We prove the inclusion $`(H^{p,q},H^{s,q})_{qts}(H^{p,t},H^{s,t})`$ with the reverse conclusion being obvious. Let $`g(H^{p,q},H^{s,q})`$. Since $`qs2,`$ we have $`(H^{p,q},H^{s,q})(H^{p,q},H^s)=(H^{p,s},H^s)`$ by the Hardy-Lorentz analog of (2.1), Lemma 3.1(i) and Theorem 7.3. Thus g is a bounded multiplier for
$$g:H^{p,s}H^s\text{ and }g:H^{p,q}H^{s,q}.$$
Therefore, by interpolation, we find g is also bounded as a multiplier
$$g:(H^{p,q},H^{p,s})_{\theta ,t}(H^{s,q},H^s)_{\theta ,t},$$
for $`0<\theta <1,`$ and $`\frac{1}{t}=\frac{1\theta }{q}+\frac{\theta }{s}`$. Then an application of Theorem 2.1(ii) implies $`g`$ is bounded as a multiplier
$$g:H^{p,t}H^{s,t}\text{ for all }qts.$$
## 8 Multipliers of $`H^{p,q}`$ and $`H_0^{p,\mathrm{}},0<p<1,0<q\mathrm{}`$ into analytic Lipschitz spaces, analytic Zygmund spaces, Bloch spaces, and $`BMOA`$
In this section we apply the results of the previous two sections to some specific target spaces belonging to the class of Bergman-Sobolev spaces. The target spaces we have in mind are the analytic Lipschitz and Zygmund spaces and the Bloch spaces. We also have some results for the case when the target space is $`BMOA`$. For the discussion that follows we assume $`fH(𝔻)`$ and that $`f`$ has non-tangential limits $`m`$-a.e. on $`𝕋`$. We denote the resulting boundary value function by the same symbol $`f`$. For $`1s\mathrm{},`$ the moduli of continuity $`\omega _s(f)(t)`$ and $`\mathrm{\Omega }_s(f)(t)`$ of $`f`$ are defined for $`t>0`$ by $`\omega _s(f)(t)=sup_{0<|h|t}T_h(f)f_s`$ and $`\mathrm{\Omega }_s(f)(t)=sup_{0<|h|t}T_h(f)2f+T_h(f)_s`$, where $`T_h`$ is the translation operator given by $`T_h(f)(e^{i\theta })=f(e^{i(\theta +h)})`$. Let $`0<\alpha 1`$. Then $`f`$ is said to belong to the analytic Lipschitz space $`\mathrm{\Lambda }_\alpha ^s(𝔻)`$ (resp. $`\lambda _\alpha ^s(𝔻)`$ if $`\omega _s(f)(t)=O(t^\alpha )`$ (resp. $`o(t^\alpha )`$) as $`t0^+`$. If the boundary value function $`fC(𝕋)`$ and $`\mathrm{\Omega }_{\mathrm{}}(f)(t)=O(t)`$ (resp. $`o(t)`$) as $`t0^+,`$ then $`f`$ is said to belong to the analytic Zygmund space $`\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻)`$ (resp. $`\lambda _{}^{\mathrm{}}(𝔻)`$). For $`1s<\mathrm{},`$ $`f`$ is said to belong to the analytic Zygmund space $`\mathrm{\Lambda }_{}^s(𝔻)`$ (resp. $`\lambda _{}^s(𝔻)`$) if $`\mathrm{\Omega }_s(f)(t)=O(t^\alpha )`$ (resp. $`o(t^\alpha ))`$ as $`t0^+`$. Theorem 8.1 below is a collection of well-known results of Hardy and Littlewood identifying various analytic Lipschitz and Zygmund spaces as Bergman-Sobolev spaces. See and . In order to cover the case $`\alpha =1`$, we recall that for $`0<s\mathrm{},`$ $`0<\beta <\mathrm{},`$ the Hardy-Sobolev space $`H_\beta ^s=\{fH(𝔻):f^{[\beta ]}H^s\}`$.
###### Theorem 8.1.
Let $`0<\alpha <1s\mathrm{}`$. Then
(i) $`\mathrm{\Lambda }_\alpha ^s(𝔻)=H(s,\mathrm{},1\alpha ,1)`$ and $`\lambda _\alpha ^s(𝔻)=H_0(s,\mathrm{},1\alpha ,1),`$
(ii) $`\mathrm{\Lambda }_{}^s(𝔻)=H(s,\mathrm{},1\alpha ,2)`$ and $`\lambda _{}^s(𝔻)=H_0(s,\mathrm{},1\alpha ,2),`$
(iii) $`\mathrm{\Lambda }_1^s(𝔻)=H_1^s.`$
We combine Theorem 8.1 with the Duren-Romberg-Shields Theorem to find the multipliers from $`H^{p,q}`$ into the analytic Lipschitz spaces $`\mathrm{\Lambda }_\alpha ^{\mathrm{}}(𝔻),`$ $`\lambda _\alpha ^{\mathrm{}}(𝔻)`$ and analytic Zygmund spaces $`\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻),`$ $`\lambda _{}^{\mathrm{}}(𝔻),`$ $`0<\alpha ,p<1,`$ $`0<q\mathrm{}`$.
###### Corollary 8.1.
Let $`0<\alpha ,p<1,`$ $`0<q\mathrm{}`$. Then
(i) $`(H^{p,q},\mathrm{\Lambda }_\alpha ^{\mathrm{}}(𝔻))=(H_0^{p,\mathrm{}},\mathrm{\Lambda }_\alpha ^{\mathrm{}}(𝔻))=H(\mathrm{},\mathrm{},1,\frac{1}{p}+\alpha )=(H^{\frac{p}{1+\alpha p}})^{},`$
(ii) For $`q\mathrm{},`$
$$(H^{p,q},\lambda _\alpha ^{\mathrm{}}(𝔻))=(H_0^{p,\mathrm{}},\lambda _\alpha ^{\mathrm{}}(𝔻))=H(\mathrm{},\mathrm{},1,1/p+\alpha )=(H^{\frac{p}{1+\alpha p}})^{},$$
(iii) $`(H^{p,\mathrm{}},\lambda _\alpha ^{\mathrm{}}(𝔻))=H_0(\mathrm{},\mathrm{},1,\frac{1}{p}+\alpha ),`$
(iv) $`(H^{p,q},\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻))=(H_0^{p,\mathrm{}},\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻))=H(\mathrm{},\mathrm{},1,\frac{1}{p}+1)=(H^{\frac{p}{1+p}})^{},`$
(v) For $`q\mathrm{},`$
$$(H^{p,q},\lambda _{}^{\mathrm{}}(𝔻))=(H_0^{p,\mathrm{}},\lambda _{}^{\mathrm{}}(𝔻))=H(\mathrm{},\mathrm{},1,1/p+1)=(H^{\frac{p}{1+p}})^{},$$
(vi) $`(H^{p,\mathrm{}},\lambda _{}^{\mathrm{}}(𝔻))=H_0(\mathrm{},\mathrm{},1,\frac{1}{p}+1).`$
Corollary 8.1 shows that for $`0<\alpha <1,`$ the secondary index $`q`$ is irrelevant with respect to the multiplier spaces $`(H^{p,q},E)`$ for the target spaces
$`E=\mathrm{\Lambda }_\alpha ^{\mathrm{}}(𝔻)`$ or $`E=\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻)`$. The same is true for the target spaces $`E=\lambda _\alpha ^{\mathrm{}}(𝔻)`$ or $`E=\lambda _{}^{\mathrm{}}(𝔻)`$ provided $`q\mathrm{}`$. A similar phenomenon occurs when the target spaces are the Bloch spaces. Let us recall that the Bloch space $``$ and the little Bloch space $`_0`$ are realized as Bergman-Sobolev spaces using the identifications $`=H(\mathrm{},\mathrm{},1,1)`$ and $`_0=H_0(\mathrm{},\mathrm{},1,1)`$. The Bloch space analog of Corollary 8.1 is the following.
###### Corollary 8.2.
Let $`0<p<1,`$ $`0<q\mathrm{}`$. Then
(i) $`(H^{p,q},)=(H_0^{p,\mathrm{}},)=H(\mathrm{},\mathrm{},1,1/p)=(H^p)^{},`$
(ii)$`(H^{p,q},_0)=(H_0^{p,\mathrm{}},_0)=H(\mathrm{},\mathrm{},1,1/p)=(H^p)^{},q\mathrm{},`$
(iii)$`(H^{p,\mathrm{}},_0)=H_0(\mathrm{},\mathrm{},1,1/p).`$
We observe here that the fractional derivative operator $`D=D^1`$ is a continuous isomorphism of the analytic Zygmund space $`\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻)`$ (resp. $`\lambda _{}^{\mathrm{}}(𝔻)`$) onto $``$ (resp. $`_0`$) Thus, Corollary 8.2 may be viewed as an isomorphic version of the Zygmund space portion of Corollary 8.1.
In contrast, the Lipschitz spaces $`\mathrm{\Lambda }_1^{\mathrm{}}(𝔻)`$ will, in general, determine different multiplier spaces $`(H^{p,q},\mathrm{\Lambda }_1^{\mathrm{}}(𝔻))`$ for different values of $`q`$. This is demonstrated in the next result which follows from Theorem 8.1 and Theorem 4.2.
###### Corollary 8.3.
Let $`0<p<1,`$ $`0<q<\mathrm{}`$. Then
(i) $`(H^{p,q},\mathrm{\Lambda }_1^{\mathrm{}}(𝔻))=(H^{\frac{p}{1+p},q})^{},`$
(ii)$`(H^{p,\mathrm{}},\mathrm{\Lambda }_1^{\mathrm{}}(𝔻))=(H_0^{\frac{p}{1+p},\mathrm{}})^{}.`$
For $`1s<\mathrm{},`$ we have the following analogs of Corollaries 8.1 and 8.3.
###### Corollary 8.4.
Let $`0<\alpha ,p<1s<\mathrm{},`$ $`0<q\mathrm{}`$. Then
(i) $`(H^{p,q},\mathrm{\Lambda }_\alpha ^s(𝔻))=(H_0^{p,\mathrm{}},\mathrm{\Lambda }_\alpha ^s(𝔻))=H(s,\mathrm{},1,\frac{1}{p}+\alpha )=(H^{\frac{p}{1+\alpha p}},H^s),`$
(ii) For $`q\mathrm{},`$
$$(H^{p,q},\lambda _\alpha ^s(𝔻))=(H_0^{p,\mathrm{}},\lambda _\alpha ^s(𝔻))=H(s,\mathrm{},1,1/p+\alpha )=(H^{\frac{p}{1+\alpha p}},H^s),$$
(iii) $`(H^{p,\mathrm{}},\lambda _\alpha ^s(𝔻))=H_0(s,\mathrm{},1,\frac{1}{p}+\alpha ),`$
(iv) $`(H^{p,q},\mathrm{\Lambda }_{}^s(𝔻))=(H_0^{p,\mathrm{}},\mathrm{\Lambda }_{}^s(𝔻))=H(s,\mathrm{},1,\frac{1}{p}+1)=(H^{\frac{p}{1+p}},H^s),`$
(v) For $`q\mathrm{},`$
$$(H^{p,q},\lambda _{}^s(𝔻))=(H_0^{p,\mathrm{}},\lambda _{}^s(𝔻))=H(s,\mathrm{},1,1/p+1)=(H^{\frac{p}{1+p}},H^s),$$
(vi) $`(H^{p,\mathrm{}},\lambda _{}^s(𝔻))=H_0(s,\mathrm{},1,\frac{1}{p}+1).`$
###### Corollary 8.5.
Let $`0<p<1s<\mathrm{},`$ $`0<q\mathrm{}`$. Then
(i)$`(H^{p,q},\mathrm{\Lambda }_1^s(𝔻))=(H^{\frac{p}{1+p},q},H^s),`$
(ii)$`(H_0^{p,\mathrm{}},\mathrm{\Lambda }_1^s(𝔻))=(H_0^{\frac{p}{1+p},\mathrm{}},H^s).`$
The space $`BMOA`$ is the space of functions $`fH(𝔻)`$ having non-tangential limits $`m`$-a.e. on $`𝕋`$ for which the resulting boundary value function $`f`$ is of bounded mean oscillation on $`𝕋`$. That is for which
$$\underset{I𝕋}{sup}[m(I)^1(ff_I)\chi _I_1]<\mathrm{}$$
( 8.1)
where the supremum in (8.1) is taken over all subintervals $`I𝕋`$ and
$`f_I=m(I)^1_If(z)𝑑m(z)`$. The space $`BMOA`$ is not a Bergman-Sobolev space. However we do have the following embedding, see .
###### Lemma 8.1.
$`H(\mathrm{},2,1,1)BMOA`$.
###### Theorem 8.2.
Let $`0<p<1,`$ $`0<q2`$. Then
$$(H^{p,q},BMOA)=(H^p)^{}=(H^{p,\mathrm{}},).$$
###### Proof.
Using the Duren-Romberg-Shields Theorem, Theorem 6.2, Lemma 8.1, Lemma 3.1(ii), and Corollary 8.2(i) we find $`(H^p)^{}=H(\mathrm{},\mathrm{},1,1/p)=(H^{p,q},H(\mathrm{},2,1,1)(H^{p,q},BMOA)(H^{p,q},)=(H^p)^{}`$ and $`(H^p)^{}=(H^{p,\mathrm{}},)`$ by Corollary 8.2(i). ∎
We have not been able to find $`(H^{p,q},BMOA)`$ for $`2<q\mathrm{}`$. However we can show
$$(H^{\frac{1}{2},\mathrm{}},BMOA)(H^{\frac{1}{2},\mathrm{}},).$$
( 8.2)
To show (8.2) let $`G`$ be a function in $``$ which is not in $`BMOA`$. The fractional integral operator $`D=D_1`$ is a continuous isomorphism of $``$ onto $`\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻)`$ and hence $`G_{[1]}\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻)`$, . But $`\mathrm{\Lambda }_{}^{\mathrm{}}(𝔻)=(H^{\frac{1}{2},\mathrm{}},)`$ by Theorem 8.1 and Corollary 8.2(i) and so $`G_{[1]}(H^{\frac{1}{2},\mathrm{}},)`$. Since the function $`f(z)=(1z)^2H^{\frac{1}{2},\mathrm{}}`$ and since $`fG_{[1]}=G`$ then $`G_{[1]}`$ fails to multiply $`H^{\frac{1}{2},\mathrm{}}`$ into $`BMOA`$. Thus $`G_{[1]}`$ belongs to $`(H^{\frac{1}{2},\mathrm{}},)`$ but not to $`(H^{\frac{1}{2},\mathrm{}},BMOA)`$ which establishes (8.2).
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# Photometric Accretion Signatures Near the Substellar Boundary
## 1 Introduction
Brown dwarfs are defined as objects that have less than the minimum mass required for stable Hydrogen fusion. In the absence of rotation and for solar metallicity this mass, referred to as the hydrogen-burning limit, is 0.075 M or 79 M<sub>Jup</sub> (Chabrier & Baraffe, 2000). Young brown dwarfs are fully convective with photospheric temperatures comparable to M stars. In the models of Baraffe et al. (1998) 2 Myr old proto-brown dwarfs of masses 0.020 M up to the hydrogen-burning limit have temperatures between 2500 K and 2900 K, corresponding to spectral types of M8.5 to M6 (Luhman, 2003). The presence of a circumstellar disk in combination with the ability of fully convective objects to generate significant magnetic fields should result in young brown dwarfs exhibiting Classical T Tauri behavior (Jayawardhana, Mohanty, & Basri, 2002).
The young brown dwarfs in this paper are found as isolated objects within the Orion OB1 star formation region. Distances and inferred ages of the Orion OB1a,b,c, and d subassociations are given in Table 1 (from Sherry (2003)). There are two possible scenarios for brown dwarf formation: (1) individual formation out of isolated low-mass molecular cores, just as stars are presumed to form from higher-massed cores, and (2) ejection of the lowest massed objects from a cluster though dynamical interactions. In the latter scenario, objects are ejected before accreting enough material to grow into stars, thereby forming brown dwarfs (Reipurth & Clarke, 2001; Kroupa et al., 2003).
Consequences of the ejection process include a broader spatial distribution, increased velocity dispersion for the lower massed objects, and a truncation of the circumstellar disk via tidal forces. The ejection velocity dispersion is of order 5 km/sec corresponding to a dispersion in proper motion of 2.6 mas/yr at the 400 pc distance of Orion. Models of the ejection process (Armitage & Clarke, 1997) predict that the disk will be truncated to an outer radius of 10 AU or less, resulting in disk lifetimes of less than 1 Myr. In order to differentiate the star-like and ejection formation scenarios we examine observational signatures, such as Classical T Tauri behavior, that are due to the presence of a circumstellar disk.
### 1.1 Accretion Signatures in T Tauris
T Tauri variables were originally identified by Joy (1945) as a class of irregular variable stars exhibiting marked changes in brightness and color on timescales of hours to days. These changes are greatest in the blue and ultraviolet (Herbst et al., 1994), and are expected to be related to disk accretion. Magnetospheric accretion models (Königl, 1991) predict that the inner circumstellar disk is truncated by the magnetic field of the star. For accretion to occur the truncation must be at or within the Keplerian co-rotation radius, at which the orbital period of the disk matches that of the star’s rotation.
Disrupted disk material flows towards the star in accretion columns that follow the magnetic field lines connecting the inner edge of the disk with the high latitude regions of the star. As evidenced by the broad velocity profiles observed in H$`\alpha `$ and other spectral features, much of the permitted line emission arises in these columns. Other emission lines, including He I and the broad component of the Ca II IR triplet, plus the veiling blue and ultraviolet continua arise in the high temperature shock fronts at the base of the accretion columns near the stellar surface (Gullbring et al., 2000). The lifetime of the classical T Tauri phase is thought to be 1 to 10 Myr (Kenyon & Hartmann, 1995).
Protostellar magnetospheric accretion is one of a dynamic process characterized by intricate magnetic field topologies and unstable field-disk interactions. Detailed MHD computations by Romanova (2002), Küker, Henning, & Rüdiger (2003), and von Rekowski & Brandenburg (2003) indicate that even for the idealized case of a dipole field aligned with the stellar rotation axis the accretion process is fundamentally unstable. An intensive photometric and spectroscopic campaign targeting the nearly edge-on Classical T Tauri AA Tau by Bouvier (2003) shows clear evidence for large scale instabilities developing in T Tauri magnetospheres as the magnetic field lines are twisted by differential rotation between the star and the inner disk. Observational support of magnetospheric accretion in AA Tau includes time delays between the H$`\alpha `$, H$`\beta `$, He II line emission and the accretion shock generated emission consistent with free-fall from $``$8 stellar radii and the presence of two rotationally modulated hot spots.
### 1.2 Accretion in Very Low Mass Stars and Brown Dwarfs
There has been considerable effort over recent years to observe and characterize the formation processes in young very low mass stars and brown stars. This has included analyses of spectroscopic signatures of accretion, veiling measurements, variability studies, and IR searches for circum(sub)stellar disk emission.
#### 1.2.1 Emission Lines
H$`\alpha `$ emission lines with widths in excess of 200 km/s have been detected in young brown dwarfs having masses as low as $``$ 20 M<sub>Jup</sub> in many star formation regions including IC 348, Taurus, Chamaeleon I, Upper Scorpius, and $`\rho `$ Oph (Jayawardhana, Mohanty, & Basri, 2002; Muzerolle et al., 2003; Natta et al., 2004; Mohanty, Jayawardhana, & Basri, 2005; Muzerolle et al., 2005). The mass accretion rate appears to scale roughly as $`M^2`$ which is steeper than the linear relation inferred by White & Basri (2003) based on studies in Taurus-Auriga. In addition, the uncertainties in the determination of stellar parameters (mass, radius) from evolutionary models of very low mass objects (Baraffe et al., 2002) during first few Myr contribute another source of uncertainty in these relations. However, Padoan et al. (2005) propose that the accretion rate is controlled by accretion from the surrounding gas in the star formation region rather than from the disk evolution. The resulting Bondi-Hoyle accretion model yields trends consistent with the Muzerolle et al. (2003) relation although the validity of this hypothesis needs to be checked through detailed observations and simulations.
Natta et al. (2004) find that the mass accretion rate scales exponentially with the width of the H$`\alpha `$ line measured at the 10% intensity level by comparing previous results in the substellar domain (Muzerolle et al. (2003), White & Basri (2003), and Barrado y Navascues et al. (2004)) with observations in Chamaeleon I and $`\rho `$ Oph and study of Classical T Tauri stars(Gullbring et al., 1998). The lowest detected accretion rates are on the order of $`10^{11}`$ M/yr.
#### 1.2.2 Optical Veiling
Optical veiling measurements of young brown dwarfs have been primarily obtained redward of 6000 Å. Barrado y Navascues, Mohanty, & Jayawardhana (2004) find $`r_{6200}`$ = 1.0 and $`r_{6750}`$ = 0.25 from the infilling of several TiO bands in the very low mass object LS-RCrA 1, where $`r_\lambda `$ is the ratio of the excess emission to the photospheric emission measured at a wavelength $`\lambda `$ ($`\AA `$). Due to the presence of \[S II\], \[O I\], \[O II\], and \[N II\] forbidden lines in the spectra they propose that this substellar object is driving an outflow analogous to that which forms the Herbig-Haro shocks seen near intermediate and low mass T Tauri stars.
Accretion shock models for very low mass stars (Muzerolle et al., 2000) suggest that accretion rates as low as $`10^9`$ M/yr would be detectable in the Johnson V band against an M6 photosphere. We expect a greater sensitivity to accretion rates when working in the Sloan Digital Sky Survey (SDSS) $`g`$ band due to increased brightness contrast between the shock and the photosphere.
#### 1.2.3 Variability
Optical photometric variability of low mass pre-main sequence stars is the result of four physical effects. These are rotational modulation of cool spots, rotational modulation of hot spots formed at the base of magnetospheric accretion columns, instabilities in the mass accretion rate, and flaring.
The study by Herbst et al. (1994) finds three distinct patterns of variability in Weak-lined T Tauris (WTTS) and Classical T Tauris (CTTS). The Type I variables are characterized by low amplitude multi-band fluctuations on timescales of $`0.5`$ to 18 days. This behavior is primarily seen in WTTS but can been detected in some WTTS. The second type of variation only occurs in the CTTS and is marked by ireegular high amplitude brightness changes, especially in the near ultraviolet and the blue, on timescales as short as a few hours. The Type III variations, like the Type II, are irregular but occur on much longer timescales of days to weeks. This last category are only seen in pre-main sequence stars having early spectral types (AO to K1).
Type I variations are believed due to rotational modulation of cool spots on the stellar surface. These spots are a consequence of the magnetic activity and can persist for over 100 rotational periods yielding periodic variations in $`VRI`$. Flaring is often seen in the $`B`$ and especially $`U`$ bands. WTTS only exhibit Type I variations. When magnetospheric accretion is present, i.e. in the CTTS stage, then Type II variations are seen. These are due to a combination of rotational modulation of the accretion hot spots together with short period fluctuations in the mass accretion rate. In most cases the variations in the accretion mask the rotation signature. Herbst et al. (1994) also identify a third class (Type III) variables such as RY Tau and SU Aur that exhibit large non-periodic variation in the $`V`$ band but do not show significant veiling. This behavior is thought to be caused by variable obscuration but is also limited to intermediate and high-mass stars (spectral types A0 to K1); we do not expect to find Type III variables in our low-mass sample.
In this sample of intermediate and low mass T Tauris Herbst et al. (1994) find that the maximum amplitudes seen in Type I variables are 0.8 magnitudes in $`V`$ and 0.5 magnitudes in $`I`$ for V410 Tau (SpT = K3). For the Type I variables of spectral type M0 and later the $`V`$ band amplitudes are typically less than 0.3 magnitudes.
Variability studies of young very low mass stars and brown dwarfs have primarily used the longer wavelength filters, most notably the Johnson-Cousins $`I`$. Herbst et al. (2002) utilized an intermediate band filter centered at 815.9 nm to study variability in the Orion Nebula Cluster and found that 46% of the stars with $`\sigma <0.1`$ were periodic, where $`\sigma `$ is the standard deviation measured over all observations for a star. Only 24% having $`\sigma >0.1`$ were periodic, the majority exhibited the irregular variation characteristic of Type II variables.
In their study of IC 348 Cohen, Herbst, & Williams (2004) found that the full range of variation for periodic stars of spectral type M0-M4 or later is $`<0.6`$ magnitudes in the $`I`$ band. No CTTS exhibited periodic variability over 5 seasons of monitoring. Cohen, Herbst, & Williams (2004) also characterized the variability in their sample by $`\sigma `$ finding that in their $`I<14.3`$ subsample the WTTS had $`\sigma <0.1`$ and the CTTS were found to have $`\sigma >0.05`$. Only the Type II variables (CTTS) had $`\sigma >0.1`$.
The survey of the Cha I star formation region in $`R`$, $`i`$, and $`J`$ by Joergens et al. (2003) found that the largest amplitudes in each band were 0.18, 0.14, and 0.13, respectively. The earliest spectral type in the sample was M5, corresponding to a maximum mass of 0.12 M. These objects showed periodic variations, presumably due to spots, whose peak-to-peak amplitudes drop in later spectral types.
$`I`$ band monitoring of the $`\sigma `$ Ori and $`ϵ`$ Ori clusters were performed by Scholz & Eislöffel (2004a) and Scholz & Eislöffel (2005). They found that periodic variables had amplitudes less than 0.2 magnitudes. The larger amplitude variables were generally non-periodic. S Ori 45, a $``$ 20 Jupiter mass member of the $`\sigma `$ Ori cluster (SpT = M8.5) was studied by Zapatero Osorio et al. (2003) in the $`IJ`$ bands who measured a periodic peak to peak amplitude $`0.2`$ magnitudes.
Caballero et al. (2004) find that nine out of 32 young brown dwarf candidates (spectral types M5.5 to L2) exhibit $`I`$ band variability. The amplitudes of mid-term (day-to-day) variations were as high as 0.36 magnitudes while shorter term variations had standard deviations $``$ 0.05 magnitudes. Correlations were found between high amplitude variability ($`>0.12`$ magnitudes) and detection of either a near-IR excess or strong H$`\alpha `$ emission indicative of a possible accretion disk.
#### 1.2.4 Near-IR Excess
The thermal emission from the inner circumstellar disk can be detected in the near-IR $`H`$ and $`K`$ bands (Meyer, Calvet, & Hillenbrand, 1997). For the less massive and cooler young brown dwarfs the reduced contrast between the disk emission and the substellar photosphere motivates disk surveys in $`L`$ band. An $`L`$ band survey by Liu, Najita, & Tokunaga (2003) finds that for low mass stars and substellar objects in IC 348 and Taurus the disk fraction is not dependent on mass or age and that the presence of disk emission is associated with accretion indicators such as H$`\alpha `$ emission.
Comparison of disk models with mid-IR photometry shows that the disks surrounding young brown dwarfs can be flared like those around Classical T Tauri stars (Mohanty et al., 2004). The model fits also suggest the presence of inner holes that are a few substellar radii in size.
### 1.3 The Orion Equatorial Region
The Orion OB1 a,b,c and d sub-associations have been the targets of several surveys. These include the optical photometric and spectroscopic study of the $`\sigma `$ Ori cluster in Orion OB1b by Walter, Wolk, & Sherry (1998), the H$`\alpha `$ Kiso object prism survey (Wiramihardja et al., 1989), and the ROSAT All-Sky Survey (Alcalá et al., 1996). With the exception of the Orion Nebula Cluster (OB1d) these are found along the celestial equator and are considered in this study. The estimated ages span 1.7 Myr to 11.4 Myr with distances between 330 and 460 parsecs corresponding to distance moduli of 7.6 to 8.3 magnitudes.
Recent imaging surveys of the Orion equatorial region include Wolk (1996), Sherry (2003), and Briceño (2003). Wolk (1996) performed $`UBVR_cI_c`$ photometric monitoring of X-ray selected stars in the Orion OB1a and OB1b associations with an estimated completeness limit of $`V`$ = 18.5. Sherry (2003) conducted a deeper $`BVR_cI_c`$ survey covering 5 square degrees in Orion OB1b around $`\sigma `$ Ori, $`ϵ`$ Ori, and $`\delta `$ Ori using the CTIO 0.9 and 1.5 meter telescopes. Their adopted faint limit was $`V`$ = 20.5. The $`B`$ band photometry was not analyzed.
The CIDA-QUEST variability survey (Briceño, 2003) is acquiring multi-epoch drift-scanned imaging in the Orion equatorial region using the 1.0/1.5 meter Schmidt telescope at The National Astronomical Observatory of Venezuela. Exposures are acquired in the VRIH$`\alpha `$ bands at a plate scale of 1.02” per pixel. The 10$`\sigma `$ limit is approximately $`V=19.7`$ although in selection of candidate pre-main sequence stars they specify $`V=1618.5`$ (Briceño, 2004), resulting in a minimum mass limit of $``$ 0.12 M.
Studies of neighboring regions include Barrado y Navascues et al. (2004) who performed a deep $`R_cI_c`$ survey of the $`\lambda `$ Orionis cluster (d = 400 pc; 5 Myr) with the CFHT 12K camera. Their completeness limits were 22.75 in both bands with the lowest mass confirmed cluster member having 0.02 M with a spectral type of M8.5. Studies of NGC 2264 including S Mon and the Cone Nebula region (Rebull et al., 2002) and of the Orion Nebula Flanking fields (Rebull et al., 2000) explicitly used ultraviolet and near-IR excess as disk indicators. Their limiting magnitudes are $`U20`$, $`V20`$, and $`I_C18`$.
### 1.4 Program
The goal of this paper is to study the trends in veiling emission, variability, and near-infrared excess in the low mass stars and brown dwarfs found in the Orion OB1a,b, and c subassociations. If brown dwarfs form in the same manner as stars we expect that these characteristics will vary smoothly across the stellar/substellar mass boundary.
The remainder of this paper is organized as follows: the SDSS and Two Micron All-Sky Survey (2MASS) observations used in this paper are presented in §2. In §3 we discuss the expected colors of young mass stars and brown dwarfs using the colors of spectroscopically verified field M dwarfs (West et al., 2004), the evolutionary models of Baraffe et al. (1998) (hereafter BCAH98), and the bolometric corrections of Girardi et al. (2004). The selection criteria for young low mass objects including an estimate on the contamination fraction are presented in §4. We discuss the empirical trends in variability and color in §5 and in §6 we provide a summary of results and upcoming work.
## 2 Observations
### 2.1 Photometry
The Sloan Digital Sky Survey (SDSS) obtains deep photometry with magnitude limits (defined by 95% detection repeatability for point sources) of $`u=22.0`$, $`g=22.2`$, $`r=22.2`$, $`i=21.2`$ and $`z=20.5`$. These five passbands, $`ugriz`$, have effective wavelengths of 3540, 4760, 6290, 7690, and 9250 Å, respectively. A technical summary of the SDSS is given by York et al. (2000). The SDSS imaging camera is described by Gunn et al. (1998). Ivezić et al. (2004) discuss the data management and photometric quality assessment system.
The Early Data Release and the Data Release One are described by Stoughton et al. (2002) and Abazajian et al. (2003). The former includes an extensive discussion of the data outputs and software. Pier et al. (2003) describe the astrometric calibration of the survey and the network of primary photometric standard stars is described by Smith et al. (2002). The photometric system itself is defined by Fukugita et al. (1996), and the system which monitors the site photometricity by Hogg et al. (2001). Abazajian et al. (2003) discuss the differences between the native SDSS 2.5m $`ugriz`$ system and the $`u^{}g^{}r^{}i^{}z^{}`$ standard star system defined on the USNO 1.0 m (Smith et al., 2002).
The SDSS low Galactic latitude data which includes the Orion equatorial imaging used in this work are described by Finkbeiner et al. (2004). 2MASS (Skrutskie et al., 1997) obtained nearly complete coverage of the sky in $`JHK_s`$. The multi-epoch data (Table 2) we use in this study were obtained under photometric conditions and cover $`80\mathrm{°}<\alpha _{2000}<90\mathrm{°}`$ and $`1.25\mathrm{°}<\delta _{2000}<1.25\mathrm{°}`$.
### 2.2 Reddening-Invariant Indices
In this work we employ reddening-invariant indices of the form $`Q_{xyz}=(xy)(yz)\times E(xy)/E(yz)`$. $`Q_{xyz}`$ is dependent upon the assumed ratio of general to selective extinction ($`R_V=A_V/E(BV)`$; Cardelli, Clayton, & Mathis (1989)). Here $`xyz`$ refer to the specific passbands, e.g. $`ugrizJHK_S`$ and $`E(xy)`$ is the color excess due to reddening in the $`xy`$ color. This definition of reddening-invariant colors follows the original Johnson & Morgan (1953) $`Q`$ that would be written as $`Q_{UBV}`$ in our notation. In the $`(xy,yz)`$ color-color diagram the $`Q_{xyz}`$ axis is perpendicular to the reddening vector.
This approach has previously been used to study pre-main sequence populations by de Geus, de Zeeuw, & Lub (1989) and Brown, de Geus, & de Zeeuw (1994) for the Scorpius-Centaurus OB and Orion OB1b associations, respectively. In both cases reddening-invariant colors were formed by $`[xy]=(xy)(VB)\times E(xy)/E(VB)`$. using the $`VBLUW`$ Walraven system. Stellar parameters were estimated using reddening-invariant two-color diagrams to compare observed stellar colors against those from Kurucz model grid (Kurucz, 1979).
While the contrast between the accretion shock and the cool photosphere is greatest in the SDSS $`u`$ band, we choose not to use the $`Q_{ugr}`$ color due to sensitivity issues and the presence of a “red leak” in the SDSS imager $`u`$ filter. M dwarfs are 2.5 magnitudes fainter in $`u`$ than in $`g`$, thus use of the $`u`$ band in Orion would not reach young brown dwarfs.
We use the extinction tables derived by D. Finkbeiner<sup>1</sup><sup>1</sup>1private communication; see http://www.astro.princeton.edu/$``$dfink/sdssfilters/ to define the coefficients used in defining reddening-invariant colors. These tables contain the $`A_X/E(BV)`$ values for the SDSS $`ugriz`$ filters for specific values of $`R_V`$ and source spectra. The values we present here are obtained using an F dwarf source spectrum and $`R_V`$ = 3.1 and 5.5. The $`R_V=3.1`$ case is the standard extinction law found in the diffuse ISM. The $`R_V=5.5`$ law is representative of that found in molecular clouds due to larger dust grains and is shown here for example. Throughout this paper we will adopt $`R_V=3.1`$ as it is characteristic of all but perhaps the most heavily extincted regions of the Orion OB1b association (Warren & Hesser, 1978).
$$Q_{gri}=\{\begin{array}{cc}(gr)1.852(ri)\hfill & R_V\text{ = 3.1,}\hfill \\ (gr)1.339(ri)\hfill & R_V\text{ = 5.5}\hfill \end{array}$$
(1)
and
$$Q_{riz}=\{\begin{array}{cc}(ri)0.987(iz)\hfill & R_V\text{ = 3.1,}\hfill \\ (ri)1.004(iz)\hfill & R_V\text{ = 5.5.}\hfill \end{array}$$
(2)
The near-IR reddening-invariant color is independent of the value of $`R_V`$, from Schlegel, Finkbeiner, & Davis (1998) we obtain
$$Q_{JHK}=(JH)1.563(HK).$$
(3)
## 3 Expected Colors of Young Brown Dwarfs
In this section we compare the expected colors of young brown dwarfs against those of field M dwarfs of equivalent temperatures. Examination of the temperature - surface gravity (Figure 1) relations based on the BCAH98 models shows for objects at the Hydrogen Burning Limit younger than 10 Myr the effective temperature ranges between 2900 K and 3000 K and the surface gravity varies from log(g) = 3.5 to 4.2.
When the locations in Hertzsprung-Russell diagrams of spectroscopically classified and presumably co-eval stars in specific star formation regions are compared against evolutionary models the inferred ages mismatch using the spectral type to $`T_{eff}`$ scale defined for disk dwarfs. This has led to the adoption of a semi-empirical temperature scale intermediate between that of the dwarfs and the giants (Luhman, 2003). This practice is motivated by the surface gravities expected in pre-main sequence stars and results in effective temperatures several 100 K warmer at a given spectral type. Thus a $`T_{eff}`$ of 2900 K corresponds to a spectral type of M4.5 for a dwarf (Reid & Hawley, 2000) but of M6.5 for a pre-main sequence star.
The photospheric colors of M dwarfs and M giants differ at the later spectral types with the M giants becoming bluer due to surface gravity effects. This color shift with lower surface gravity is due to the change in the optical depth for the formation of molecules, most importantly TiO, that dominate the stellar spectrum. Luhman (1999) notes that for objects of later spectral types, e.g. M8, the trend is reversed with the lower surface gravity stars becoming redder longward of 8500 Å. Due to these deviations from disk dwarf colors UV-excess and blue-excess techniques based on the ZAMS become ambiguous for pre-main sequence stars later than M4 (Rebull et al., 2000).
### 3.1 Disk Dwarf Colors
Figure 2, based on the sample of West et al. (2004), shows that the $`Q_{riz}`$ color index increases towards cooler spectral types among M dwarfs, reaching a maximum value $`1.2`$ at M8. At the ages in which we expect to see accretion activity this turnover occurs for objects near 0.02 $`M_{}`$ which are too faint to be included in this survey.
To study the change in $`Q_{gri}`$ relative to the disk dwarf locus we obtain linear fits between $`Q_{gri}`$ and $`Q_{riz}`$ using photometry of M dwarfs in the SDSS Third Data Release (DR3; Abazajian et al. (2005)) obtaining
$$Q_{gri}=\{\begin{array}{cc}3.371Q_{riz}+1.119\hfill & R_V\text{ = 3.1,}\hfill \\ 2.241Q_{riz}+1.163\hfill & R_V\text{ = 5.5}\hfill \end{array}$$
(4)
using a minimum threshold to $`Q_{riz}`$ of 0.5 to avoid the slope change in the color-color diagram at the earliest M dwarfs. The fit residuals have RMS values of 0.148 and 0.117 for $`R_V`$ = 3.1 and 5.5, respectively.
The $`R_V`$ reddening-invariant two-color diagram is shown in Figure 3 with the theoretical isochrones of Girardi et al. (2004) overplotted. These isochrones are based on the “AMES” (4000 K $`T_{eff}`$ 2800 K) and “AMES-dusty” (2800 K $`>T_{eff}`$ 500 K) model atmospheres of Allard, Hauschildt, & Schwenke (2000) for cool stars and brown dwarfs. While the latter cover log(g) between 3.5 and 6.0, the $`ugriz`$ bolometric corrections were only computed for the warmer “AMES” models having log(g) of 5.5 and 6.0.
The fit residuals ($`\mathrm{\Delta }Q_{gri}`$) are plotted against $`Q_{riz}`$ for the model isochrones and the spectroscopic sample of West et al. (2004). For the purposes of illustration this M dwarf sample is subdivided at M4.5 with the two resulting histograms in color-color space scaled to match peak counts. As is seen in Figure 4 the observed reddening-invariant colors are well-described by the log(g) = 5.0 and 5.5 models appropriate for the disk dwarfs. The blueward trends in both $`Q_{riz}`$ and $`\mathrm{\Delta }Q_{gri}`$ increasing with decreasing surface gravity as shown by the log(g) = 4.5, 4.0, and 3.5 models. The dotted lines connect models of $`T_{eff}=`$ 2500 K, 2700 K, 2900 K, and 3100 K. We refer to $`\mathrm{\Delta }Q_{gri}`$ as the observed $`Q_{gri}`$ offset since it is measured relative to the $`(Q_{riz},Q_{gri})`$ linear fit. As described below, understanding the effect of low surface gravity on $`\mathrm{\Delta }Q_{gri}`$ is crucial for estimating the $`g`$ band veiling from reddening-invariant colors.
### 3.2 Veiling Models
The identification of accretion shocks by photometry involves detection of the additional continuum flux generated by the reprocessing of the high temperature shock back along the accretion column. The veiling produced by this excess continuum is defined as the ratio of the shock flux to that of the stellar photosphere within a specified wavelength range or filter passband.
The observed increase in brightness in an arbitrary filter $`x`$ is
$$\mathrm{\Delta }m_x=2.5\mathrm{log}(1+r_x)$$
(5)
where $`r_x`$ is the veiling in $`x`$ and $`m_x`$ is the observed magnitude. The observed veilings in low mass T Tauris are significantly higher in the near-ultraviolet and blue regions of the spectrum than in the red. For all but the mostly heavily veiled sources the contribution of the accretion continuum to the photosphere is minimal at wavelengths above 7000Å (Muzerolle et al., 2003) and therefore should be minimal in the SDSS $`i`$ and $`z`$ bands. This is a contrast effect between the $`10^4`$ K shock emission and the $``$ 3000 K star.
The change in a reddening-invariant color due to veiling is then
$`\mathrm{\Delta }Q_{xyz}`$ $`=`$ $`2.5\{\mathrm{log}[(1+r_x)/(1+r_y)]+`$ (6)
$`\mathrm{log}[(1+r_y)/(1+r_z)]E(xy)/E(yz)\}.`$
Our veiling models are characterized by the veilings in the SDSS $`g`$ and $`r`$ band, $`r_g`$ and $`r_r`$. We explicitly assume that $`r_i`$ and $`r_z`$ are zero (cf. supra) and obtain
$$\mathrm{\Delta }Q_{gri}=2.5\{\mathrm{log}[(1+r_g)/(1+r_r)]\mathrm{log}[1+r_r]E(gr)/E(ri)\}$$
(7)
and
$$\mathrm{\Delta }Q_{riz}=2.5\{\mathrm{log}[1+r_r]\}.$$
(8)
In the general case where $`r_r>0`$ we can not deduce the photospheric spectral type from $`Q_{riz}`$. However, since veiling makes stars look bluer (earlier, more massive), we can use the observed $`Q_{riz}`$ to deduce the earliest possible spectral type. For veiling models consisting of a non-zero $`r_g`$ and $`r_r`$ we can examine the variations in the $`Q_{riz}`$ and $`Q_{gri}`$ color. Figure 5 shows the color shifts for $`r_g`$ varying continuously from -0.5 to 1 and $`r_r/r_g`$ = 0.00, 0.25, 0.50, 0.75, and 1.00. The negative veiling values are used to model the cool spots seen in Weak-lined T Tauris.
When an unknown veiling is present in the $`r`$ band, the key diagnostic feature in the $`(Q_{riz},Q_{gri})`$ color-color diagram is the vertical displacement of the star from the stellar locus along the $`Q_{gri}`$ reddening-invariant index. This is shown in Figure 6 as a function of $`r_g`$ assuming the linear fits to the locus given above. In the case of weak to moderate veiling, $`r_g<1`$, we see that differing the $`r_r`$ to $`r_g`$ ratio only results in a minimal spread of derived $`r_g`$ for a given $`\mathrm{\Delta }Q_{gri}`$. We can therefore use the observed $`\mathrm{\Delta }Q_{gri}`$ as a proxy for continuum veiling even when the spectral type and $`Q_{riz}`$ color are not known a priori assuming that the intrinsic $`Q_{gri}`$ offset due to low surface gravity can be inferred.
## 4 Candidate Selection
Possible low mass and very low mass accreting stars were selected on the basis of color, location in color-magnitude diagrams, and variability. An estimate of contamination was obtained by applying these criteria to a neighboring field.
### 4.1 Color and Magnitude Criteria
We adopted selection criteria $`(ri)>0.6`$ and $`Q_{riz}>0.35`$ to target M dwarf colored stars. This corresponds to $`T_{eff}<`$ 3700 K or spectral types later than M0.5 in the field and M1 on the Luhman (2003) scale. The maximum stellar mass inferred from the 2 Myr Baraffe et al. (1998) isochrone is 0.8 M.
To reduce the scatter in the reddening invariant colors we required photometric errors less than 0.05 magnitudes in $`grizJH`$. This was relaxed for the $`K_S`$ band to extend the detection limit. The resulting magnitude and error limits are summarized in Table 3. The limiting $`J`$ band magnitude of 15.5 corresponds to a lower mass limit of 0.05 M ($``$ M7) based on the BCAH98 2 Myr isochrone. Propagation of these errors into the expressions for the reddening-invariant colors yields maximum errors for $`Q_{gri}`$ of 0.088 and 0.075 and for $`Q_{riz}`$ of 0.045 and 0.046, each in the cases of $`R_V`$ = 3.1 and 5.5, respectively. The maximum error in $`Q_{JHK}`$ is 0.159 which is dominated by the $`K_S`$ error of 0.1 at the faint limit.
In order to remove background stars we required that candidates were brighter than the main sequence placed at the distance of Orion in both the $`(ri,r)`$ amd $`(iz,i)`$ color-magnitude diagrams. The color-absolute magnitude relations in the SDSS passbands for late K and M dwarfs are $`M_r=(6.110\pm 0.418)+(3.800\pm 0.320)(ri)`$ and $`M_i=(5.936\pm 0.496)+(6.308\pm 0.803)(iz)`$ (Golimowski et al. 2005 in prep.). The errors in these relations are dominated not by photometric uncertainly (which is typically 1-2%), but by the intrinsic cosmic scatter of the sample.
### 4.2 Variability
Due to the non-periodic variations of low mass CTTS and the irregular sampling by the SDSS we expect the resulting lightcurves to be fairly random. A truly random sequence has a standard deviation ($`\sigma `$) equal to 1.0/$`\sqrt{12}`$ of the peak-to-peak amplitude. Given a finite number of samples and a probability that the true value of $`\sigma `$ exceeds a specific threshold we need to find the corresponding minimum threshold for the observed $`\sigma `$ that indicates significant variability. We generated a series of Monte Carlo simulations to determine the maximum $`\sigma `$ in a random sequence given an intrinsic $`\sigma `$ of the parent population and a limited number of observations. The scaling factors presented in Table 4 reflect the maximum $`\sigma `$ seen in 99% of the simulated data sets using 10,000 realizations per point. In the limit of large number of observations the scaling factors tend towards unity.
We identify variable stars on the basis of the $`g`$ band $`\sigma `$ exceeding both a threshold of 0.05 magnitudes (scaled by the correction factor defined above), and three times the photometric error in $`g`$. Of these we also select a subset that meet the same conditions applied to the $`z`$ band. This $`z`$ band variability criterion matches the typical boundary between Type I and Type II $`I`$ band variations in very low mass stars and young brown dwarfs.
### 4.3 Results
On the basis of these criteria we identify 507 stars that are significantly variable in the $`g`$ band of which 215 also meet the criteria for $`z`$ band variability. In Figures 7 and 8 we see that the majority of the variable stars form a locus parallel to and roughly 2 magnitudes above the main sequence placed at the $`mM=8.2`$ distance of the Orion star formation region. The spatial distribution of these low mass candidates is concentrated in the Orion OB1a and OB1b associations and the NGC 2068/2071 protocluster (Figure 9). The Orion OB1c subassociation members comprise a dispersed population in this region.
### 4.4 Comparison Field
In order to assess the possible contamination of our Orion sample by variable field stars we apply the criteria defined above to a neighboring equatorial field bounded by $`40\mathrm{°}<\alpha <70\mathrm{°}`$ and $`1.25\mathrm{°}<\delta <1.25\mathrm{°}`$ which encompasses three times the area of our survey. This is a complex region containing at least three populations of young stars which we now describe.
On the basis of ROSAT All-Sky Survey (Trümpler, 1983) detections Neuhaüser (1997) subsequently identified a population of $``$ 30 Myr low mass WTTS south of the Taurus-Auriga cloud. While we expect to see Class I variability in these stars signs of active accretion or the presence of inner circumstellar disks should not be evident due to the typical 10 Myr disk lifetime. We also expect to detect the background 8 Myr old Gould Belt stars (Guillout et al., 1998) which may include CTTS. As the Gould Belt midplane crosses our equatorial survey area near $`\alpha =70`$ we suspect some Gould Belt stars might be found in our Orion sample.
The third star formation complex is that associated with the nearby (80 $`\pm `$ 20 pc) translucent cloud MBM 18 (Larson & Reed, 2003). This high latitude cloud contains compact regions of CO(1-0) emission with line intensities above 3 K. If MBM 18 is an active star formation region then due to its proximity we expect all but the lowest massed young stars to be too bright for the SDSS,
Application of our Orion CTTS criteria results in 133 $`g`$ band variables of which 19 are also variable in the $`z`$ band. After scaling by the ratio of the survey areas these imply possible contamination fractions of 8.7% and 2.9% for the Orion $`g`$ band variables and its $`z`$ band variable subset, respectively.
## 5 Empirical Results For The Orion Population
### 5.1 Veiling Signatures
In Figure 10 we show the reddening-invariant $`Q_{riz}`$ color shift between the faint and the bright states defined by the $`g`$ band brightness against the 2MASS $`J`$ magnitude. The observed trend is for $`Q_{riz}`$ to decrease when the star is brighter in $`g`$ which is consistent with the presence of veiling in the $`r`$ band. In this figure the horizontal dotted line marks the substellar boundary based on the BCAH98 2 Myr isochrone placed at m-M = 8.22. We convert the BCAH98 $`JHK`$ magnitudes from the CIT to the 2MASS system following Carpenter (2001). The derived effective temperature from $`Q_{riz}`$ changes by approximately 100 K for every 0.1 change in the color (see Figure 2). Due to this uncertainty in determination of the photospheric $`Q_{riz}`$ color we will present all color and variability trends against the 2MASS $`J`$ magnitude which should be relatively unaffected by both shock emission and disk thermal emission.
In Figure 11 we see that the bright state $`\mathrm{\Delta }Q_{gri}`$ during both faint and bright states slightly decreases in the less massive (fainter) objects. However, when the observed $`\mathrm{\Delta }Q_{gri}`$ are compared against those predicted from the cooler ($``$ 2700 K) low surface gravity models (dot-dash line) it is clear that the bulk of veiling signature is probably due to the shifts in photospheric colors. Thus, at the substellar boundary we only see evidence for weak $`g`$ band veiling which consistent with the results of Muzerolle et al. (2003) who find a deficit of continuum veiling at 5500 Å for accreting stars of spectral types M5 to M7.
### 5.2 Variability
If the veiling in a specific band $`x`$ ranges from $`r_x^{low}`$ to $`r_x^{high}`$ then the observed magnitude change is
$$\mathrm{\Delta }m_x=2.5log\left(\frac{1+r_x^{high}}{1+r_x^{low}}\right).$$
(9)
This scales montonically, albeit non-linearly, with $`r_x^{high}`$, if
$$\frac{dr_x^{low}}{dr_x^{high}}\frac{1+r_x^{low}}{1+r_x^{high}}$$
(10)
which in the case where $`r_x^{low}`$ scales with $`r_x^{high}`$, i.e. $`r_x^{low}=\alpha r_x^{high}`$ requires $`\alpha 1`$, which is implied by $`r_x^{low}r_x^{high}`$. In general we then expect the amplitude and RMS of the variability to be related to the veiling and hence with the shock flux in that band.
In Figure 12 we see that the maximum RMS of the $`g`$ band variations peaks at $`J13`$ and that no high-amplitude ($`\sigma _g>0.5)`$ variables are seen below $`J=14`$. The triangles mark the median $`\sigma _g`$ values computed in for one-magnitude bins spanning $`11<J<16`$. We also see that the maximum ratio of the $`g`$ band to $`z`$ band RMS ($`\sigma g`$/$`\sigma z`$) peaks at the same magnitude and then decreases in the fainter stars (Figure 13). Only five faint ($`J>14`$) candidate accreting stars have $`\sigma _g/\sigma _z>3`$. This trend is perhaps indicative of cooling of the accretion shock temperature with diminished mass.
For the intermediate mass Classical T Tauris studied by Gullbring et al. (1998), the Balmer jump diminishes as the veiling and shock bolometric luminosity increase. Conversely, for the weakly accreting systems expected in very low mass stars and young brown dwarfs the Balmer jump should be relatively large suggesting that the bulk of the shock emission will be in the near-ultraviolet rather than the optical bands. The shock models of Muzerolle et al. (2003) suggest that shocks with lower energy flux carried by the accreting column emit more at redder wavelengths due to quasi-blackbody emission from the heated photosphere. The trend we see is consistent with a decrease in the shock energy flux with decreasing mass.
### 5.3 Near-Infrared Excess
Near-Infrared ($`JHK`$) photometry has often been used to detect the presence of excess emission primarily in the $`K`$ band due to thermal emission from the inner circumstellar disk. Meyer, Calvet, & Hillenbrand (1997) find a Classical T Tauri locus in $`(JH,HK)`$ based on the dereddened colors of K7/M0 intermediate massed stars in Taurus. In Figure 14 we see that the $`K_S`$ excess indicating either a disk signatures or significant extinction drops for objects having $`J<14`$. For the 2 Myr BCAH98 isochrone this corresponds to 0.11 M and $`T_{eff}`$ = 3000 K ($``$ M5.5). This lack of signature in $`JHK`$ is expected for substellar objects since they are less luminous and so the inner circumstellar disk will be cooler. As discussed by Liu, Najita, & Tokunaga (2003) this contrast problem between the disk and the stellar photosphere motivates searching for young brown dwarf disks in the $`L`$ and longer wavelength bands.
## 6 Summary
Analysis of 2MASS and multi-epoch SDSS photometry of low mass pre-main sequence objects stars within the Orion OB1 assocation shows that:
1. $`g`$ band variations having a standard deviation greater than 0.05 magnitudes persist for objects having reddening-invariant colors redder than, and $`J`$ magnitudes fainter than, that expected for the stellar/substellar boundary.
2. While identification of $`g`$ band veiling for spectral types later than M4 is hindered by our uncertain knowledge of surface gravity effects on photospheric colors, we see evidence for accretion signatures even in these very low mass objects. The general trend is for the veiling signature to diminish with decreasing temperature (and mass).
3. No disk signatures based on $`JK`$ color excess are seen for spectral types later than M4. We infer based on variability and weak veiling that the accretion process and hence the disks are present, and that relative contribution to the $`K`$ band flux is weak. This is consistent with earlier work on young brown dwarf disks.
We have demonstrated that multi-epoch and multi-band imaging surveys enable identification and classification of pre-main sequence objects based on variability and reddening-invariant colors. While use of the SDSS $`g`$ band was mandated by the sensitivity issues and the red-leak in the imager’s filter system, inclusion of a broad-band filter that is shortward of the Balmer jump such the SDSS $`u`$ in the next generation of surveys, i.e. the Large Synoptic Survey Telescope or LSST, will increase the sensitivity to the veiling emission from the accretion shock for the later and less massive objects.
We are supplementing this imaging survey with an SDSS spectroscopic program targeting low mass stars in Orion and the Taurus-Auriga star formation region. Results from this study in addition to analysis of the veiling trends presented here and their constraints on shock luminosities and mass accretion rates will be discussed in subsequent work.
We thank the anonymous referee for comments that have greatly improved this paper. PMM and JAS acknowledge support from LANL Laboratory Directed Research and Development (LDRD) program 20030486DR. Funding for the Sloan Digital Sky Survey (SDSS) has been provided by the Alfred P. Sloan Foundation, the Participating Institutions, the National Aeronautics and Space Administration, the National Science Foundation, the U.S. Department of Energy, the Japanese Monbukagakusho, and the Max Planck Society. The SDSS is a joint project of the University of Chicago, Fermilab, the Institute for Advanced Study, the Japan Participation Group, The Johns Hopkins University, the Korean Scientist Group, Los Alamos National Laboratory, the Max-Planck-Institute for Astronomy (MPIA), the Max-Planck-Institute for Astrophysics (MPA), New Mexico State University, the University of Pittsburgh, Princeton University, University of Portsmouth, the United States Naval Observatory, and the University of Washington. Facilities: SDSS, 2MASS
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# Medium modifications of nucleon electromagnetic form factors
## 1 Introduction
The structure of the nucleon and its modifications in the nuclear medium is a very active field of experimental and theoretical research. The basic quantities, which reflect the charge and current distributions in the nucleon, are the electromagnetic form factors , which are currently investigated in elastic electron- nucleon scattering experiments from intermediate to very high energies . The knowledge of the nucleon form factors is also inevitable to understand the electromagnetic structure of nuclei. Electron-nucleus scattering experiments under quasielastic kinematic conditions, like the measurement of inclusive response functions in the intermediate energy region and recent measurements of polarization transfer in semi-exclusive knock-out processes , are ideal places to study the form factors of a nucleon bound in the nuclear medium. Because the structure of the quark core and the surrounding meson cloud may be different for a bound nucleon and a free nucleon, one expects medium modifications of the nucleon form factors, and the exploration of these effects is an important subject at current electron accelerator facilities .
On the theoretical side, effective quark theories are the ideal tools to describe the electromagnetic form factors of the nucleon. Much progress for the case of the free nucleon has been made in Faddeev type descriptions based on the Schwinger-Dyson method . An important point which still has to be implemented in these calculations is the role of the pion cloud around the nucleon, and the recently developed method of chiral extrapolations of lattice results provides important hints. On the other hand, the calculation of form factors at finite nucleon density requires also a description of the equation of state of the many-nucleon system, and here progress has been made by using the Nambu-Jona-Lasinio (NJL) model as an effective quark theory: Recent works have shown how to account for the saturation properties of nuclear matter in this model , and when combined with the quark-diquark description of the single nucleon this provides a successful description of both nucleon and nuclear structure functions for deep inelastic scattering <sup>1</sup><sup>1</sup>1Recently the model has been extended to describe the equation of state at high densities..
The purpose of this paper is to discuss the results for the nucleon form factors obtained in the simple quark-scalar diquark description of the nucleon at finite density in the NJL model. We have to note from the beginning that this can only be a first step toward a realistic description, because it is known that axial vector diquarks are important for spin-dependent quantities, and the pion cloud is important for magnetic moments and the size of the nucleon. While the axial diquarks could be included in a further step like it was done for the structure functions, a reliable description of pion cloud effects makes it necessary to go beyond the standard ladder approximation scheme. However, like the simple quark-scalar diquark model of Ref. served as a basis for the more elaborate description of structure functions , it will also be the basis of a more realistic description of form factors including axial diquarks and the pion cloud. To provide this basis is the main intention of the present paper.
In Sect. 2 we will briefly review the model for the nucleon and the nuclear matter equation of state. Sect. 3 is devoted to the nucleon form factors at finite density, and in Sect.4 we discuss the numerical results. As an application, we discuss the response function for quasielastic electron scattering in Sect.5. For this purpose we will also elucidate the nucleon-nucleon interaction in our model in order to include the correlations within the relativistic RPA. A summary will be presented in Sect. 6.
## 2 The model
In this work, we use the NJL model as an effective quark theory to describe the nucleon as a quark-diquark bound state, and nuclear matter (NM) in the mean field approximation. The details are explained in Refs. , and here we will only briefly summarize those points which will be needed for our calculations.
The NJL model is characterized by a chirally symmetric 4-fermi interaction between the quarks. Any such interaction can be Fierz symmetrized and decomposed into various $`q\overline{q}`$ channels . Writing out explicitly only those channels which are relevant for our present discussion, we have
$$=\overline{\psi }\left(i\overline{)}m\right)\psi +G_\pi \left(\left(\overline{\psi }\psi \right)^2\left(\overline{\psi }(\gamma _5𝝉)\psi \right)^2\right)G_\omega \left(\overline{\psi }\gamma ^\mu \psi \right)^2+\mathrm{}$$
(2.1)
where $`m`$ is the current quark mass. In a mean field description of the isospin symmetric nuclear matter ground state $`|\rho `$, the Lagrangian can be expressed as
$`=\overline{\psi }\left(i\overline{)}M\overline{)}V\right)\psi {\displaystyle \frac{(Mm)^2}{4G_\pi }}+{\displaystyle \frac{V_\mu V^\mu }{4G_\omega }}+_I,`$ (2.2)
where $`M=m2G_\pi \rho |\overline{\psi }\psi |\rho `$ and $`V^\mu =2G_\omega \rho |\overline{\psi }\gamma ^\mu \psi |\rho `$, and $`_I`$ is the normal ordered interaction Lagrangian. The effect of the mean scalar field is thus included in the density-dependent constituent quark mass $`M`$, and the effect of the mean vector field is to shift the quark momentum according to $`p^\mu =p_Q^\mu +V^\mu `$, where $`p_Q^\mu `$ is the kinetic momentum. The propagator of the constituent quark therefore has the following dependence on the mean vector field <sup>2</sup><sup>2</sup>2In this section, Green functions in the presence of the mean vector field are denoted by a tilde, and those without the vector field have no tilde. In the loop integrals for the electromagnetic form factors in Sect.3, however, it is always possible to eliminate the vector field by a shift of the integration variable, and therefore the tilde-Green functions do not appear in later sections. : $`\stackrel{~}{S}(k)=S(k_\mathrm{Q})`$.
One can use a further Fierz transformation to decompose $`_I`$ into a sum of $`qq`$ channel interaction terms . For our purposes we need only the interaction in the scalar diquark ($`J^\pi =0^+,T=0`$, color $`\overline{3}`$) channel:
$`_{I,s}=G_s\left(\overline{\psi }\left(\gamma _5C\right)\tau _2\beta ^A\overline{\psi }^T\right)\left(\psi ^T\left(C^1\gamma _5\right)\tau _2\beta ^A\psi \right),`$ (2.3)
where $`\beta ^A=\sqrt{3/2}\lambda ^A(A=2,5,7)`$ are the color $`\overline{3}`$ matrices and $`C=i\gamma _2\gamma _0`$. The coupling constant $`G_s`$ will be determined so as to reproduce the free nucleon mass.
The reduced t-matrix in the scalar diquark channel is given by
$`\stackrel{~}{\tau }_s(q)={\displaystyle \frac{4iG_s}{1+2G_s\stackrel{~}{\mathrm{\Pi }}_s(q)}}=\tau _s(q_\mathrm{D})`$ (2.4)
with the scalar $`qq`$ bubble graph
$`\stackrel{~}{\mathrm{\Pi }}_s(q)=6i{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\mathrm{tr}_D\left[\gamma _5S(k)\gamma _5S\left((qk)\right)\right]}=\mathrm{\Pi }_s(q_\mathrm{D}).`$ (2.5)
Here $`q_\mathrm{D}^\mu =q^\mu 2V^\mu `$ is the kinetic momentum of the diquark.
The relativistic Faddeev equation in the NJL model can been solved numerically for the free nucleon, but here we restrict ourselves to the static approximation, where the momentum dependence of the quark exchange kernel is neglected. The solution for the quark-diquark t-matrix then takes the simple analytic form
$`\stackrel{~}{T}_N(p)={\displaystyle \frac{3}{M}}{\displaystyle \frac{1}{1+\frac{3}{M}\stackrel{~}{\mathrm{\Pi }}_N(p)}}=T_N(p_\mathrm{N}),`$ (2.6)
with the quark-diquark bubble graph given by
$`\stackrel{~}{\mathrm{\Pi }}_N(p)={\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\stackrel{~}{S}(k)\stackrel{~}{\tau }_s(pk)}=\mathrm{\Pi }_N(p_\mathrm{N}),`$ (2.7)
where $`p_\mathrm{N}^\mu =p^\mu 3V^\mu `$ is the kinetic momentum of the nucleon. The nucleon mass $`M_N`$ is defined as the pole of (2.6) at $`\overline{)}p_\mathrm{N}=M_N`$, and the positive energy spectrum has the form $`p_0=ϵ_pE_{Np}+3V_0`$ with $`E_{Np}=\sqrt{M_N^2+𝒑^2}`$. The nucleon vertex function in the non-covariant normalization is defined by the pole behavior of the quark-diquark t-matrix:
$`T_N(p){\displaystyle \frac{\mathrm{\Gamma }_N(p)\overline{\mathrm{\Gamma }}_N(p)}{p_0ϵ_p}}\mathrm{as}p_0ϵ_p.`$ (2.8)
From this definition and Eq.(2.6) one obtains
$`\mathrm{\Gamma }_N(p)=\sqrt{Z_N{\displaystyle \frac{M_N}{E_{Np}}}}u_N(p_\mathrm{N}),`$ (2.9)
where $`u_N`$ is a free Dirac spinor for mass $`M_N`$ normalized as $`\overline{u}_Nu_N=1`$. The normalization factor $`Z_N`$ is easily obtained from this definition and will be given in Eq.(LABEL:zn) below. We note that with this normalization the vertex function satisfies the relation
$`\overline{\mathrm{\Gamma }}_N(p)\left({\displaystyle \frac{\mathrm{\Pi }_N(p)}{p_\mu }}\right)\mathrm{\Gamma }_N(p)={\displaystyle \frac{p^\mu }{E_{Np}}}`$ (2.10)
In the numerical calculations of this paper, we will approximate the quantity $`\tau _s`$ by a “contact+pole” form:
$`\tau _s(q)4iG_s{\displaystyle \frac{ig_s}{q^2M_s^2}}4iG_sig_s\mathrm{\Delta }_{Fs}(q).`$ (2.11)
Here $`\mathrm{\Delta }_{Fs}(q)`$ is the Feynman propagator for a scalar particle of mass $`M_s`$, which is defined as the pole of $`\tau _s`$ of Eq.(2.4). The residue at the pole ($`g_s`$) will be given in Eq. (LABEL:gs) below.
In the calculation of the nucleon form factors, we will also consider the effects of the pion cloud around the constituent quarks. In this case, the propagator $`S(p)`$ of the quark involves an additional self energy correction from the pion cloud ($`\mathrm{\Sigma }_Q`$). Here we will use a simple pole approximation for $`S(p)`$:
$`S(p)=Z_QS_F(p),Z_Q^1=1\left({\displaystyle \frac{\mathrm{\Sigma }_Q}{\overline{)}p}}\right)_{\overline{)}p=M},`$ (2.12)
where $`S_F`$ is the Feynman propagator of a constituent quark with mass $`M`$. In this approximation, the pion effects can be renormalized by $`\psi \sqrt{Z_Q}\psi `$ and a redefinition of the four fermi coupling constants $`G_\alpha G_\alpha /Z_Q^2`$, see Ref.(). For the calculation of the form factors, however, we will keep the factor $`Z_Q`$ explicitly for clarity <sup>3</sup><sup>3</sup>3We also note that such a renormalization procedure is no longer possible when one considers the pion cloud effects around the nucleon, which goes beyond the simple ladder approximation on the quark level.. Introducing (2.12) and (2.11) into the expressions (2.5) and (2.7) for the bubble graphs shows that the diquark and nucleon normalization factors can be written as
$`g_s^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{\Pi }_s(q)}{q^2}}\right)_{q^2=M_s^2}={\displaystyle \frac{1}{2}}Z_Q^2\left({\displaystyle \frac{\widehat{\mathrm{\Pi }}_s(q)}{q^2}}\right)_{q^2=M_s^2}Z_Q^2\widehat{g}_s^1`$
$`Z_N^1`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Pi }_N(p)}{\overline{)}p}}\right)_{\overline{)}p=M_N}=Z_Qg_s\left({\displaystyle \frac{\widehat{\mathrm{\Pi }}_N(p)}{\overline{)}p}}\right)_{\overline{)}p=M_N}Z_Qg_s\widehat{Z}_N^1,`$
where $`\widehat{\mathrm{\Pi }}_s`$ and $`\widehat{\mathrm{\Pi }}_N`$ are defined in terms of the pole parts only:
$`\widehat{\mathrm{\Pi }}_s(q)`$ $`=`$ $`6i{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\mathrm{tr}_D\left[\gamma _5S_F(k)\gamma _5S_F\left((qk)\right)\right]}`$ (2.15)
$`\widehat{\mathrm{\Pi }}_N(p)`$ $`=`$ $`i{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}S_F(k)\mathrm{\Delta }_{Fs}(pk)}`$ (2.16)
The equation of state of NM in the NJL model can be derived in a formal way from the quark Lagrangian (2.2) by using hadronization techniques, but in the mean field approximation the resulting energy density of isospin symmetric NM has the simple form
$`=_V{\displaystyle \frac{V_0^2}{4G_\omega }}+4{\displaystyle \frac{d^3p}{(2\pi )^3}\mathrm{\Theta }\left(p_F|𝒑|\right)ϵ_p},`$ (2.17)
where the vacuum contribution (quark loop) is
$`_V=12i{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\mathrm{ln}\frac{k^2M^2+iϵ}{k^2M_0^2+iϵ}}+{\displaystyle \frac{(Mm)^2}{4G_\pi }}{\displaystyle \frac{(M_0m)^2}{4G_\pi }}.`$ (2.18)
Here $`M_0`$ the constituent quark mass for zero nucleon density. The condition $`/V_0=0`$ leads to $`V_0=6G_\omega \rho `$, and we can eliminate the vector field in (2.17) in favor of the baryon density. The resulting expression has then to be minimized with respect to the constituent quark mass $`M`$ for fixed density. For zero density this condition becomes identical to the familiar gap equation of the NJL model, and for finite density the nonlinear $`M`$-dependence of the nucleon mass $`M_N`$ is essential to obtain saturation of the binding energy per nucleon .
In order to fully define the model one has to specify a cut-off procedure. In the calculations in this paper we will use the proper time regularization scheme , where one evaluates loop integrals over a product of propagators by introducing Feynman parameters, performing a Wick rotation and replacing the denominator ($`A`$) of the loop integral according to
$`{\displaystyle \frac{1}{A^n}}{\displaystyle \frac{1}{(n1)!}}{\displaystyle _{1/\mathrm{\Lambda }_{\mathrm{UV}}^2}^{1/\mathrm{\Lambda }_{\mathrm{IR}}^2}}d\tau \tau ^{n1}e^{\tau A}(n1),`$ (2.19)
where $`\mathrm{\Lambda }_{\mathrm{IR}}`$ and $`\mathrm{\Lambda }_{\mathrm{UV}}`$ are the infrared and ultraviolet cut-offs, respectively. The infrared cut-off plays the important role of eliminating the unphysical thresholds for the decay of the nucleon and mesons into quarks , thereby taking into consideration a particular aspect of confinement physics.
## 3 Nucleon electromagnetic form factors
The electromagnetic current of the nucleon in the quark-diquark model is represented by the Feynman diagrams of Fig.1 and given by
$`j_N^\mu (q)=\overline{\mathrm{\Gamma }}_N(p^{}){\displaystyle }{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}}[\left(S(p^{}k)\mathrm{\Lambda }_Q^\mu S(pk)\right)\tau _s(k)`$
$`+i\left(\tau _s(p^{}k)\mathrm{\Lambda }_D^\mu \tau _s(pk)\right)S(k)]\mathrm{\Gamma }_N(p).`$ (3.1)
Here $`\mathrm{\Lambda }_Q^\mu `$ and $`\mathrm{\Lambda }_D^\mu `$ are the electromagnetic vertices of the quark and the scalar diquark, both depending on the final and initial particle momenta. It is an easy task to use the Ward-Takahashi identities for the quark and diquark vertices
$`q_\mu \left(S(\mathrm{}^{})\mathrm{\Lambda }_Q^\mu S(\mathrm{})\right)`$ $`=`$ $`Q_Q\left(S(\mathrm{}^{})S(\mathrm{})\right)`$ (3.2)
$`q_\mu \left(\tau _s(\mathrm{}^{})\mathrm{\Lambda }_D^\mu \tau _s(\mathrm{})\right)`$ $`=`$ $`iQ_D\left(\tau _s(\mathrm{}^{})\tau _s(\mathrm{})\right)`$ (3.3)
to show the current and charge conservation for the nucleon:
$`q_\mu j_N^\mu (q)`$ $`=`$ $`Q_N\overline{\mathrm{\Gamma }}_N(p^{})\left(\mathrm{\Pi }_N(p^{})\mathrm{\Pi }_N(p)\right)\mathrm{\Gamma }_N(p)=0`$ (3.4)
$`j^\mu (0)`$ $`=`$ $`Q_N{\displaystyle \frac{p^\mu }{E_{Np}}},`$ (3.5)
where the electric charge of the nucleon $`Q_N=Q_Q+Q_D`$ is the sum of the quark and diquark electric charges.
The electromagnetic vertices in (3.1) describe the finite extension of the constituent quarks and the diquark. In general, they should be calculated off-shell consistently with the propagators $`\tau _s`$ and $`S`$, using Feynman diagrams or some ansatz which satisfies the Ward identities. In the present work, however, our principal aim is to investigate the medium effects in a simple model calculation. For this purpose, we limit the complications caused by the quark and diquark sizes to a minimum, and use an on-shell (or pole) approximation for the quark and diquark currents appearing in Eq.(3.1):
$`\left(S(\mathrm{}^{})\mathrm{\Lambda }_Q^\mu S(\mathrm{})\right)`$ $``$ $`Z_Q\left(S_F(\mathrm{}^{})\widehat{\mathrm{\Lambda }}_Q^\mu S_F(\mathrm{})\right)`$ (3.6)
$`\left(\tau _s(\mathrm{}^{})\mathrm{\Lambda }_D^\mu \tau _s(\mathrm{})\right)`$ $``$ $`g_s\left(\mathrm{\Delta }_{Fs}(\mathrm{}^{})\widehat{\mathrm{\Lambda }}_D^\mu \mathrm{\Delta }_{Fs}(\mathrm{})\right),`$ (3.7)
where the on-shell (o.s.) vertices are denoted by a hat and given by the pole residues of the full quantities by
$`\widehat{\mathrm{\Lambda }}_Q^\mu `$ $``$ $`Z_Q\left(\mathrm{\Lambda }_Q^\mu \right)_{\mathrm{o}.\mathrm{s}.}\gamma ^\mu F_{1Q}(q^2)+{\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M}}F_{2Q}(q^2)`$ (3.8)
$`\widehat{\mathrm{\Lambda }}_D^\mu `$ $``$ $`g_s\left(\mathrm{\Lambda }_D^\mu \right)_{\mathrm{o}.\mathrm{s}.}(\mathrm{}^{}+\mathrm{})^\mu F_D(q^2).`$ (3.9)
Here we introduced the quark and diquark form factors which satisfy $`F_{1Q}(0)=Q_Q,F_D(0)=Q_D`$.
To understand (3.6) and (3.7), we note that in general the on-shell approximation for a vertex function $`\mathrm{\Lambda }^\mu (\mathrm{}^{},\mathrm{})`$ can be formulated only if it appears between pole parts of Green functions, because only in this case one can approximate it by its value for on-shell momenta $`\mathrm{}^{},\mathrm{}`$ <sup>4</sup><sup>4</sup>4This approximation is the basis of the standard convolution formalism to calculate quark light-cone momentum distributions and structure functions.. This is why in Eq.(3.6) and (3.7) we have replaced the propagators left and right to the vertex functions by their pole parts (see (2.11) and (2.12)), which is also essential in order to have charge conservation with the vertices (3.8) and (3.9).
We now can write down the form of the nucleon current (3.1) which will be used in the further calculations:
$`j_N^\mu (q)=\sqrt{{\displaystyle \frac{M_N}{E_{Np}}}{\displaystyle \frac{M_N}{E_{Np^{}}}}}\overline{u}_N(p^{})\left(𝒪_C^\mu +𝒪_Q^\mu +𝒪_D^\mu \right)u_N(p).`$ (3.10)
Here the first and second terms denote the contributions of the contact term and the pole term of the diquark t-matrix to the quark diagram (first diagram of Fig.1), and the third term is the contribution from the diquark diagram:
$`𝒪_C^\mu `$ $`=`$ $`{\displaystyle \frac{4iG_s}{\widehat{g}_s}}\widehat{Z}_N{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}S_F(p^{}k)\left(\gamma ^\mu F_{1Q}(q^2)+\frac{i\sigma ^{\mu \nu }q_\nu }{2M}F_{2Q}(q^2)\right)S_F(pk)}`$
$`𝒪_Q^\mu `$ $`=`$ $`i\widehat{Z}_N{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}S_F(p^{}k)\left(\gamma ^\mu F_{1Q}(q^2)+\frac{i\sigma ^{\mu \nu }q_\nu }{2M}F_{2Q}(q^2)\right)S_F(pk)\mathrm{\Delta }_{Fs}(k)}`$
(3.12)
$`𝒪_D^\mu `$ $`=`$ $`i\widehat{Z}_NF_D(q^2){\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\mathrm{\Delta }_{Fs}(p^{}k)\left(p+p^{}2k\right)^\mu \mathrm{\Delta }_{Fs}(pk)S_F(k)}.`$ (3.13)
(For the contact term, we replaced $`G_sG_s/Z_Q^2`$, so that $`G_s`$ in (LABEL:contact) is the renormalized coupling in the sense explained in Sect.2.) By using the elementary Ward-Takahashi identities $`\overline{)}q=S_F^1(\mathrm{}^{})S_F^1(\mathrm{})`$ and $`\mathrm{}^2\mathrm{}^2=\mathrm{\Delta }_{Fs}^1(\mathrm{}^{})\mathrm{\Delta }_{Fs}^1(\mathrm{})`$ and the fact that on the nucleon mass shell $`\mathrm{\Pi }_N(p)=\mathrm{\Pi }_N(p^{})`$, it is easy to check that the 3 parts in (3.10) satisfy current conservation separately<sup>5</sup><sup>5</sup>5These formal manipulations involve shifts of the integration variables. In the actual calculations based on our regularization scheme, however, we always checked that current and charge conservation are satisfied exactly. (Therefore, the explicit expressions given in the Appendices all satisfy the charge conservation.). Similarly, charge conservation can be checked by using the elementary Ward identities $`\gamma ^\mu =S_F^1(\mathrm{})/\mathrm{}^\mu `$ and $`2\mathrm{}^\mu =\mathrm{\Delta }_{Fs}^1(\mathrm{})/\mathrm{}^\mu `$, as well as $`F_{1Q}(0)=Q_Q`$, $`F_D(0)=Q_D`$. It has to be noted, however, the general Ward-Takahashi identity for an off-shell nucleon (the first equality in Eq.(3.4) without the nucleon spinor) is not valid in this approximation scheme.
It is not very difficult to evaluate the three loop integrals in (LABEL:contact)-(3.13), and in Appendix A the results are given in terms of the Dirac-Pauli form factors $`F_{1N}`$ and $`F_{2N}`$, which are defined by
$`j_N^\mu `$ $`=`$ $`\sqrt{{\displaystyle \frac{M_N}{E_{Np}}}{\displaystyle \frac{M_N}{E_{Np^{}}}}}\overline{u}_N(p^{})\left[\gamma ^\mu F_{1N}(q^2)+{\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M_N}}F_{2N}(q^2)\right]u_N(p).`$
In the following we will discuss various steps in the calculation of the nucleon form factors.
### 3.1 Naive quark-diquark model
The simplest approximation consists in assuming point couplings of the quarks and diquarks to the photon, i.e., replacing $`F_{1Q}Q_Q={\displaystyle \frac{1}{6}}+{\displaystyle \frac{\tau _3}{2}}`$, $`F_{2Q}0`$, $`F_DQ_D={\displaystyle \frac{1}{3}}`$ in Eqs. (LABEL:contact)-(3.13). This approximation will be called the “naive quark-diquark model”, and the detailed expressions can be found in Appendix A.
### 3.2 Effects of finite diquark size
Here we consider the effect of the diquark form factor $`F_D`$, which has been defined in (3.9). The vertex $`\mathrm{\Lambda }_D^\mu `$ is shown graphically in Fig.2, where the quark-diquark vertex functions are those appearing in the Lagrangian(2.3). We obtain
$`\widehat{\mathrm{\Lambda }}_D^\mu =ig_s{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\left\{\mathrm{Tr}\left[\gamma _5S(p^{}+k)\mathrm{\Lambda }_Q^\mu S(p+k)\gamma _5S(k)\right]\right\}},`$ (3.15)
where the trace refers to color, isospin and Dirac indices. Using the on-shell approximation (3.6), the definition of quark form factors (3.8), and also (2.12) and (LABEL:gs), we obtain
$`\widehat{\mathrm{\Lambda }}_D^\mu `$ $`=`$ $`6i\widehat{g}_s{\displaystyle }{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}}\{F_{1Q}^{(0)}(Q^2)\mathrm{Tr}_D\left[\gamma _5S_F(p^{}+k)\gamma ^\mu S_F(p+k)\gamma _5S_F(k)\right]`$ (3.16)
$`+`$ $`F_{2Q}^{(0)}(Q^2)\mathrm{Tr}_D\left[\gamma _5S_F(p^{}+k){\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M}}S_F(p+k)\gamma _5S_F(k)\right]\}.`$
Here $`F_{1Q}^{(0)}`$ and $`F_{2Q}^{(0)}`$ are the isoscalar parts of the quark form factors $`F_{1Q}`$ and $`F_{2Q}`$. The resulting diquark form factor $`F_D`$ is given in Appendix A.
### 3.3 Effects of meson cloud around constituent quarks
Here we consider the quark form factors arising from the pion cloud and vector mesons. For the pion cloud, we obtain from the definition (3.8) and the Feynman diagrams of Fig.3,
$`\widehat{\mathrm{\Lambda }}_Q^\mu `$ $`=`$ $`Z_Q\gamma ^\mu Q_Q+Z_Q{\displaystyle }{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}}[\gamma _5\tau _i\left(S(p^{}k)\mathrm{\Lambda }_{Q0}^\mu S(pk)\right)\tau _i\gamma _5\tau _\pi (k)`$ (3.17)
$`+`$ $`i\left(\tau _\pi (p^{}k)\mathrm{\Lambda }_\pi ^\mu \tau _\pi (pk)\right)\gamma _5S(k)\gamma _5].`$
Here $`\tau _\pi `$ is the reduced $`q\overline{q}`$ t-matrix is the pion channel, which can be approximated in the same way as the diquark t-matrix (2.11):
$`\tau _\pi (k)`$ $`=`$ $`{\displaystyle \frac{2iG_\pi }{1+2G_\pi \mathrm{\Pi }_\pi (k^2)}}2iG_\pi +{\displaystyle \frac{ig_\pi }{k^2M_\pi ^2}}2iG_\pi +ig_\pi \mathrm{\Delta }_{F\pi }(k).`$
(3.18)
The bubble graph in the pion channel is $`\mathrm{\Pi }_\pi =\mathrm{\Pi }_s`$ (see (2.5)), and the residue $`g_\pi `$ will be given in Eq.(3.24) below.
The expression (3.17) is formally similar to the nucleon current (3.1). The presence of the first term simply expresses the fact that in the NJL model the (bare) quarks are present from the beginning, and the factor $`Z_Q`$ gives the probability of having a constituent quark without its pion cloud. It has been expressed in Eq.(2.12) in terms of the self energy
$`\mathrm{\Sigma }_Q(p)=3{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\left(\gamma _5S(pk)\gamma _5\right)\tau _\pi (k)}.`$ (3.19)
The quark electromagnetic vertex $`\mathrm{\Lambda }_{Q0}^\mu `$ in (3.17) will be approximated by its point form after processing the renormalization factor $`Z_Q`$, and the pion electromagnetic vertex $`\mathrm{\Lambda }_\pi ^\mu \tau _i\mathrm{\Lambda }_{\pi ,ij}^\mu \tau _j`$ is similar to the diquark vertex of Fig.2, but with point quark-photon couplings, as will be specified below.
We now follow the same steps as for the calculation of the nucleon current: We use the on-shell approximation for the quark and pion vertices
$`\left(S(\mathrm{}^{})\mathrm{\Lambda }_{Q0}^\mu S(\mathrm{})\right)`$ $``$ $`Z_Q\left(S_F(\mathrm{}^{})\widehat{\mathrm{\Lambda }}_{Q0}^\mu S_F(\mathrm{})\right)`$ (3.20)
$`\left(\tau _\pi (\mathrm{}^{})\mathrm{\Lambda }_\pi ^\mu \tau _\pi (\mathrm{})\right)`$ $``$ $`g_\pi \left(\mathrm{\Delta }_{F\pi }(\mathrm{}^{})\widehat{\mathrm{\Lambda }}_\pi ^\mu \mathrm{\Delta }_{F\pi }(\mathrm{})\right),`$ (3.21)
where the on-shell vertices are defined by
$`\widehat{\mathrm{\Lambda }}_{Q0}^\mu `$ $``$ $`Z_Q\left(\mathrm{\Lambda }_{Q0}^\mu \right)_{\mathrm{o}.\mathrm{s}.}\gamma ^\mu Q_Q`$ (3.22)
$`\widehat{\mathrm{\Lambda }}_{\pi ,ij}^\mu `$ $``$ $`g_\pi \left(\mathrm{\Lambda }_{\pi ,ij}^\mu \right)_{\mathrm{o}.\mathrm{s}.}(\mathrm{}^{}+\mathrm{})^\mu \left(iϵ_{ij3}\right)F_\pi (q^2).`$ (3.23)
Using $`S(k)=Z_QS_F(k)`$ in the expression for $`\mathrm{\Pi }_\pi `$, we get
$`g_\pi ^1=\left({\displaystyle \frac{\mathrm{\Pi }_\pi (q)}{q^2}}\right)_{q^2=M_\pi ^2}=Z_Q^2\left({\displaystyle \frac{\widehat{\mathrm{\Pi }}_\pi (q)}{q^2}}\right)_{q^2=M_\pi ^2}Z_Q^2\widehat{g}_\pi ^1,`$ (3.24)
with the renormalized bubble graph $`\widehat{\mathrm{\Pi }}_\pi =\widehat{\mathrm{\Pi }}_s`$, see (2.15). Then Eq. (3.17) becomes
$`\widehat{\mathrm{\Lambda }}_Q^\mu `$ $`=`$ $`Z_Q\gamma ^\mu Q_Qi\widehat{g}_\pi {\displaystyle \frac{1}{2}}(1\tau _3){\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\gamma _5\left(S_F(p^{}k)\gamma ^\mu S_F(pk)\right)\gamma _5\mathrm{\Delta }_{F\pi }(k)}`$
$``$ $`2i\tau _3\widehat{g}_\pi F_\pi (q^2){\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\left(\mathrm{\Delta }_{F\pi }(p^{}k)\left(p^{}+p2k\right)^\mu \mathrm{\Delta }_{F\pi }(pk)\right)\gamma _5S_F(k)\gamma _5}.`$
Here we note that the contribution of the contact term ($`2iG_\pi `$) to the quark diagram has been dropped in order to avoid double counting: Because we always assume that our interaction Lagrangians are Fierz symmetric, it is easy to see that this contribution can be incorporated into the vector meson channel, which will be separately considered below. By using $`S=Z_QS_F`$ in the self energy (3.19) and in the expression for $`Z_Q`$ of Eq.(2.12), we see that
$`Z_Q=1+\left({\displaystyle \frac{\widehat{\mathrm{\Sigma }}_Q}{\overline{)}p}}\right)_{\overline{)}p=M},`$ (3.26)
where the renormalized quark self energy is given by
$`\widehat{\mathrm{\Sigma }}_Q(p)=3i\widehat{g}_\pi {\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\left(\gamma _5S_F(pk)\gamma _5\right)\mathrm{\Delta }_{F\pi }(k)}.`$ (3.27)
The further evaluation of the loop integral (LABEL:jq1) is left to Appendix B, where the contributions to the quark form factors $`F_{1Q}`$ and $`F_{2Q}`$ are given. The pion electromagnetic vertex is evaluated from the definition (3.23) and a Feynman diagram similar to Fig.2 for an external pion, but with a point quark-photon coupling:
$`\widehat{\mathrm{\Lambda }}_{ij}^\mu =\left(iϵ_{ij3}\right)\mathrm{\hspace{0.17em}6}i\widehat{g}_\pi {\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}\mathrm{Tr}_D\left[\gamma _5S_F(p^{}+k)\gamma ^\mu S_F(p+k)\gamma _5S_F(k)\right]}.`$ (3.28)
The explicit form of $`F_\pi (q^2)`$ is given in Appendix B.
Finally, we consider the corrections of the quark electromagnetic vertex arising from vector mesons, similar to vector meson dominance (VMD) models. If our original interaction Lagrangian contains terms of the form
$`_{I,v}=G_\omega \left(\overline{\psi }\gamma ^\mu \psi \right)^2G_\rho \left(\overline{\psi }\gamma ^\mu 𝝉\psi \right)^2,`$ (3.29)
then the point-like quark-photon vertices in the diagrams of Fig.3 and the pion vertex $`\mathrm{\Lambda }_\pi ^\mu `$ are replaced by the VMD vertex shown in Fig.4.
Because of the transverse structure of the bubble graphs in the vector channel, this leads to the following renormalization:
$`\gamma ^\mu Q_Q`$ $``$ $`{\displaystyle \frac{1}{6}}\left[\gamma ^\mu {\displaystyle \frac{2G_\omega \widehat{\mathrm{\Pi }}_V(q^2)}{1+2G_\omega \widehat{\mathrm{\Pi }}_V(q^2)}}\left(\gamma ^\mu {\displaystyle \frac{q^\mu \overline{)}q}{q^2}}\right)\right]`$ (3.30)
$`+`$ $`{\displaystyle \frac{\tau _3}{2}}\left[\gamma ^\mu {\displaystyle \frac{2G_\rho \widehat{\mathrm{\Pi }}_V(q^2)}{1+2G_\rho \widehat{\mathrm{\Pi }}_V(q^2)}}\left(\gamma ^\mu {\displaystyle \frac{q^\mu \overline{)}q}{q^2}}\right)\right],`$ (3.31)
where the form of $`\widehat{\mathrm{\Pi }}_V`$ is given in Appendix B. Because our quark-photon vertex in (3.8) is defined for on-shell quarks, the terms $`{\displaystyle \frac{q^\mu \overline{)}q}{q^2}}`$ do not contribute. Therefore the isoscalar (or isovector) parts in the quark electromagnetic vertex should be multiplied by a form factor <sup>6</sup><sup>6</sup>6Because the scalar diquark has isospin zero, this eventually also holds for the nucleon electromagnetic vertex. $`F_\omega (q^2)`$ (or $`F_\rho (q^2)`$), where
$`F_\alpha (q^2)={\displaystyle \frac{1}{1+2G_\alpha \widehat{\mathrm{\Pi }}_V(q^2)}}(\alpha =\omega ,\rho ).`$ (3.32)
## 4 Results for the nucleon form factors
In this section we will show our results for the nucleon form factors. First we discuss our model parameters and particle masses. The parameters are the same as in Refs.: The IR cut-off $`\mathrm{\Lambda }_{\mathrm{IR}}`$ is fixed as $`0.2`$ GeV, the constituent quark mass at zero nucleon density is $`M_0=0.4`$ GeV, and $`G_\pi `$, $`\mathrm{\Lambda }_{\mathrm{UV}}`$ are determined so as to reproduce $`M_\pi =0.14`$ GeV and $`f_\pi =93`$ MeV. This gives $`G_\pi =19.60`$ GeV<sup>-2</sup> and $`\mathrm{\Lambda }_{\mathrm{UV}}=0.6385`$ GeV. The coupling constant $`G_s`$ is determined so as to reproduce $`M_{N0}=0.94`$ GeV, which gives the ratio $`G_s/G_\pi =0.508`$. The coupling constant $`G_\omega `$ is determined so that the curve for the NM binding energy per nucleon $`(E_B/A`$) as a function of the density passes through the empirical saturation point<sup>7</sup><sup>7</sup>7We recall from Ref. that, in this simple NJL model, we cannot adjust both the empirical binding energy and saturation density at the same time. Therefore, although the binding energy curve passes through the empirical saturation point, its minimum is at a different point, ($`\rho ,E_B/A)=(0.22`$ fm$`{}_{}{}^{3},\mathrm{\hspace{0.17em}17.3}`$ MeV). ($`\rho ,E_B/A)=(0.16`$ fm$`{}_{}{}^{3},\mathrm{\hspace{0.17em}15}`$ MeV), which gives the ratio $`G_\omega /G_\pi =0.37`$. Finally, for the VMD form factors (3.32) we also need the coupling constant $`G_\rho `$, which is determined by reproducing the empirical symmetry energy coefficient $`a_4=35`$ MeV at the density $`\rho =0.16`$ fm<sup>-3</sup>, which gives $`G_\rho /G_\pi =0.091`$
In Table 1 we list the effective quark, diquark, nucleon and pion masses for the densities $`\rho =0,\mathrm{\hspace{0.17em}0.08},\mathrm{\hspace{0.17em}0.16},`$ and $`0.24`$ fm<sup>-3</sup>. Concerning the pion mass in the medium, we use a general result based on chiral symmetry , which for the NJL model implies that the product $`M_\pi ^2M`$ is a constant independent of density, see Eq.(2.58) of Ref.. We therefore use <sup>8</sup><sup>8</sup>8More precisely, this pion mass is defined at zero momentum, and includes nucleonic (Z-graph and contact) terms, which are important in order to guarantee the Goldstone nature of the pion in the medium. We note that these nucleonic contributions to the scalar diquark (or sigma) mass at normal densities are numerically small compared to the $`qq`$ (or $`q\overline{q}`$) polarizations, although they become important for small $`M`$ and guarantee the stability of the system w.r.t. variations in $`M`$, see Ref.() for details. $`M_\pi ^2=M_{\pi 0}^2M_0/M`$. Also listed in Table 1 are the values of $`\widehat{g}_s`$, $`\widehat{g}_\pi `$, $`\widehat{Z}_N`$ and $`Z_Q`$.
We will discuss our results for the nucleon form factor in terms of the Dirac-Pauli form factors defined by Eq.(LABEL:par1). In the discussion of medium effects, in particular the effects of the reduced nucleon mass ($`M_N<M_{N0}`$), we will also refer to an equivalent parametrization in terms of the “orbital form factor” $`G_L`$ and the “spin form factor” $`G_S`$:
$`j_N^\mu `$ $`=`$ $`\sqrt{{\displaystyle \frac{M_N}{E_{Np}}}{\displaystyle \frac{M_N}{E_{Np^{}}}}}\overline{u}_N(p^{})\left[{\displaystyle \frac{\left(p^{}+p\right)^\mu }{2M_{N0}}}G_L(q^2)+{\displaystyle \frac{i\sigma ^{\mu \nu }q_\nu }{2M_{N0}}}G_S(q^2)\right]u_N(p)`$
The relations to the Dirac-Pauli form factors are
$`G_L={\displaystyle \frac{M_{N0}}{M_N}}F_{1N},G_S={\displaystyle \frac{M_{N0}}{M_N}}(F_{1N}+F_{2N}).`$ (4.2)
We introduce these form factors here, because $`F_1`$ and $`F_2`$ do not directly reflect the enhancement of the nucleon orbital current ($`𝒋_{\mathrm{N},\mathrm{orb}}`$) arising from the reduced nucleon mass (enhanced nuclear magneton) <sup>9</sup><sup>9</sup>9Note that the appearance of the medium modified nucleon mass ($`M_N`$) in the Pauli term of (LABEL:par1) is a mere definition of $`F_2`$.. Moreover, the parametrization (LABEL:par2) for the space part of the current has more connection to the traditional calculations of nuclear magnetic properties , because the values of these form factors at $`q=0`$ reduce to the orbital and spin g-factors: $`g_L=G_L(0)`$, $`g_S=G_S(0)`$.
Here we would like to point out that these different ways to discuss the medium modifications of the nucleon current remind us that the form factors of a nucleon in the medium are not directly observable quantities: Ultimately the current $`j_N^\mu `$ has to be used in a nuclear structure calculation of observable cross sections. Our current $`j_N^\mu `$ reflects only those effects which are not taken into account in nuclear structure calculations, i.e., the effects of the nuclear mean fields on the internal motion of quarks in the nucleon. Other effects, which explicitly depend on the density and have their origin in the Pauli principle on the level of nucleons, must be considered in the nuclear part of the calculation <sup>10</sup><sup>10</sup>10This is also evident from the fact that the full current of a nucleon in the medium, including the explicitly density dependent parts, cannot be parametrized in the Lorentz invariant forms (LABEL:par1) or (LABEL:par2).. As an example of such a calculation for the case of nuclear matter, we will consider the response function for quasielastic electron scattering in Sect.5.
For the zero density (single nucleon) case, it is possible to define a Breit frame where $`q_0=0`$, and in this frame the nucleon current can be expressed in terms of the familiar electric and magnetic form factors
$`G_E(q^2)`$ $`=`$ $`F_1(q^2)+{\displaystyle \frac{q^2}{4M_N^2}}F_2(q^2)`$ (4.3)
$`G_M(q^2)`$ $`=`$ $`F_1(q^2)+F_2(q^2).`$ (4.4)
We will compare our calculated form factors for zero density with the empirical dipole form factors, defined by
$`G_{Ep}`$ $`=`$ $`G_{Mp}/\mu _p=G_{Mn}/\mu _n={\displaystyle \frac{1}{\left(1q^2/0.71\mathrm{GeV}^2\right)^2}}`$
$`F_{1n}`$ $`=`$ $`0.`$ (4.5)
For a nucleon moving in the medium, however, one cannot define a Breit frame, and consequently the combinations (4.3) and (4.4) do not enter naturally in the expressions for nuclear observables, like response functions or elastic form factors. For the finite density case we will therefore discuss our results in terms of the form factors $`F_1`$ and $`F_2`$, or $`G_L`$ and $`G_S`$.
The results for the form factors at zero density (free nucleon case) are shown in Figs.5-8. There we plot (i) the results of the naive quark-diquark model (see Sect.3.1) without (dotted lines) and with (dashed lines) the contact term contribution to the quark diagram of Fig.1, (ii) the results obtained by including in addition the effects of the diquark form factor (dash-dotted lines), and (iii) the total result including also the pion and VMD effects.
Fig. 5 shows that in the naive quark-diquark model the electric size of the proton is too small and the form factors $`F_{1p}`$ and $`G_{Ep}`$ fall off too slowly. The situation improves when the diquark form factor is included, and also the pion cloud gives some positive contribution to the electric size of the proton <sup>11</sup><sup>11</sup>11We obtain $`<r_E^2>_p=0.421`$fm$`{}_{}{}^{2}+0.062`$fm$`{}_{}{}^{2}=0.483`$ fm<sup>2</sup>, where the two terms come from the Dirac ($`F_1`$) part and the anomalous ($`F_2`$) part. The fact that this is too small compared to the experimental value of $`0.74`$fm<sup>2</sup> is partially because the magnetic moment is too small, but also because the slope of $`F_1`$ is too small.. The total result for $`F_{1p}`$ still lies above the empirical dipole form factor, but we can expect that the further inclusion of axial vector diquarks and pion cloud effects around the nucleon will lead to a satisfactory description.
Fig. 6 shows that the proton magnetic moment in the naive quark-diquark picture is too small, which is expected and also well known. The finite size effects of the scalar diquark do not contribute much in this case. The inclusion of the contact term in the quark diagram of Fig.1 improves the situation somewhat. By using a Fierz transformation, this term is actually seen to be equivalent to a vector meson contribution (like shown in Fig.4 for the quark), but with a tensor coupling ($`\sigma ^{\mu \nu }q_\nu `$) to the nucleon. Also the pion cloud, which leads to anomalous magnetic moments of the constituent up and down quarks<sup>12</sup><sup>12</sup>12We obtain $`\mu _u=\frac{2}{3}+0.061`$, $`\mu _d=\frac{1}{3}0.123`$ for the magnetic moments of u,d quarks in the free nucleon., gives a positive contribution to the proton magnetic moment, but the total result is still too small. It is, however, known that the axial vector diquark and the pion cloud around the nucleon give large contributions to the magnetic moment and the associated form factors, and Fig.6 shows how far one can go in the simple quark - scalar diquark description.
The importance of the diquark form factor is also seen for the neutron form factors $`F_{1n}`$ and $`G_{En}`$, which are shown in Fig. 7. The naive quark-diquark model gives an electric form factor which is too large in comparison to the experimental one (note that the experimental $`G_{En}`$ is smaller than the “dipole form factor” shown in Fig. 7), and the diquark form factor, which suppresses the (positive) contribution from the second diagram of Fig.1 relative to the (negative) first one, is essential to obtain reasonable values. The electric size of the neutron is somewhat too small in magnitude, but this is because the absolute value of the magnetic moment is too small <sup>13</sup><sup>13</sup>13We obtain $`<r_E^2>_n=0.003`$fm$`{}_{}{}^{2}0.072`$fm$`{}_{}{}^{2}=0.069`$ fm<sup>2</sup>, where the two terms come from the Dirac ($`F_1`$) part and the anomalous ($`F_2`$) part. The fact that this is too small in magnitude compared to the experimental value of $`0.12`$fm<sup>2</sup> is because the absolute value of the magnetic moment is too small. For discussions on the role of these two contributions to the neutron electric radius, see for example .. In this connection, it is interesting to observe that the result for $`F_{1n}`$, and in particular its contribution to the electric radius, is very small, and therefore the electric size of the neutron is almost entirely due to the “Foldy term” .
The situation for the form factors $`F_{2n}`$ and $`G_{Mn}`$ shown in Fig. 8 is similar to the case of the proton (Fig. 6), i.e., the contact term and the pion cloud around the constituent quarks give some improvements of the magnetic moment, but the total result is still too small. This, and the fact that the form factors fall off too slowly, again points out the necessity to include the axial vector diquark channel.
The medium modifications of the nucleon form factors are shown in Figs. 9-12, where we plot the results for $`\rho =0`$ (dotted lines), $`\rho =0.08`$ fm<sup>-3</sup> (dashed lines), $`\rho =0.16`$ fm<sup>-3</sup> (solid lines), and $`\rho =0.24`$ fm<sup>-3</sup> (dash-dotted lines).
The result for $`F_{1p}`$ of Fig. 9 indicates that the electric size of the proton in the medium is somewhat enhanced. The orbital form factor $`G_{Lp}`$ shown in the insert of Fig.9 demonstrates the enhancement of the orbital current ($`𝒋_{\mathrm{N},\mathrm{orb}}`$) due to the reduced effective nucleon mass, see Eq.(4.2). We have to remind, however, that the isoscalar part of this enhancement is in a sense spurious, because in an actual nuclear calculation it is canceled by the “backflow” effect, which in our language arises from Z-graphs, i.e., the Pauli blocking part of the $`N\overline{N}`$ excitation piece (see the detailed discussions in Ref. on the backflow in relativistic meson-nucleon theories). Namely, the proton orbital g-factor, which is roughly $`M_{N0}/M_N`$ in a Hartree calculation, becomes approximately $`\frac{1}{2}(1+M_{N0}/M_N)`$ after the inclusion of the backflow, where the first term is the isoscalar and the second one the isovector piece<sup>14</sup><sup>14</sup>14It is well known that the pion effects further enhance the isovector piece..
Fig. 10 shows that the medium effects tend to decrease the “intrinsic” anomalous magnetic moment of the proton, but when combined with the enhancement of the nuclear magneton, the spin g-factor is enhanced, as shown in the insert of Fig.10. It is interesting that a very similar result has been obtained also in hadronic models. Therefore, the quark effects considered here do not lead to a quenching of the spin g-factor, as would be desirable to explain the missing quenching of isovector nuclear magnetic moments , but rather to an enhancement <sup>15</sup><sup>15</sup>15This is in contrast to the quenching of the axial vector coupling constant observed in the quark-diquark calculations of Ref. including the axial vector diquark, and in hadronic models .. The figure also shows that the magnetic size of the proton becomes somewhat larger in the medium.
The results for the neutron form factor $`F_{1n}`$ in Fig. 11 show that the effect of finite diquark size, which was very important for the zero density case (Fig.7) to reduce $`F_{1n}`$ to reasonable values, increases with increasing density. That is, the diquark form factor at finite density further suppresses the positive contribution of the diquark diagram in Fig.1. The orbital form factor $`G_{Ln}`$ shown in the insert of Fig.11 again demonstrates the enhancement of the nuclear magneton, but we have to keep in mind that the backflow effect will change the neutron orbital g-factor from the present value $`0`$ to roughly $`\frac{1}{2}\left(1M_{N0}/M_N\right)<0`$, and that effects of the pion cloud around the nucleon are known to further enhance the magnitude of the neutron orbital g-factor.
The medium effects on the neutron form factor $`F_{2n}`$ shown in Fig.12 are qualitatively similar to the proton case of Fig.10: They decrease the absolute value of the “intrinsic” anomalous magnetic moment of the neutron, but when combined with the enhancement of the nuclear magneton the neutron spin g-factor becomes slightly enhanced in magnitude, as shown in the insert of Fig.12. The magnetic size of the neutron is not changed much, and the total medium effects on the neutron spin form factor are very small.
An interesting common feature of the results shown in Fig. 9-12 is that the medium modifications of the orbital and spin form factors ($`G_L(Q^2)`$ and $`G_S(Q^2)`$) always decrease with increasing $`Q^2`$. This is consistent with the intuitive expectation that the internal structure of the nucleon at short distances is not influenced much by the mean nuclear fields, and again indicates that these form factors reflect the change of the nucleon current in the medium more directly and transparently than the form factors $`F_1`$, $`F_2`$ themselves, or any other combination of them.
## 5 Application: The longitudinal nuclear response function
As an application of the medium effects discussed in the previous section to the calculation of nuclear quantities, we consider the longitudinal response function for quasielastic electron scattering. Generally, quasielastic processes are the ideal places to investigate medium modifications of nucleon form factors, because the on-shell kinematics, which is used in the derivation of the nucleon form factors, is justified. Recently, experiments on the polarization transfer in proton knock-out reactions have been carried out in the region of quasielastic kinematics at large momentum transfers, and the results were discussed in connection to the predicted medium modifications of nucleon form factors. In this work, however, we will consider the simpler case of the inclusive quasielastic response function in the region of lower momentum transfers. In previous works, based on purely hadronic models, it was shown that this quantity is quite sensitive to medium modifications of form factors, and here we wish to apply our more microscopic quark description to the longitudinal response function in nuclear matter.
The longitudinal response function in isospin symmetric nuclear matter is expressed as
$`S_L(\omega ,𝒒)={\displaystyle \frac{2Z}{\pi \rho }}\mathrm{Im}\mathrm{\Pi }_L(\omega ,𝒒),`$ (5.1)
where $`\mathrm{\Pi }_L(\omega ,𝒒)`$ is the 2-point (correlation) function for two external operators $`j_N^0`$. In the mean field (Hartree) approximation, this is expressed by the first diagram of Fig.13, and if we include the direct terms of the NN interaction in the ladder approximation, this gives the familiar RPA series.
The relativistic calculation of $`\mathrm{\Pi }_L`$ follows the lines given in Ref., i.e., the bubble graphs in Fig.13 consist of the nucleon particle- nucleon hole excitations and the Pauli blocking part of the nucleon particle- antinucleon excitations (Z-graphs). The other effects, which are not taken into account explicitly in Fig.13, are summarized in the density dependent electromagnetic vertices. Previous calculations based on hadronic models incorporated the vacuum fluctuations on the level of nucleons (that is, the change of the nucleon-antinucleon vacuum polarization graphs in the presence of the nuclear mean fields), but it is more appropriate to describe these vacuum fluctuations on the level of quarks <sup>16</sup><sup>16</sup>16Although this follows naturally from the derivation of the nucleon lagrangian from the quark lagrangian in the path integral approach , or from the derivation of the effective NN interaction in quark theories following the Landau-Migdal approach, it remains to be demonstrated explicitly for electromagnetic quantities.. Therefore we use our nucleon current (3.1) for $`\mu =0`$ at the electromagnetic vertices of Fig.13.
For the RPA calculation, we need the NN interaction kernel in our effective quark theory, which is shown graphically in Fig.14 and expressed as <sup>17</sup><sup>17</sup>17The quantities $`(\mathrm{𝟏})_i`$ and $`(\gamma ^\mu )_i`$ in (5.3) and (5.4) express the Dirac matrices acting between the spinors of nucleon i=1,2, i.e., in order to get the NN interaction one has to multiply the spinors $`\sqrt{M_N/E_N}u_N`$ of the initial and final states.
$`V_{NN}(k)V_\sigma (k)(\mathrm{𝟏})_1(\mathrm{𝟏})_2+V_\omega (k)(\gamma _\mu )_1(\gamma ^\mu )_2`$ (5.3)
$`={\displaystyle \frac{2G_\pi }{12G_\pi \widehat{\mathrm{\Pi }}_\sigma (k^2)+2G_\pi \delta M_s^2}}\left({\displaystyle \frac{\mathrm{d}M_N}{\mathrm{d}M}}\right)^2F_\sigma ^2(k^2)(\mathrm{𝟏})_1(\mathrm{𝟏})_2`$
$`+`$ $`{\displaystyle \frac{2G_\omega }{1+2G_\omega \widehat{\mathrm{\Pi }}_V(k^2)}}\mathrm{\hspace{0.17em}9}F_\omega ^2(k^2)(\gamma _\mu )_1(\gamma ^\mu )_2`$ (5.4)
Here the first factor in (5.3) is the reduced $`q\overline{q}`$ t-matrix in the sigma meson channel, and the corresponding bubble graph $`\widehat{\mathrm{\Pi }}_\sigma `$ is given in Appendix C. The derivation of the effective NN interaction in the NJL model, however, has shown that in addition to the part $`\widehat{\mathrm{\Pi }}_\sigma `$, which describes the $`q\overline{q}`$ exchange (see Fig.14), there is also a nuclear part which consists of (i) the Z-graph, and (ii) a contact term arising from an induced $`\overline{N}\sigma ^2N`$ interaction. Concerning the Z-graph, we note that this is just the Pauli blocking part to the $`N\overline{N}`$ bubble graph, which is taken into account explicitly in the RPA series of Fig.13 and therefore should not be included in the interaction. (Numerically the Z-graph contribution is small compared to $`\widehat{\mathrm{\Pi }}_\sigma `$ because of the reduced $`\sigma `$N coupling in the medium, see Ref..) The self energy correction for the sigma meson arising from the induced $`\overline{N}\sigma ^2N`$ contact interaction is included in (5.3) as the density-dependent constant
$`\delta M_\sigma ^2=4\left({\displaystyle \frac{\mathrm{d}^2M_N}{\mathrm{d}M^2}}\right){\displaystyle \frac{\mathrm{d}^3k}{(2\pi )^3}\frac{M_N}{E_N(k)}\mathrm{\Theta }(p_F|𝒌|)}.`$ (5.5)
The $`NN\sigma `$ coupling constant at zero momentum is proportional to $`\mathrm{d}M_N/\mathrm{d}M`$, and its square appears in (5.3). The $`NN\sigma `$ vertex form factor $`F_\sigma (k^2)`$ is normalized to $`F_\sigma (0)=1`$. The first factor in (5.4) is the reduced $`q\overline{q}`$ t-matrix in the $`\omega `$-meson channel, and the corresponding bubble graph $`\widehat{\mathrm{\Pi }}_V`$ is the same as in the VMD correction to the quark form factors (Fig.4). The $`NN\omega `$ coupling constant at zero momentum is proportional to the number of quarks in the nucleon, and the vertex form factor $`F_\omega (k^2)`$ is normalized as $`F_\omega (0)=1`$.
By taking matrix elements of (5.3), (5.4) between the nucleon spinors, it is easy to see that for $`k=0`$ and for nucleons at the Fermi surface we just get the Landau-Migdal interaction derived more generally in Ref.() (except for the Z-graph contributions as explained above). The form factor $`F_\omega `$ is equal to $`(F_{1p}+F_{1n})`$, which was calculated in the previous sections. The scalar form factor $`F_\sigma `$ should in principle be calculated independently from the Feynman diagrams of Fig.1 by using the external operator $`\mathrm{𝟏}`$. However, because the calculations discussed below show that the RPA effects are numerically not very important, we will simply assume the same form factor as for the $`NN\omega `$ coupling ($`F_\sigma =F_\omega `$).
The two parts of the NN interaction, $`V_\sigma `$ of (5.3) and $`V_\omega `$ of (5.4) , are shown by the lower and upper dashed lines, respectively, in Fig.15 for the kinematics needed in the calculation of the response function, $`|𝒒|=0.5`$ GeV and $`0<\omega <0.4`$ GeV. As in the case of relativistic hadronic models, we see large cancellations between the attractive scalar and repulsive vector parts. The solid line shows the naive sum of the two dashed lines, while the dash-dotted line shows the combination $`V_\sigma +V_\omega (1\omega ^2/|𝒒|^2)`$, which includes the effect of the longitudinal space component of the $`\omega `$ exchange and is more relevant for the longitudinal response function <sup>18</sup><sup>18</sup>18We have to note, however, that the true interaction in the nuclear medium in the longitudinal channel is more repulsive than shown by the dash-dotted line in Fig.15, because of the difference between the Dirac matrices $`\mathrm{𝟏}`$ and $`\gamma ^0`$ in Eq.(5.3). On the average it is repulsive on the small $`\omega `$ side and becomes attractive on the large $`\omega `$ side, as the RPA result of Fig.16 shows.. The results shown in Fig. 15 can be reproduced almost exactly by an approximate form in terms of Yukawa potentials, if the coupling constants and meson masses are defined at $`k^2=0`$. This is discussed in Appendix C, where also numerical values are given.
The result for the longitudinal response function of NM ($`\rho =0.16`$ fm<sup>-3</sup>) for $`|𝒒|=0.5`$ GeV is shown in Fig. 16. The dashed line is the Hartree response with the free dipole form factors, and the dash-dotted line is obtained by adding our calculated medium corrections $`\mathrm{\Delta }F(q^2;\rho )=F(q^2,\rho )F(q^2,\rho =0)`$ to the free dipole form factors. (Here $`F`$ denotes any of the Dirac-Pauli form factors.) We see that, even in this region of relatively low momentum transfers, the medium effects are appreciable. Finally, the solid line shows the result obtained by further adding the RPA corrections with the NN interaction derived above and using the density dependent meson-nucleon form factors.
We do not show a comparison to experimental data in Fig.16 because of two reasons: First, our calculation refers to NM and cannot be applied directly to finite nuclei, although for the case of <sup>40</sup>Ca the results are qualitatively very similar to the NM results, see Ref.. Second, the analysis of the experimental data is still controversial, mainly because of the model dependence of the Coulomb corrections <sup>19</sup><sup>19</sup>19The analysis of Ref., which did not take into account the Coulomb corrections, was refined in Ref. by using a particular theoretical model for the Coulomb corrections, but it has been pointed out recently that the results depend on which theoretical prescriptions are used..
## 6 Summary
In this paper we used a simple quark-scalar diquark picture for the single nucleon to describe the electromagnetic form factors of a nucleon bound in the nuclear medium. We used the nuclear matter equation of state derived within the same effective quark theory to assess the effect of the mean nuclear fields on the internal quark structure and the form factors of a bound nucleon, taking into account also the meson cloud around the constituent quarks. We have shown that this simple model gives reasonable results for the Dirac form factors of the free nucleon if the finite size of the diquark is taken into account. This is particularly important for the neutron in order to obtain a small $`F_1`$, consistent with observations. Concerning the Pauli form factors, in particular the anomalous magnetic moments and the behavior for large $`Q^2`$, one would need to further add the effect of the axial vector diquark and its finite size to achieve a reasonable description.
The medium modifications of the form factors associated with the orbital current are significant in the region of low and intermediate $`Q^2`$, and partially associated with an increase of the electric size in the medium. The form factors associated with the spin current are enhanced because of the reduced nucleon mass, but due to a simultaneous decrease of the intrinsic anomalous magnetic moment, the total changes are not very large; in particular for the neutron they are almost zero. For both kinds of form factors - orbital and spin -, the medium modifications decrease with increasing $`Q^2`$. This is consistent with the intuitive expectation that the mean fields, which reflect the long range nuclear correlations, should not influence the structure of the nucleon at short distances.
As an application, we considered the longitudinal response function of nuclear matter for inclusive quasielastic electron scattering. Even in this region of relatively low momentum transfer, the medium effects on the nucleon form factors give rise to an appreciable quenching of the response function. In order to take into account the RPA-type correlations between the nucleons, we derived the NN interaction in the present quark model for the isoscalar channel, which is most important for the longitudinal response function. In the limit of zero momentum transfer (Landau limit), this interaction agrees with the more general Landau-Migdal effective interaction, which was derived in earlier works for this particular quark theory, and has many features which are similar to relativistic meson-nucleon theories. In particular, it is on the average repulsive in the region below the quasielastic peak, and becomes attractive for higher energy transfers.
We finally would like to remark that the language of effective quark theories is very appropriate to address the problem of medium modifications of nucleon properties. This work represents one further step toward the goal of describing relativistic nuclear systems by taking into account the quark substructure of the constituent hadrons.
Acknowledgment
W.B. wishes to thank A.W. Thomas for many helpful discussions. This work was supported by the Grant in Aid for Scientific Research of the Japanese Ministry of Education, Culture, Sports, Science and Technology, Project No. C-16540267.
Appendices
## Appendix A Expressions for the nucleon form factors
Here we give the explicit expressions for the contributions of the 3 terms in Eq.(LABEL:contact)-(3.13) to the nucleon Dirac-Pauli form factors defined in (LABEL:par1). They can be derived by invariant integration in the usual way, i.e., by (i) introducing Feynman parameters, (ii) performing a shift so that the denominator depends only on $`k^2`$, (iii) expressing the result in the form (LABEL:par1) by using the Dirac equation for the nucleon spinors, (iv) performing a Wick rotation according to
$`i{\displaystyle \frac{\mathrm{d}^4k}{(2\pi )^4}f(k^2)}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}f(t)`$
where $`t=k_0^2+𝒌^2`$ is the square of the Euclidean length, and finally (v) introducing a Lorentz invariant regularization scheme, in our case the proper time regularization (2.19). Below we give the expressions which are obtained after performing the steps (i)-(iv), using $`Q^2=q^2>0`$ as the variable.
In the expressions given below there enter the diquark and nucleon wave function normalization factors, which are defined by (LABEL:gs) and (LABEL:zn) in terms of the renormalized bubble graphs (2.15) and (2.16), for which we have the expressions
$`\widehat{\mathrm{\Pi }}_s(p^2)`$ $`=`$ $`24{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}\left({\displaystyle \frac{1}{t+M^2}}+{\displaystyle \frac{p^2}{2}}{\displaystyle _0^1}dx{\displaystyle \frac{1}{\left(t+M^2p^2x(1x)\right)^2}}\right)`$
$`\widehat{\mathrm{\Pi }}_N(p)`$ $`=`$ $`\widehat{g}_s{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle \frac{\overline{)}px+M}{\left(t+M^2(1x)+M_s^2xp^2x(1x)\right)^2}}.`$
If we denote
$`D_1(Q^2,x)=t+M^2+Q^2x(1x),`$ (A.3)
the contributions of the contact term (LABEL:contact) to the form factors are as follows:
$`F_{1N}^{(C)}(Q^2)`$ $`=`$ $`{\displaystyle \frac{4G_s}{\widehat{g}_s}}\widehat{Z}_N{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx`$ (A.4)
$`\times `$ $`[F_{1Q}(Q^2)({\displaystyle \frac{1}{D_1(Q^2,x)}}{\displaystyle \frac{1}{t+M^2}}+{\displaystyle \frac{Q^2}{2}}{\displaystyle \frac{1}{D_1(Q^2,x)^2}})`$
$`+`$ $`F_{2Q}(Q^2){\displaystyle \frac{Q^2}{2}}{\displaystyle \frac{(1+x)(1x\frac{M_N}{M})}{D_1(Q^2,x)^2}}]`$
$`F_{2N}^{(C)}(Q^2)`$ $`=`$ $`{\displaystyle \frac{4G_s}{\widehat{g}_s}}\widehat{Z}_NM_NM{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx`$
$`\times `$ $`\left[F_{1Q}(Q^2){\displaystyle \frac{2}{D_1(Q^2,x)^2}}+F_{2Q}(Q^2){\displaystyle \frac{(1x\frac{M_N}{M})^2\frac{Q^2x}{M^2}}{D_1(Q^2,x)^2}}\right].`$
If we denote
$`D_2(Q^2,x,y)`$ $`=`$ $`t+M_s^2(1x)+M^2xM_N^2x(1x)+{\displaystyle \frac{Q^2}{4}}(x^2y^2)`$
$`N_1(Q^2,x)`$ $`=`$ $`(M_N+M)^2M_s^2{\displaystyle \frac{t}{2}}2M_N^2x(1x)2M_NMx`$
$`N_2(Q^2,x)`$ $`=`$ $`\left(1+{\displaystyle \frac{M_N}{M}}(1x)\right)^2{\displaystyle \frac{Q^2}{4M^2}}(x^2y^2),`$
the contributions of the quark term (3.12) to the form factors are as follows:
$`F_{1N}^{(Q)}`$ $`=`$ $`\widehat{Z}_N{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx`$ (A.7)
$`\times `$ $`[F_{1Q}(Q^2)({\displaystyle \frac{1}{D_1(Q^2,x)^2}}+{\displaystyle _x^x}\mathrm{d}y{\displaystyle \frac{N_1(Q^2,x)}{D_2(Q^2,x,y)^3}})`$
$``$ $`F_{2Q}(Q^2){\displaystyle \frac{Q^2}{2}}{\displaystyle _x^x}\mathrm{d}y{\displaystyle \frac{1+\frac{M_N}{M}(1x)}{D_2(Q^2,x,y)^3}}]`$
$`F_{2N}^{(Q)}`$ $`=`$ $`\widehat{Z}_N{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle _x^x}dy`$
$`\times `$ $`\left[F_{1Q}(Q^2)\mathrm{\hspace{0.17em}2}M_Nx{\displaystyle \frac{M_N(1x)+M}{D_2(Q^2,x,y)^3}}+F_{2Q}(Q^2)M_NM{\displaystyle \frac{N_2(Q^2,x)}{D_2(Q^2,x,y)^3}}\right].`$
If we denote
$`D_3(Q^2,x,y)=t+M^2(1x)+M_s^2xM_N^2x(1x)+{\displaystyle \frac{Q^2}{4}}(x^2y^2),`$
(A.9)
the contributions of the diquark term (3.13) to the form factors are as follows:
$`F_{1N}^{(D)}`$ $`=`$ $`\widehat{Z}_NF_D(Q^2){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle _x^x}dy{\displaystyle \frac{2M_N(1x)(M_Nx+M)+\frac{t}{2}}{D_3(Q^2,x,y)^3}}`$
$`F_{2N}^{(D)}`$ $`=`$ $`\widehat{Z}_NF_D(Q^2){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle _x^x}dy\mathrm{\hspace{0.17em}2}M_N(1x){\displaystyle \frac{M_Nx+M}{D_3(Q^2,x,y)^3}}.`$
The diquark form factor $`F_D`$ is calculated from the expression (3.16) as
$`F_D`$ $`=`$ $`12\widehat{g}_s{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}\mathrm{d}x\{F_{1Q}^{(0)}(Q^2)[{\displaystyle \frac{1}{D_1(Q^2,x)^2}}+M_s^2{\displaystyle _x^x}\mathrm{d}y{\displaystyle \frac{x}{D_4(Q^2,x,y)^3}}]`$ (A.12)
$``$ $`{\displaystyle \frac{Q^2}{2}}F_{2Q}^{(0)}(Q^2){\displaystyle _x^x}\mathrm{d}y{\displaystyle \frac{1}{D_4(Q^2,x,y)^3}}\}`$
where we defined
$`D_4(Q^2,x,y)=t+M^2M_s^2x(1x)+{\displaystyle \frac{Q^2}{4}}(x^2y^2),`$ (A.13)
and $`F_{1Q}^{(0)}`$, $`F_{2Q}^{(0)}`$ are the isoscalar parts of the quark form factors.
## Appendix B Expressions for the quark form factors
Here we give the expressions for the quark electromagnetic vertex (LABEL:jq1) in terms of the quark form factors defined in Eq.(3.8). In the expressions given below there enter the pion and quark wave function normalizations $`\widehat{g}_\pi `$ and $`\widehat{Z}_Q`$, which are defined in terms of the renormalized self energies $`\widehat{\mathrm{\Pi }}_\pi `$ and $`\widehat{\mathrm{\Sigma }}_Q`$ by Eqs. (3.24) and (3.26). The expression for $`\widehat{\mathrm{\Pi }}_\pi =\widehat{\mathrm{\Pi }}_s`$ has been given in Eq.(LABEL:pisc), and
$`\widehat{\mathrm{\Sigma }}_Q(p)`$ $`=`$ $`3\widehat{g}_\pi {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle \frac{\overline{)}pxM}{\left(t+M^2(1x)+M_\pi ^2xp^2x(1x)\right)^2}}`$
As explained in the main text, the isoscalar (or isovector) parts of the quark vertices given below should eventually be further multiplied by the VMD form factors (3.32), where the expression for $`\widehat{\mathrm{\Pi }}_V`$ is
$`\widehat{\mathrm{\Pi }}_V(Q^2)=48Q^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle \frac{x(1x)}{D_1(Q^2,x)^2}}.`$ (B.2)
The quark diagram (second term in (LABEL:jq1)) gives the following contributions to the quark form factors:
$`F_{1Q}^{(Q)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1\tau _3\right)\widehat{g}_\pi {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx\left[{\displaystyle \frac{1}{D_1(Q^2,x)^2}}+{\displaystyle _x^x}dy{\displaystyle \frac{2M^2x^2M_\pi ^2\frac{t}{2}}{D_5(Q^2,x,y)^3}}\right]`$
$`F_{2Q}^{(Q)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1\tau _3\right)\widehat{g}_\pi {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle _x^x}dy{\displaystyle \frac{2M^2x^2}{D_5(Q^2,x,y)^3}}`$ (B.4)
where
$`D_5(Q^2,x,y)=t+M_\pi ^2(1x)+M^2x^2+{\displaystyle \frac{Q^2}{4}}(x^2y^2).`$ (B.5)
The pion diagram (third term in (LABEL:jq1)) gives the following contributions to the quark form factors:
$`F_{1Q}^{(\pi )}`$ $`=`$ $`2\tau _3\widehat{g}_\pi F_\pi (q^2){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle _x^x}dy{\displaystyle \frac{2M^2(1x)^2\frac{t}{2}}{D_6(Q^2,x,y)^3}}`$
$`F_{2Q}^{(\pi )}`$ $`=`$ $`4\tau _3M^2\widehat{g}_\pi F_\pi (q^2){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx{\displaystyle _x^x}dy{\displaystyle \frac{(1x)^2}{D_6(Q^2,x,y)^3}},`$
where
$`D_6(Q^2,x,y)=t+M^2(1x)^2+M_\pi ^2x+{\displaystyle \frac{Q^2}{4}}(x^2y^2).`$ (B.8)
The pion form factor is calculated from (3.23) and (3.28) as
$`F_\pi `$ $`=`$ $`12\widehat{g}_\pi {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}{\displaystyle _0^1}dx\left[{\displaystyle \frac{1}{D_1(Q^2,x)^2}}+M_\pi ^2{\displaystyle _x^x}dy{\displaystyle \frac{x}{D_7(Q^2,x,y)^3}}\right],`$
where
$`D_7(Q^2,x,y)=t+M^2M_\pi ^2x(1x)+{\displaystyle \frac{Q^2}{4}}(x^2y^2).`$ (B.10)
## Appendix C The NN interaction
Here we provide some details on the NN interaction Eqs.(5.3).
The bubble graph in the sigma channel, which appears in (5.3), has the form
$`\widehat{\mathrm{\Pi }}_\sigma (k^2)`$ $`=`$ $`6i{\displaystyle \frac{\mathrm{d}^4q}{(2\pi )^4}\mathrm{Tr}_D\left(S_F(q)S_F(k+q)\right)}`$
$`=`$ $`12{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{t\mathrm{d}t}{16\pi ^2}}\left[{\displaystyle \frac{2}{t+M^2}}+\left(k^24M^2\right){\displaystyle _0^1}dx{\displaystyle \frac{1}{\left(t+M^2k^2x(1x)\right)^2}}\right].`$
The expression for $`\widehat{\mathrm{\Pi }}_V`$ has been given in (B.2).
In the calculations of Sect.5, the NN interaction (5.3) is used without further approximations. Here we wish to give an approximate form in terms of Yukawa potentials: Expanding $`\widehat{\mathrm{\Pi }}_\sigma `$ and $`\widehat{\mathrm{\Pi }}_V`$ around $`k^2=0`$, we obtain
$`V_{\mathrm{NN}}(k)={\displaystyle \frac{g_{\sigma N}}{k^2M_\sigma ^2}}(\mathrm{𝟏})(\mathrm{𝟏}){\displaystyle \frac{g_{\omega N}}{k^2M_\omega ^2}}(\gamma _\mu )(\gamma ^\mu ),`$ (C.2)
where
$`g_{\sigma N}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{d}M_N}{\mathrm{d}M}}\right)^2g_\sigma =(2.34)^2\times 21.75=119.4`$ (C.3)
$`g_{\omega N}`$ $`=`$ $`9g_\omega =9\times 18.43=165.9,`$ (C.4)
where the quark-meson couplings are defined by
$`g_\sigma =\left[\left(\widehat{\mathrm{\Pi }}_\sigma /k^2\right)_{k^2=0}\right]^1`$ and $`g_\omega =\left[\left(\widehat{\mathrm{\Pi }}_V/k^2\right)_{k^2=0}\right]^1`$, and the numerical values given above are obtained from our nuclear matter EOS for $`\rho =0.16`$ fm<sup>-3</sup>. The meson masses defined at zero momentum are
$`M_\sigma ^2`$ $`=`$ $`g_\sigma \left({\displaystyle \frac{1}{2G_\pi }}\widehat{\mathrm{\Pi }}_\sigma (0)+\delta M_\sigma ^2\right)=\left(0.81\mathrm{GeV}\right)^2`$ (C.5)
$`M_\omega ^2`$ $`=`$ $`{\displaystyle \frac{g_\omega }{2G_\omega }}=\left(1.13\mathrm{GeV}\right)^2.`$ (C.6)
Note that these masses are different from the pole positions.
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