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# 1 Introduction
## 1 Introduction
High-energy circularly polarized photons, if available, would help in solving several important problems in high energy physics like that of spin-crisis (see e.g. and literature cited there). If we start with linearly polarized multi-GeV photons, then they can be converted into circularly polarized ones provided that the suggestion of Cabibbo and collaborators to use a specially chosen crystalline plate for this purpose is true. This idea is now under experimental investigation at CERN within NA59 project, where the linearly polarized photons are also produced by means of a crystalline target from unpolarized electrons passing through it. From theoretical point of view, the conversion process itself is very interesting since for multi-GeV photons it is caused by the polarization of vacuum in the presence of the periodical electric field of a crystal. This phenomenon has not been investigated so far.
To develop a description of the photon propagation, we start in Sec.2 from the Maxwell equations taking into account a current induced in a medium by the incident wave in a rather general form. When the wavelength of a photon is much shorter than any other characteristic distance scale of the problem, the parameters of a wave packet change slowly while it propagates in a medium. Using this fact,we obtain the solution to the Maxwell equations in the short-wavelength approximation, which is the analog of the eikonal approximation in the fast particle scattering theory. The approximation used becomes valid starting with a relatively small photon energy $`\omega _{}^>m`$ , where $`m`$ is the electron mass ( we use $`\mathrm{}=c=\mathrm{\hspace{0.17em}1}`$ units). Comparing our solution with that of obtained for an amorphous medium, we find that a matrix, describing a change of the polarization and intensity of a wave packet is deeply connected with the forward scattering amplitude. At least, it is the case for large thicknesses, when this change becomes really prominent. In the multi-GeV energy region of interest, the photon scattering via virtual $`e^+e^{}`$ pairs is the only process relevant to the problem. Its amplitude for a separate atom is well known ( see for the forward scattering amplitude ). In crystals, along with this incoherent ( amorphous like ) contribution to the amplitude, the coherent one, caused by the periodicity of a lattice is present. The latter is calculated in Sec.3 by means of the so called quasiclassical operator method. The details of this method along with many applications can be found in . We use the results obtained to consider the circular - to - linear polarization conversion process. Optimal ( according to the criterion formulated in ) thicknesses and orientations are found for diamond, silicon and germanium crystals for photon energy $`\omega \mathrm{\hspace{0.17em}100}GeV`$.
## 2 Propagation of short-wavelength photons
For the electric field $`𝐄`$ of a wave we obtain from the Maxwell equations
$$\left(\frac{^2}{t^2}^2\right)E_i(x)=\mathrm{\hspace{0.17em}4}\pi d^4x^{}R_{im}(x)\chi _{mj}(x,x^{})E_j(x^{}),$$
(1)
where the operator
$$R_{im}(x)=\frac{^2}{x_ix_m}\delta _{im}\frac{^2}{t^2}.$$
When deriving eq.(1), the following relationship of the electric $`𝐄`$ and the induced $`𝐃`$ fields was used:
$$D_i(x)=E_i(x)+4\pi d^4x^{}\chi _{ij}(x,x^{})E_j(x^{}),$$
(2)
as well as the condition $`D_i(x)/x_i=\mathrm{\hspace{0.17em}0}`$. By definition, differentiating the integral in eq.(2) over $`t`$ , we obtain $`𝐣(x)`$ being a density of the induced current. When the function $`\chi _{ij}(x,x^{})`$ depends only on the difference $`tt^{}`$ , the equation for the vector potential $`𝐀`$ is essentially the same as eq.(1) for the electric field.
In the vacuum the solution to eq.(1), satisfying the condition $`E_i(x)/x_i=\mathrm{\hspace{0.17em}0}`$, reads
$$𝐄_0(x)=𝑑𝐤g(𝐤𝐤_0)𝐞^{}(𝐤)e^{ikx},$$
(3)
where $`x`$ and $`k`$ are 4-vectors: $`x(t,𝐫)`$, $`k(\omega ,𝐤)`$ with $`𝐤=\omega 𝝂,𝝂^2=\mathrm{\hspace{0.17em}1}`$, and $`𝐞^{}(𝐤)`$ being an arbitrary vector, perpendicular to $`𝐤`$. If we introduce $`\delta _{ij}^{}=\delta _{ij}\nu _i\nu _j`$, then $`e_i^{}(𝐤)=\delta _{ij}^{}e_j`$. We assume that the function $`g`$ in eq.(3) vanishes except of the narrow region where $`𝐤𝐤_0𝐤_0,`$. So, the incident wave is a wave packet propagating along $`𝐤_0`$. Let it encounter a medium at some boundary . We try to satisfy eq.(1) in the medium using the field $`𝐄`$ in the form of
$$E_i(x)=𝑑𝐤g(𝐤𝐤_0)F_{ij}(𝐫,𝐤)e_j^{}(𝐤)e^{ikx},$$
(4)
with $`F_{ij}(𝐫,𝐤)=\delta _{ij}`$ on the boundary. Assuming that $`\omega _0`$ is sufficiently large, we expect that the function $`F_{ij}(𝐫,𝐤)`$ varies very slowly with respect to $`𝐫`$. Substituting eq.(4) into eq.(1) and neglecting the term $`^2F`$ ( keeping only $`𝐤/𝐫F`$), we obtain
$$\begin{array}{cc}& 0=d𝐤g(𝐤𝐤_0)\{\mathrm{\hspace{0.17em}2}ie^{ikx}k_l\frac{F_{ij}(𝐫,𝐤)}{r_l}+\hfill \\ & \\ & +\mathrm{\hspace{0.17em}4}\pi d^4x^{}R_{im}(x)\chi _{ml}(x,x^{})F_{lj}(𝐫^{},𝐤)e^{ikx^{}}\}e_j^{}(𝐤).\hfill \end{array}$$
(5)
The integral $`d^4x^{}`$ in eq.(5) can be easily taken if we use the Fourier-transforms:
$$\begin{array}{cc}& \chi _{ml}(x,x^{})=\frac{d^4k_1d^4k_2}{(2\pi )^8}e^{i(k_2x^{}k_1x)}\chi _{ml}(k_1,k_2),\hfill \\ & \\ & F_{lj}(𝐫^{},𝐤)=\frac{d𝐪}{(2\pi )^3}e^{i\mathrm{𝐪𝐫}^{}}F_{lj}(𝐪,𝐤).\hfill \end{array}$$
(6)
The most general form of the function $`\chi _{ml}(k_1,k_2)`$ in crystals reads
$$\chi _{ml}(k_1,k_2)=(2\pi )^4\underset{𝐐}{}c_{ml}(𝐐,k_1)\delta (k_2k_1Q),$$
(7)
where $`Q(0,𝐐)`$ is a reciprocal lattice vector. Using this representation, we obtain from eq.(5)
$$\begin{array}{cc}& 0=d𝐤g(𝐤𝐤_0)e^{ikx}\{\mathrm{\hspace{0.17em}2}ik_l\frac{F_{ij}(𝐫,𝐤)}{r_l}+\hfill \\ & \\ & +\mathrm{\hspace{0.17em}4}\pi \omega ^2\delta _{im}^{}\frac{d𝐪}{(2\pi )^3}\underset{𝐐}{}c_{ml}(𝐐,k+qQ)e^{i\mathrm{𝐐𝐫}}F_{lj}(𝐪,𝐤)e^{i\mathrm{𝐪𝐫}}\}e_j^{}(𝐤).\hfill \end{array}$$
(8)
The second argument of the function $`c_{ml}`$ in eq.(8) is $`k_1=k+qQ(\omega ,𝐤+𝐪𝐐)`$. This function depends actually on $`\omega _1=\omega `$ and $`𝐧_1=𝐤_1/𝐤_1𝝂+𝐬^{}`$, where $`𝐬=(𝐪𝐐)/\omega `$. We have already used the fact that $`𝐬\mathrm{\hspace{0.17em}1}`$ neglecting it as compared to unity. In particularly, that is why we have $`\delta _{im}^{}`$ in eq.(8) instead of $`\delta _{im}k_{1i}k_{1m}/\omega ^2`$. We shall see below that owing to this replacement the longitudinal components of $`F_{lj}(𝐫,𝐤)`$ do not appear at any depth if they are absent on the boundary. However, when $`𝝂`$ is almost parallel to some crystal axis, the transverse ( with respect to this axis ) component of $`𝝂`$ should be compared with $`𝐬`$. We can rule out $`𝐪𝐐`$ from the second argument of the function $`c_{ml}`$ when the angle of incidence $`\vartheta _0`$ with respect to this axis ( transverse component of $`𝝂`$ ) is sufficiently large: $`\vartheta _0𝐬`$. In what follows, we assume this condition to be fulfilled, bearing in mind that really it does not lead to serious limitations as a typical magnitude of $`𝐬`$ at $`\omega \mathrm{\hspace{0.17em}1}GeV`$ is $`\mathrm{\hspace{0.17em}10}^6`$. The integral over $`𝐪`$ turns into $`F_{lj}(𝐫,𝐤)`$ according to the definition eq.(6) , when we change $`c_{ml}(𝐐,k+qQ)c_{ml}(𝐐,k)`$ in eq.(8). This means that the integral over $`𝐫^{}`$ in eq.(5) converges in the domain $`𝐫𝐫^{}R`$ where $`R`$ is a characteristic scale for the noticeable change of the function $`F_{lj}(𝐫^{})`$. So that, we could substitute $`F_{lj}(𝐫^{})F_{lj}(𝐫)`$ yet in eq.(5). Now eq.(8) goes over into
$$\begin{array}{cc}& 0=d𝐤g(𝐤𝐤_0)e^{ikx}\{\mathrm{\hspace{0.17em}2}ik_l\frac{F_{ij}(𝐫,𝐤)}{r_l}+\hfill \\ & \\ & +\mathrm{\hspace{0.17em}4}\pi \omega ^2\delta _{im}^{}\underset{𝐐}{}c_{ml}(𝐐,k)e^{i\mathrm{𝐐𝐫}}F_{lj}(𝐫,𝐤)\}e_j^{}(𝐤).\hfill \end{array}$$
(9)
Here we are interested in the transverse ( with respect to $`𝝂`$ ) tensor $`F_{ij}`$. Let us analyze, whether such a form of the sought tensor is consistent with eq.(9) and boundary conditions. According to eq.(9), the left longitudinal components of this tensor $`F_j^{}=\nu _iF_{ij}(𝐫,𝐤)`$ are independent of the penetration depth. If we suppose that $`F_j^{}=\mathrm{\hspace{0.17em}0}`$ on the boundary, then $`F_{lj}\delta _{ln}F_{nj}=\delta _{ln}^{}F_{nj}`$ at any depth. In particularly, this implies that only transverse components of the tensor $`c_{ml}(𝐐,k)`$ are present in eq.(9). The right longitudinal components of the sought tensor do not enter eq.(9) owing to the factor $`e_j^{}(𝐤)`$. So, all the tensors below are self-consistently assumed to be transverse. They can be presented as two-dimensional matrices. Remembering this, we obtain from eq.(9)
$$F(𝐫,𝐤)=\mathrm{exp}\left\{\mathrm{\hspace{0.17em}2}\pi i\omega z\underset{𝐐}{}c(𝐐,k)e^{i𝐐𝝆}\underset{0}{\overset{1}{}}𝑑xe^{ixzQ_{}}\right\},$$
(10)
where $`𝝆=𝐫z𝝂`$,$`z=𝝂𝐫`$,$`Q_{}=𝝂𝐐`$. For $`Q_{}=\mathrm{\hspace{0.17em}0}`$ the integral in eq.(10) equals unity, while for $`Q_{}\mathrm{\hspace{0.17em}0}`$ it is of the order of $`(Q_{}z)^1`$ . Noticeable effects appear when the main ( $`Q_{}=\mathrm{\hspace{0.17em}0}`$ ) term in the phase of eq.(10) is of the order of unity. As shown in the next Section, it happens for $`z`$ of several centimeters when the contribution of terms with $`Q_{}\mathrm{\hspace{0.17em}0}`$ to this phase $`(Q_{}z)^1\mathrm{\hspace{0.17em}1}`$ and can be neglected. Since matrices $`_𝐐c_{ml}(𝐐,k)e^{i\mathrm{𝐐𝐫}}`$ are in general non-commutative at different $`z`$ , the exact solution to eq.(9) is the $`z`$ -ordered exponential function. However, corrections to eq.(10) due to $`z`$ -ordering are connected with terms $`Q_{}\mathrm{\hspace{0.17em}0}`$ being as small ($`(Q_{}z)^1`$ ) as terms already neglected. It must be remembered also that ( for the coherent yield ) $`𝐐`$ is a discrete vector and that large values of $`𝐐`$ do not contribute to the sum in eq.(10). These means that for sufficiently thick crystals we should keep only terms with $`Q_{}=\mathrm{\hspace{0.17em}0}`$ in the sum of eq.(10):
$$F(𝐫,𝐤)=\mathrm{exp}\left\{\mathrm{\hspace{0.17em}2}\pi i\omega z\underset{𝐐_{}}{}c(𝐐_{},k)e^{i𝐐_{}𝝆}\right\}.$$
(11)
The condition $`Q_{}=\mathrm{\hspace{0.17em}0}`$ implies that $`𝐐_{}=\mathrm{\hspace{0.17em}0}`$ too, if $`𝝂`$ is not aligned on some crystal plane. More precisely, let $`\psi `$ be the angle of $`𝝂`$ with respect to the plane. For $`\psi \psi _0=d_{pl}/z`$ , with $`d_{pl}`$ being the inter-planar distance, we can retain the single term $`c(𝐐=\mathrm{\hspace{0.17em}0},k)`$ of the sum in eq.(11). The angle $`\psi _0`$ is extremely small, even for very thin films. The integration over $`𝝂`$ in eq.(4) should smear out a dependence of the general result on $`𝝆`$. We are not interested here whether such a dependence is observable or not. In what follows, the sum $`_𝐐_{}c(𝐐_{},k)e^{i𝐐_{}𝝆}`$ in eq.(11) will be replaced by the matrix $`c(𝐐=\mathrm{\hspace{0.17em}0},k)c(0,k)`$ which coincides with the average of this sum over $`𝝆`$ . The matrix $`c(0,k)`$ has no singularities and is practically constant in the region $`\psi <\psi _0`$. So that, any directions are allowed for $`𝝂`$ .
Keeping in the sum $`_𝐐_{}c(𝐐_{},k)e^{i𝐐_{}𝝆}`$ the only term with $`𝐐_{}=0`$ , we disregard possible small-angle photon scattering. In this approximation there is no exchange between fractions of the wave-packet having different $`𝝂`$ ,i.e., these photons propagate independently. So that, for given direction, we can go over to the conventional idealization of a monochromatic plane wave. . Then the electric field is obtained from eq.(11) and eq.(4) with the change in the latter $`g(𝐤𝐤_0)\delta (𝐤𝐤_0)`$ and then $`𝐤_0𝐤`$
$$E(x)=e^{i\omega t}e^{i\omega zn}e^{}(𝐤),$$
(12)
where the two-dimensional vectors $`E(x)`$ and $`e^{}(𝐤)`$ are correspondingly the electric field in a medium and the polarization vector of the incident wave. The quantity $`n=\delta _{ij}^{}+2\pi c(0,k)_{ij}`$ is a complex matrix representing, by definition, the index of refraction. It depends on $`\omega `$ and ( in crystals ) on $`𝝂`$ . We introduce a matrix $`\mathrm{\Pi }=\mathrm{\hspace{0.17em}4}\pi \omega ^2c(0,k)`$, then $`n=I+\mathrm{\Pi }/(2\omega ^2)`$. The photon density matrix at depth $`L`$ reads
$$\rho (L)=\mathrm{exp}\left\{iL\frac{\mathrm{\Pi }}{2\omega }\right\}\rho (0)\mathrm{exp}\left\{iL\frac{\mathrm{\Pi }^{}}{2\omega }\right\},$$
(13)
where $`\rho (0)`$ is the density matrix of the incident photon.
Comparing eq.(12) with the corresponding result of obtained for an amorphous medium by the direct summation of fields coming from individual scatterers, we find that in this case $`\mathrm{\Pi }=\mathrm{\hspace{0.17em}4}\pi f(0)N`$. In this formula $`f(0)`$ is the forward scattering amplitude for an individual particle ( the cross section $`d\sigma /d\mathrm{\Omega }=f(𝚫)^2`$ , $`𝚫`$$`=𝐤_2𝐤_1`$ ) and $`N`$ is a number of particles per unit volume. So, to find the polarization operator $`\mathrm{\Pi }`$, we have to calculate the forward scattering amplitude normalized in such a way, that it reproduces $`\mathrm{\hspace{0.17em}4}\pi f(0)`$ for an individual particle. Moreover, using the amplitude ( see Sec.3 for its normalization ) $`T(k_1,k_2)`$ for an arbitrary momentum transfer $`𝚫`$ , we obtain for $`\chi (k_1,k_2)`$ defined in eq.(6)
$$\chi (k_1,k_2)=T(k_1,k_2)/(2\omega ^2).$$
(14)
From this relation and eq.(7), we can extract $`c_{ml}(𝐐,k_1)`$ and then, if needed, take into account terms with $`𝐐\mathrm{\hspace{0.17em}0}`$ omitted at the transition from eq.(10) to eq.(12).
As any other $`\mathrm{\hspace{0.17em}2}\times 2`$ matrix, $`\mathrm{\Pi }/(2\omega )`$ in eq.(13) can be presented in the form of $`\mathrm{\Pi }/(2\omega )a+𝐛𝝈`$ , where $`𝝈`$ are the Pauli matrices. Note that the eigenvalues of the matrix $`a+𝐛𝝈`$ are $`a\pm \sqrt{𝐛^2}`$. Using a similar representation for the initial density matrix $`\rho (0)=(I+𝜼𝝈)/2`$ , with $`𝜼`$ being the initial Stokes vector, we find from eq.(13)
$$\begin{array}{cc}& \rho (L)=\frac{1}{2}e^{\gamma L}\{\frac{1+𝜼𝝈}{2}(\mathrm{cosh}\alpha +\mathrm{cos}\beta )+(\mathrm{cosh}\alpha \mathrm{cos}\beta )[\frac{1𝜼𝝈}{2}(𝐠_1^2+𝐠_2^2)+\hfill \\ & \\ & +(𝐠_1\times 𝐠_2,𝝈𝜼)+\left(𝐠_1𝜼\right)\left(𝐠_1𝝈\right)+\left(𝐠_2𝜼\right)\left(𝐠_2𝝈\right)]\hfill \\ & \\ & [(𝐠_2,𝝈+𝜼)+(𝐠_1\times 𝜼,𝝈)]\mathrm{sin}\beta [(𝐠_1,𝝈+𝜼)(𝐠_2\times 𝜼,𝝈)]\mathrm{sinh}\alpha \},\hfill \end{array}$$
(15)
where $`\gamma =\mathrm{\hspace{0.17em}2}Im(a)`$ , $`\alpha =\mathrm{\hspace{0.17em}2}LIm(\sqrt{𝐛^2})`$ , $`\beta =\mathrm{\hspace{0.17em}2}LRe(\sqrt{𝐛^2})`$ , $`𝐠𝐠_1+i𝐠_2=𝐛/\sqrt{𝐛^2}`$. The real vectors $`𝐠_1`$ and $`𝐠_2`$ satisfy the conditions $`𝐠_1𝐠_2=\mathrm{\hspace{0.17em}0}`$ , $`𝐠_1^2𝐠_2^2=\mathrm{\hspace{0.17em}1}`$ , since $`𝐠^2=\mathrm{\hspace{0.17em}1}`$. As $`tr\rho (0)=\mathrm{\hspace{0.17em}1}`$ , the fraction of outgoing photons is given by $`FRAC=tr\rho (L)`$. The Stokes vector at the depth $`L`$ is $`𝜼(L)=tr𝝈\rho (L)/FRAC`$. In the explicit form:
$$𝜼(L)=\frac{𝐒}{P},FRAC=e^{\gamma L}P,𝜼^2(L)=\mathrm{\hspace{0.17em}1}\frac{1𝜼^2}{P^2},$$
(16)
where
$$\begin{array}{cc}& P=𝐠_1^2\mathrm{cosh}\alpha 𝐠_2^2\mathrm{cos}\beta (𝐠_1,𝜼)\mathrm{sinh}\alpha (𝐠_2,𝜼)\mathrm{sin}\beta \hfill \\ & \\ & (\mathrm{cosh}\alpha \mathrm{cos}\beta )(𝐠_1\times 𝐠_2,𝜼);\hfill \\ & \\ & 𝐒=𝜼\left(𝐠_1^2\mathrm{cos}\beta 𝐠_2^2\mathrm{cosh}\alpha \right)\left[𝐠_1𝐠_2\times 𝜼\right]\mathrm{sinh}\alpha \left[𝐠_2+𝐠_1\times 𝜼\right]\mathrm{sin}\beta +\hfill \\ & \\ & [𝐠_1\times 𝐠_2+𝐠_1(𝐠_1,𝜼)+𝐠_2(𝐠_2,𝜼)](\mathrm{cosh}\alpha \mathrm{cos}\beta ).\hfill \end{array}$$
(17)
We emphasize that, according to eq.(16), the polarization degree $`𝜼(L)`$ is an increasing function of depth $`L`$ except for $`𝜼^2=\mathrm{\hspace{0.17em}1}`$ or $`P^2=\mathrm{\hspace{0.17em}1}`$.
## 3 Scattering of photons in a crystal
We now assume that $`\omega m`$. Generally speaking, the amplitude of the forward Compton scattering off electrons $`f_C(0)`$ may be of comparable size with that due to virtual $`e^+e^{}`$ pairs ( VP ) even in $`GeV`$ energy region ( see discussion in ). However, we can neglect $`f_C(0)`$ in what follows since it is, first of all, the unit matrix multiplied by some scalar factor and, second, real . That is why its contributions to $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }^{}`$ in eq.(13) are cancelled. For the same reason we can omit the real part of the incoherent contribution to the forward scattering VP amplitude. Nevertheless, its imaginary part, which does not cancel in eq.(13), will be taken into account. It is the unit matrix multiplied by $`\omega W_p`$, with $`W_p`$ being the probability per unit length of the $`e^+e^{}`$ pair production by a photon. Remember that the incoherent probability $`W_p`$ in crystals is smaller than in corresponding amorphous media (see ,e.g. ). So, we must calculate here only the coherent contribution to the VP amplitude.
The properly normalized transition amplitude of the photon with 4-momentum $`k_{1\mu }=(\omega _1,𝐤_1)`$ and polarization vector $`e_{1\mu }`$ into the photon with $`k_{2\mu }=(\omega _2,𝐤_2)`$, $`e_{2\mu }`$ reads
$$T(k_1,k_2)=\mathrm{\hspace{0.17em}2}i\alpha d^4x_1d^4x_2Tr\left[G(x_2,x_1)\widehat{e}_1e^{ik_1x_1}G(x_1,x_2)\widehat{e}_2^{}e^{ik_2x_2}\right],$$
(18)
where $`\alpha =\mathrm{\hspace{0.17em}1}/137`$ is the fine structure constant. The electron Green function $`G(x_1,x_2)`$ can be expressed via solutions $`\mathrm{\Psi }_n^{(\pm )}(x)`$ to the Dirac equation in the corresponding external field
$$iG(x_2,x_1)=\theta (t_2t_1)\underset{n}{}\mathrm{\Psi }_n^{(+)}(x_2)\overline{\mathrm{\Psi }}_n^{(+)}(x_1)\theta (t_1t_2)\underset{n}{}\mathrm{\Psi }_n^{()}(x_2)\overline{\mathrm{\Psi }}_n^{()}(x_1).$$
(19)
Recollect that ( see discussion in ) for $`\omega m`$ the contribution to the amplitude , in terms of the non-covariant perturbation theory, is given by the diagram where the pair production by the initial photon precedes the annihilation of this pair into the final photon . We find, keeping only terms proportional to $`\theta (t_2t_1)`$
$$\begin{array}{cc}& T(k_1,k_2)=\mathrm{\hspace{0.17em}2}i\alpha \underset{n,m}{}d^4x_1d^4x_2\theta (t_2t_1)V_{nm}(x_1,e_1,k_1)V_{nm}^+(x_2,e_2,k_2),\hfill \\ & \\ & V_{nm}(x,e,k)=\overline{\mathrm{\Psi }}_n^{(+)}(x)\widehat{e}e^{ikx}\mathrm{\Psi }_m^{()}(x).\hfill \end{array}$$
(20)
It is really important that in the Fourier-transform of a crystal potential
$$U(𝐪)=𝑑𝐫e^{i\mathrm{𝐪𝐫}}U(𝐫)$$
only small momenta $`𝐪m`$ are present. For such potentials, the quasiclassical operator method can be applied in its standard form to calculate the amplitude $`T(k_1,k_2)`$. Then further transformations in eq.(20) are completely identical to those leading to the quasiclassical expression for the probability of $`e^+e^{}`$ pair production by a photon in an external field ( see Sec.3 in ), because $`V_{nm}`$ in eq.(20) is essentially the matrix element for this process. Notice that the amplitude $`T(k_1,k_2)`$ is the contraction of $`e_{1\mu }e_{2\nu }^{}`$ with the sought tensor $`T^{\mu \nu }(k_1,k_2)`$. For the transverse ( with respect to $`𝝂_1=𝐤_1/\omega _1`$ ) components of this tensor we obtain from eq.(20)
$$\begin{array}{cc}& T^{ij}=\frac{i\alpha }{(2\pi )^2}\delta (\omega _2\omega _1)𝑑𝐫exp(i𝚫𝐫)\underset{0}{\overset{\mathrm{}}{}}𝑑\tau \underset{0}{\overset{\omega }{}}𝑑\epsilon \left[\frac{\omega }{\epsilon (\omega \epsilon )}\right]^2𝑑𝐩_{}(0)e^{iA}B^{ij},\hfill \\ & \\ & A=\frac{\omega m^2}{2\epsilon (\omega \epsilon )}\left[\tau +\underset{\tau /2}{\overset{\tau /2}{}}\frac{ds}{m^2}𝐩_{}^2(s)\right];\hfill \\ & \\ & B^{ij}=\delta _{}^{ij}\left[m^2+\left(𝐩_1𝐩_2\right)\right]+\left(\frac{2\epsilon \omega }{\omega }\right)^2p_1^ip_2^jp_2^ip_1^j,\hfill \end{array}$$
(21)
where $`𝐩_{1,2}𝐩(\tau /2)`$ being the momentum of an electron on a classical trajectory at the corresponding time, $`𝚫=𝐤_2𝐤_1`$. When the external field vanishes, the amplitude $`T`$ must vanish as well. It assumes a subtraction which really affects only the term proportional to $`m^2`$ in $`B^{ij}`$. This subtraction will be performed in the explicit form below.
Generally speaking, for further calculations we have to know the dependence of $`𝐩_{}(t)`$ on time ,i.e. the corresponding mechanical problem should be solved. For arbitrary potential it cannot be done analytically. Fortunately, it is sufficient to know the alteration of the quantity $`𝐩_{}(t)`$ during the process formation time $`\tau _f`$. In eq.(21) it is the characteristic size of the variable $`\tau `$ contributing to the integral. In the rectilinear trajectory approximation (RTA), we substitute into equations of motion $`𝐫+𝝂_1t`$ instead of the exact solution $`𝐫(t)`$. We refer for details to stating here that in crystals the amplitude $`T`$ can be calculated by the use of RTA for arbitrary photon energy ($`\omega m`$) and crystal orientation.
Representing the crystal potential $`U(𝐫)`$ as a sum over vectors $`𝐪`$ of the reciprocal lattice
$$U(𝐫)=\underset{𝐪}{}G(𝐪)e^{i\mathrm{𝐪𝐫}},$$
we find within RTA for the transverse momentum
$$𝐩_{}(t)=𝐩_{}(0)+𝜹(t)𝐩_{}(0)\underset{𝐪}{}G(𝐪)e^{i\mathrm{𝐪𝐫}}\frac{𝐪_{}}{q_{}}\left(e^{iq_{}t}\mathrm{\hspace{0.17em}1}\right),$$
(22)
where $`q_{}=(𝝂_1𝐪)`$. Substituting eq.(22) into eq.(21) , we can take the Gaussian-type integral over $`𝐩_{}(0)`$. Performing the subtraction at vanishing external field and going over to the variable $`y=\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}2}\epsilon /\omega `$ , we find
$$\begin{array}{cc}& T^{ij}=\frac{\alpha }{\pi }\delta (\omega _2\omega _1)𝑑𝐫exp(i𝚫𝐫)\underset{1}{\overset{1}{}}\frac{dy}{1y^2}\underset{0}{\overset{\mathrm{}}{}}\frac{d\tau }{\tau }\mathrm{exp}\left[i\beta \tau \left(1+\psi (\tau )\right)\right]B_1^{ij};\hfill \\ & \\ & \psi (\tau )=\underset{𝐪,𝐪^{}}{}\frac{G(𝐪)G(𝐪^{})}{m^2}\frac{(𝐪_{}𝐪_{}^{})}{q_{}q_{}^{}}e^{i(𝐪+𝐪^{},𝐫)}\left[\frac{sin(q_{}+q_{}^{})\tau /2}{(q_{}+q_{}^{})\tau /2}\frac{sinq_{}\tau /2}{q_{}\tau /2}\frac{sinq_{}^{}\tau /2}{q_{}^{}\tau /2}\right]\hfill \\ & \\ & B_1^{ij}=y^2a_1^ia_2^ja_2^ia_1^j+\delta _{}^{ij}\left\{\left(𝐚_1𝐚_2\right)+m^2\left[(y^21)\tau \frac{d\psi }{d\tau }\psi \right]\right\};\hfill \\ & \\ & 𝐚_{1,2}=\underset{𝐪}{}G(𝐪)\frac{𝐪_{}}{q_{}}e^{i\mathrm{𝐪𝐫}}\left[e^{iq_{}\tau /2}\frac{sinq_{}\tau /2}{q_{}\tau /2}\right],\beta =\frac{2m^2}{\omega (1y^2)}.\hfill \end{array}$$
(23)
This expression is still rather complicated for the numerical calculations. It would be drastically simplified if we neglect the function $`\psi (\tau )`$ as compared to $`\mathrm{\hspace{0.17em}1}`$ in the phase of eq.(23). We suppose that the angle of incidence $`\vartheta _0`$ with respect to some major crystal axis $`\vartheta _0\mathrm{\hspace{0.17em}1}`$, since precisely for such angles the strengthening of electromagnetic effects happens in crystals as compared with amorphous media. If $`𝝂_1`$ is far in azimuth from any major crystal plane, then $`q_{}𝐪_{}\vartheta _0`$ and $`\psi _{off}(\tau )(V_0/m/\vartheta _0)^2`$. In this estimate $`V_0`$ denotes a typical magnitude of $`G(𝐪)`$ being of the order of the axial potential well depth. So, in this case we can omit $`\psi (\tau )`$ in the phase of eq.(23) for $`\vartheta _0V_0/m`$. If now $`𝝂_1`$ is aligned on some crystal plane, there is a subset of $`𝐪(𝐪^{})`$ for which $`q_{}(q_{}^{})`$ are extremely small or vanish. Those $`𝐪`$ are perpendicular to the plane. Expanding in $`q_{}`$ and $`q_{}^{}`$ , we obtain for the corresponding contribution to the phase $`\psi _{pl}(\tau )=(\tau /m)^2(dU_{pl}(x)/dx)^2/12`$ . Here $`x`$ is the distance from the plane, $`U_{pl}(x)`$ being the planar potential. As long as $`\psi (\tau )\mathrm{\hspace{0.17em}1}`$ , we can substitute $`\tau \beta ^1`$ into our estimate which turns into $`\psi _{pl}(\tau )\kappa ^2(x)`$. The magnitude of the strong field parameter
$$\kappa (x)=\frac{E_{pl}(x)}{E_0}\frac{\omega }{m},E_0=\frac{m^2}{e}=\mathrm{\hspace{0.17em}1.32}10^{16}V/cm$$
can be estimated using the Table 1 of , where, in particularly, averaged over $`x`$ values of $`(E_{pl}(x)/E_0)^2`$ are presented for $`(110)`$ plane of several crystals. As a result, for commonly used crystals, $`\psi _{pl}(\tau )`$ is sufficiently small when $`\omega `$ is less than several $`TeV`$. Assuming that $`\omega `$ and $`\vartheta _0`$ satisfy the formulated conditions, we, finally, rule out $`\psi (\tau )`$ from the phase of eq.(23) .
Now the integration over r in eq.(23) can be performed
$$𝑑𝐫exp\left[i(𝐪+𝐪^{}𝚫,𝐫)\right]=(2\pi )^3\delta (𝐪+𝐪^{}𝚫).$$
As a result, the tensor $`T^{ij}`$ acquires the form
$$T^{ij}(k_1,k_2)=(2\pi )^3\underset{𝐐}{}\pi ^{ij}(𝐐,k_1)\delta (k_2k_1Q),$$
where $`Q=𝐪+𝐪^{}`$. Using this presentation of $`T^{ij}(k_1,k_2)`$ and eq.(7),eq.(14) we find that $`c^{ij}(𝐐,k_1)`$ in eq.(7) are $`c^{ij}(𝐐,k_1)=\pi ^{ij}(𝐐,k_1)/(4\pi \omega ^2)`$. Here we are interested only in $`\pi ^{ij}(0,k_1)`$ being just the polarization operator $`\mathrm{\Pi }`$ which determines a development of the photon density matrix at large thicknesses according to eq.(13) . Taking elementary integrals over $`\tau `$ and $`y`$, we obtain from eq.(23)
$$\mathrm{\Pi }^{ij}=\frac{\alpha \omega ^2}{8m^2}\underset{𝐪}{}\frac{G(𝐪)^2}{m^2}\left\{𝐪_{}^2\delta _{}^{ij}\left[id_1(\mu )+d_2(\mu )\right]+\left(\frac{1}{2}𝐪_{}^2\delta _{}^{ij}q_{}^iq_{}^j\right)\left[id_3(\mu )d_4(\mu )\right]\right\},$$
(24)
where $`\mu =\mathrm{\hspace{0.17em}2}m^2/(\omega q_{})`$ and functions $`d_l`$ are
$$\begin{array}{cc}& d_1(x)=x^2\left[\left(1+x\frac{1}{2}x^2\right)g(\sqrt{1x})(1+x)\sqrt{1x}\right]\theta (1x),\hfill \\ & \\ & d_3(x)=x^3(\frac{x}{2}g(\sqrt{1x})+\sqrt{1x})\theta (1x),g(x)=\mathrm{ln}\frac{1+x}{1x},\hfill \\ & \\ & d_4(x)=\frac{x^2}{\pi }\{(\frac{x}{2}g(\sqrt{1+x})\sqrt{1+x})^2(x\mathrm{arctan}\frac{1}{\sqrt{x1}}\sqrt{x1})^2\theta (x1)+\hfill \\ & \\ & +[(\frac{x}{2}g(\sqrt{1x})+\sqrt{1x})^2\left(\frac{\pi x}{2}\right)^2]\theta (1x)\},\hfill \\ & \\ & d_2(x)=d_4(x)+\frac{2x^2}{\pi }\{\frac{1}{2}g(\sqrt{1+x})(\frac{x1}{2}g(\sqrt{1+x})+\sqrt{1+x})3\hfill \\ & \\ & \left[\frac{1}{2}g(\sqrt{1x})\left(\frac{x+1}{2}g(\sqrt{1+x})\sqrt{1x}\right)+1\frac{\pi ^2(1+x)}{4}\right]\theta (1x)+\hfill \\ & \\ & +[\mathrm{arctan}\frac{1}{\sqrt{x1}}((x+1)\mathrm{arctan}\frac{1}{\sqrt{x1}}+\sqrt{x1})1]\theta (x1)\},\hfill \end{array}$$
(25)
where $`\theta (x)=\mathrm{\hspace{0.17em}1}`$ for $`x>\mathrm{\hspace{0.17em}0}`$ and vanishes for $`x<\mathrm{\hspace{0.17em}0}`$.
The photon energy $`\omega `$ enters eq.(24) in the combination $`\mu `$ except of the factor $`\omega ^2`$ in front of the sum. This sum is almost independent of $`\omega `$ near its maximum since the latter is given by $`\mu \mathrm{\hspace{0.17em}1}`$. As a result, the $`\omega `$-dependence of the quantities $`\alpha `$ and $`\beta `$ in eq.(15) and eq.(17) describing the behavior of the photon polarization is reduced to the factor $`\omega L`$ for optimal orientations. Correspondingly, the optimal thickness $`L_{opt}`$ is roughly proportional to $`\omega ^1`$. As mentioned above, there is a subset of $`𝐪`$ at perfect planar alignment for which $`q_{}`$ vanishes ($`\mu \mathrm{}`$). The contribution of this subset to $`\mathrm{\Pi }^{ij}`$ , i.e. the yield of the corresponding planar potential obtained from eq.(24) reproduces well known results ( see, e.g. and literature cited there ) derived within the Born approximation:
$$\frac{\alpha m^2}{45\pi }\kappa ^2(x)\left(4e_1^ie_1^j+\mathrm{\hspace{0.17em}7}e_2^ie_2^j\right),$$
where $`\mathrm{}`$ means averaging over the coordinate $`x`$ , $`𝐞_1`$ is the unit vector perpendicular to the plane and $`𝐞_2=𝝂_1\times 𝐞_1`$ . What is lost when $`\psi (\tau )`$ has been omitted in the phase of eq.(23), are higher order corrections in crystal potential .
Using the explicit form of $`\mathrm{\Pi }^{ij}`$ eq.(24) and adding the incoherent yield,we can find the quantities presented in eq.(15), which describes the properties of a photon beam for any initial conditions dependent on the crystal thickness. As explained above, the incoherent yield is present in eq.(15) only in the absorption coefficient $`\gamma `$. We can use any basis to calculate $`\mathrm{\Pi }^{ij}`$. Let this basis be formed by two real unit vectors $`𝐞_1`$ and $`𝐞_2`$ satisfying $`𝐞_1𝐞_2=𝝂_1𝐞_1=𝝂_1𝐞_2=\mathrm{\hspace{0.17em}0}`$. We should use the same basis to obtain the initial Stokes vector $`𝜼`$ . Supposing that $`𝝂_1`$ is near some axis direction $`𝝂_3`$, we choose $`𝐞_2`$ in the plane ( reaction plane ) containing $`𝝂_1`$ and $`𝝂_3`$. For the basis chosen, the circular polarization degree is given by the magnitude of the second Stokes parameter $`\xi (L)=\eta _2(L)`$. If we define the angle of incidence $`\vartheta _0`$ as the angle between $`𝝂_1`$ and $`𝝂_3`$, and $`\phi _0`$ which is the angle between the reaction plane and the $`(1\overline{1}0)`$ plane, then the explicit form of the basis vectors reads
$$𝐞_1=𝐞_x\mathrm{sin}\phi _0+𝐞_y\mathrm{cos}\phi _0,𝐞_2=𝝂_3\mathrm{sin}\vartheta _0(𝐞_x\mathrm{cos}\phi _0+𝐞_y\mathrm{sin}\phi _0)\mathrm{cos}\vartheta _0,$$
where $`𝐞_x𝐞_y=𝝂_3𝐞_x=𝝂_3𝐞_y=\mathrm{\hspace{0.17em}0}`$ and $`𝐞_x`$ is in the $`(1\overline{1}0)`$ plane. In this basis the quantities $`a`$ and $`𝐛`$ presented in eq.(15) take the form
$$\begin{array}{cc}& (a,𝐛)=\frac{\alpha \omega }{16m^2}\underset{𝐪}{}\frac{G(𝐪)^2}{m^2}(A,𝐁),\hfill \\ & \\ & A=𝐪_{}^2[id_1(\mu )+d_2(\mu )],B_1=(𝐞_1𝐪_{})(𝐞_2𝐪_{})[id_3(\mu )d_4(\mu )],\hfill \\ & \\ & B_2=\mathrm{\hspace{0.17em}0},B_3=\frac{1}{2}[(𝐞_2𝐪_{})^2(𝐞_1𝐪_{})^2][id_3(\mu )d_4(\mu )].\hfill \end{array}$$
(26)
Remember that the eigenvalues of the propagation matrix $`\mathrm{\Pi }/(2\omega )`$ in eq.(13) are $`a\pm \sqrt{𝐛^2}`$. The corresponding eigenvectors $`𝐞_\pm `$ satisfying the normalization condition $`𝐞_\pm ^2=\mathrm{\hspace{0.17em}1}`$ are
$$𝐞_+=\frac{𝐞_1+r𝐞_2}{\sqrt{1+r^2}},𝐞_{}=\frac{𝐞_2r𝐞_2}{\sqrt{1+r^2}},r=\frac{b_1}{b_3+\sqrt{𝐛^2}}.$$
(27)
In general, the eigenvectors $`𝐞_\pm `$ are complex. However, they become real when the quantity $`r`$ does so. In particularly, when $`𝝂_1`$ is in the symmetry plane of a crystal like $`(1\overline{1}0)`$ plane ($`\phi _0=\mathrm{\hspace{0.17em}0}`$), $`r`$ vanishes and $`𝐞_\pm `$ coincide with $`𝐞_{1,2}`$. Just such a case ($`\phi _0=\pi /2`$,$`𝝂_3`$ along $`<110>`$ axis) was the only orientation considered in where a quantity accounting for the polarization conversion was calculated by the use of dispersion relations. In our notation it corresponds to the term in eq.(24) proportional to $`d_4(\mu )`$ . We were unable to reproduce the factor in eq.(5) of , nevertheless, we emphasize that the expression in braces of the cited equation coincides with $`d_4(\mu )/\mu ^2`$.
Consider now, as an example, fully linearly polarized ($`\eta _2=\mathrm{\hspace{0.17em}0},𝜼^2=\mathrm{\hspace{0.17em}1}`$) initial photon beam. An efficiency of the polarization conversion process is determined not only by $`\xi (L)`$ but also by the fraction of surviving photons $`FRAC(L)`$. We use the criterion of such an efficiency suggested in
$$FOM(L)=\mathrm{\hspace{0.17em}10}\xi (L)\sqrt{FRAC(L)},$$
$`L`$ being the crystal thickness. Recollect that $`\xi (L)`$ and $`FRAC(L)`$ depend also on $`\omega `$ and $`𝝂_1`$. For given orientation of a crystal , we still have a free parameter which is the angle $`\varphi `$ of the initial polarization vector with respect to the basis vectors $`𝐞_1`$ and $`𝐞_2`$. We do not claim here to the final analysis of the polarization conversion process, so that the yield for different orientations will be compared at the same initial condition. Namely, we set $`\eta _1=1`$ which corresponds to $`\varphi =\pi /4`$ with respect to $`𝐞_2`$ ( $`\varphi =\mathrm{\hspace{0.17em}3}\pi /4`$ with respect to $`𝐞_1`$). Evidently, this is the best choice for the alignment on the symmetry plane of a crystal when the basis vectors are the eigenvectors of the matrix $`\mathrm{\Pi }^{ij}`$ as well.
In Fig.1 the maximum values of the figure of merit ( $`FOM`$ ) are shown near $`<110>`$ axis of a diamond crystal for $`\omega =\mathrm{\hspace{0.17em}100}GeV`$ as a function of the orientation (angles $`\vartheta _0`$ and $`\phi _0`$) . For each direction $`(\vartheta _0,\phi _0)`$, the figure of merit was calculated first as a function of $`L`$ , then its maximum value was found. Just these maximum values of FOM are plotted, so that different directions correspond usually to different optimal thicknesses $`L_{opt}`$. Owing to the crystal symmetry, there is no need to perform calculations for $`\phi _0`$ beyond the interval chosen, since they will simply reproduce the results already obtained for corresponding $`\phi _0`$ within the interval. For diamond crystal, the largest effect is achieved for $`\vartheta _0=\mathrm{\hspace{0.17em}1.5}mrad`$ off the $`<110>`$ axis on the $`(1\overline{1}0`$ plane ($`\phi _0=\mathrm{\hspace{0.17em}0}`$), where ( see Fig.1 ) $`FOM\mathrm{\hspace{0.17em}3.7}`$. However, in this case the size of $`L_{opt}\mathrm{\hspace{0.17em}4.7}cm`$ seems to be too large for practical use. Generally, a typical size of the optimal thickness is of a few $`cm`$. For the same axis of a silicon crystal ( see Fig.2 ) , the peak is at $`\vartheta _0=\mathrm{\hspace{0.17em}2.3}mrad`$ and $`\phi _0=\mathrm{\hspace{0.17em}0}`$ with $`FOM\mathrm{\hspace{0.17em}2.6}`$. We emphasize rather narrow angular with of the peak in both directions. The optimal thickness in this case is $`L_{opt}\mathrm{\hspace{0.17em}8}cm`$. We have performed the same kind of calculations for three major axes of diamond, silicon and germanium crystals. Comparing types of crystals, we see that the effect is the largest for diamond and the smallest for germanium, where $`FOM\mathrm{\hspace{0.17em}2.0}`$ can be obtained at $`L_{opt}\mathrm{\hspace{0.17em}2.8}cm`$. What about orientations, the $`<110>`$ axis is the most preferable, while two others give a comparable but smaller yield.
For $`\phi _0=\mathrm{\hspace{0.17em}0}`$ the position of a peak is determined by the condition $`\mu =\mathrm{\hspace{0.17em}2}m^2/(\omega q_{})=\mathrm{\hspace{0.17em}1}`$ for the smallest non-zero $`q_{}`$. This is connected with the threshold behavior of the functions $`d_{1,3}(\mu )`$ in eq.(24) at $`\mu =\mathrm{\hspace{0.17em}1}`$. From this condition we obtain for $`fcc(d)`$ structure near $`<110>`$ axis $`\vartheta _0^{max}=m^2l_c/(\pi \omega )`$ , where $`l_c`$ is a lattice constant. In particularly, for $`Si`$ we have $`\vartheta _0^{max}(mrad)\mathrm{\hspace{0.17em}229}/\omega (GeV)`$. This fact is illustrated by Fig.3 where $`FOM`$ (upper curves), $`FRAC`$ (lower curves), and the degree of circular polarization $`\xi `$ are shown for $`\phi _0=\mathrm{\hspace{0.17em}0}`$ near $`<110>`$ axis of a $`\mathrm{\hspace{0.17em}10}cm`$ thick silicon crystal as functions of $`\omega `$. Three sets of curves ( from the left to the right ) in Fig.3 correspond to the angles of incidence $`\vartheta _0(mrad)=`$ 2.29 , 2.08 and 1.91 respectively. All curves in Fig.3 have peaks exactly at the positions prescribed by the condition obtained above. For a given orientation ( at fixed $`\vartheta _0`$ ), rather narrow shape of these peaks does not allow us to handle with the same efficiency a photon beam having the wide energy spread.
The magnitude of $`FOM(L)`$ being proportional to $`\xi (L)`$ diminishes for partially polarized initial photon beam roughly proportionally to $`𝜼<\mathrm{\hspace{0.17em}1}`$ as compared to the fully polarized case ($`𝜼=\mathrm{\hspace{0.17em}1}`$ ). Depending on the orientation, more or less noticeable change of the polarization degree $`𝜼(L)`$ occurs for $`𝜼<\mathrm{\hspace{0.17em}1}`$. It can be seen in Figs.4,5 where the absolute values of three Stokes parameters and the polarization degree are presented as functions of the silicon crystal thickness $`L`$ at $`\omega =\mathrm{\hspace{0.17em}100}GeV`$. The calculations were carried out using eq.(16), eq.(17), and eq.(26) for $`\eta _1=0.5`$ , $`\eta _2=\eta _3=\mathrm{\hspace{0.17em}0}`$. The angle of incidence $`\vartheta _0=\mathrm{\hspace{0.17em}2.29}mrad`$ is the same for both figures. At $`\phi _0=\mathrm{\hspace{0.17em}0}`$ ( Fig.4 ), $`\eta _3(L)`$ is small and $`𝜼(L)`$ ( curve ”tot” ) practically does not change due to the smallness of the parameter $`\alpha `$ ( see eq.(15)) within the whole interval of $`L`$ presented. The fracture of the curve (1) means only that $`\eta _1(L)`$ changes its sign at $`L=\mathrm{\hspace{0.17em}21}cm`$. Since the polarization degree practically conserves and $`\eta _3(L)`$ can be neglected for this orientation, we could measure the linear polarization $`\eta _1(L)`$ to determine the circular polarization $`\eta _2(L)`$ appeared. However, the situation drastically changes already for a small alteration of $`\phi _0`$. It is seen in Fig.5 calculated at $`\phi _0=\mathrm{\hspace{0.17em}0.015}`$. Now the measurement of the linear polarization does not help in the determination of the circular one since the polarization degree is no more constant.
In conclusion, formulas derived describe the propagation of hard polarized photons through crystals. They are valid in a wide photon energy range for any orientation and any crystal type as long as the approximations used are correct. Our calculations show that the linear polarization of multi-GeV photons can be converted with an appropriate efficiency into the circular one using properly chosen single crystals. However, if we do not use theoretical results for some quantities involved ( e.g. for the polarization degree ), the only way to determine the circular polarization appeared in a crystal is the direct measurement of it.
Acknowledgements
The author is grateful to V.M.Katkov and A.I.Milstein for many fruitful discussions.
Figure captions
1. Fig.1 The maximum values of $`FOM`$ ( the figure of merit , see text ) near $`<110>`$ axis of a diamond crystal for $`\omega =\mathrm{\hspace{0.17em}100}GeV`$ dependent on the angles of incidence $`\vartheta _0`$ and $`\phi _0`$.
2. Fig.2 The same as in Fig.1 but for a silicon crystal in the neighborhood of the point $`\vartheta _0=\mathrm{\hspace{0.17em}2.3}mrad`$ , $`\phi _0=\mathrm{\hspace{0.17em}0}`$.
3. Fig.3 $`FOM`$ (upper curves), $`FRAC`$ (lower curves), and the degree of circular polarization $`\xi `$ for $`\phi _0=\mathrm{\hspace{0.17em}0}`$ near $`<110>`$ axis of a $`\mathrm{\hspace{0.17em}10}cm`$ thick silicon crystal as functions of $`\omega (GeV)`$. Three sets of curves ( from the left to the right ) correspond to the angles of incidence $`\vartheta _0(mrad)=`$ 2.29 , 2.08 and 1.91 respectively.
4. Fig.4 Absolute values of the Stokes parameters $`\eta _1(L)`$ (curve 1), $`\eta _2(L)`$ (curve 2), $`\eta _3(L)`$ (curve 3), and the polarization degree $`𝜼(L)`$ (curve ”tot” ) depending on the thickness $`L`$ of a silicon crystal at $`\omega =\mathrm{\hspace{0.17em}100}GeV`$, $`\eta _1=0.5`$ , $`\eta _2=\eta _3=\mathrm{\hspace{0.17em}0}`$, $`\vartheta _0=\mathrm{\hspace{0.17em}2.29}mrad`$, $`\phi _0=\mathrm{\hspace{0.17em}0}`$ .
5. Fig.5 The same as in Fig.4 but for $`\phi _0=\mathrm{\hspace{0.17em}0.015}`$.
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# Astrophysical Reaction Rates From Statistical Model Calculations
## 1 Introduction
Nuclear reaction rates are an essential ingredient for all investigations of nucleosynthetic or energy generating processes in astrophysics. Highly unstable nuclei are produced in such processes which again can become targets for subsequent reactions. Cross sections and astrophysical reaction rates for a large number of nuclei are required to perform complete network calculations which take into account all possible reaction links and do not postulate a priori simplifications. Despite concerted experimental efforts, most of the involved nuclei are currently not accessible in the laboratory and therefore theoretical models have to be invoked in order to predict reaction rates.
In predictions of cross sections and reaction rates for astrophysical applications, slightly different points are emphasized than in pure nuclear physics investigations. Firstly, one is confined to the very low-energy region, from thermal energies up to a few MeV. Secondly, since most of the ingredients for the calculations are experimentally undetermined, one has to develop reliable phenomenological models to predict these properties with an acceptable accuracy across the nuclear chart. This task is made even harder by the lack of information on specific properties, such as optical potentials for $`\alpha `$ particles at the astrophysically relevant energies even for stable nuclei. Therefore, one has to be satisfied with a somewhat more limited accuracy, as compared to usual nuclear physics standards. The accuracy of the rates is estimated to be within a factor of $`1.52`$, with an even better average deviation, e.g. of 1.4 for neutron capture. Considering the substantially larger uncertainties in many astrophysical scenarios, this seems to be acceptable. Thus, the real challenge is not the application of well-established models, but rather to provide all the necessary ingredients in as reliable a way as possible, also for nuclei where no such information is available. Many efforts have been directed at addressing those problems and the current status of the investigations make it worthwhile to publish a full set of theoretical rates, intended to supersede early reaction rate tabulations .
For the majority of nuclear reactions in astrophysics, the statistical model (Wolfenstein-Hauser-Feshbach approach) can be applied. This is appropriate provided the level density in the contributing energy window around the peak of the projectile energy distribution is sufficiently high to justify a statistical treatment. The critical level density is usually estimated between 5 and 10 MeV<sup>-1</sup> . Furthermore, the compound nucleus picture will only dominate when the energy of the incident particle is low enough ($`<20`$ MeV). While the latter point is practically always satisfied in astrophysical environments, the level density may fall below the critical value in certain nuclei lighter than Fe, at shell closures, and for very neutron-rich or proton-rich isotopes near the drip lines with correspondingly low separation energies. In these cases, single resonances or direct capture contributions will become significant and have to be treated individually. In this tabulation, we are not concerned with such effects but rather give a full set of rates calculated in the statistical model. However, the limits of its applicability will be discussed (Sec. 2.4).
In the following section, we concisely summarize the theoretical background for easy reference as well as the nuclear properties used as input in the calculations. This is followed by a section defining the reaction rates, explaining the fitting procedure and giving more details on the tabulated values. The paper is concluded by a summary and the rate tables.
## 2 The Statistical Model
### 2.1 Theory
The averaged transmission coefficients $`T`$ comprise the central quantities in statistical model calculations. They do not reflect a resonance behavior, but rather describe absorption via an imaginary part in the (optical) nucleon-nucleus potential . This leads to the well-known expression
$`\sigma ^{\mu \nu }(E_{ij})`$ $`=`$ $`{\displaystyle \frac{\pi \mathrm{}^2/(2\mu _{ij}E_{ij})}{(2J_i^\mu +1)(2J_j+1)}}`$
$`\times {\displaystyle \underset{J,\pi }{}}(2J+1){\displaystyle \frac{T_j^\mu (E,J,\pi ,E_i^\mu ,J_i^\mu ,\pi _i^\mu )T_o^\nu (E,J,\pi ,E_m^\nu ,J_m^\nu ,\pi _m^\nu )}{T_{\mathrm{tot}}(E,J,\pi )}}`$
for the cross section $`\sigma ^{\mu \nu }`$ of the reaction $`i^\mu (j,o)m^\nu `$ from the target state $`i^\mu `$ to the excited state $`m^\nu `$ of the final nucleus, with a center of mass energy $`E`$<sub>ij</sub> and reduced mass $`\mu _{ij}`$. $`J`$ denotes the spin, $`E`$ the corresponding excitation energy, and $`\pi `$ the parity of excited states. When these properties are used without subscripts they describe the compound nucleus; subscripts refer to states of the participating nuclei in the reaction $`i^\mu (j,o)m^\nu `$ and superscripts indicate the specific excited states. The total transmission coefficient $`T_{\mathrm{tot}}=_{\nu ,o}T_o^\nu `$ describes the transmission into all possible bound and unbound states $`\nu `$ in all energetically accessible exit channels $`o`$ (including the entrance channel $`i`$). Experiments measure $`\sigma ^{\mathrm{lab}}=_\nu \sigma ^{0\nu }(E_{ij})`$, summed over all excited states of the final nucleus, with the target in the ground state. Target states $`\mu `$ in an astrophysical plasma of temperature $`T^{}`$ are thermally populated and the astrophysical cross section $`\sigma ^{}`$ is given by
$$\sigma ^{}(E_{ij})=\frac{\underset{\mu }{}(2J_i^\mu +1)\mathrm{exp}(E_i^\mu /kT^{})\underset{\nu }{}\sigma ^{\mu \nu }(E_{ij})}{_\mu (2J_i^\mu +1)\mathrm{exp}(E_i^\mu /kT^{})},$$
(2)
$`k`$ being the Boltzmann constant. The summation over $`\nu `$ replaces $`T_o^\nu (E,J,\pi )`$ in Eq. (2.1) by the total transmission coefficient
$`T_o(E,J,\pi )`$ $`=`$ $`{\displaystyle \underset{\nu =0}{\overset{\nu _m}{}}}T_o^\nu (E,J,\pi ,E_m^\nu ,J_m^\nu ,\pi _m^\nu )`$
$`+{\displaystyle _{E_m^{\nu _m}}^{ES_{m,o}}}{\displaystyle \underset{J_m,\pi _m}{}}T_o(E,J,\pi ,E_m,J_m,\pi _m)\rho (E_m,J_m,\pi _m)dE_m.`$
Here $`S_{m,o}`$ is the channel separation energy, and the summation over excited states above the highest experimentally known state $`\nu _m`$ is changed to an integration over the level density $`\rho `$. The summation over target states $`\mu `$ in Eq. (2) has to be generalized accordingly.
The important ingredients of statistical model calculations as indicated in Eqs. (2.1) through (2.1) are the particle and $`\gamma `$-transmission coefficients $`T`$ and the level density of excited states $`\rho `$. Therefore, the reliability of such calculations is determined by the accuracy with which these components can be evaluated (often for unstable nuclei). It is in these quantities that various statistical model calculations differ. The reaction rates given in this paper are calculated with the code NON-SMOKER , derived from the well-known SMOKER code . (The code MOST is another code derived from SMOKER.) In the following we want to briefly outline the methods utilized in that code to estimate those nuclear properties. The challenge is in the goal to provide them in as reliable a way as possible, also for unstable nuclei for which no experimental information is available. Thus, global descriptions are employed which minimize the overall error and are trusted to be reliable also far from stability.
### 2.2 Transmission Coefficients
The transition from an excited state in the compound nucleus $`(E,J,\pi )`$ to the state $`(E_i^\mu ,J_i^\mu ,\pi _i^\mu )`$ in nucleus $`i`$ via the emission of a particle $`j`$ is given by a summation over all quantum mechanically allowed partial waves
$$T_j^\mu (E,J,\pi ,E_i^\mu ,J_i^\mu ,\pi _i^\mu )=\underset{l=|Js|}{\overset{J+s}{}}\underset{s=|J_i^\mu J_j|}{\overset{J_i^\mu +J_j}{}}T_{j_{ls}}(E_{ij}^\mu ).$$
(4)
Here the angular momentum $`\stackrel{}{l}`$ and the channel spin $`\stackrel{}{s}=\stackrel{}{J}_j+\stackrel{}{J}_i^\mu `$ couple to $`\stackrel{}{J}=\stackrel{}{l}+\stackrel{}{s}`$. The transition energy in channel $`j`$ is $`E_{ij}^\mu `$=$`ES_jE_i^\mu `$, where $`S_j`$ is the channel separation energy.
The total transmission coefficients for this tabulation are then calculated by applying Eq. (2.1) and utilizing up to 19 experimentally known excited states. The data are taken from , up to the first level for which the spin assignment was not known. Ground state spin and parities are known for many unstable nuclei. Far off stability, ground state spins and parities are taken from , if experimental values are not available.
#### 2.2.1 Particle Transmission Coefficients
The individual particle transmission coefficients $`T_{j_{ls}}`$ are calculated by solving the Schrödinger equation with an optical potential for the particle-nucleus interaction. We employ the optical potential for neutrons and protons given by , based on microscopic infinite nuclear matter calculations for a given density, applied with a local density approximation. It includes corrections of the imaginary part .
The optical potential for $`\alpha `$ particles from was shown to be quite accurate for a wide range of nuclei and is used in this work. However, it was realized that for heavily charged nuclei a more sophisticated potential had to be adopted at the comparatively low energies of astrophysical interest. Promising is the folding approach , with a parameterized mass- and energy-dependence of the real volume integral . Microscopic and deformation information should be considered in the parametrization of the imaginary potential . However, due to the scarcity of experimental data, the potential parameters can as yet only be extracted for a limited mass and energy range. The optical $`\alpha `$+nucleus potential is likely to introduce the largest uncertainties in the charged particle rates presented here. Further experimental work is clearly necessary and most welcome.
Deformed nuclei are treated by using an effective spherical potential of equal volume, based on averaging the deformed potential over all possible angles between the incoming particle and the orientation of the deformed nucleus.
#### 2.2.2 Radiative transmission coefficients
At least the dominant $`\gamma `$-transitions (E1 and M1) have to be included in the calculation of the total photon width. The smaller, and therefore less important, M1 transitions are treated, as usual, in the simple single particle approach ($`TE^3`$ ), as also discussed in . The E1 transitions are calculated on the basis of the Lorentzian representation of the Giant Dipole Resonance (GDR). Within this model, the E1 transmission coefficient for the transition emitting a photon of energy $`E_\gamma `$ in a compound nucleus $`{}_{N}{}^{A}Z`$ is given by
$$T_{E1}(E_\gamma )=\frac{8}{3}\frac{NZ}{A}\frac{e^2}{\mathrm{}c}\frac{1+\chi }{Mc^2}\underset{i=1}{\overset{2}{}}\frac{i}{3}\frac{\mathrm{\Gamma }_{G,i}E_\gamma ^4}{(E_\gamma ^2E_{G,i}^2)^2+\mathrm{\Gamma }_{G,i}^2E_\gamma ^2}.$$
(5)
Here, $`M`$ is the proton mass, $`\chi (=0.2)`$ accounts for the neutron-proton exchange contribution , and the summation over $`i`$ includes two terms which correspond to the split of the GDR in statically deformed nuclei, with oscillations along ($`i=1`$) and perpendicular ($`i=2`$) to the axis of rotational symmetry. Many microscopic and macroscopic models have been devoted to the calculation of the GDR energies ($`E_G`$) and widths ($`\mathrm{\Gamma }_G`$). Here, the (hydrodynamic) droplet model approach is used for $`E_G`$, which gives an excellent fit to the GDR energies and can also predict the split of the resonance for deformed nuclei, when making use of the deformation, calculated within the droplet model. In that case, the two resonance energies are related to the mean value calculated by the relations $`E_{G,1}+2E_{G,2}=3E_G`$, $`E_{G,2}/E_{G,1}=0.911\eta +0.089`$. $`\eta `$ is the ratio of the diameter along the nuclear symmetry axis to the diameter perpendicular to it, and is obtained from the experimentally known deformation or mass model predictions. For the width $`\mathrm{\Gamma }_G`$ of the GDR the description of is used, which applies to spherical and deformed nuclei and can be described as a superposition of a macroscopic width due to the viscosity of the nuclear fluid and a coupling to quadrupole surface vibrations of the nucleus (see also ).
Direct application of Eq. (5) would overestimate the radiation width by about 30% (see e.g. ). This is due to the fact that, for low energy $`\gamma `$-transitions, the Lorentz curve is suppressed and the GDR width increases with excitation energy (e.g. ). To account for these deficiencies, various treatments of an energy-dependent width have been suggested. We use the form
$$\mathrm{\Gamma }_G\left(E_\gamma \right)=\mathrm{\Gamma }_G\sqrt{\frac{E_\gamma }{E_G}.}$$
(6)
Another effect has to be taken into account for certain $`\alpha `$-capture reactions. Because of isospin selection rules, $`\gamma `$-transitions between levels with isospin $`I=0`$ are forbidden. This leads to a suppression of the cross section for ($`\alpha `$,$`\gamma `$) reactions on self-conjugate ($`N=Z`$) targets, due to isospin conservation. A suppression could also be found for capture reactions leading into self-conjugate nuclei, although somewhat less pronounced because $`I=1`$ states can be populated according to the isospin coupling coefficients. The suppression is usually treated as a suppression of the $`\gamma `$-width. In previous rate tabulations it was either neglected or accounted for in a phenomenological way by dividing the $`\gamma `$-width by quite arbitrary factors . In the code NON-SMOKER the appropriate $`\gamma `$-widths are obtained by explicitly accounting for isospin mixing and suppression of the appropriate $`\gamma `$-transitions . A detailed account of the procedure can be found in .
#### 2.2.3 Width fluctuation corrections
In addition to the ingredients required for Eq. (2.1), like the transmission coefficients for particles and photons, width fluctuation corrections $`W(j,o,J,\pi )`$ have to be employed. They define the correlation factors with which all partial channels for an incoming particle $`j`$ and outgoing particle $`o`$, passing through the excited state $`(E,J,\pi )`$, have to be multiplied. This takes into account that the decay of the state is not fully statistical, but some memory of the way of formation is retained and influences the available decay choices. The major effect is elastic scattering, the incoming particle can be immediately re-emitted before the nucleus equilibrates. Once the particle is absorbed and not re-emitted in the very first (pre-compound) step, the equilibration is very likely. This corresponds to enhancing the elastic channel by a factor $`W_j`$. In order to conserve the total cross section, the individual transmission coefficients in the outgoing channels have to be renormalized to $`T_j^{}`$. The total cross section is proportional to $`T_j`$ and, when summing over the elastic channel ($`W_jT_j^{}`$) and all outgoing channels ($`T_{tot}^{}T_j^{}`$), one obtains the condition $`T_j`$=$`T_j^{}(W_jT_j^{}/T_{tot}^{})+T_j^{}(T_{tot}^{}T_j^{})/T_{tot}^{}`$ . We can (almost) solve for $`T_j^{}`$
$$T_j^{}=\frac{T_j}{1+T_j^{}(W_j1)/T_{tot}^{}}.$$
(7)
This requires an iterative solution for $`T^{}`$ (starting in the first iteration with $`T_j`$ and $`T_{tot}`$), which converges rapidly. The enhancement factor $`W_j`$ has to be known in order to apply Eq. (7). A general expression in closed form was derived , but is computationally expensive to use. A fit to results from Monte Carlo calculations gave
$$W_j=1+\frac{2}{1+T_j^{1/2}}.$$
(8)
For a general discussion of approximation methods see . Eqs. (7) and (8) redefine the transmission coefficients of Eq. (2.1) in such a manner that the total width is redistributed by enhancing the elastic channel and weak channels over the dominant one. Cross sections near threshold energies of new channel openings, where very different channel strengths exist, can only be described correctly when taking width fluctuation corrections into account. The width fluctuation corrections of are only an approximation to the correct treatment. However, it was shown that they are quite adequate .
### 2.3 Level Densities
Until recently, the nuclear level density has given rise to the largest uncertainties in the description of nuclear reactions in the statistical model . Implemented in the NON-SMOKER code is a recently improved treatment . It is based on a shifted Fermi-gas formalism with an energy-dependent level density parameter $`a`$ together with microscopic corrections from nuclear mass models. This leads to improved fits to known level densities in a wide range of masses . More sophisticated Monte Carlo shell model calculations , as well as combinatorial approaches (see e.g. ), have shown excellent agreement with this phenomenological approach and justified the application of the Fermi-gas description at and above the neutron separation energy.
An in-depth description of the model and its application to astrophysical problems can be found in . Here, we only want to briefly summarize the inputs used for calculating the rates presented in this tabulation. It should be noted that we applied our description throughout the nuclear chart, without relying on experimental level density parameters in specific cases as has been done before . This may lead to locally slightly larger deviations from experiment but it improves the reliability when extrapolating to unknown isotopes.
The microscopic correction and the pairing corrections comprise crucial inputs for the level density formalism used here (see for details). They can be extracted from mass models. There is a choice of several mass models in the NON-SMOKER code. The Finite Range Droplet Model (FRDM) and an extended Thomas-Fermi approach with Strutinski Integral (ETFSI-Q) have been chosen for the reaction rate calculations in this work. It has to be emphasized that experimental mass values were included where possible. This is straightforward for the separation energies which were calculated from the mass differences; it was ensured that either only experimental or theoretical values were used in the differences, thus avoiding unphysical breaks at transition points from experiment to theory. The microscopic corrections were obtained by subtracting the well-defined spherical macroscopic (droplet) term of the FRDM from the total mass energy derived from experiment, from the FRDM or from ETFSI-Q, respectively (cf. Eq. (17) in ). The validity of the resulting rates is discussed in Sec. 3.3. Rates based on other mass models can be obtained from the authors on request or on-line (see Sec. 4).
The shifted Fermi-gas approach diverges for $`U=E\delta =0`$ (i.e. $`E=\delta `$, if $`\delta `$ is a positive backshift, with $`E`$ being the excitation energy and $`\delta `$ being an energy shift due to pairing corrections). In order to obtain the correct behavior at very low excitation energies, the Fermi-gas description can be combined with the constant temperature formula (; and references therein)
$$\rho (U)\frac{\mathrm{exp}(U/T_{\mathrm{nucl}})}{T_{\mathrm{nucl}}}.$$
(9)
The two formulations are matched by a tangential fit determining the nuclear temperature $`T_{\mathrm{nucl}}`$.
### 2.4 Applicability of the Statistical Model
The statistical model can be applied provided that the use of averaged transmission coefficients (Eq. (4)) is permitted. This will be the case for high level densities with completely overlapping resonances, typical for the compound nucleus reaction mechanism. For light nuclei, decreasing particle separation energies or at shell closures, level densities will eventually become too low for the application of the statistical model at astrophysical temperatures. In those cases, single resonances and contributions from the direct reaction mechanism have to be taken into account . Based on the level density description outlined in Sec. 2.3, a quantitative criterion for the applicability was derived recently . In the present work we give tables of all reaction rates regardless of applicability but specify the allowed temperature range in the tables. The estimate is quite conservative and thus the rates can still be accurate slightly below the given lower limits of the temperature.
## 3 Astrophysical Reaction Rates
### 3.1 Definition
The nuclear reaction rate per particle pair at a given stellar temperature $`T^{}`$ is determined by folding the reaction cross section $`\sigma `$$`{}_{}{}^{}(E)`$ from Eq. (2) with the Maxwell-Boltzmann velocity distribution of the projectiles
$$\sigma ^{}v=\sigma v^{}=\left(\frac{8}{\pi \mu }\right)^{1/2}\frac{1}{\left(kT^{}\right)^{3/2}}\underset{0}{\overset{\mathrm{}}{}}\sigma ^{}(E)E\mathrm{exp}\left(\frac{E}{kT^{}}\right)𝑑E.$$
(10)
It has to be emphasized that only the use of the stellar cross section $`\sigma `$ (Eq. (2)) yields a reaction rate with the desired behavior that the inverse reaction can be calculated by using detailed balance. Therefore, laboratory rates – which only measure $`\sigma `$$`{}_{}{}^{\mathrm{lab}}=_\nu \sigma ^{0\nu }`$, i.e. the cross section with the target being in the ground state – should always be measured in the direction that is least affected by excited target states. This is usually the reaction with positive $`Q`$-value (exoergic reaction). For astrophysical applications, such rates have to be corrected for the stellar enhancement effect due to the thermal excitation of the target . The stellar enhancement factors (SEF) $`f^{}`$ are defined by
$$f^{}=\frac{\sigma ^{}}{\sigma ^{\mathrm{lab}}}.$$
(11)
The values of $`f^{}`$for a range of temperatures for nuclei close to stability are given in the tables. Stellar enhancement factors for neutron capture reactions and a discussion of the involved uncertainties can also be found in a recent compilation of neutron cross sections for the $`s`$ process .
### 3.2 Partition functions and reverse rates
The temperature-dependent partition function $`G(T^{})`$ normalized to the ground state target spin $`J_i^0`$ is defined as
$`(2J_i^0+1)G(T^{})=`$ $`{\displaystyle \underset{\mu =0}{\overset{\mu _m}{}}}(2J_i^\mu +1)e^{E_i^\mu /kT^{}}`$
$`+{\displaystyle \underset{E_i^{\mu _m}}{\overset{E_i^{\mathrm{max}}}{}}}{\displaystyle \underset{J^\mu ,\pi ^\mu }{}}(2J^\mu +1)e^{ϵ/kT^{}}\rho (ϵ,J^\mu ,\pi ^\mu )dϵ,`$
with $`\rho `$ being the level density and $`\mu `$<sub>m</sub> the last included experimentally known state. The included experimental levels were the same as for the calculation of the transmission coefficients (Eq. (2.1)). For the temperature range considered here, the maximum energy $`E_i^{\mathrm{max}}`$ above which there are no more significant contributions to the partition function is of the order of $`2030`$ MeV. With that definition, the stellar reaction rate $`\sigma _iv^{}`$ for a reaction with particles in all channels is related to the rate of the reverse reaction $`\sigma _mv^{}`$ by
$$N_A\sigma _mv^{}=\left(\frac{A_iA_j}{A_oA_m}\right)^{3/2}\frac{(2J_i+1)(2J_j+1)}{(2J_o+1)(2J_m+1)}\frac{G_i(T^{})}{G_m(T^{})}e^{Q/kT^{}}N_A\sigma _iv^{},$$
(13)
where $`N_A`$ is Avogadro’s number, $`J`$ and $`A`$ are spins and masses $`A`$ (in atomic mass units $`u`$) of the particles involved in the reaction $`i`$($`j`$,$`o`$)$`m`$, $`Q`$ is the reaction $`Q`$-value. To calculate photodisintegration rates from capture rates the appropriate relation is
$$\lambda _\gamma =\left(\frac{A_iA_j}{A_m}\right)^{3/2}\frac{(2J_i+1)(2J_j+1)}{(2J_m+1)}\frac{G_i(T^{})}{G_m(T^{})}\left(T^{}\right)^{3/2}Fe^{Q/kT^{}}N_A\sigma _iv^{}.$$
(14)
For $`\lambda _\gamma `$ in s<sup>-1</sup> and using the usual practical units, i.e. temperatures $`T_9=T^{}/10^9`$ K and $`N_A\sigma v^{}`$ in cm<sup>3</sup> s<sup>-1</sup> mole<sup>-1</sup>, one obtains
$$\left(T^{}\right)^{3/2}F=\left(\frac{ukT^{}}{2\pi \mathrm{}^2}\right)^{3/2}\frac{1}{N_A}=T_9^{3/2}9.8685\times 10^9\mathrm{mole}\mathrm{cm}^3.$$
(15)
The numerical factor $`F`$ as well as the spin and mass factors are already accounted for in the parameter $`a_0^{\mathrm{rev}}`$ which is tabulated. See Sec. 3.3.2 for further details.
In the tabulated rates, the thermal effects are already considered. Statistical model calculations not including full photon cascades may also be prone to some error arising from the decay of unbound particle states in reactions with negative $`Q`$-values (endoergic reactions). Furthermore, due to the exponential dependence of the inverse rate on the $`Q`$-value (Eqs. (13) and (14)), inaccuracies in the rate would be strongly enhanced when computing the exoergic rate from the endoergic one. In order to minimize the error, reactions are always calculated in the exoergic direction (with the exceptions of capture reactions, where photodisintegration is always treated as the reverse reaction regardless of $`Q`$-value) and detailed balance is applied to obtain the values for the endoergic reaction. This treatment has the additional advantage that it ensures consistent values for forward and reverse reactions, which is essential for application in astrophysical nuclear reaction networks. To calculate the actual (endoergic) reaction rate, those fits have to be multiplied by the ratio of the partition functions of the final nucleus and the target $`G_i/G_m`$ at the appropriate temperature (see Sec. 3.3.2 and Sec. 3.3.6). For that purpose, partition functions are tabulated separately.
### 3.3 Analytic reaction rate fits
#### 3.3.1 Parametrization
Reaction rates have been calculated for a temperature grid of 24 temperatures: $`T_9`$=0.1, 0.15, 0.2, 0.3, 0.4, 0.5, 0.6, 0.7, 0.8, 0.9, 1.0, 1.5, 2.0, 2.5, 3.0, 3.5, 4.0, 4.5, 5.0, 6.0, 7.0, 8.0, 9.0, 10.0. For easy application in astrophysical investigations, all reaction types ((n,$`\gamma `$), (n,p), (n,$`\alpha `$), (p,$`\gamma `$), (p,n), (p,$`\alpha `$), ($`\alpha `$,$`\gamma `$), ($`\alpha `$,n), ($`\alpha `$,p), ($`\gamma `$,n), ($`\gamma `$,p), ($`\gamma `$,$`\alpha `$)) are fitted with the same parametrization
$`\begin{array}{c}N_A\sigma v^{}\\ \lambda _\gamma \end{array}\}`$ $`=`$ $`\mathrm{exp}(a_0+a_1T_9^1+a_2T_9^{1/3}+a_3T_9^{1/3}+a_4T_9`$ (19)
$`+a_5T_9^{5/3}+a_6\mathrm{ln}T_9),`$
with the seven open parameters $`a_0a_6`$ and the stellar temperature $`T_9`$ given in 10<sup>9</sup> K. This parametrization proves to be flexible enough to accommodate the different temperature dependencies of the various reaction types across the fitted temperature range of $`0.01T_910.`$ Parametrizations of the present rates in the form used in and others can be obtained from the authors on request.
#### 3.3.2 Parameters for the reverse rates
The parameters for the reverse rates are not given explicitly but can easily be computed from the information in the tables. To calculate the reverse rate of the reaction $`i(j,o)m`$, i.e. the reaction $`m(o,j)i`$, Eq. (19) is employed and the seven parameters $`a_0^{\mathrm{rev}}a_6^{\mathrm{rev}}`$ for the reverse rate are determined as follows:
$`a_0^{\mathrm{rev}}`$ $`=`$ $`a_0^{\mathrm{rev}},\mathrm{as}\mathrm{tabulated}`$
$`a_1^{\mathrm{rev}}`$ $`=`$ $`a_111.6045Q`$
$`a_2^{\mathrm{rev}}`$ $`=`$ $`a_2`$
$`a_3^{\mathrm{rev}}`$ $`=`$ $`a_3`$
$`a_4^{\mathrm{rev}}`$ $`=`$ $`a_4`$ (20)
$`a_5^{\mathrm{rev}}`$ $`=`$ $`a_5`$
$`a_6^{\mathrm{rev}}`$ $`=`$ $`\{\begin{array}{cc}a_6+1.5\hfill & \hfill (\mathrm{i})\\ a_6\hfill & \hfill (\mathrm{ii})\end{array}`$ (23)
The above relations are derived from Eqs. (13) and (14), using Eq. (19) and taking the logarithms on both sides. For the coefficient $`a_6^{\mathrm{rev}}`$, case (i) applies when calculating a photodisintegration rate from a capture rate, case (ii) for all other rates. Finally, for the reverse reaction case, the value found by application of Eqs. (19) and (20) has to be multiplied by the ratio of the partition functions for residual and target nucleus $`G_i/G_m`$. Examples are shown in Sec. 3.3.6.
#### 3.3.3 Fit accuracy
The flexibility of the fitting function makes it prone to numerical problems outside the calculated range at low temperatures, where the rates should be close to zero. In some cases they tend to diverge strongly. This difficulty can be avoided by providing fit data at low temperatures additionally to the calculated values by appropriately extrapolating the rates to lower temperatures. This is achieved by either assuming s-wave capture for $`T_9<0.1`$ for exoergic neutron capture reactions (Maxwellian averaged capture cross sections in the energy range $`5E100`$ keV for targets along the line of stability can be found in another compilation ) or by considering proper Coulomb barrier penetration factors in the charged particle channels. Thus, both accuracy and flexibility can be ensured within a single parametrization. However, it has to be emphasized that *the given parameterization is only valid within the temperature range of* $`0.01T_910.`$, although many fits will show “proper” behavior down to lower temperature. Caution is advised when using derived ($`\gamma `$,p) and ($`\gamma `$,n) rates at the proton dripline (see below). For all cases, it is recommended to use the fits only down to the temperature $`T_{\mathrm{low}}^{\mathrm{fit}}`$ given in the table. The temperatures of the validity of the fits are given in the tables for each reaction, to emphasize the importance of the given fit ranges.
As a measure of the accuracy of a given fit, the quantity $`\zeta `$ is quoted in the tables. It is defined by
$$\zeta =\frac{1}{24}\underset{i=1}{\overset{24}{}}\left(\frac{r_if_i}{f_i}\right)^2,$$
(24)
with $`r`$ being the original rate value as calculated at each of the 24 temperatures $`T_9`$= 0.1, 0.15 …10.0, and $`f`$ the rate calculated from the fit at these temperatures. Contributions with $`r<10^{20}`$ cm<sup>3</sup> s<sup>-1</sup> mole<sup>-1</sup> are neglected as lower accuracy at at low rates is inconsequential. Note that while a small value of $`\zeta `$ is indicative of an accurate fit over the entire temperature range, large $`\zeta `$ generally signify deviations of the calculated from the fitted rate at the lowest temperatures only.
The fit parameters are tabulated regardless of the validity of the statistical model of nuclear reactions in the given temperature range (see Sec. 2.4). The estimated lower temperature limit of the validity of the statistical model, $`T_{\mathrm{low}}^{\mathrm{HF}}`$ is given separately for each rate in the tables. Below that limit the calculation of the rate by means of the statistical model may not be justified, although the fit to the calculated rate will still be accurate. At temperatures below the applicability limit, rates may be over-estimated and should be compared to calculations considering single resonance and direct reaction contributions. Especially close to the driplines, fits of reactions with low $`Q`$-value cannot be applied at low temperatures. Although the fit may be valid, it should not be used at low temperature because the statistical model will not be applicable anymore.
#### 3.3.4 Computed rate sets
Two different sets of rates have been calculated. They differ in the mass model used, which enters into the computation of the separation energies and $`Q`$-values as well as into the microscopic input to the level density calculation (see Sec. 2.3). One set was calculated employing the well-known FRDM mass model , which excellently reproduces masses and other ground state properties of nuclei close to stability. It is also the most comprehensive set across the nuclear chart.
Recently, it was suggested that so-called shell quenching effects may arise for neutron-rich nuclei far off stability . Fully microscopic calculations and experimental data indicate the weakening of nuclear shell gaps for neutron-rich nuclei. In the absence of microscopic calculations of the required nuclear properties for the whole nuclear chart, it is possible to phenomenologically include such quenching into existing mass formulae. This has been done in the ETFSI-Q model , based on the ETFSI-1 mass model . However, it does not cover the full range of isotopes. Therefore, we provide the alternative sets of reaction rates obtained with the two mass models, so that the rates based on the FRDM can be used for large-scale studies and close to stability, and the rates based on ETFSI-Q for investigations concerning neutron-rich unstable isotopes and the $`r`$ process. It has to be noted that one should refrain from mixing rates from the two sets as this will lead to inconsistencies and artificial effects in the results.
#### 3.3.5 Mass ranges of the tabulated fits
Due to the extensive number of nuclear reactions in the considered mass range, we have to limit the printed version of our reaction rate fits. Full rate libraries both for reactions calculated with the FRDM and with the ETFSI-Q (as well as ETFSI-1) mass model can be obtained on-line or from the authors on request (see Sec. 4). In the printed version, only the FRDM set is given and no capture rates for reactions with negative $`Q`$ value are shown.
The full electronic versions of the tables available on-line include all reactions in the range $`10Z83`$ (FRDM) and $`24Z83`$ (ETFSI-Q). This amounts to 5369 (FRDM) and 4628 (ETFSI-Q) involved nuclei. The isotope ranges for which rate fits are available are given in Table A; for the FRDM, the mass range is also indicated by the heavy lines in Tables IA-IC. Rate fits are given for all n-, p-, and $`\alpha `$-capture reactions and for those (n,p), (n,$`\alpha `$), (p,n), (p,$`\alpha `$), ($`\alpha `$,n), and ($`\alpha `$,p) reactions having positive $`Q`$-value. Reverse rates are not given explicitly but can be computed by a two-step procedure as described in Secs. 3.3.2 and 3.3.6. The stellar enhancement factors close to stability as well as the partition functions for all isotopes are given for a temperature grid of 24 temperatures: $`T_9`$ = 0.1, 0.15, 0.2, 0.3, 0.4, 0.5, 0.6, 0.7, 0.8, 0.9, 1.0, 1.5, 2.0, 2.5, 3.0, 3.5, 4.0, 4.5, 5.0, 6.0, 7.0, 8.0, 9.0, 10.0 .
The printed Tables II$``$IV contain the following (calculated with the FRDM mass model):
* (n,$`\gamma `$) rate fits from stability to close to the neutron dripline in the range $`10Z49`$, and to $`N=Z+59`$ for $`50Z83`$ (Table II).
* (n,p) and (n,$`\alpha `$) rate fits around stability in the range $`10Z83`$ (Table II).
* (p,$`\gamma `$), (p,n), (p,$`\alpha `$), ($`\alpha `$,$`\gamma `$), ($`\alpha `$,n), and ($`\alpha `$,p) rate fits from proton-rich nuclei to stability in the range $`10Z50`$ (Tables III and IV).
* Reverse (endoergic) rates are not given explicitly but can be computed with the help of partition functions from the information given in the tables (see Secs. 3.3.2, 3.3.6).
* Stellar enhancement factors $`f^{}`$ at selected 16 temperatures are only quoted close to the valley of stability (Tables II$``$IV).
* Partition functions for all involved isotopes are given at selected 20 temperatures in the range $`0.1T_910`$ (Table V).
An overview of the provided rates is given in Tables IA$``$IC, which show in detail for which n-, p-, and $`\alpha `$-induced reactions rate fit parameters are available in Tables II$``$IV, respectively.
#### 3.3.6 Examples of use of tables
This section is intended to help with interpreting the information given in the tables. We give two examples for calculating the reaction rate for a given reaction and its inverse reaction at a temperature of $`T_9=2.0`$.
The first example is the capture reaction <sup>35</sup>Ar(p,$`\gamma `$)<sup>36</sup>K. From Table III one finds a $`Q`$-value of $`Q=1.666`$ MeV and the parameters $`a_0=`$128.39, $`a_1=4.0033`$, $`a_2=137.67`$, $`a_3=276.87`$, $`a_4=`$17.691, $`a_5=1.0728`$, $`a_6=`$123.68. With the help of Eq. (19) one calculates $`N_A\sigma v^{}=92.5`$ cm<sup>3</sup>s<sup>-1</sup>mole<sup>-1</sup> at $`T_9=2.0`$. Because both $`T_{\mathrm{low}}^{\mathrm{HF}}`$ and $`T_{\mathrm{low}}^{\mathrm{fit}}`$ are considerably smaller than our temperature $`T_9`$, it is safe to assume that the statistical model is applicable and the fit to the rate is valid. In order to obtain the value for the reverse rate, one first has to determine the parameter values in the given parametrization. The parameter $`a_0^{\mathrm{rev}}=151.83`$ is given in the table. The remaining parameters are derived according to Eq. (20) for a photodisintegration rate. This yields $`a_1^{\mathrm{rev}}=23.3364`$, $`a_6^{\mathrm{rev}}=125.18`$; all other parameters assume the same value as for the forward reaction. Using those, Eq. (19) gives a value of $`\lambda _\gamma `$$`{}_{}{}^{}=`$2.5$`\times `$10<sup>8</sup> s<sup>-1</sup>. This has to be multiplied by the ratio of the partition functions in order to obtain the valid rate factor for <sup>36</sup>K($`\gamma `$,p)<sup>35</sup>Ar:
$$\lambda _\gamma =\lambda _\gamma ^{}\frac{G_{{}_{}{}^{35}\mathrm{Ar}}}{G_{{}_{}{}^{36}\mathrm{K}}}=\lambda _\gamma ^{}\frac{1.001}{1.203}=2.1\times 10^8\mathrm{s}^1.$$
The values of the partition functions at $`T_9=2.0`$ were taken from Table V. Note that for capture rates the procedure is always the same as described above regardless of whether it is an exoergic or an endoergic reaction.
The second example we consider is the reaction <sup>34</sup>S($`\alpha `$,n)<sup>37</sup>Ar, again at $`T_9=2.0`$. It is not to be found in Table IV because of its negative $`Q`$-value. Therefore, one has to rely on Table II to calculate this reaction as the inverse reaction of <sup>37</sup>Ar(n,$`\alpha `$)<sup>34</sup>S. The parameters found in Table II are $`a_0=`$20.072, $`a_1=0.019613`$, $`a_2=1.8224`$, $`a_3=4.759`$, $`a_4=`$0.56437, $`a_5=0.033893`$, $`a_6=`$1.7801, and $`Q=4.63`$ MeV. Again, both $`T_{\mathrm{low}}^{\mathrm{HF}}`$ and $`T_{\mathrm{low}}^{\mathrm{fit}}`$ are lower than the temperature of interest. The rate $`N_A\sigma v^{}=5.2\times 10^7`$ cm<sup>3</sup>s<sup>-1</sup>mole<sup>-1</sup> for <sup>37</sup>Ar(n,$`\alpha `$)<sup>34</sup>S results from the direct application of Eq. (19). In order to calculate the rate of <sup>34</sup>S($`\alpha `$,n)<sup>37</sup>Ar, the parameters are determined by application of Eq. (20). This yields the value $`a_1^{\mathrm{rev}}=53.748448`$. The value of $`a_0^{\mathrm{rev}}=20.199`$ is taken from the table and all other parameters remain the same as for the forward reaction. With Eq. (19) one arrives at the rate value $`r^{}=1.272\times 10^4`$ cm<sup>3</sup>s<sup>-1</sup>mole<sup>-1</sup> at $`T_9=2.0`$. Multiplying this by the appropriate ratio of partition functions taken from Table V yields the final result
$$N_A\sigma v^{}=r^{}\frac{G_{{}_{}{}^{37}\mathrm{Ar}}}{G_{{}_{}{}^{34}S}}=r^{}\frac{1.0}{1.0}=1.27\times 10^4\mathrm{cm}^3\mathrm{s}^1\mathrm{mole}^1.$$
## 4 Summary
Thermonuclear reaction rates for neutron-, proton- and $`\alpha `$-induced reactions and their inverses have been calculated in the statistical model. All rates from the proton dripline to the neutron dripline for $`10Z83`$ (Ne to Bi) have been fitted to a unique function with seven free parameters. Tables of these parameters are provided on-line for two sets of rates, calculated with input from two different mass models. Furthermore, the stellar enhancement factors are given in order to facilitate comparison with experimental ground state rates. A printed subset of the on-line tables for the FRDM presented here shows fit parameters for (n,$`\gamma `$) and (p,$`\gamma `$) reactions from close to their respective driplines to stability, and for other n-, p-, and $`\alpha `$-induced reactions with positive $`Q`$-values near stability. A prescription on deriving rates for inverse reactions with negative $`Q`$-values is given, as is a listing of the necessary partition functions.
It should further be noted that only purely theoretical rates are given here which do not use any direct experimental information (except for nuclear masses and excited state information where available). The methods to predict nuclear properties needed in the statistical model calculations are chosen to be as reliable as possible in order to retain predictive power. This is a compromise which may lead to locally enhanced inaccuracies but it emphasizes the importance of reliable predictions of rates far off stability.
In real applications, these rates should be supplemented or replaced with experimental rates as they become available. Such a combination of theoretical and experimental rates is provided, e.g., in the REACLIB compilation. Latest information on the current version of REACLIB can be found on the WWW at *http://ie.lbl.gov/astro.html*. Further details on the NON-SMOKER code and the cross section and reaction rate calculations are presented at *http://quasar.physik.unibas.ch/~tommy/reaclib.html*. Rates including further mass models can also be obtained from the authors on request or directly at the latter URL.
## Acknowledgements
This work was supported in part by the Swiss National Science Foundation (grant 2000-053798.98) and the Austrian Academy of Sciences (APART). T. R. is a PROFIL fellow of the Swiss National Science Foundation (grant 2124-055832.98). We thank S.E. Woosley and R.D. Hoffman for discussions. We also want to thank the Consulting Editor Wilfried Scholz for helpful comments and help in the preparation of the final manuscript.
## 5 Explanation of Tables
### Table IA: Neutron-Induced Reaction Rates Available in Table II
This is an overview of which neutron-induced reaction rates are available in the printed and the online versions. The full lines delimit the range of rates in the electronic version as given in Table A for the FRDM. The entries at a single neutron and proton number specify the reactions on the given target nucleus listed in the printed Table II. Only reactions with positive $`Q`$-value are shown. In addition to the marked rates, their reverse rates (with negative $`Q`$-value) can be inferred from the information in Table II as explained in Sec. 3.3.2. The reactions are denoted as follows:
* (n,$`\gamma `$)
* (n,p)
* (n,$`\alpha `$)
The box at the lower left corner gives the location in the $`Z`$, $`N`$ plane of the final nucleus relative to the target nucleus for (n,$`\gamma `$), (n,p), and (n,$`\alpha `$), thereby specifying also the inverse reaction fits derivable from Table II.
### Table IB: Proton-Induced Reaction Rates Available in Table III
Same as Table IA but for proton-induced reactions. The marked reactions correspond to the entries in Table III. The reactions are denoted as follows:
* (p,$`\gamma `$)
* (p,n)
* (p,$`\alpha `$)
The box at the lower left corner gives the location in the $`Z`$, $`N`$ plane of the final nucleus relative to the target nucleus for (p,$`\gamma `$), (p,n), and (p,$`\alpha `$), thereby specifying also the inverse reaction fits derivable from Table III.
### Table IC: Alpha Particle-Induced Reaction Rates Available in Table IV
Same as Table IA but for $`\alpha `$-particle induced reactions. The marked reactions correspond to the entries in Table IV. The reactions are denoted as follows:
* ($`\alpha `$,$`\gamma `$)
* ($`\alpha `$,n)
* ($`\alpha `$,p)
The box at the lower left corner gives the location in the $`Z`$, $`N`$ plane of the final nucleus relative to the target nucleus for ($`\alpha `$,$`\gamma `$), ($`\alpha `$,n), and ($`\alpha `$,n), thereby specifying also the inverse reaction fits derivable from Table IV.
### Table II: Neutron-Induced Reaction Rates
Fits to stellar rates $`N_A\sigma v^{}`$ for (n,$`\gamma `$), (n,p), (n,$`\alpha `$) reactions, calculated including masses from the FRDM. The rates in cm<sup>3</sup> mole<sup>-1</sup> s<sup>-1</sup> are computed by the use of Eq. (19), with the temperature given in units of 10<sup>9</sup> K. The fits are valid in the temperature range $`T_{\mathrm{low}}^{\mathrm{fit}}<T_910`$, with $`T_{\mathrm{low}}^{\mathrm{fit}}`$ given in the table. It should be noted that while the fit may still be formally valid and accurate, the application of the statistical model may not be justified at low temperatures. An estimate for the applicability of the statistical model is given by $`T_{\mathrm{low}}^{\mathrm{HF}}`$. The following information is provided:
* Reaction target
* Reaction type and final nucleus
* Reaction $`Q`$-value
* Target ground state spin (same as in Table V)
* Final nucleus ground state spin (same as in Table V)
* Estimate of the lower temperature limit for the applicability of the (Hauser-Feshbach) statistical model; “n.c.” indicates that the limit was not calculated for the given reaction.
* Lower temperature limit for the fit; usually 0.01. Note that the fits give extrapolated rates below $`T_9=0.1`$, which may be less accurate, especially if they are very small.
* Fit accuracy $`\zeta `$ (Eq. (24))
* If the field is blank, the seven fit parameters below are followed by the stellar enhancement factors (SEF) $`f^{}`$ (Eq. (11)) at the 16 temperatures given in the head of the table.
A value of 1 indicates that all SEF are unity; no SEF are printed.
A value of 0 indicates that no SEF were calculated.
* Seven fit parameters for the forward rate
* First parameter for the reverse rate fit (see Sec. 3.3.2)
* Temperatures at which the SEF were calculated
### Table III: Proton-Induced Reaction Rates
Same as Table I for the reaction types (p,$`\gamma `$), (p,n), (p,$`\alpha `$).
### Table IV: Alpha Particle-Induced Reaction Rates
Same as Table I for the reaction types ($`\alpha `$,$`\gamma `$), ($`\alpha `$,n), ($`\alpha `$,p).
### Table V: Partition functions
Partition functions of isotopes for various temperatures calculated with a level density making use of FRDM input. Included are only those partition functions for nuclei involved in the reactions given in Tables II$``$IV.
* Isotope for which the partition functions are tabulated.
* Temperature (in 10<sup>9</sup> K) at which the partition functions have been calculated.
* A value of 1 indicates that all partition function are unity; no partition functions are then printed explicitly.
* Ground state spin of nucleus, either from experiment or from theory .
* Partition functions normalized to the ground state (Eq. (3.2)) for the 20 temperatures specified in the table header.
## The remaining tables can be found on the ADNDT server
or at http://quasar.physik.unibas.ch/~tommy/adndt.html.
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# 1 Introduction
## 1 Introduction
While in the context of local quantum field theories CPT has to be conserved, CPT violating effects may show up in the framework of quantum gravity. As an example, Hawking radiation of black holes can be understood as a pair creation process near the event horizon, with one particle falling into the black hole and the other one escaping. Since with the particle falling into the black hole some phase information of the quantum state is lost, the thermic final state is a mixed state rather than a pure one. As Hawking has pointed out , such an evolution of a pure state into a mixed state violates the laws of conventional quantum mechanics (QMV). If the space time possesses a foamy structure at the Planck scale, including the creation and annihilation of black holes with Planck radius and Planck lifetime, such effects also may influence microscopical processes in the vacuum . In the following Page showed that such processes violate also CPT and the possibility of experimental tests in the $`K_0\overline{K_0}`$ sector was discussed by Eberhard . Ellis, Hagellin, Nanopoulos and Srednicki independently developed an evolution equation formalism in the space of density matrices containing three CPT violating (EHNS) parameters $`\alpha ,\beta ,\gamma `$ which have a dimension of mass and which might be expected to be of order $`m_K^2/M_{Pl}10^{20}`$ GeV in the Kaon sector. Recently the topic has been reconsidered by Ellis, Mavromatos and Nanopoulos and Huet and Peskin . CPT violating processes in the neutrino sector have been discussed for the first time by Liu et al. and in the following in , where neutrino oscillations due to CPT violation has been discussed as a solution to the solar neutrino problem. Recently another paper explored the possibility of explaining the atmospheric neutrino anomaly with quantum foam effects and came to a negative conclusion. In this note we extend the discussion of quantum foam effects in the neutrino sector to the cases of neutrinoless double beta decay and oscillations of neutrinos from astrophysical sources, supernovae as well as active galactic nuclei. New, extremely stringent bounds are found improving constraints found in the literature by several orders of magnitude.
## 2 Density matrix formalism
For mixed states it is useful to work in the framework of the density matrix formalism, following the methodology as presented in ref. . We start with the Schrödinger equation for the density matrix,
$$i\frac{d}{dt}\rho =[H,\rho ].$$
(1)
Here $`\rho `$ is the density matrix of the system, which can be expanded in the Pauli matrix basis,
$$\rho =\rho ^0I+\rho ^i\sigma ^i,$$
(2)
where $`I`$ is the unity matrix and $`\sigma ^i`$ are the Pauli matrices. In a lepton number violating parametrization for the evolution equation of the components of the density matrix has been assumed:
$$\frac{d}{dt}\left(\begin{array}{c}\rho ^0\\ \rho ^1\\ \rho ^2\\ \rho ^3\end{array}\right)=2\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& \mathrm{\Delta }m^2/(4E)& 0\\ 0& \mathrm{\Delta }m^2/(4E)& \alpha & \beta \\ 0& 0& \beta & \gamma \end{array}\right)\left(\begin{array}{c}\rho ^0\\ \rho ^1\\ \rho ^2\\ \rho ^3\end{array}\right).$$
(3)
Here $`\beta \alpha ,\gamma `$ . In also an alternative, lepton-number conserving parametrization has been discussed. However this parametrization will not influence neither the double beta decay observable nor the oscillation probability in the asymptotics of large propagation distances compared to the standard case of neutrino masses . Thus we will concentrate on the lepton-number violating case in the following. It should be mentioned however that a full analysis of the generalized dynamics requires six parameters .
Moreover, it should be stressed that this non-relativistic ansatz may not be suitable to describe ultrarelativisic particles such as neutrinos. However, while the covariant treatment of open quantum systems is still an unsolved problem, the density matrix ansatz has been successfully used in previous works to derive the “standard” mass mechanism neutrino oscillation probability also, see . Thus, while future works should improve the present ansatz, this approach seems to be suitable to provide at least a possibility for the comparison of the sensitivity of different experiments and a rough estimation for the order of magnitude of the obtained bounds.
## 3 Neutrinoless double beta decay
Neutrinoless double beta decay is one of the most sensitive tools in neutrino physics. It corresponds to two single beta decays occuring simultaneous in one nucleus, with a virtual neutrino propagating between the vertices. Important impact of this process has been derived on the reconstruction of the neutrino mass spectrum, physics beyond the standard model as well as more exotic phenomena such as violations of the equivalence principle or Lorentz invariance (for an overview see ). In the following we will study the potential of neutrinoless double beta decay for searches for CPT violations due to quantum foam interactions in the neutrino sector. The observable measured in neutrinoless double beta decay is the $`ee`$ entry of the neutrino mass matrix in the flavor space,
$$m_{ee}=\overline{m}\frac{\delta m}{2}cos(2\theta )$$
(4)
in a two neutrino scenario with $`\overline{m}=(m_1+m_2)/2`$ and $`\delta m=(m_2m_1)`$ and $`m_{1,2}`$ being the mass eigenstates. This quantity will be modified in the presence of QMV. The recent experimental constraint is $`m_{ee}<0.3`$ eV, obtained from the Heidelberg–Moscow experiment searching for double beta decays of <sup>76</sup>Ge . The GENIUS project will be sensitive to $`m_{ee}=10^210^3`$ eV . In the density matrix formalism the double beta decay observable can be expressed as follows:
$`Tr(\rho _{\nu _e}𝒪)`$ $`=`$ $`Tr\left(\begin{array}{cc}\rho _0+\rho _3& \rho _1i\rho _2\\ \rho _1+i\rho _2& \rho _0+\rho _3\end{array}\right).\left(\begin{array}{cc}m_1& 0\\ 0& m_2\end{array}\right)`$ (9)
$`=`$ $`(m_1+m_2)\rho _0+(m_2m_1)\rho _3.`$ (10)
The propagation time of the neutrino
$$t=\frac{1}{4\pi \mathrm{\Delta }E}610^{24}s$$
(11)
can be estimated by taking its energy to be of the size of the nuclear Fermi momentum $`p_F100`$ MeV for <sup>76</sup>Ge. Assuming $`\beta \alpha ,\gamma `$, eq. (5) yields
$`{\displaystyle \frac{d}{dt}}\rho _0`$ $`=`$ $`0`$ (12)
$`{\displaystyle \frac{d}{dt}}\rho _3`$ $`=`$ $`2\gamma \rho _3`$ (13)
and thus, using eq. (4)
$`\rho _0`$ $`=`$ $`{\displaystyle \frac{1}{2}}`$ (14)
$`\rho _3`$ $`=`$ $`e^{2\gamma t}{\displaystyle \frac{cos(2\theta )}{2}}.`$ (15)
This implies
$$m_{ee}^{QMV}=\overline{m}+e^{2\gamma t}\frac{\mathrm{\Delta }m}{2}\mathrm{cos}2\theta .$$
(16)
Due to the tiny propagation time (6) no significant variation of the double beta decay observable is obtained. However, from this analysis we realize that the distance plays a crucial role in constraining the QMV parameters, so we shall consider the bounds on the neutrino oscillation probability where neutrinos are propagating over large distances.
## 4 Oscillations of neutrinos from astrophysical sources
In the following we study the effect of quantum mechanics violation in neutrino oscillations from astrophysical sources. The most distant sources that have been discussed in the context of neutrino oscillations are supernovae (SN) and active galactic nuclei (AGN). While astrophysical sources have been discussed in the context of QMV effects on life time measurements , they have not been considered for the case of QMV induced neutrino oscillations so far.
For the neutrino oscillation case we get the survival and disappearance oscillation probabilities
$`P(\nu _x\nu _x)`$ $`=`$ $`Tr[\rho _{\nu _x}(t)\rho _{\nu _x}]`$ (17)
$`P(\nu _x\nu _x^{})`$ $`=`$ $`Tr[\rho _{\nu _x}(t)\rho _{\nu _x^{}}],`$ (18)
respectively. Here the density matrices can be parametrized as
$`\rho _{\nu _x}`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}^2\theta & \mathrm{cos}\theta \mathrm{sin}\theta \\ \mathrm{cos}\theta \mathrm{sin}\theta & \mathrm{sin}^2\theta \end{array}\right),`$ (21)
$`\rho _{\nu _x^{}}`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{sin}^2\theta & \mathrm{cos}\theta \mathrm{sin}\theta \\ \mathrm{cos}\theta \mathrm{sin}\theta & \mathrm{cos}^2\theta \end{array}\right).`$ (24)
As initial condition we assume
$$\rho (t=0)=\rho (\nu _e)$$
(25)
and thus :
$`\rho _0`$ $`=`$ $`{\displaystyle \frac{1}{2}}`$ (26)
$`\rho _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}(2\theta )`$ (27)
$`\rho _2`$ $`=`$ $`0`$ (28)
$`\rho _3`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{cos}(2\theta ).`$ (29)
The interesting observable is the oscillation propability
$$P_{\nu _x\nu _x^{}}^{QMV}=Tr[\rho (t)\rho _x]=\frac{1}{2}\frac{1}{2}e^{\gamma L}\mathrm{cos}^22\theta \frac{1}{2}e^{\alpha L}\mathrm{sin}^22\theta \mathrm{cos}(\frac{\mathrm{\Delta }m^2}{2E_\nu }L),$$
(30)
where $`\beta \alpha ,\gamma `$ has been assumed. For the n-flavour case the oscillation probability for large propagation distances is given by ,
$$P_{\nu _x\nu _x^{}}^{QMV}=\frac{1}{n}\frac{1}{n}e^{\gamma L},$$
(31)
where L is the propagation distance of the neutrinos.
This QMV oscillation probability can easily be distinguished from the asymptotics of the “standard” mass induced oscillation probability:
$$P_{\nu _x\nu _x^{}}^{mass}=\frac{\mathrm{sin}^22\theta }{2}.$$
(32)
The quantity $`P^{mass}`$ is fixed experimentally to $`P_{\nu _\mu \nu _\tau }^{mass}0.5`$ due to the maximal mixing in atmospheric neutrinos and Pνeντmass
<
0.05
<
subscriptsuperscript𝑃𝑚𝑎𝑠𝑠subscript𝜈𝑒subscript𝜈𝜏0.05P^{mass}_{\nu_{e}\to\nu_{\tau}}\mathrel{\vbox{\hbox{$<$}\nointerlineskip\hbox{$\sim$}}}0.05 due to the CHOOZ bound .
Supernovae 1987a: In supernovae strong neutrino oscillations will significantly distort the $`\nu _e`$ spectra at the earth, since the $`\nu _e`$ will aquire the spectra of the more energetic $`\nu _\mu `$ and $`\nu _\tau `$. The distance is very large. As a result, the condition that QMV should satisfy the bound on the oscillation probability gives a very strong bound. In the case of supernova 1987a, $`L50Mpc710^{39}GeV`$, so that the observed constraint on the oscillation probability $`P_{\nu _e\nu _{\mu ,\tau }}^{exp}<0.2`$ is satisfied for the three neutrino case when
$$\gamma <\frac{0.6}{L}10^{40}GeV.$$
(33)
We assumed here that $`P^{exp}`$ is the accuracy with which the deviations from the asymptotics $`1/n=1/3`$ can be measured. Due to the unknown energy dependence of the EHNS parameters and the Lorentz non-invariant ansatz it is difficult to compare these bounds with the bound coming from K-physics. Following we assume $`\gamma `$ to be of the order $`E_\nu ^2/M_{Pl}`$ and scale the obtained bound by the neutrino energy to the kaon mass squared,
$$\gamma _\nu \frac{E_\nu ^2}{m_K^2},$$
(34)
implying $`\gamma _K<10^{37}GeV`$, which is an improvement of about 16 orders of magnitude. This disfavors strongly any solution of the solar or atmospheric neutrino problem by lepton number violating QMV effects. If one assumes that the same QMV parametrization is valid for the $`K`$system, then any observational possibility in the $`K`$system will also be excluded by the present constraints from the supernovae analysis. A relativistic treatment of the problem can modify this bound to some extent, but it is most unlikely that the modification is by several orders of magnitude.
Active Galactic Nuclei (AGN) : AGN can be intense sources of high energy neutrinos ($`E_\nu 𝒪`$ (1 PeV)) . According to representative models the flux of these neutrinos is flavor dependent and the $`\nu _\tau `$ flux is reduced by at least two orders of magnitude compared to the $`\nu _e`$, $`\nu _\mu `$ fluxes. An unique appearance signal of high energy $`\nu _\tau `$ neutrinos can be a double bang signal of the produced $`\tau `$ leptons: The first bang originates from the CC interaction of the $`\tau `$ neutrino and the second one from the hadronic decay of the $`\tau `$ lepton. Deep underwater or ice neutrino detectors have been estimated to be sensitive on neutrino oscillation probabilities of
$$P_{\nu _{e,\mu }\nu _\tau }^{exp}<5\times 10^3.$$
(35)
Since the QMV effects become strong for large distances and higher energies, it is likely that when we have data from the active galactic nuclei on neutrino oscillations, these bounds will be modified by several orders of magnitude. Considering the distance to be $`L100Mpc`$ and the average energy of the neutrinos to be around 1 PeV, a bound on the neutrino oscillation of $`P_{\nu _e\nu _\mu }<5\times 10^3`$ will imply a corresponding bound on the QMV parameter
$$\gamma _\nu <10^{42}GeV.$$
(36)
Translation to the kaon mass scale yields
$$\gamma _K<10^{55}GeV,$$
(37)
which would imply the by far strongest bound on QMV parameters. This will provide a decisive test for any contribution of lepton number violating QMV effects in the neutrino sector.
## 5 Conclusions
We studied the effects of violation of quantum mechanics due to quantum space time foam interactions in neutrino experiments. While the non-observation of neutrinoless double beta decay does not give any significant constraint, the supernova 1987a implies a constraint being 16 orders of magnitude more stringent than the bounds known from the literature. This disfavors strongly any possibility of observable effects of lepton number violating QMV in any other experiments. The non-observation of QMV induced neutrino oscillations from active galactic nuclei will be able to improve this bound by many orders of magnitude. While the chosen non-relativistic ansatz might not be totally suitable for neutrinos, it should be at least useful to compare the sensitivity of different neutrino sources. Moreover the bounds obtained are that stringent, that, even in view of this ambiguity, they should be considered as the most restrictive ones.
## Acknowledgements
We thank J. Ellis, E. Lisi, S. Pakvasa, A.Y. Smirnov, the referee and especially N.E. Mavromatos for comments and useful discussions.
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# Matrix Theory over the Complex Quaternion Algebra
Abstract.We present in this paper some fundamental tools for developing matrix analysis over the complex quaternion algebra. As applications, we consider generalized inverses, eigenvalues and eigenvectors, similarity, determinants of complex quaternion matrices, and so on.
AMS Mathematics Subject Classification: 15A06; 15A24; 15A33
Key words: complex quaternion; matrix representation; universal similarity factorization; generalized inverse; eigenvalues and eigenvectors; determinant.
1. Introduction
The complex quaternion algebra (biquaternion algebra) $``$ is well known as a four dimensional vector space over the complex number field $``$ with its basis $`1,e_1,e_2,e_3`$ satisfying the multiplication laws
$$e_1^2=e_2^2=e_3^2=1,e_1e_2e_3=1.$$
$`(1.1)`$
$$e_1e_2=e_2e_1=e_3,e_2e_3=e_3e_2=e_1,e_3e_1=e_3e_3=e_2,$$
$`(1.2)`$
and 1 acting as unity element. In that case, any element in $``$ can be written as
$$a=a_0+a_1e_1+a_2e_2+a_3e_3,$$
$`(1.3)`$
where $`a_0`$$`a_3`$. According to this definition, real numbers, complex numbers, and real quaternions all can be regarded as the special cases of complex quaternions. A well-known fundamental fact on the complex quaternion algebra $``$ (see, e. g., ) is that it is algebraically isomorphic to the $`2\times 2`$ total matrix algebra $`^{2\times 2}`$ through the bijective map $`\psi :^{2\times 2}`$ satisfying
$$\psi (1)=\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right],\psi (e_1)=\left[\begin{array}{cc}i& 0\\ 0& i\end{array}\right],\psi (e_2)=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],\psi (e_3)=\left[\begin{array}{cc}0& i\\ i& 0\end{array}\right].$$
These four matrices are well-known as Pauli matrices. Based on this map, every element $`a=a_0+a_1e_1+a_2e_2+a_3e_3`$ has a faithful complex matrix representation as follows
$$\psi (a):=\left[\begin{array}{cc}a_0+a_1i& (a_2+a_3i)\\ a_2a_3i& a_0a_1i\end{array}\right]^{2\times 2},$$
$`(1.4)`$
In this article, we shall reveal a deeper relationship between $`a`$ and $`\psi (a)`$, which can simply be stated that there is an independent invertible matrix $`Q`$ of size 2 over $``$ such that all $`a`$ satisfy the following universal similarity factorization equality
$$Q^1\mathrm{diag}(a,a)Q=\psi (a),$$
where $`Q`$ has no relation with the expression of $`a`$. Moreover we also extend this equality to all $`m\times n`$ matrices over $``$. On the basis of these results, we shall consider several basic problems related to complex quaternion matrices, such as, generalized inverse, eigenvalues and eigenvectors, similarity, and determinant of complex quaternion matrices.
Some known terminology on complex quaternions are listed below (see, e.g., ). For $`a=a_0+a_1e_1+a_2e_2+a_3e_3,`$ the dual quaternion of $`a`$ is
$$\overline{a}=a_0a_1e_1a_2e_2a_3e_3;$$
$`(1.5)`$
the complex conjugate of $`a`$ is
$$a^{}=\overline{a_0}+\overline{a_1}e_1+\overline{a_2}e_2+\overline{a_3}e_3;$$
$`(1.6)`$
the Hermitian conjugate of $`a`$ is
$$a^{}=(\overline{a})^{}=\overline{a_0}\overline{a_1}e_1\overline{a_2}e_2\overline{a_3}e_3;$$
$`(1.7)`$
the weak norm of $`a`$ is
$$n(a)=a_0^2+a_1^2+a_2^2+a_3^2.$$
$`(1.8)`$
A quaternion $`a`$ is said to be real if $`a^{}=a`$, to be pure imaginary if $`a^{}=a`$, to be scalar if $`\overline{a}=a`$, to be Hermitian if $`a^{}=a`$.
For any $`A=(a_{st})^{m\times n}`$, the dual of $`A`$ is $`\overline{A}=(\overline{a_{ts}})^{n\times m}`$; the Hermitian conjugate of $`A`$ is $`A^{}=(a_{ts}^{})^{n\times m}`$. A square matrix $`A`$ is said to be self-dual if $`\overline{A}=A`$, it is Hermitian if $`A^{}=A`$, it is unitary if $`AA^{}=A^{}A=I`$, the indentity matrix, it is invertible if there is a matrix $`B`$ over $`𝒬`$ such that $`AB=BA=I`$.
Some known basic properties on complex quaternions and matrices of complex quaternions are listed below.
Lemma 1.1. Let $`a,b`$ be given. Then
(a) $`\overline{\overline{a}}=a,(a^{})^{}=a,(a^{})^{}=a.`$
(b) $`\overline{a+b}=\overline{a}+\overline{b},(a+b)^{}=a^{}+b^{},(a+b)^{}=a^{}+b^{}.`$
(c) $`\overline{ab}=\overline{b}\overline{a},(ab)^{}=a^{}b^{},(ab)^{}=b^{}a^{}.`$
(d) $`a\overline{a}=\overline{a}a=n(a)=n(\overline{a});`$
(e) $`a`$ is invertible if and only if $`n(a)0,`$ in that case $`a^1=n^1(a)\overline{a}.`$
Lemma 1.2. Let $`A^{m\times n},B^{n\times p}`$ be given. Then
(a) $`\overline{\overline{A}}=A,(A^{})^{}=A.`$
(b) $`\overline{AB}=\overline{B}\overline{A},(AB)^{}=B^{}A^{}.`$
(c) $`(AB)^1=B^1A^1,`$ if $`A`$ and $`B`$ are invertible.
(d) $`(\overline{A})^1=\overline{(A^1)},(A^{})^1=(A^1)^{},`$ if $`A`$ is invertible.
2. A universal similarity factorization equality over complex quaternion algebra
We first present a general result on the universal similarity factorization of elements over $`2\times 2`$ total matrix algebra.
Lemma 2.1.Let $`M_2(𝔽)`$ be the $`2\times 2`$ total matrix algebra over an arbitrary field $`𝔽`$ with its basis $`e_{11},e_{12},e_{21}`$ and $`e_{22}`$ satisfying the following multiplication rules
$$e_{st}e_{pq}=\{\begin{array}{c}e_{sq},t=p\\ 0,tp\end{array},s,t,p,q=1,\mathrm{\hspace{0.17em}2}.$$
$`(2.1)`$
Then for any $`a=a_{11}e_{11}+a_{12}e_{12}+a_{21}e_{21}+a_{22}e_{22}M_2(),`$ where $`a_{st}`$, the corresponding diagonal matrix $`\mathrm{diag}(a,a)`$ satisfies the following universal similarity factorization equality
$$Q\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill a\end{array}\right]Q^1=\left[\begin{array}{cc}a_{11}& a_{12}\\ a_{21}& a_{22}\end{array}\right]^{2\times 2},$$
$`(2.2)`$
where $`Q`$ has the independent form
$$Q=Q^1=\left[\begin{array}{cc}e_{11}& e_{21}\\ e_{12}& e_{22}\end{array}\right].$$
$`(2.3)`$
Proof. According to Eq.(2.1), it is easy to verify that the unity element in $`M_2()`$ is $`e=e_{11}+e_{22}`$. In that case, the matrix $`Q`$ in Eq.(2.3) satisfies
$`Q^2=\left[\begin{array}{cc}e_{11}& e_{21}\\ e_{12}& e_{22}\end{array}\right]=\left[\begin{array}{cc}e_{11}^2+e_{21}e_{12}& e_{11}e_{21}+e_{21}e_{22}\\ e_{12}e_{11}+e_{22}e_{12}& e_{12}e_{21}+e_{22}^2\end{array}\right]=\left[\begin{array}{cc}e_{11}+e_{22}& 0\\ 0& e_{11}+e_{22}\end{array}\right]=eI_2,`$
which implies that $`Q`$ is invertible over $`M_2()`$ and $`Q=Q^1`$. Next multiplying the three matrices in the left-hand side of Eq.(2.2) yields the right-hand side of Eq.(2.2). $`\mathrm{}`$
Theorem 2.2.Let $`a=a_0+a_1e_1+a_2e_2+a_3e_3`$ be given. Then the diagonal matrix $`\mathrm{diag}(a,a)`$ satisfies the following universal factorization similarity equality
$$Q\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill a\end{array}\right]Q=\left[\begin{array}{cc}a_0+a_1i& (a_2+a_3i)\\ a_2a_3i& a_0a_1i\end{array}\right]=\psi (a)^{2\times 2},$$
$`(2.4)`$
where $`Q`$ is an unitary matrix over $``$
$$Q=A^1=Q^{}=\frac{1}{2}\left[\begin{array}{cc}1ie_1& e_2+ie_3\\ e_2+ie_3& 1+ie_1\end{array}\right].$$
$`(2.5)`$
Proof. According to Eq.(2.4), we choose a new basis for $``$ as follows
$$e_{11}=\frac{1}{2}\left(\mathrm{\hspace{0.17em}1}ie_1\right),e_{12}=\frac{1}{2}\left(e_2+ie_3\right),e_{21}=\frac{1}{2}\left(e_2+ie_3\right),e_{22}=\frac{1}{2}\left(\mathrm{\hspace{0.17em}1}+ie_2\right).$$
$`(2.6)`$
Then it is not difficult to verify that the above basis satisfies the multiplication rules in Eq.(2.1). Under this new basis, any element $`a=a_0+a_1e_1+a_2e_2+a_3e_3`$ can be expressed as
$$a=(a_0+ia_1)e_{11}+(a_2ia_3)e_{12}+(a_2ia_3)e_{21}+(a_0ia_1)e_{22}.$$
$`(2.7)`$
Substituting Eqs.(2.6) and (2.7) into Eq.(2.2), we obtain Eqs.(2.4) and (2.5). $`\mathrm{}`$
The equality in Eq.(2.4) can also equivalently be expressed as
$$Q\left[\begin{array}{cc}\hfill a_{11}& \hfill a_{12}\\ \hfill a_{21}& \hfill a_{22}\end{array}\right]Q=\left[\begin{array}{cc}\hfill a& \hfill 0\\ \hfill 0& \hfill a\end{array}\right]^{2\times 2},$$
$`(2.8)`$
where $`a_{st}`$ is arbitrary, $`Q`$ is as in Eq.(2.5), and $`a`$ has the form
$$a=\frac{1}{2}(a_{11}+a_{22})+\frac{1}{2}(a_{22}a_{11})ie_1+\frac{1}{2}(a_{21}a_{12})e_2+\frac{1}{2}(a_{12}+a_{21})ie_3.$$
$`(2.9)`$
This equality shows that every $`2\times 2`$ complex matrix is uniformly similar to an diagonal matrix with the form $`aI_2`$ over the complex quaternion algebra $``$.
The complex quaternions and their complex matrix representations satisfy the following operation properties.
Theorem 2.3Let $`a=a_0+a_1e_1+a_2e_2+a_3e_3,b,\lambda `$ be given. Then
(a) $`a=b\psi (a)=\psi (b).`$
(b) $`\psi (a+b)=\psi (a)+\psi (b),\psi (ab)=\psi (a)\psi (b),\psi (\lambda a)=\psi (a\lambda )=\lambda \psi (a),\psi (1)=I_2.`$
(c) $`\psi (\overline{a})=\left[\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right]\psi ^T(a)\left[\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right].`$
(d) $`\psi (a^{})=\left[\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right]\overline{\psi (a)}\left[\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right].`$
(e) $`\psi (a^{})=\overline{\psi (a)}^T=\psi ^{}(a),`$ the conjugate transpose of $`\psi (a).`$
(f) $`det\psi (a)=n(a)=a_0^2+a_1^2+a_2^2+a_3^2;`$
(g) $`a=\frac{1}{4}E_2\psi (a)E_2^{},`$ where $`E_2=[\mathrm{\hspace{0.17em}1}ie_1,e_2+ie_3].`$
(h) $`a`$ is invertible if and only if $`\psi (a)`$ is invertible, in that case, $`\psi (a^1)=\psi ^1(a)`$ and $`a^1=\frac{1}{4}E_2\psi ^1(a)E_2^{}.`$
For a noninvertible element over $``$, we can define its Moore-Penrose inverse as follows.
Definition. Let $`a`$ be given. If the following four equations
$$axa=a,xax=x,(ax)^{}=ax,(xa)^{}=xa$$
$`(2.10)`$
have a common solution $`x`$, then this solution is called the Moore-Penrose inverse of $`a`$, and denoted by $`x=a^+`$.
The existence and the uniqueness of the Moore-Penrose inverse of a complex quaternion $`a`$ can be determined by its matrix representation $`\psi (a)`$ over $``$. In fact, according to Theorem 2.3(a), (b) and (e), the four equations in Eq.(2.10) are equivalent to the following four equations over $``$
$$\psi (a)\psi (x)\psi (a)=\psi (a),\psi (x)\psi (a)\psi (x)=\psi (x),$$
$$[\psi (a)\psi (x)]^{}=\psi (a)\psi (x),[\psi (x)\psi (a)]^{}=\psi (x)\psi (a).$$
According to the complex matrix theory, the following four equations
$$\psi (a)Y\psi (a)=\psi (a),Y\psi (a)Y=Y,[\psi (a)Y]^{}=\psi (a)Y,[Y\psi (a)]^{}=Y\psi (a).$$
has a unique solution $`Y=\psi ^+(a)`$, the Moore-Penrose inverse of $`\psi (a)`$. Then by Eq.(2.8), it follows that there must be a unique $`x`$ over $``$ such that $`\psi (x)=Y=\psi ^+(a)`$, in which case, this $`x`$ can be expressed as
$$x=\frac{1}{4}E_2YE_2^{}=\frac{1}{4}E_2\psi ^+(a)E_2^{}.$$
Correspondingly this $`x`$ is the unique solution to Eq.(2.10). In summary, we have the following.
Theorem 2.4.Let $`a`$ be given. Then its Moore-Penrose inverse $`a^+`$ exists uniquely, and satisfies the following four equalities
$$\psi (a^+)=\psi ^+(a),a^+=\frac{1}{4}E_2\psi ^+(a)E_2^{},$$
where $`E_2=[\mathrm{\hspace{0.17em}1}ie_1,e_2+ie_3].`$
One of the basic problems related to complex quaternions is concerned with similarity of two complex quaternions. As usual, two complex quaternions $`a`$ and $`b`$ are said to be similar if there is an invertible complex quaternion $`x`$ such that $`x^1ax=b`$, and this is written as $`ab`$. It is easy to verify that the similarity mentioned here is an equivalence relation on complex quaternions.
A simple result follows immediately from the above definition, Eqs.(2.4) and (2.8).
Theorem 2.5.Let $`a,b`$ be given. Then
$$ab\psi (a)\psi (b).$$
$`(2.11)`$
From Eq.(2.11), we easily find the following.
Theorem 2.6.Let $`a=a_0+a_1e_1+a_2e_2+a_3e_3`$ given with $`a.`$
(a) If $`a_1^2+a_2^2+a_3^20,`$ then $`aa_0+\tau (a)e_1,`$ where $`\tau (a)`$ is complex number satisfying $`\tau ^2(a)=a_1^2+a_2^2+a_3^2.`$
(b) If $`a_1^2+a_2^2+a_3^2=0,`$ then $`aa_0\frac{1}{2}e_2+\frac{1}{2}ie_3.`$
Proof. For any $`a`$, the characteristic polynomial of its complex matrix representation $`\psi (a)`$ is
$$|\lambda I_2\psi (a)|=\left|\begin{array}{cc}\lambda (a_0+a_1i)& a_2+a_3i\\ a_2+a_3i& \lambda (a_0a_1i)\end{array}\right|=(\lambda a_0)^2+a_1^2+a_2^2+a_3^2.$$
From it we immediately know that if $`a_1^2+a_2^2+a_3^20`$, then
$$\psi (a)\left[\begin{array}{cc}a_0+\tau (a)i& 0\\ 0& a_0\tau (a)i\end{array}\right]=\psi [a_0+\tau (a)e_1],$$
$`(2.12)`$
and if $`a_1^2+a_2^2+a_3^2=0`$, then
$$\psi (a)\left[\begin{array}{cc}a_0& 1\\ 0& a_0\end{array}\right]=\psi \left(a_0\frac{1}{2}e_2+\frac{1}{2}ie_3\right).$$
$`(2.13)`$
Correspondingly applying Eq.(2.11) to Eqs.(2.12) and (2.13) may lead to Part (a) and Part (b) of this theorem. $`\mathrm{}`$
3. Two universal factorization equalities on complex quaternion matrices
In this section, we extend the universal similarity factorization equality in (2.4) to any $`m\times n`$ matrix over $``$, and give some of its consequences.
Theorem 3.1.Let $`A=A_0+A_1e_1+A_2e_2+A_3e_3^{m\times n}`$ be given. Then $`A`$ satisfies the following universal factorization equality
$$Q_{2m}\left[\begin{array}{cc}\hfill A& \hfill 0\\ \hfill 0& \hfill A\end{array}\right]Q_{2n}=\left[\begin{array}{cc}A_0+A_1i& (A_2+A_3i)\\ A_2A_3i& A_0A_1i\end{array}\right]:=\mathrm{\Psi }(A)^{2m\times 2n},$$
$`(3.1)`$
where $`Q_{2t}`$ has the independent form
$$Q_{2t}=Q_{2t}^1=Q_{2t}^{}=\frac{1}{2}\left[\begin{array}{cc}(\mathrm{\hspace{0.17em}1}ie_1)I_t& (e_2+ie_3)I_t\\ (e_2+ie_3)I_t& (\mathrm{\hspace{0.17em}1}+ie_1)I_t\end{array}\right],t=m,n.$$
$`(3.2)`$
In particular, when $`m=n,`$ Eq.(3.1)becomes a universal similarity factorization equality over $``$ .
Proof. It follows directly from multiplying out the three block matrices in the left-hand side of Eq.(3.1). $`\mathrm{}`$
The matrix $`\mathrm{\Psi }(A)`$ in Eq.(3.1) is called the complex representation of $`A`$. If setting
$$A_{11}=A_0+A_1i,A_{12}=(A_2+A_3i),A_{21}=A_2A_3i,A_{11}=A_0A_1i$$
in Eq.(3.1), then it can equivalently be expressed as
$$Q_{2m}\left[\begin{array}{cc}\hfill A_{11}& \hfill A_{12}\\ \hfill A_{21}& \hfill A_{22}\end{array}\right]Q_{2n}=\left[\begin{array}{cc}\hfill A& \hfill O\\ \hfill 0& \hfill A\end{array}\right]^{2m\times 2n},$$
$`(3.3)`$
where $`A_{st}^{m\times n}`$ are arbitrary, $`Q_{2t}`$ is as in Eq.(3.2), and $`A`$ has the following form
$$A=\frac{1}{2}(A_{11}+A_{22})+\frac{1}{2}(A_{22}A_{11})ie_1+\frac{1}{2}(A_{21}A_{12})e_2+\frac{1}{2}(A_{12}+A_{21})ie_3^{m\times n}.$$
$`(3.4)`$
which can also be stated that for any matrix $`M^{2m\times 2n}`$, there must be a unique matrix $`A^{m\times n}`$ such that
$$\mathrm{\Psi }(A)=M.$$
$`(3.5)`$
Various operation properties on complex representation of complex quaternion matrices can easily be derived from Eq.(3.1).
Theorem 3.2.Let $`A,B^{m\times n},C^{n\times p},`$ and $`\lambda `$ be given. Then
(a) $`A=B\mathrm{\Psi }(A)=\mathrm{\Psi }(B).`$
(b) $`\mathrm{\Psi }(A+B)=\mathrm{\Psi }(A)+\mathrm{\Psi }(B),\mathrm{\Psi }(AC)=\mathrm{\Psi }(A)\mathrm{\Psi }(C),\mathrm{\Psi }(\lambda A)=\mathrm{\Psi }(A\lambda )=\lambda \mathrm{\Psi }(A).`$
(c) $`\mathrm{\Psi }(A^{})=\left[\begin{array}{cc}O& I_n\\ I_n& O\end{array}\right]\mathrm{\Psi }^T(A)\left[\begin{array}{cc}O& I_m\\ I_m& O\end{array}\right].`$
(d) $`\mathrm{\Psi }_1(A^{})=\mathrm{\Psi }^{}(A),`$ the conjugate transpose of $`\mathrm{\Psi }(A).`$
(e) $`A=\frac{1}{4}E_{2m}\mathrm{\Psi }(A)E_{2n}^{},`$ where $`E_{2t}=[(\mathrm{\hspace{0.17em}1}ie_1)I_t,(e_2+ie_3)I_t],t=m,n.`$
(f) $`A`$ is invertible if and only if $`\mathrm{\Psi }(A)`$ is invertible, in that case, $`\mathrm{\Psi }(A^1)=\mathrm{\Psi }^1(A)`$ and $`A^1=\frac{1}{4}E_{2m}\mathrm{\Psi }^1(A)E_{2m}^{}.`$
(g) $`\mathrm{\Psi }(A)E_{2n}^{}E_{2n}=E_{2m}^{}E_{2m}\mathrm{\Psi }(A).`$
(h) $`A`$ is Hermitian if and only if $`\mathrm{\Psi }(A)`$ is Hermitian over $``$.
(i) $`A`$ is unitary if and only if $`\mathrm{\Psi }(A)`$ is unitary over $``$.
Another universal factorization equality on complex quaternion matrices is established as follows.
Theorem 3.3.Let $`A=(a_{st})^{m\times n},`$ and denote $`D_A=(a_{st}I_2)^{2m\times 2n}`$. Then $`D_A`$ satisfies the following universal factorization equality
$$P_{2m}D_AP_{2n}=\left[\begin{array}{ccc}\psi (a_{11})& \mathrm{}& \psi (a_{1n})\\ \mathrm{}& & \mathrm{}\\ \psi (a_{m1})& \mathrm{}& \psi (a_{mn})\end{array}\right]:=\psi (A)^{2m\times 2n},$$
$`(3.6)`$
where $`P_{2t}`$ has the form
$$P_{2t}=P_{2t}^1=P_{2t}^{}=\mathrm{diag}(Q,\mathrm{},Q),Q=Q^1=Q^{}=\frac{1}{2}\left[\begin{array}{cc}1ie_1& e_2+ie_3\\ e_2+ie_3& 1+ie_1\end{array}\right].t=m,n.$$
In particular, when $`m=n,`$ Eq.(3.6) becomes a universal similarity factorization equality over $`.`$
Proof. Observe from Eq.(2.1) that $`Q(a_{st}I_2)Q=\psi (a_{st})`$. We immediately obtain
$$P_{2m}D_AP_{2n}=[Q(a_{st}I_2)Q]_{m\times n}=[\psi (a_{st})]_{m\times n}=\psi (A).$$
which is exactly Eq.(3.6). $`\mathrm{}`$
The complex matrix $`\psi (A)`$ in Eq.(3.6) is also called the complex representation of $`A`$. Clearly the two complex matrices $`\mathrm{\Psi }(A)`$ in Eq.(3.1) and $`\psi (A)`$ Eq.(3.6) are permutationally equivalent, that is, there are two permutation matrices $`G`$ and $`H`$ such that $`G\mathrm{\Psi }(A)H=\psi (A)`$.
For convenience of application, Eq.(3.6) can be simply stated that for any $`M^{2m\times 2n},`$ there must exist an $`A^{m\times n}`$ such that
$$\psi (A)=M.$$
$`(3.7)`$
Theorem 3.4.Let $`A,B^{m\times n},C^{n\times p},\lambda `$ be given. Then
(a) $`A=B\psi (A)=\psi (B).`$
(b) $`\psi (A+B)=\psi (A)+\psi (B),\psi (AC)=\psi (A)\psi (C),\psi (\lambda A)=\psi (A\lambda )=\lambda \psi (A).`$
(c) $`\psi (I_m)=I_{2m}.`$
(d) $`\psi (A^{})=\psi ^{}(A),`$ the conjugate transpose of $`\psi (A).`$
(e) $`A`$ is invertible if and only if $`\psi (A)`$ is invertible, in that case, $`\psi (A^1)=\psi ^1(A).`$
(h) $`A`$ is Hermitian if and only if $`\psi (A)`$ is Hermitian.
(i) $`A`$ is unitary if and only if $`\psi (A)`$ is unitary.
Just as for complex quaternions, we can define the Moore-Penrose inverse of any $`m\times n`$ complex quaternion matrix as follows.
Definition. Let $`A^{m\times n}`$ be given. If the following four equations
$$AXA=A,XAX=X,(AX)^{}=AX,(XA)^{}=XA$$
$`(3.8)`$
have a common solution for $`X`$, then this solution is called the Moore-Penrose inverse of $`A`$, and denoted by $`X=A^+`$.
The existence and the uniqueness of the Moore-Penrose inverse of a complex quaternion matrix $`A`$ can be determined by its complex representation $`\mathrm{\Psi }(A)`$. In fact, according to Theorem 3.2(a), (c) and (d), the four equations in Eq.(3.8) are equivalent to the following four complex matrix equations
$$\mathrm{\Psi }(A)\mathrm{\Psi }(X)\mathrm{\Psi }(A)=\mathrm{\Psi }(A),\mathrm{\Psi }(X)\mathrm{\Psi }(A)\mathrm{\Psi }(X)=\mathrm{\Psi }(X),$$
$$[\mathrm{\Psi }(A)\mathrm{\Psi }(X)]^{}=\mathrm{\Psi }(A)\mathrm{\Psi }(X),[\mathrm{\Psi }(X)\mathrm{\Psi }(A)]^{}=\mathrm{\Psi }(X)\mathrm{\Psi }(A).$$
According to the complex matrix theory, the following four equations
$$\mathrm{\Psi }(A)Y\mathrm{\Psi }(A)=\mathrm{\Psi }(A),Y\mathrm{\Psi }(A)Y=Y,[\mathrm{\Psi }(A)Y]^{}=\mathrm{\Psi }(A)Y,[Y\mathrm{\Psi }(A)]^{}=Y\mathrm{\Psi }(A)$$
have a unique common solution $`Y=\mathrm{\Psi }^+(A)`$, the Moore-Penrose of $`\mathrm{\Psi }(A)`$. Then by Eq.(3.5), there is a unique matrix $`X`$ over $``$ such that $`\mathrm{\Psi }(X)=Y=\mathrm{\Psi }^+(A)`$, in which case, this $`X`$ can be expressed as
$$X=\frac{1}{4}E_{2n}YE_{2m}^{}=\frac{1}{4}E_{2n}\mathrm{\Psi }^+(A)E_{2m}^{}.$$
Correspondingly this $`X`$ is the unique solution to Eq.(3.8). Hence we have the following.
Theorem 3.5.Let $`A^{m\times n}`$ be given. Then its Moore-Penrose inverse $`A^+`$ of $`A`$ exists uniquely, and satisfies the following two equalities
$$\mathrm{\Psi }(A^+)=\mathrm{\Psi }^+(A),A^+=\frac{1}{4}E_{2n}\mathrm{\Psi }^+(A)E_{2m}^{},$$
where $`E_{2t}=[(\mathrm{\hspace{0.17em}1}ie_1)I_t,(e_2+ie_3)I_t],t=m,n.`$
Through the complex representation in Eq.(3.1), we can define the rank of any $`A^{m\times n}`$ as follows $`\mathrm{rank}(A)=\frac{1}{2}\mathrm{rank}[\mathrm{\Psi }(A)].`$ Obviously the rank of a complex quaternion matrix $`A`$ is a fraction if the rank of $`\mathrm{\Psi }(A)`$ is an odd number. In particular, $`A^{m\times m}`$ is invertible if and only if $`\mathrm{rank}(A)=m.`$
On the basis of the above results we now are able to investigate various kinds of problems related to complex quaternion matrices. In the next three sections, we consider three basic problems—right eigenvalues and eigenvectors, similarity, as well as determinants of square complex quaternion matrices.
4. Right eigenvalues and eigenvectors of complex quaternion matrices
As usual, right eigenvalue equation for a complex quaternion matrix $`A^{n\times n}`$ is defined by
$$AX=X\lambda ,X^{n\times 1},\lambda .$$
$`(4.1)`$
If a $`\lambda `$ and a nonzero $`X^{n\times 1}`$ satisfy Eq.(4.1), then $`\lambda `$ is called a right eigenvalue of $`A`$ and $`X`$ is called an eigenvector associated with $`\lambda `$. In particular, if a $`\lambda `$ and an $`X^{n\times 1}`$ with rank$`(X)=1`$ satisfy Eq.(4.1), then $`\lambda `$ is called a regular right eigenvalue of $`A`$ and $`X`$ is called an regular eigenvector associated with $`\lambda `$. In this section, we shall prove that any square complex quaternion matrix has at least one complex right eigenvalue, and also has at least one regular right eigenvalue. To do so, we need some preparation.
Definition. For any $`X=X_0+X_1e_1+X_2e_2+X_3e_3^{n\times 1}`$, we call the following complex matrix
$$\left[\begin{array}{c}\hfill X_0+X_1i\\ \hfill X_2X_3i\end{array}\right]:=\stackrel{}{X}^{2n\times 1},$$
$`(4.2)`$
uniquely determined by $`X`$, the complex adjoint vector of $`X`$.
Lemma 4.1.Let $`A^{m\times n},X^{n\times 1},`$ and $`\lambda `$ be given. Then $`\stackrel{}{AX}=\mathrm{\Psi }(A)\stackrel{}{X}`$ and $`\stackrel{}{X\lambda }=\stackrel{}{X}\lambda .`$
Proof. It is easy to see that for any column matrix $`Y^{n\times 1}`$, the corresponding $`\mathrm{\Psi }(Y)`$ and $`\stackrel{}{Y}`$ satisfy the relation $`\stackrel{}{Y}=\mathrm{\Psi }(Y)[\mathrm{\hspace{0.17em}1},0]^T.`$ Now applying it to $`AX`$ and $`X\lambda `$ we find
$$\stackrel{}{AX}=\mathrm{\Psi }(AX)[\mathrm{\hspace{0.17em}1},0]^T=\mathrm{\Psi }(A)\mathrm{\Psi }(X)[\mathrm{\hspace{0.17em}1},0]^T=\mathrm{\Psi }(A)\stackrel{}{X},$$
$$\stackrel{}{X\lambda }=\mathrm{\Psi }(X\lambda )[\mathrm{\hspace{0.17em}1},0]^T=\mathrm{\Psi }(X)\lambda [\mathrm{\hspace{0.17em}1},0]^T=\stackrel{}{X}\lambda .\mathrm{}$$
Based on the above notation, we can deduce the following several results.
Theorem 4.2.Let $`A^{n\times n}`$ be given. Then all the eigenvalues of $`\mathrm{\Psi }(A)`$ are right eigenvalues of $`A,`$ and all the eigenvectors of $`\mathrm{\Psi }(A)`$ can be used for constructing eigenvectors of $`A.`$
Proof. Assume that $`\lambda `$ and $`OY^{2n\times 1}`$ satisfy
$$\mathrm{\Psi }(A)Y=Y\lambda .$$
Then we set $`X=E_{2n}Y^{n\times 1},`$ where $`E_{2n}=[(\mathrm{\hspace{0.17em}1}ie_1)I_n,(e_2+ie_3)I_n]`$, and it satisfies $`E_{2n}E_{2n}^{}=4I_n.`$ Combining the above results with Theorem 3.2(e) and (g), we obtain
$$AX=\frac{1}{4}E_{2n}\mathrm{\Psi }(A)E_{2n}^{}E_{2n}Y=\frac{1}{4}E_{2n}E_{2n}^{}E_{2n}\mathrm{\Psi }(A)Y=E_{2n}Y\lambda =X\lambda ,$$
which shows that $`\lambda `$ is a right eigenvalue of $`A`$ and $`X=E_{2n}Y`$ is an eigenvector of $`A`$ associated with this $`\lambda `$. So the conclusion of the theorem is true. $`\mathrm{}`$
Conversely, we have the following result.
Theorem 4.3.Let $`A^{n\times n}`$ be given. Then all the complex right eigenvalues of $`A`$ are eigenvalues of $`\mathrm{\Psi }(A),`$ too.
Proof. Assume that $`\lambda `$ and $`OX^{n\times 1}`$ satisfy $`AX=X\lambda .`$ Then applying Lemma 4.1 to the both sides of this equality, we find $`\mathrm{\Psi }(A)\stackrel{}{X}=\stackrel{}{X}\lambda ,`$ which shows that $`\lambda `$ and $`\stackrel{}{X}`$ are a pair of eigenvalue and eigenvector of $`\mathrm{\Psi }(A)`$. $`\mathrm{}`$
Next are two results on the regular right eigenvalues and eigenvectors of complex quaternion matrices.
Theorem 4.4.Let $`A^{n\times n}`$ be given. Then $`A`$ has at least one regular right eigenvalue.
Proof. We first assume that the complex representation $`\mathrm{\Psi }(A)`$ has at least two linearly independent eigenvectors $`Y_1`$ and $`Y_2`$, i.e.,
$$\mathrm{\Psi }(A)Y_1=Y_1\lambda _1,\mathrm{\Psi }(A)Y_2=Y_2\lambda _2,$$
$`(4.3)`$
where $`\mathrm{rank}[Y_1,Y_2]=2,`$ $`\lambda _1`$ and $`\lambda _2`$ are two complex numbers with $`\lambda _1=\lambda _2`$, or $`\lambda _1\lambda _2`$. Then
$$\mathrm{\Psi }(A)[Y_1,Y_2]=[Y_1,Y_2]\left[\begin{array}{cc}\hfill \lambda _1& \hfill 0\\ \hfill 0& \hfill \lambda _2\end{array}\right].$$
$`(4.4)`$
According to Eq.(3.5), there must exist an $`X^{n\times 1}`$ such that
$$\mathrm{\Psi }(X)=[Y_1,Y_2],\mathrm{rank}(X)=\frac{1}{2}\mathrm{rank}\mathrm{\Psi }(X)=1,$$
$`(4.5)`$
meanwhile there must exist a $`\lambda `$ such that
$$\mathrm{\Psi }(\lambda )=\left[\begin{array}{cc}\hfill \lambda _1& \hfill 0\\ \hfill 0& \hfill \lambda _2\end{array}\right].$$
$`(4.6)`$
Putting Eqs.(4.5) and (4.6) into Eq.(4.4) and then applying Theorem 3.2(a) to it, we obtain
$$\mathrm{\Psi }(A)\mathrm{\Psi }(X)=\mathrm{\Psi }(X)\mathrm{\Psi }(\lambda )AX=X\lambda ,and\mathrm{rank}(X)=1,$$
which shows that $`\lambda `$ and $`X`$ are a pair of regular right eigenvalue and eigenvector of $`A`$.
Next assume that the complex matrix $`\mathrm{\Psi }(A)`$ has only one eigenvalue $`\lambda _1`$ and only one linearly independent eigenvector $`Y_1`$ associated with $`\lambda _1`$. Then according to complex matrix theory, there exists another vector $`Y_2`$ over $``$ such that
$$\mathrm{\Psi }(A)[Y_1,Y_2]=[Y_1,Y_2]\left[\begin{array}{cc}\lambda _1& 1\\ 0& \lambda _1\end{array}\right],$$
$`(4.7)`$
where rank$`[Y_1,Y_2]=2.`$ Then it follows by Eq.(3.5) that there must exist a $`\lambda `$ and an $`X^{n\times 1}`$ such that
$$\mathrm{\Psi }(\lambda )=\left[\begin{array}{cc}\lambda _1& 1\\ 0& \lambda _1\end{array}\right],\mathrm{\Psi }(X)=[Y_1,Y_2],\mathrm{rank}(X)=\frac{1}{2}\mathrm{rank}\mathrm{\Psi }(X)=1.$$
$`(4.8)`$
In that case, combining Eqs.(4.7) and (4.8) and Theorem 3.2(a), we get
$$\mathrm{\Psi }(A)\mathrm{\Psi }(X)=\mathrm{\Psi }(X)\mathrm{\Psi }(\lambda )AX=X\lambda ,r(X)=1,$$
which shows that $`\lambda `$ and $`X`$ are a pair of regular right eigenvalue and eigenvector of $`A`$. $`\mathrm{}`$
Theorem 4.5.Let $`A^{n\times n}`$ be given. Then from any regular right eigenvalue of $`A`$ we can derive eigenvalues of $`\mathrm{\Psi }(A).`$
Proof. Assume that $`\lambda =\lambda _0+\lambda _1e_1+\lambda _2e_2+\lambda _3e_3`$ and $`X^{n\times 1}`$ with $`\mathrm{rank}(X)=1`$ satisfy
$$AX=X\lambda .$$
$`(4.9)`$
If $`\lambda `$ in Eq.(4.9), then from Theorem 3.3(a) and (b) we get
$$\mathrm{\Psi }(A)\mathrm{\Psi }(X)=\mathrm{\Psi }(X)\mathrm{\Psi }(\lambda )=\mathrm{\Psi }(X)\left[\begin{array}{cc}\hfill \lambda & \hfill 0\\ \hfill 0& \hfill \lambda \end{array}\right],$$
which clearly shows that $`\lambda `$ is an eigenvalue of $`\mathrm{\Psi }(A)`$, and $`\mathrm{\Psi }(X)`$ are two eigenvectors of $`\mathrm{\Psi }(A)`$. If $`\lambda `$ in Eq.(4.9), then by Theorem 2.6 we know that this $`\lambda `$ can be expressed as
$$\lambda =p[\lambda _0+\tau (\lambda )e_1]p^1,$$
$`(4.10)`$
when $`\tau ^2(\lambda )=\lambda _1^2+\lambda _2^2+\lambda _3^20`$, or
$$\lambda =q\left(\lambda _0\frac{1}{2}e_2+\frac{1}{2}ie_3\right)q^1,$$
$`(4.11)`$
when $`\lambda _1^2+\lambda _2^2+\lambda _3^2=0`$. Under the condition in (4.10), the equality in Eq.(4.9) can be equivalently expressed as
$$A(Xp)=(Xp)[\lambda _0+\tau (\lambda )e_1].$$
$`(4.12)`$
Applying Theorem 3.2(b) to the both sides of Eq.(4.12) yields
$$\mathrm{\Psi }(A)\mathrm{\Psi }(Xp)=\mathrm{\Psi }_1(Xp)\mathrm{\Psi }[\lambda _0+\tau (\lambda )e_1]=\mathrm{\Psi }(Xp)\left[\begin{array}{cc}\lambda _0+\tau (\lambda )i& 0\\ 0& \lambda _0\tau (\lambda )i\end{array}\right],$$
which shows that $`\lambda _0\pm \tau (\lambda )i`$ are two eigenvalues of $`\mathrm{\Psi }(A)`$, and $`\mathrm{\Psi }(Xp)`$ are two eigenvectors of $`\mathrm{\Psi }(A)`$. Under the condition in Eq.(4.11), the equality in (4.9) can equivalently be expressed as
$$A(Xq)=(Xq)\left(\lambda _0\frac{1}{2}e_2+\frac{1}{2}ie_3\right).$$
$`(4.13)`$
Then applying Theorem 3.2(b) to the both sides of Eq.(4.13) yields
$$\mathrm{\Psi }(A)\mathrm{\Psi }(Xq)=\mathrm{\Psi }(Xq)\mathrm{\Psi }\left(\lambda _0\frac{1}{2}e_2+\frac{1}{2}ie_3\right)=\mathrm{\Psi }(Xq)\left[\begin{array}{cc}\lambda _0& 1\\ 0& \lambda _0\end{array}\right],$$
which shows that $`\lambda _0`$ is an eigenvalue of $`\mathrm{\Psi }(A)`$, and the first column of $`\mathrm{\Psi }(Xp)`$ is an eigenvector of $`\mathrm{\Psi }(A)`$. $`\mathrm{}`$
5. Similarity of complex quaternion matrices
Two square matrices $`A`$ and $`B`$ of the same size over $``$ are said to be similar if there exists an invertible matrix $`X`$ over $``$ such that $`X^1AX=B`$. Based on the results in the preceding two subsections, we can easily find a simple result charactering the similarity of two complex quaternion matrices.
Theorem 5.1.Let $`A,B^{n\times n}`$ be given. Then the following three statements are equivalent:
(a) $`A`$ and $`B`$ are similar over $`.`$
(b) $`\mathrm{\Psi }(A)`$ and $`\mathrm{\Psi }(B)`$ are similar over $``$.
(c) $`\psi (A)`$ and $`\psi (B)`$ are similar over $``$.
Proof. Assume first that $`AB`$ over $``$. Then there is an invertible matrix $`X`$ such that $`X^1AX=B`$. Now applying Theorem 3.2(b) and (f) to its both sides yields $`\mathrm{\Psi }^1(X)\mathrm{\Psi }(A)\mathrm{\Psi }(X)=\mathrm{\Psi }(B)`$, which shows that $`\mathrm{\Psi }(A)\mathrm{\Psi }(B)`$ over $``$. Conversely, assume that $`\mathrm{\Psi }(A)\mathrm{\Psi }(B)`$ over $``$. Then there is an invertible matrix $`Y^{2n\times 2n}`$ such that $`Y^1\mathrm{\Psi }(A)Y=\mathrm{\Psi }(B)`$. For this $`Y`$, by (3.5) there must be an invertible matrix $`X^{n\times n}`$ such that $`\mathrm{\Psi }(X)=Y`$. Thus, $`\mathrm{\Psi }(X^1)\mathrm{\Psi }(A)\mathrm{\Psi }(X)=\mathrm{\Psi }(B)`$, and consequently $`X^1AX=B`$, which shows that $`AB`$. The equivalence of (a) and (c) can also be shown in the same manner. $`\mathrm{}`$
Based on this result, we can extend various results on similarity of complex matrices to complex quaternion matrices.
Theorem 5.2.Let $`A^{n\times n}`$ be given. Then $`A`$ is similar to a diagonal matrix over $``$ if and only if the sizes of the Jordan blocks in the Jordan canonical form of the complex representation $`\psi (A)`$ of $`A`$ are not greater than 2.
Proof. Assume that $`A`$ is diagonalizable over $``$, i.e., there is an invertible matrix $`P`$ over $``$ such that
$$P^1AP=\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _n),$$
$`(5.1)`$
where $`\lambda _1`$$`\lambda _n`$. Then by Theorem 5.1, we obtain
$$\psi ^1(P)\psi (A)\psi (P)=\mathrm{diag}(\psi (\lambda _1),\psi (\lambda _2),\mathrm{},\psi (\lambda _n)),$$
$`(5.2)`$
where $`\psi (\lambda _t)`$ is a $`2\times 2`$ complex matrix. This equality implies that the sizes of Jordan block in the Jordan canonical form of $`\psi (A)`$ are not greater than 2. Conversely assume that there is an invertible matrix $`G`$ over $``$ such that
$$G^1\psi (A)G=\mathrm{diag}(J(\mu _1),\mathrm{},J(\mu _k),\mu _{k+1},\mathrm{},\mu _r),$$
$`(5.3)`$
where $`J(\mu _t)(t=1`$$`k`$) is a $`2\times 2`$ complex Jordan block, $`\mu _{k+1}`$$`\mu _r`$ are complex eigenvalues of $`\psi (A)`$ with $`2k+(rk)=2n`$. In that case, by Eq.(3.7) we know that there is an invertible matrix $`P`$ over $``$ such that $`\varphi (P)=G`$, there is a $`\lambda _t`$ over $``$ such that $`\psi (\lambda _t)=J(\mu _t)(t=1`$$`k`$), and there are $`\lambda _{k+1},\mathrm{},\lambda _r`$ such that
$$\psi (\lambda _{k+1})=\left[\begin{array}{cc}\mu _{k+1}& 0\\ 0& \mu _{k+2}\end{array}\right],\psi (\lambda _{k+2})=\left[\begin{array}{cc}\mu _{k+3}& 0\\ 0& \mu _{k+4}\end{array}\right],\mathrm{},\psi (\lambda _n)=\left[\begin{array}{cc}\mu _{r1}& 0\\ 0& \mu _r\end{array}\right],$$
$`(5.4)`$
In that case, let $`D=\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _n)`$. Then Eq.(5.3) becomes
$$\psi ^1(P)\psi (A)\psi (P)=\psi (D)P^1AP=D.$$
Hence $`A`$ is diagonalizable over $``$. $`\mathrm{}`$.
The result in the above theorem alternatively implies that if the size of a Jordan block in the Jordan canonical form of the complex representation $`\psi _(A)`$ of $`A`$ is greater than 2, then $`A`$ can not be diagonalizable over $``$.
Theorem 5.3.Let $`A^{n\times n}`$ be given. Then $`A`$ is similar to a complex matrix over $``$ if and only if the complex representation $`\mathrm{\Psi }(A)`$ of $`A`$ is similar to a block diagonal matrix $`\mathrm{diag}(J,J)`$ over $`,`$ where $`J^{n\times n}`$.
Proof. If $`AJ^{n\times n}`$, then it follows by Theorem 5.1 and Eq.(3.1) that $`\mathrm{\Psi }(A)\mathrm{\Psi }(J)=\mathrm{diag}(J,J)`$. Conversely if $`\mathrm{\Psi }(A)\mathrm{diag}(J,J)`$ with $`J^{n\times n}`$, then it follows by (3.1) that $`\mathrm{\Psi }(J)=\mathrm{diag}(J,J)`$. Consequently $`AJ`$ follows by Theorem 5.1. $`\mathrm{}`$
6. Determinants of complex quaternion matrices
As one the most fundamental problems in the theory of quaternion matrices, the determinants of complex quaternion matrices, including the determinants of real quaternion matrices, have been considered by lots of authors from different aspects, see, e.g., , , and . Among all of them, a direct and simple method for defining determinants of quaternion matrices is through their representations in the central fields of the corresponding quaternion algebras. In , Zhang presents this kind of definition for determinants of real quaternion matrices, and demonstrates the consistency of his definition with some other classic definitions on determinants based on the products and sums of entries in matrices.
Following the introduction of the universal similarity equalities for generalized quaternion matrices, we shall easily find that it is a most reasonable method to define determinants of quaternion matrices through their representations in the central field of corresponding quaternion algebra. In fact, from Eq.(3.1), we know that for any $`A^{n\times n}`$
$$Q_{2n}\mathrm{diag}(A,A)Q_{2n}^1=\mathrm{\Psi }(A)^{2n\times 2n}.$$
$`(6.1)`$
Therefore if we hope that determinants of quaternion matrices over $``$ satisfy the following two basic properties
$$\mathrm{det}(MN)=\mathrm{det}M\mathrm{det}Nand\mathrm{det}I_n=1$$
for any $`M,N^{n\times n}`$, then as a natural consequence of these two properties, the determinant of the diagonal matrix $`\mathrm{diag}(A,A)`$ satisfies the following equality
$$\mathrm{det}[\mathrm{diag}(A,A)]=\mathrm{det}[Q_{2n}\mathrm{diag}(A,A)Q_{2n}^1]=\mathrm{det}\mathrm{\Psi }(A).$$
$`(6.2)`$
From this equality we see that there seem only two kinds of natural choices for the definition of determinant of $`A^{n\times n}`$, namely, define
$$\mathrm{det}A:=|\mathrm{\Psi }(A)|,$$
$`(6.3)`$
or
$$\mathrm{det}A:=|\mathrm{\Psi }(A)|^{\frac{1}{2}},$$
$`(6.4)`$
where $`|\mathrm{\Psi }(A)|`$ is the ordinary determinant of the complex matrix $`\mathrm{\Psi }(A)`$. For simplicity, we prefer Eq.(6.3) as the definition of determinant of matrix over $``$ and call it the central determinant of matrix $`A`$, and denote it by $`|A|_c`$.
Some basic operation properties on the central determinants of complex quaternion matrices are listed below without proofs.
Theorem 6.1.Let $`A,B^{n\times n},\lambda ,`$ and $`\mu `$ be given.
(a) If $`A^{n\times n},`$ then $`|A|_c=|A|^2.`$
(b) $`A`$ is invertible $``$ $`|A|_c0`$.
(c) $`|AB|_c=|A|_c|B|_c.`$
(d) $`|\lambda A|_c=|A\lambda |_c=\lambda ^{2n}|A|_c.`$
(e) $`|\mu A|_c=|A\mu |_c=n^{2n}(\mu )|A|_c,`$ where $`n(\mu )`$ is the weak norm of $`\mu `$.
(f) $`|A^1|_c=|A|_c^1.`$
(g) $`|A^{}|_c=\overline{|A|}_c.`$
(h) If $`A=\left[\begin{array}{ccc}\hfill a_{11}& \hfill \mathrm{}& \\ & \hfill \mathrm{}& \mathrm{}\\ & & a_{nn}\end{array}\right],`$ then $`|A|_c=n(a_{11})n(a_{22})\mathrm{}n(a_{nn}).`$
(i) If $`A=\left[\begin{array}{cc}A_1& \\ 0& A_2\end{array}\right]`$, where $`A_1`$ and $`A_2`$ are square, then $`|A|_c=|A_1|_c|A_2|_c`$.
(j) If $`AB,`$ then $`|A|_c=|B|_c.`$
By means of the determinants of complex quaternion matrices defined above, we can easily define the central characteristic polynomial of any $`A^{n\times n}`$ as follows
$$p_A(\lambda )=|\lambda I_{2n}\mathrm{\Psi }(A)|,$$
$`(6.5)`$
which is a complex polynomial of degree $`2n`$. From it we easily get the following.
Theorem 6.2(Cayley-Hamilton theorem). Let $`A^{n\times n}`$ be given. Then $`p_A(A)=0`$.
Proof. Since $`p_A(\lambda )=|\lambda I_{2n}\mathrm{\Psi }(A)|`$, so $`p_A[\mathrm{\Psi }(A)]=0.`$ Putting Eq.(6.1) in it and simplifying the equality yields
$$p_A[\mathrm{diag}(A,A)]=\mathrm{diag}(p_A(A),p_A(A))=0,$$
which implies that $`p_A(A)=0`$. $`\mathrm{}`$
7. Conclusions
In the article, we have established a fundamental universal similarity factorization equality over the complex quaternion algebra $``$. This equality clearly reveals the intrinsic relationship between the complex quaternion algebra $``$ and the $`2\times 2`$ total matrix algebra, and could serve as a powerful tool for investigating various problems related to complex quaternions and their applications. In addition to the results in Sections 3—6, we can also apply Eqs.(2.4), (3.1) and (3.6) to investigate some other basic topics related to complex quaternion matrices, such as, singular value decompositions, norms, numerical ranges of complex quaternion matrices and so on. Finally we point out that the equality (2.4) can also extended to the complex Clifford algebra $`_n`$, and a set of perfect matrix representation theory on the complex Clifford algebra $`_n`$ can explicitly established. We shall examine this topic in another paper.
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# The Near Infrared and Multiwavelength Afterglow of GRB 000301c
## 1 Introduction
Infrared observations can be used to improve our understanding of gamma ray burst afterglows in several ways. First, they can be combined with optical measurements to obtain spectral slope measurements with a much wider wavelength baseline, and hence yield a more accurate spectral slope than optical data alone. Second, they can be used to test observed light curve breaks for wavelength dependence, which is an important discriminant between breaks due to ejecta collimation and other possible causes. Finally, it has been suggested that bursts are preferentially located in dusty regions, e.g., under “hypernova” scenarios where bursters are a final evolutionary stage of some class of massive stars. If so, and if the dusty region extends beyond the expected dust destruction distance (Waxman & Draine 1999), then near-infrared observations will detect some afterglows that are obscured at optical wavelengths. At less extreme dust optical depths, near-infrared data help to characterize the host galaxy extinction and so infer both extinction corrected fluxes and properties of dust and gas in high redshift GRB host galaxies.
We present here near-infrared (NIR) photometric observations of the afterglow of the gamma ray burst GRB 000301c. These constitute the best-sampled near infrared light curve for any afterglow to date.
GRB 000301c was detected independently by the All-Sky Monitor on the Rossi X-Ray Timing Explorer and by two spacecraft (Ulysses and NEAR) of the current Interplanetary Network on 2000 March 1.4108 UT. The event was a single peaked GRB lasting approximately 10 seconds (Smith, Hurley, & Cline 2000) at low energies ($`<10\text{keV}`$) and 2 seconds at higher energies ($`>25\text{keV}`$) (Jensen et al 2000). Coordinates were available approximately 36 hours after the burst, and an optical counterpart was reported by Fynbo et al (2000) based on observations at March 3.21 UT. The redshift of the burst was first reported as $`z=1.95\pm 0.1`$ (Smette et al 2000a) and subsequently refined to $`z=2.028\pm 0.025`$ (Smette et al 2000b) using the observed Lyman break in a near-UV spectrum of the GRB afterglow. Weak metal lines in absorption yielded further improvements to $`z=2.0335\pm 0.0003`$ (Castro et al 2000) or $`z=2.0404\pm 0.0008`$ (Jensen et al 2000).
## 2 Photometric Data
## 3 Infrared Telescope Facility Images
We present data obtained at the NASA Infrared Telescope Facility (IRTF) on Mauna Kea as part of a service mode Target of Opportunity program for broadband near-IR followup of gamma ray bursts using the NSFCam imager. Observing conditions were good throughout the period, with clear skies and subarcsecond seeing on all nights. All data were taken with a plate scale of $`0.3^{\prime \prime }`$ per pixel. A log of the observations is given in table 1.
We reduced and analyzed the data following standard near-IR procedures. Raw sky flats were generated for each filter and night using the mean of all available frames after outlier rejection to eliminate the influence of objects on the flatfield. Final flats were generated by subtracting stacked dark frames (taken with the same exposure time, number of coadds, and number of nondestructive reads as the data) from the raw flats and normalizing the result to have mean $`=1`$. Individual frames were then sky-subtracted (using the median of three to six frames taken immediately before and/or after the object frame) and flatfielded. Frames were aligned using bilinear interpolation to implement a simple shift based on the measured centroid of a bright star. Finally, the aligned frames were combined using a clipped mean after removing any sky subtraction residuals through subtraction of the modal pixel value.
We measured the GRB fluxes on all nights using aperture photometry with an aperture diameter of $`1.8^{\prime \prime }`$. All afterglow photometry was taken relative to the bright star $`5.7^{\prime \prime }`$ west and $`1^{\prime \prime }`$ south of the optical transient. (This is star A of Garnavich et al 2000a.) The J and K’ band magnitudes of this star were $`J_{\text{CIT}}=16.64\pm 0.01`$ and $`K^{}=15.97\pm 0.02`$. These were calibrated using observations of UKIRT faint standard 27 (FS27, Casali & Hawarden 1992) taken immediately before the first IRTF observations of GRB 000301c. The reference magnitudes used for FS 27 were $`K^{}=13.14`$ (based on photometric transformation equations from Wainscoat & Cowie 1992) and $`J_{\text{CIT}}=13.45`$ (based on equations from the NSFCam documentation). The standard star observations were processed in the same way as the GRB images, and a larger photometry aperture ($`5.4^{\prime \prime }`$ diameter) was used for flux calibration, to reduce sensitivity to any centering errors or seeing variations between the standard star and GRB frames. No correction was made for atmospheric extinction; however, any such correction would be small ($`0.01\text{mag}`$), because the standard was observed at airmass 1.09, and the GRB field at airmass 1.02. The reference star showed no evidence for variability either in our near-IR data or in optical data from two groups (Garnavich, private communication; Halpern, private communication). The aperture photometry included a local sky estimation and subtraction using the mode of pixel values in an annulus around each point source. This step should control any residual sky level or first order gradient in the sky. Moreover, by selecting the annulus to avoid bright objects and using the mode for sky level estimation, we also control the possible influence of other sources on sky level estimates.
Our counts always remained below the nominal linearity limit for NSFCam (Leggett & Denault 1996); we therefore did not apply nonlinearity corrections to our data. We placed the J band data on the CIT magnitude system using $`J_{\text{CIT}}=J_{\text{MK}}0.01(JK)`$, where $`J_{\text{MK}}`$ is the natural magnitude system for the current Mauna Kea J filter in NSFCam. This gives $`J_{\text{CIT}}=19.11`$ for the afterglow on March 4.652. The K’ magnitudes were left on the NSFCam instrumental system.
Errors due to photon counting statistics were computed based on an iteratively clipped variance of each night’s final stacked images, suitably corrected for the correlated noise introduced by bilinear interpolation. Sky subtraction errors arise only from the difference between sky level in the photometry aperture and in the sky annulus. These errors are separated into statistical errors (due to photon noise in the sky annulus) and systematic (due to objects in the sky annulus or any other source of bias in estimating the true background under the transient source). The statistical part is always small compared to statistical errors from the object flux measurement, due to the large number of pixels in the sky annulus ($`900`$) and much smaller number ($`28`$) in the photometric aperture. Systematic errors in sky subtraction are potentially larger, but we believe they are under reasonable control in our data set because the sky level was removed in two steps (globally, during data reduction, and locally, during aperture photometry) accounting for up to first order gradients in both time and space; and because the weather was cooperative, with relatively little temporal variation in sky background ($`<10\%`$ minimum to maximum in any night and filter, and usually less) and subarcsecond seeing (allowing small photometry apertures). Further possible error sources include centroiding of the afterglow and reference star (expected to be systematic-error limited at the $`0.1`$ pixel level) and residual flatfielding difficulties. Neither will be large compared to our photon counting noise. To be conservative, we estimate that sky subtraction and flatfielding errors combined may affect all of our photometry at up to the $`4\%`$ level. Table 1 lists both pure photon counting errors and error bars including this systematic error added in quadrature.
### 3.1 Data from Other Sources
In addition to the IRTF data, we include in our analysis data from other observatories presented in the literature (both the GRB Coordinate Network Circulars and preprints by Masetti et al 2000 and Jensen et al 2000). These data are summarized in table 2. For data points reported multiple times, we use the most recently reported value. In particular, for the Uttar Pradesh State Observatory data we take values from Masetti et al (2000) rather than Sagar et al (2000).
Two early time K’ data points come from Calar Alto data of Stecklum et al (2000) and Subaru data of Kobayashi et al (2000a,b). Flux values from both were measured relative to the Garnavich et al (2000a) star A. The use of a uniform reference star allows the data from multiple observatories to be compared with reasonable confidence. Residual differences in color terms should be small since all three observatories used the same photometric bandpass. To allow for color terms, we have used $`\sigma =0.03\text{mag}`$ as the effective error on the Subaru data, rather than the photometric error $`0.01\text{mag}`$ reported by Kobayashi et al (2000a,b).
Similarly, the optical data reported in the literature and used in this paper has all been calibrated to either stars A-D of Garnavich et al (2000a) (for some R band data) or to the photometry of Henden et al (2000) (for other optical filters and the rest of the R band data). The R band fluxes measured by Henden et al for stars A-D agree with those from Garnavich et al within the uncertainties, which are $`5\%`$. The mean and median magnitude differences between the two calibrations are $`0.050`$ and $`0.036`$ magnitude respectively, in the sense that Garnavich et al report brighter magnitudes than do Henden et al. The possibility of inconsistent photometric zero points in different data sets therefore amounts to about $`0.05`$ magnitude between authors using calibrations from the two alternative sources, and less for authors using the same calibration. The two largest sets of uniformly reduced optical data currently available are Masetti et al (2000) and Jensen et al (2000). Both groups use the Henden et al calibration. We have increased all R band fluxes reported on the Garnavich et al (2000a) zero point by $`0.04`$ magnitude to adjust all R band photometry to the Henden et al (2000) zero point.
Another possible systematic difference between photometry from different observatories is color term differences. To estimate the importance of these effects, we compare the magnitude difference between the GRB afterglow and a comparison star for thick and thinned CCD chips behind a standard Cousins R filter (Bessell 1990). The quantum efficiency of a thick chip was approximated by $`\text{QE}=0.45+(\lambda /\mu m0.65)`$, and that of a thin chip by $`\text{QE}=0.850.8(\lambda /\mu m0.65)`$. (Both these relations are approximations based on plots from the Steward Observatory CCD Lab web pages and valid for the R band spectral region.) We take the afterglow spectrum to be a power law with $`f_\nu \nu ^1`$ (a characteristic value for this afterglow before any reddening corrections; see below), and take stellar spectral energy distributions from Gunn and Stryker (1983). The result depends on the color of the comparison star. For stars with the colors of unreddened G0 to K4 dwarfs (which includes stars A and D of Garnavich et al 2000a), the difference in inferred magnitude for the two $`\text{QE}(\lambda )`$ models is $`0.01`$ magnitude. For redder or bluer stars, the differences can be substantial: Up to $`0.05`$ magnitude for O stars, $`0.07`$ magnitude for M0 stars (approximately matching Garnavich et al’s stars B and C), and $`0.20`$ magnitude for the coolest M stars. However, only $`9\%`$ of stars in the Henden et al list are redder than stars B and C. Overall, then, color terms will be a modest source of error (no greater than the random photometric error and the zero point error) except for pathological choices of comparison stars. Jensen et al (2000) explicitly address this issue by using comparison stars of color similar to the GRB afterglow in their analysis.
A final possible source of error arises from possible “crowded field” effects in the photometry. Star A is located $`6^{\prime \prime }`$ from the GRB afterglow. With sufficiently poor seeing, aperture photometry of the afterglow might be affected by the wings of this star’s point spread function (PSF). We have estimated the magnitude of this error as a function of aperture size and seeing through numerical calculations with a Gaussian PSF model. We find that photometric errors from this source are negligible ($`0.01\text{mag})`$ for seeing better than $`2^{\prime \prime }`$ (FWHM). They remain $`0.015\text{mag}`$ even for $`3^{\prime \prime }`$ seeing provided that the photometric aperture radius used is $`1.8^{\prime \prime }`$. PSF-fitting photometry, as done by Masetti et al (2000) and Jensen et al (2000), is more robust to such errors.
The ultraviolet flux at $`0.28\mu m\lambda 0.33\mu m`$ is taken from the continuum level in the Hubble Space Telescope STIS spectrum of the afterglow (Smette et al 2000b). This measurement was $`f_\lambda =(7.3_{1.8}^{+0.8}\pm 0.6)\times 10^{18}\text{erg}\text{cm}^2\text{s}^1\text{Å}^1`$. The dominant (first) error bar comes from uncertainties in wavelength calibration coupled with strongly wavelength-dependent sensitivity; the second error bar is a random error. This corresponds to $`f_\nu (0.305\mu m)=(2.26_{0.56}^{+0.25}\pm 0.19)\mu \text{Jy}`$. Over this wavelength range, $`f_\lambda `$ is approximately constant, so $`f_\nu \nu ^2`$. This slope is steeper than the optical-IR slopes derived in section 6 below, but is consistent with a wavelength independent intrinsic slope if intergalactic Lyman line absorption is considered (Madau 1995).
## 4 Light curve fitting
Both the K’ and R light curves are shown in figure 1. The K’ band light curve shows a smooth rollover from an initially constant flux to a rapid decay at late time. The R band light curve shows a qualitatively similar behavior, but the early time R flux decays more steeply than the early K’ flux, and there are additional irregularities in the R band light curve that are arguably significant (e.g., Masetti et al 2000; Sagar et al 2000).
Broken power laws can be empirically fitted by functions of the form $`f=f_0\left[(t/t_b)^{\alpha _1\beta }+(t/t_b)^{\alpha _2\beta }\right]^{1/\beta }`$. With $`\alpha _1<\alpha _2`$ and $`\beta >0`$, this function describes a light curve falling as $`t^{\alpha _1}`$ at $`tt_b`$ and $`t^{\alpha _2}`$ at $`tt_b`$. $`\beta `$ controls the sharpness of the break, with larger $`\beta `$ implying a sharper break. The function that Stanek et al (1999) used to fit the light curve of GRB 990510 is the special case $`\beta =1`$ of this function. Rewritten in magnitudes, the fitting function becomes
$$m=m_b+\frac{2.5}{\beta }\left\{\mathrm{log}_{10}\left[(t/t_b)^{\alpha _1\beta }+(t/t_b)^{\alpha _2\beta }\right]\mathrm{log}_{10}(2)\right\}$$
(1)
where $`m_b`$ is the magnitude at time $`t_b`$.
The K’ light curve can be well fitted by this model. The best fit values (with $`\chi ^2=0.422`$ for one degree of freedom \[ d.o.f\]) are $`t_b=3.57`$ days, $`K_b=17.76`$, $`\alpha _1=0.09`$, $`\alpha _2=2.26`$, and $`\beta =4.23`$.
The R band data show significant departures from such a light curve. The first R band data point is clearly fainter than the extrapolation of the other early R data. If we exclude this point, the best fit to the R data has $`\chi ^2/\text{d.o.f}=2.8`$ for 41 degrees of freedom. This fit has $`t_b=5.1`$ days, $`R_b=21.18`$, $`\alpha _1=0.69`$, $`\alpha _2=2.77`$, and $`\beta =3.2`$. (Including the earliest point gives instead $`\chi ^2/\text{d.o.f}4.7`$, with $`t_b=4.97`$ days, $`R_b=21.15`$, $`\alpha _1=0.56`$, $`\alpha _2=2.81`$, and $`\beta =2.25`$.) The largest deviations from the smoothe light curve (aside from the first data point) are from the bump at $`t3.75`$ days, the plateau at $`t7`$ days, and the steep decline between the two. We believe all three are likely real features. This behavior may indicate variations in the external medium density (see Kumar & Panaitescu 2000b) or refreshed shock effects (Panaitescu, Meszaros, & Rees 1998; Dai & Lu 2000). Either explanation suggests that the K’ data should show similar effects. To test this, we examine the behavior of $`RK^{}`$ below.
If we assume a fixed $`RK^{}`$ color and compare the $`K^{}`$ data to the R band light curve model, the best fit color is $`RK^{}=2.80`$ and $`\chi ^2/\text{d.o.f}=6.6`$ for 5 degrees of freedom (6 data points - 1 color). The same time interval includes about half the R data (31 points), and the data $``$ model residuals in R yield $`\chi ^2/\text{d.o.f}=3.9`$. This relatively large $`\chi ^2`$ is due to the presence of the bump, drop, and plateau in this time interval. In short, the smooth R band fit does not describe the R data very well in this time period, but does an even worse job describing the K’ data.
## 5 $`RK^{}`$ color variations
To assess the significance of the light curve fit differences above, we have examined in detail the $`RK^{}`$ colors of the afterglow for the six $`K^{}`$ data points. At each epoch, we use selected R data points to obtain the most reliable interpolated or extrapolated R flux at the time of the K’ measurement. All measurements are adjusted to the zero point of Henden et al (2000). Details for each epoch follow.
March 3.215: We interpolate between the March 3.21 and March 3.25 Nordic Optical Telescope measurements by Jensen et al (2000), obtaining $`R=20.145\pm 0.04`$. Hence, $`RK^{}=2.62\pm 0.07`$.
March 3.55: We take a weighted average of the three measurements from March 3.51 (Garnavich et al 2000a; Halpern et al 2000a; Veillet et al 2000a,b). We apply corrections to place the first two epochs on the Henden et al zero point, and further add $`0.02`$ magnitudes to adjust the R flux to the $`K^{}`$ epoch, one hour later. (This corresponds to a local $`t^1`$ decay.) We obtain $`R=20.32\pm 0.04`$ and $`RK^{}=2.79\pm 0.05`$.
March 4.64: We combine four R measurements from March 4.48 to 4.50 (three from Jensen et al 2000; one from Halpern et al 2000b, adjusted by 0.05 mag to the Henden et al zero point) and interpolate to the March 4.909 point of Bhargavi & Cowsik to obtain $`R=20.59\pm 0.04`$. Hence, $`RK^{}=2.93\pm 0.06`$.
March 5.61: We extrapolate the March 5.63 R band point from Veillet et al (2000) to the K’ epoch, obtaining $`R=20.85\pm 0.04`$ and $`RK^{}=2.84\pm 0.06`$.
March 6.595: This epoch is during the “plateau” in the R band light curve (Bernabei et al 2000). Fitting a single power law through the R data from March 6.135 to March 7.22, evaluating it at March 6.595, and estimating the flux error through Monte Carlo simulations, we find $`R=21.57\pm 0.07`$. The resulting color is $`RK^{}=3.00\pm 0.14`$. This color does depend on our assessment that the “plateau” is a real feature. If we extend the fit to earlier times (starting at March 5.93), the color becomes $`0.14`$ mag bluer, though a single power law fit does not describe the data well over the full period from March 5.93 to March 7.22.
March 8.590: We use the R data from March 8.146 through 9.52. We again fit a single power law decay through all the R data, evaluate it at the epoch of the K’ measurement, and determine the error bar from simulations. The result is $`R=21.95\pm 0.04`$, so that $`RK^{}=2.66\pm 0.10`$.
These values of $`RK^{}`$ are plotted as a function of time in figure 2. If we assume a single $`RK^{}`$ color throughout the afterglow, the best fit is $`RK^{}=2.79`$. With $`\chi ^2/\text{d.o.f}=2.87`$ for $`5`$ degrees of freedom, this fit is not especially good.
We consider three possible explanations for the observed variation of $`RK^{}`$ with time. First, it is possible that the color variations are real. Second, as suggested by Masetti et al (2000), it may be that there are no true color variations, but that the afterglow exhibits achromatic fluctuations on time scales short compared to the interval between observations. If so, the fluctuations need to have amplitude $`10\%`$ and to occur on timescales $`\delta t3.3\text{hours}=0.04t`$ to explain the March 4.64 data, and $`\delta t7\text{minutes}=0.003t`$ to explain the March 3.21 data. The latter in particular is physically implausible for an afterglow expanding into a realistic external medium. Finally, if the error bars on the photometry are systematically underestimated, the data may be reconciled with no color evolution. In order to reduce $`\chi ^2`$ to $`5`$ (so that $`\chi ^2/\text{dof}=1`$), we need to add a systematic error contribution of $`0.078`$ magnitudes in quadrature to all of the $`K^{}`$ and interpolated $`R`$ fluxes. Given the discussion of systematic errors above and our effort to adjust all R band fluxes to the Henden et al calibration, this large an additional error is unlikely. We therefore believe that at least some of the measured color variation is real.
## 6 Spectral Energy Distributions
The long spectral baseline between the K’ band ($`2.1\mu m`$) and the optical/UV bands allows one to obtain accurate spectral slope measurements. We have calculated the burst spectral energy distribution at selected times based on the availability of multiwavelength data. Where necessary, flux measurements were interpolated between adjacent data points at one wavelength in order to determine a contemporaneous flux with another wavelength. For this operation, we always interpolated the light curves with good sampling (R, K’) to match the time of a sparsely sampled wavelength (UV $`2825\text{Å}`$, B). Photometric zero points for the conversion of magnitudes to flux density units were taken from Fukugita, Shimasaku, & Ichikawa (1995) for the optical, and from Campins, Rieke, & Lebofsky (1985) for the near-IR.
The resulting spectral energy distributions are shown in figure 3. The best-fit spectral slopes assuming unbroken power law spectra are given in table 3, first uncorrected for extinction; then with a correction for foreground Galactic extinction ($`E_{BV}=0.053`$, Schlegel et al 1998) assuming an $`R_VA_V/E_{BV}=3.1`$ extinction law; and finally with an additional correction for possible extinction in the GRB host galaxy (see below). Also included are the slopes derived from R and K’ data alone, corrected only for Galactic extinction.
The broad band spectral energy distributions plotted in figure 3 show no compelling evidence for extinction at the redshift of the GRB. (This is a change from the first preprint of this paper, largely due to revised calibration of the March 6.375 UV data point at $`3050\text{Å}`$.) However, a modest amount of host galaxy extinction remains consistent with the data.
The $`2175`$Å dust absorption feature falls into the observed R band at the redshift $`z=2.03`$ of this burst. Assuming that there is no intrinsic emission feature at this wavelength in the GRB afterglow spectrum, we can place an upper limit of $`A_V0.1`$ magnitude for Milky Way type dust at $`z=2.03`$. Dust with a different reddening law is less strongly constrained. A Small Magellanic Cloud reddening law is plausible for $`A_V0.12`$, while a Large Magellanic Cloud (LMC) law is plausible for $`A_V0.2`$. The measure of plausibility here is that a single power law should approximately fit the spectrum at any given epoch. As an illustrative case, we have shown in figure 3 the spectral energy distributions corrected for $`A_V=0.09`$ of SMC type dust. This correction gives the minimum $`\chi ^2/\text{d.o.f}`$ ($`1.45`$ for 4 degrees of freedom) for the residuals of the plotted spectral energy distributions relative to the spectral slopes reported in table 3. LMC extinction does marginally worse than SMC extinction, while any amount of Milky Way extinction in the host degrades $`\chi ^2`$. We have used the analytic extinction law fitting forms of Pei (1992) in deriving these estimates. Jensen et al (2000) have applied a similar analysis incorporating optical spectra as well as broadband colors. They also find a significantly better fit for SMC extinction than for either MW or LMC extinction, and derive $`A_V=0.14\pm 0.01`$ magnitudes for the SMC model. Using our multiple epoch SEDs, we find a somewhat worse $`\chi ^2/\text{d.o.f}=1.98`$ for their extinction estimate than for ours, but the two results are probably consistent within the errors, especially if there are unidentified systematic errors of $`0.08`$ magnitude in the photometry (see section 5).
The apparent weakness of the $`2175`$Å feature in the extinction curve of the host galaxy is reminiscent of dust attenuation laws for actively star forming galaxies (e.g., the Magellanic Clouds \[Pei 1992 and references therein\] and starburst galaxies \[Gordon, Calzetti, & Witt 1997\]). This may be further circumstantial evidence linking GRBs to actively star forming galaxies. Alternatively, such an extinction law might be observed if GRBs preferentially destroy the small carbonaceous particles thought to be carriers of the $`2175`$ Å feature, but this explanation would only work if much of the dust optical depth arises near the maximum radius where the burst can destroy grains (cf. Waxman & Draine 2000; Fruchter, Krolik, & Rhoads 2000).
In order to determine physical parameters of the afterglow, we need to measure the peak flux density and the locations of breaks in the afterglow spectrum (Wijers & Galama 1999). We now do this (insofar as possible) by combining our optical-IR spectral slope measurements with the submillimeter and radio data. We are looking for four numbers: The frequency $`\nu _{hboxmax}`$ and flux density $`f_{\nu ,\text{max}}`$ at the peak in $`f_\nu `$; the cooling frequency $`\nu _c`$, and the self-absorption frequency $`\nu _{\text{abs}}`$. The spectral slope is expected to be $`p/2`$ for $`\nu >\nu _c`$, $`(p1)/2`$ for $`\nu _{\text{max}}<\nu <\nu _c`$, $`+1/3`$ for $`\nu _{\text{abs}}<\nu <\nu _{\text{max}}`$, and $`+2`$ for $`\nu <\nu _{\text{abs}}`$ (Sari, Piran, & Narayan 1998). Here $`p`$ is the power law index of electrons recently accelerated at the external shock of the expanding GRB remnant.
Extrapolating the optical-IR spectral slopes to lower frequencies, we see that a strong spectral break is required near or above the $`250\text{GHz}`$ measurement by Bertoldi (2000) on March 4.385. The radio data from March 5.67 (Berger et al 2000) are compatible with $`f_\nu \nu ^{1/3}`$ between $`22`$ and $`250\text{GHz}`$, so we determine $`\nu _{\text{max}}`$ and $`f_{\nu ,\text{max}}`$ by extrapolating this behavior until it intersects the extrapolation from optical-IR data (see figure 4). Using fluxes corrected for both Galactic dust and $`A_V=0.09`$ magnitude of SMC type extinction at $`z=2.03`$, we obtain $`\mathrm{log}(\nu _{\text{max}}/\text{Hz})=11.81\pm 0.10`$ and $`\mathrm{log}(f_{\nu ,\text{max}}/\text{mJy})=0.46\pm 0.05`$, where the error bars account only for photometric errors on the data. If we do not apply any correction for host galaxy extinction, we instead obtain $`\mathrm{log}(\nu _{\text{max}}/\text{Hz})=12.16\pm 0.08`$ and $`\mathrm{log}(f_{\nu ,\text{max}}/\text{mJy})=0.58\pm 0.05`$. We can also estimate $`\nu _{\text{abs}}`$ using the radio data from Berger et al (2000) that is shown in figure 4. We conservatively estimate $`3\text{GHz}\nu _{\text{abs}}15\text{GHz}`$ at March 5.66. Frail (personal communication) reports a best fit value of about $`6.8\text{GHz}`$.
The location of the cooling frequency $`\nu _c`$ is harder to constrain, not least because it is a relatively modest break of $`0.5`$ in spectral index and because this burst is not very well observed at X-ray wavelengths. If we take the afterglow behavior to be reasonably described by the “standard” model at the earliest observed times, then the expected behaviors are $`f_\nu \nu ^{(p1)/2}`$ and $`f_\nu t^{3(p1)/4}`$ for $`\nu <\nu _c`$, or $`f_\nu \nu ^{p/2}`$ and $`f_\nu t^{1/23p/4}`$ for $`\nu >\nu +c`$. While this model is based on a uniform ambient medium and spherically symmetric burst, it is also valid for GRBs collimated into an angle $`\zeta `$ at early times, while $`\mathrm{\Gamma }>1/\zeta `$ (Rhoads 1997, 1999). The early time R band data (excluding the discrepant March 2.906 data point) gives $`f_\nu t^{0.71}`$. The spectral slope is $`f_\nu \nu ^{0.9}`$ if we omit host galaxy extinction corrections, but could be as blue as $`\nu ^{0.6}`$ for moderate host galaxy extinction. Both theoretical models of electron acceleration in relativistic shocks (Bednarz & Ostrowski 1998; Gallant, Achterberg, & Kirk 1999) and experience with other afterglows give typical values $`p2.3`$. Comparable values of $`p`$ are marginally consistent with the early time data if we adopt $`A_V0.12`$ in the host galaxy and $`\nu _c>\nu 10^{15}\text{Hz}`$ at $`t=3`$ days. If we require $`p2`$, then $`\nu _c<10^{15}\text{Hz}`$ yields much poorer agreement. Only if we allow $`1.4p1.6`$ can we find a viable model having $`\nu _c`$ between optical and radio wavelengths. However, a new parameter (the upper cutoff in the electron energy spectrum) must be added to the “standard” model to accomodate $`p<2`$ and still keep a finite energy in relativistic electrons. Additionally, it is very difficult to fit the late time decay $`f_\nu t^{2.8}`$ with any $`p<2`$ model; the only viable possibility we presently know of is the “naked afterglow” model of Kumar & Panaitescu (2000b). If we consider the early time K’ light curve, $`f_\nu t^{0.1}`$, and no regime under the standard model offers a self-consistent estimate of $`p`$. Overall, we regard $`p2`$ and $`\nu _c10^{15}\text{Hz}`$ around UT 2000 March 3.5 as the most plausible solution.
## 7 Physical Parameters of the Afterglow
We can use the observed spectral energy distribution of GRB 000301C to place interesting limits on some of the afterglow’s physical parameters, using the method of Wijers & Galama (1999). To do so, we assume that the afterglow is reasonably approximated by a spherical burst (or a section thereof) expanding into a uniform ambient medium up to the time of our SED measurement. Even if the light curve break is attributed to a jet, the method should work for times before the break. We take the spectrum to peak at either of the values derived above ($`\mathrm{log}(\nu _{\text{max}}/\text{Hz})=11.81`$ and $`f_{\nu ,\text{max}}=2.9\text{mJy}`$ with the host galaxy extinction correction, or $`\mathrm{log}(\nu _{\text{max}}/\text{Hz})=12.16`$ and $`f_{\nu ,\text{max}}=3.8\text{mJy}`$ with only Milky Way extinction corrections). In addition, we take $`\mathrm{log}(\nu _c/\text{Hz})15`$. Using $`\nu _{\text{abs}}=6.8\text{GHz}`$ then gives “best estimate” limits of $`E>3\times 10^{53}\text{erg}\times \mathrm{\Omega }/(4\pi )`$, $`\xi _e0.11`$, $`\xi _B1.6\times 10^3`$, and $`n1cm^3`$. Here $`E`$ is the kinetic energy of the ejecta, $`\xi _e`$ and $`\xi _B`$ are the fractions of the local energy that go into relativistic electrons and magnetic fields immediately behind the expanding GRB remnant blast wave, and $`n`$ is the number density of the ambient medium. We obtain limits rather than measurements because of the relatively weak constraints on $`\nu _c`$ and $`\nu _{\text{abs}}`$. We have taken the most conservative pair of ($`\nu _{\text{max}}`$, $`f_{\nu ,\text{max}}`$) in deriving these limits on physical quantities. The physical quantities scale with $`\nu _{\text{abs}}`$ as $`E\nu _{\text{abs}}^{5/6}`$, $`\xi _e\nu _{\text{abs}}^{5/6}`$, $`\xi _B\nu _{\text{abs}}^{5/2}`$, and $`n\nu _{\text{abs}}^{25/6}`$ (Wijers & Galama 1999). If we allow $`3\text{GHz}\nu _{\text{abs}}15\text{GHz}`$, and further take the one-sigma variations on $`\nu _{\text{max}}`$ and $`f_{\nu ,\text{max}}`$ to weaken the limits, we obtain more conservative limits of $`E>1.3\times 10^{53}\text{erg}\times \mathrm{\Omega }/(4\pi )`$, $`\xi _e0.04`$, $`\xi _B0.02`$, and $`n0.014cm^3`$. We have used $`p=2.3`$, but other values affect the results only weakly except for $`\xi _e`$, whose behavior approaches $`(p2)^1`$ as $`p2`$. If the ill-constrained cooling frequency is substantially above $`10^{15}\text{Hz}`$ at 3 days, then the bounds on $`E`$ and $`\xi _e`$ rise as $`\nu _c^{1/4}`$, $`\xi _B`$ falls as $`\nu _c^{5/4}`$, and $`n`$ rises as $`\nu _c^{3/4}`$ (Wijers & Galama 1999). On the other hand, $`\nu _c10^{15}\text{Hz}`$ requires $`p1.5`$. In this case the equations describing the afterglow and the inversion of measured quantities to obtain physical parameters would require new formulae.
It is interesting to compare our estimate of $`E`$ with the gamma ray fluence of the burst. Jensen et al (2000) estimate a fluence of $`2.1\times 10^6\text{erg}\text{cm}^2`$ in the $`25`$ to $`100\text{keV}`$ band and a corresponding energy of $`2.3\times 10^{52}\text{erg}\times \mathrm{\Omega }/(4\pi )`$. They also report a similar fluence in the $`150`$ to $`1000\text{keV}`$ band, so the total energy might be $`5\times 10^{52}\text{erg}\times \mathrm{\Omega }/(4\pi )`$. Thus, our most conservative limits $`E>1.3\times 10^{53}\text{erg}\times \mathrm{\Omega }/(4\pi )`$ implies that $`0.4`$ of the blast wave energy was emitted in gamma rays, while for our best estimate $`E>3\times 10^{53}\text{erg}\times \mathrm{\Omega }/(4\pi )`$, this fraction is reduced to $`0.15`$.
## 8 Discussion
GRB 000301c is the third burst for which a strong break in the light curve is clearly observed. Several classes of breaks are predicted by fireball models. The most basic of these are due to features in the synchrotron spectrum moving through the observed bandpass (e.g., Paczyński & Rhoads 1993; Sari, Piran, & Narayan 1998). However, this class of features predicts relatively modest changes in light curve slope, with the break occurring first at short wavelengths and evolving to longer ones. Jetlike burst ejecta, on the other hand, are expected to give strong breaks that are essentially independent of wavelength (Rhoads 1997, 1999; Sari, Piran, & Halpern 1999), and the observed breaks have generally been interpreted as evidence for collimation of the GRB ejecta (e.g., in GRB 990123, by Castro-Tirado et al 1999, Kulkarni et al 1999, Fruchter et al 1999, and Galama et al 1999; and in GRB 990510, by Stanek et al 1999 and Harrison et al 1999). A difficulty with this model is that the predicted break is quite gradual (Rhoads 1999; Panaitescu & Meszaros 1999; Moderski, Sikora, & Bulik 2000; Kumar & Panaitescu 2000a), while observed breaks are rather sharp.
The current burst is no exception. The prediction of Rhoads (1999) for the light curve around the break time for a collimated jet is
$$f_\nu \frac{[t/t_b]^{3(p1)/4}}{\left(1+3.72[t/t_b]^{2/5}\right)^{5/2}\left(1+2.07[t/t_b]^{5/12}\right)^{3(p1)/5}}$$
(2)
where $`t_b`$ is the fiducial break time defined in Rhoads (1999). The break in this predicted light curve is extremely broad, and would give $`\chi ^2`$ little better than a single power law in fitting the observed break in either $`K^{}`$ or $`R`$ band. The model curve is based on numerical integration of the remnant’s dynamical equations, and ignores differences in light travel time between the center and edge of the remnant, which will only smooth the break further (e.g., Moderski et al 2000; Panaitescu & Meszaros 1999).
If we ignore the issue of break sharpness and fit a collimated jet model to the observed R band light curve, we can infer the opening angle from the measured break time. To do so, we need a reasonable measurement of $`t_b`$ and crude estimates of $`\mathrm{\Omega }/E`$ and $`n`$, since the inferred opening angle scales as $`(t_b^3n\mathrm{\Omega }/E)^{1/8}`$ (Rhoads 1999). We use $`t_b=5.1`$ days and $`E=3\times 10^{53}\text{erg}\times \mathrm{\Omega }/(4\pi )`$. To estimate $`n`$ more precisely, we use the column density $`N(HI)10^{21.2\pm 0.5}`$ inferred from Lyman $`\alpha `$ absorption (Jensen et al 2000) and estimate the linear size of the source as $`0.2^{\prime \prime }`$ based on its nondetection in late HST images. This implies a number density $`n0.4`$, consistent with our earlier estimate. Using $`n1`$, the inferred opening angle becomes $`2.5^{}`$, or $`10^3`$ of the sky if the jet is bipolar.
The transition to the nonrelativistic regime has been proposed as another mechanism for light curve breaks both in this burst (Dai & Lu 2000) and others (Dai & Lu 1999). However, we do not know of a detailed calculation of the sharpness of this break, making a fair evaluation of this possibility difficult. Light travel time effects seem likely to broaden this feature to $`\delta t/t1`$, as with most other features.
A final possible cause for sharp breaks in GRB afterglow light curves is discontinuities in the ambient density distribution. Assuming that the density is a function of radius alone, a minimum timescale for breaks due to such discontinuities is $`\mathrm{\Delta }tt`$, where $`t`$ is the time elapsed in the observer’s frame since the burst and $`\mathrm{\Delta }t`$ the characteristic duration of a light curve feature. This duration is set by differential light travel time effects between material moving along the line of sight and off-axis material moving in direction $`1/\mathrm{\Gamma }`$, and is a rough minimum for any afterglow light curve feature provided the ambient medium density is approximately independent of angle from the line of sight. The time required for material already in the expanding blast wave to cool by adiabatic losses is also relevant for determining the sharpness of a density discontinuity break, as is the emission from angles $`>1/\mathrm{\Gamma }`$ off the line of sight (Kumar & Panaitescu 2000b).
If we believe that the observed $`RK^{}`$ variations are real, then the greatest difficulty posed by the observations of GRB 000301c is in finding a model whose light curve steepens in $`K^{}`$ before it steepens in R. For most mechanisms, breaks will occur either first at short wavelengths (e.g. the cooling break), or simultaneously at all wavelengths (e.g. “beaming” breaks). One speculative way out is to suppose that a discontinuity in the ambient density is encountered while the cooling break is between the R and K’ filters. The predicted appearance of an afterglow at frequencies above and below this break is expected to differ qualitatively: A high frequency image would show an annular structure and a low frequency image a more nearly filled disk. This is caused by the difference in the apparent dynamical age of the remnant along the line of sight (where we see things changing quickly) and near the edge of the observed afterglow (where light travel time is larger, and we see material at an earlier and hotter stage of its evolution). (E.g., Granot, Piran, & Sari 1999.) Now, the same variation of “lookback time” from the center to edge of the afterglow implies that we see the effect of an ambient density drop first in the middle of the afterglow, and that the fractional effect of such a discontinuity on the afterglow flux will initially be larger at long wavelengths than short ones. This mechanism can reproduce the sign of the observed effect. Detailed calculations would be necessary to see if it can approach the observed magnitude, given that the two filters are separated by only a factor of $`3`$ in wavelength.
## 9 Summary
We present a K’ band ($`2.1\mu m`$) light curve of the GRB 000301c afterglow, combining four epochs from a Target of Opportunity program at the NASA Infrared Telescope Facility with two additional measurements from Calar Alto (Stecklum et al 2000) and Subaru (Kobayashi et al). This light curve can be well fitted by a broken power law evolution, with a very flat early time slope ($`t^{0.1}`$), a steep late time slope ($`t^{2.3}`$), and a rather sharp break. A similar fit for R band ($`0.7\mu m`$) data from the literature yields steeper slopes ($`0.7`$ and $`2.8`$) and a later break time. A detailed analysis of the $`RK^{}`$ colors shows modest deviations from constant color, which are significant at the $`2\sigma `$ level. These may indicate the presence of systematic errors $`0.08`$ magnitude due to inhomogeneous data sets. However, we believe that the likely level of such systematic errors is $`0.05`$ magnitude, and that at least part of the color variability is likely real.
The strong break and steep late time slope in the light curves are reminiscent of GRB 990510, which has been interpreted as a collimated burst (Stanek et al 1999; Harrison et al 1999; Kumar & Panaitescu 2000a). If we fit a collimated burst model to the R band light curve, the estimated opening angle becomes $`2.5^{}`$.
However, while a jet model can reproduce the size of the break, it does not provide a natural explanation for the observed rapidity of the break. Nor, however, do other models. Indeed, there is considerable variability in the R band light curve of GRB 000301C on time scales $`\delta t/t0.4`$. This highlights the potential for unusual temporal variability in GRB afterglows.
Fitting a standard synchrotron spectral energy distribution to the burst, we place the peak of $`f_\nu `$ around $`3.4\text{mJy}`$ at about $`10^{12}\text{Hz}`$ ($`300\mu m`$). The dominant uncertainty in this measurement is the correction for extinction in the GRB host galaxy. Random errors are $`0.10`$ dex in $`\nu _{\text{max}}`$ and $`0.05`$ dex in $`f_{\nu ,\text{max}}`$, while the correction for host galaxy extinction is uncertain at perhaps twice this level. These strong constraints on the $`f_\nu `$ peak are possible because of the precise optical-IR spectral slope measurements afforded by K’ observations.
Combining the measured spectral peak with constraints on the cooling frequency and self-absorption frequency, we infer that the blast wave energy required to power this afterglow was $`E>3\times 10^{53}\text{erg}`$ if isotropic. The corresponding efficiency of gamma ray production in the burst was $`0.15`$.
It is a pleasure to thank the NASA IRTF director and staff for making this observing program possible. Special thanks are due to Bob Joseph, for arranging the target of opportunity mechanisms; Bill Vacca, for help with observing strategies; and Bill Golisch, Paul Fukumura-Sawada, and Dave Griep for observing. We also thank Sylvio Klose, Dale Frail, and Naoto Kobayashi for useful communications. JER’s work is supported by an Institute Fellowship at STScI.
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# Untitled Document
SUPERHEATED DROP AS A NEUTRON SPECTROMETER
Mala Das, B. K. Chatterjee, B. Roy<sup>1</sup><sup>1</sup>1Author for communication : Tel.no.-350 2402/03, Ext.-305, Fax no.- 91 033 350 6790, e-mail-biva@boseinst.ernet.in and S. C. Roy
Department of Physics, Bose Institute
93 / 1 A. P. C. Road, Calcutta 700009, India
## Abstract
Superheated drops are known to vaporise when exposed to energetic nuclear radiation since the discovery of bubble chamber. The application of superheated drops in neutron research specially in neutron dosimetry is a subject of intense research for quite sometime. As the degree of superheat increases in a given liquid, less and less energetic neutrons are required to cause nucleation. This property of superheated liquid makes it possible to use it as a neutron spectrometer. Neutron detection efficiency of superheated drops made of R12 exposed to Am-Be neutron source have been measured over a wide range of temperature -17<sup>o</sup>C to 60<sup>o</sup>C and the results have been utilized to construct the energy spectrum of the neutron source. This paper demonstrates that a suitable neutron spectrometer may be constructed by using a single liquid and varying the temperature of the liquid suitably at a closer grid.
PACS No. : 29.30 Hs, 29.40.-n
Keywards : SDD, R12, temperature dependence, neutron, efficiency, spectrometry.
1. INTRODUCTION
A fluid kept in the liquid state above its boiling temperature is called superheated. Application of superheated liquid to detect ionizing radiation is well known from the times of bubble chamber. The resurgence of its use has been observed from the late seventies and the investigations on the subjects in the last two decades turned into an almost maturing technology especially to detect neutrons and more recently to detect photons. The suitability of using superheated drops as a neutron dosimeter has already been established. The superheated drops are now commercially available as neutron and photon dosimeter in the trade name superheated drop detector (SDD) <sup>2</sup><sup>2</sup>2Superheated Drop Detector is the registered trademark of Apfel Enterprises Inc, NewHaven, CT, USA and bubble detector (BD). <sup>3</sup><sup>3</sup>3Bubble Detector is the registered trademark of Bubble Technology Industries Ltd. In a superheated liquid minimum energy (threshold) required to nucleate decreases as the degree of superheat increases. The degree of superheat of a liquid could be defined simply by the difference of ambient temperature above the boiling point of the liquid. Therefore, liquids with lower boiling points possess higher degree of superheats at a given ambient temperature (above their boiling points) and as the ambient temperature increases the given liquid becomes more and more superheated. This property of the superheated liquid are being utilised to develop neutron dosimetry and neutron spectrometry. There are two distinct types of methodologies used in developing neutron spectrometer. In one, a collection of superheated samples made of liquids with different boiling points (i.e. with different threshold neutron energies) are utilised, while in the other, two liquids are chosen and the temperature of the liquids are varied at four different temperatures to obtain eight sets of threshold energies (equivalent to eight different samples with different boiling points). The temperature variation method is superior than using samples with different boiling points. By controlling the temperature of the sample one can, in principle, change the threshold neutron energy at any desired level (equivalent to using ’finer’ windows to scan the spectrum), while in the other method one is limited by the availability of liquids with lower boiling points (equivalent to using ’coarser’ windows). In the present work we measured the detection efficiency of a single sample made of R-12(Dichlorodifluoromethane : C$`Cl_2`$$`F_2`$), which is known to be sensitive to neutrons of energies from thermal to tens of MeVs, by (almost) continuously changing temperatures over a wide range. The response of the sample to Am-Be neutrons have been measured at about thirty different temperatures in the range which is equivalent to using thirty different samples with thirty different boiling points (the experiment has been actually performed more than 50 different temperatures in the range -17$`{}_{}{}^{o}C`$ to about 60$`{}_{}{}^{o}C`$ but it has been observed that the sample started responding from about 0.5$`{}_{}{}^{o}C`$ to Am-Be neutrons). In addition to the advantage of using single liquid, controlling the temperature enables one to scan the energy spectrum by finer ’windows’ thereby improving the inherent energy resolution of the spectrometer when compared with other such spectrometers using superheated liquid.
2. PRINCIPLE OF OPERATION
The superheated state of the liquid is a metastable state and the nucleation in this state can be initiated by the presence of heterogeneous nucleation sites such as air bubbles, solid impurites, gas pockets etc. or by radiation interactions. The nucleation in superheated state starts with the formation of a critical sized vapour embroy. The free energy required to form a spherical vapour bubble of radius r in a liquid is given by
$$G=4\pi r^2\gamma \left(T\right)\frac{4}{3}\pi r^3\left(p_vp_o\right)$$
(1)
where $`\gamma (T)`$ is the liquid-vapour interfacial tension, $`p_v`$ is vapour pressure of the superheated liquid and $`p_o`$ is the ambient pressure. The difference $`p_v`$ -$`p_o`$ is called the degree of superheat of a given liquid. One can see from equation (1) that G is maximum at
$$r=2\gamma \left(T\right)/\left(p_vp_o\right)=r_c$$
(2)
where $`r_c`$ is called the critical radius. When a bubble grows to the size of the critical radius it becomes thermodynamically unstable and grows very fast till the entire liquid droplet vaporises.
The minimum amount of energy (W) needed to form a vapour bubble of critical size $`r_c`$ as given by Gibbs from reversible thermodynamics is
$$W=16\pi \gamma {}_{}{}^{3}\left(T\right)/3(p_vp_o)^2$$
(3)
where $`\gamma \left(T\right)`$ = C $`\left(T_cTd\right)`$ with $`T_c`$ the critical temperature of the liquid and C and d are constants .
With increase in temperature, since the degree of superheat ($`p_v`$-$`p_o`$) increases and $`\gamma (T)`$ decreases, the minimum energy ($`W`$) required for vapour bubble nucleation will be less. The variation of $`W`$ with temperature for superheated drops of R-12 is shown in figure 1. Therefore $`W`$, the threshold energy for nucleation depends on the type and the temperature of the liquid. When a neutron of energy $`E_n`$ interacts with a nucleus of atomic weight A, the maximum energy that can be transferred to the nucleus from the neutron is through the elastic head on collision and is given by,
$$E_i=4AE_n/\left(A+1\right)^2$$
(4)
After receiving the energy, the nucleus is scattered from its atom and moves through the liquid losing its energy through Coulombic interaction until it comes to rest. For a given neutron energy, different nuclei of the liquid will receive different amount of energy, depending on their atomic weight. The ion with the highest value of linear energy transfer (LET) or ($`dE/dx`$) in the liquid, will play the major role in vapour nucleation. The energy deposited along that part of the ion’s path (L) corresponding to about twice the critical radius contributes significantly to bubble formation. For nucleation to occur this deposited must exceed $`W`$ the minimum energy required for bubble formation.
Usually most of the energy is lost into heat and a very small fraction of the deposited energy is utilised in nucleation and W/$`E_c`$ is called the thermodynamic efficiency ($`\eta _T`$ of nucleation .
$$W=kr_cdE/dx=2\eta _Tr_cdE/dx$$
(5)
where $`k`$ (constant) equals twice the thermodynamic efficiency ($`\eta `$). Hence, in the equation
$$W/r_c=kdE/dx$$
(6)
relates the threshold energy (corresponding to the $`dE/dx`$) for nucleation to the ambient temperature (corresponding to $`W/r_c`$). This enables us to convert the temperature scale of superheated drops to the (threshold) energy scale of incident neutrons. Therefore by varying the ambient temperature of the superheated drops, one can observe the variation of the SDD response at different neutron energies which has been used in neutron spectrometry.
3. EXPERIMENT
The experiment was performed with superheated drops of R12 at different temperatures by using the volumetric method, described by Das et al.. The vial containing the sample was connected to a graduated horizontal glass tube with a small coloured water column as marker. The vapourization of a liquid drop displaces the water column by the distance corresponding to the volume of vapour formed. The details of the preparation of the sample is given elsewhere.
If neutrons of flux $`\psi `$ are incident on superheated drops of volume $`V`$, liquid of density$`\rho _L`$ and molecular weight $`M`$ the vaporization rate is given by
$$\frac{dV}{dt}=V\psi \frac{N_A\rho _L}{M}\eta dn_i\sigma _i$$
(7)
| where | N<sub>A</sub> | = | Avogadro Number |
| --- | --- | --- | --- |
| | d | = | average droplet volume |
| | n<sub>i</sub> | = | weight factor of the ith element in the molecule whose |
| | | | neutron nucleus elastic scattering cross section is $`\sigma _i`$ |
| | $`\eta `$ | = | efficiency of neutron detection. |
Due to nucleation by neutrons, the displacement of the water column along a horizontal glass tube was measured as a function of time. The procedure of calculating $`\eta `$ from the measured displacement of water column has been explained in detail in one of our recent publications.
The temperature of the sample was controlled by an indigeneously made temperature controller. For low temperature measurements the sample was placed in an alcohol bath sitting on the top of a cold finger dipped in liquid nitrogen. The upper part of the finger was wrapped with heating tape and by applying different voltages to the tape, different steady temperatures of the bath can be achieved. For measurements of higher temperatures the same setup was used without the liquid nitrogen. The fluctuation of temperature in these measurements was found to be within $`\pm 0.1^oC`$.
In the experiment the ambient temperature of the sample was increased slowly from low to any desired higher value and the nucleation rate was measured at each temperature. The measurement was performed from in a temperature range of -17$`{}_{}{}^{o}C`$ to about 60$`{}_{}{}^{o}C`$ in a close grid necessary to obtain the energy spectrum of the source from temperature. The nucleation due to background radiation and due to other fluctuations has been subtracted. The liquid was observed to become unstable due to spontaneous nucleatation at about 60$`{}_{}{}^{o}C`$.
4. EXPERIMENTAL RESULTS AND COMPUTATION OF THE NEUTRON ENERGY
The variation of neutron detection efficiency ($`\eta `$) for R12 with temperature in presence of neutrons from Am-Be neutron source is presented in figure 2. The solid line in figure 2 is the spline smoothing of the efficiency data at different temperatures. The uncertainties presented in the figure are the total experimental uncertainties of estimating $`\eta `$. The derivative of efficiency, d$`\eta `$/dT against temperature is shown in figure 3. Now one has to estimate the equivalence of the energy of the detected neutrons with the temperature of the SDD sample. One way to do it is to expose the sample at different temperatures to different monoenergetic neutron sources and to note the threshold neutron energies for nucleation . A novel approach has been used in this work.
As has been presented in Section 2, the different nuclei of the superheated liquid would receive different amount of energy and they must have different dE/dx. In case of R-12 containing C, Cl and F, dE/dx of these ions with different neutron energy are presented in Figure 4. From Figure 4, it is clear that the dE/dx values of C, Cl and F are comparable in the neutron energy of our interest and we take the average value of dE/dx of all the ions using the equation below
$$\left(dE/dx\right)_{average}=n_i\sigma _i\left(dE/dx\right)_i/n_i\sigma _i$$
(8)
where $`n_i`$ is the number of ions of the i-th element in a R12 molecule, $`\sigma _i`$ is the neutron-nucleus elastic scattering cross section and $`(dE/dx)_i`$ is the LET of the i-th ion. The variation of average dE/dx as a function of neutron energy is shown in figure 5.
From equations (3) and (5), we obtain
$$(dE/dx)=8\pi \gamma {}_{}{}^{2}\left(T\right)/3k(p_vp_o)$$
(9)
From the equation above the dE/dx has been plotted against temperature for different arbitrary values of $`k`$ of which only four such plots ($`k`$=1, 0.1, 0.05, 0.0195) are presented in figure 6 using the equivalence between the temperature of the sample and the incident neutron energy from figure 5. The variation of threshold neutron energy for nucleation with temperature of the sample for different $`k`$ has been studied of which four such variations are shown in figure 7. With the optimum value of $`k`$ the temperature axis of the figure 3 has been converted to neutron energy and the resulting spectrum was fitted with the peak neutron energy of the <sup>241</sup>Am-Be neutron spectrum. The best fit is obtained for k equals 0.0195. The analysis has been performed upto a maximum temperature of 42.5$`{}_{}{}^{o}C`$. The final neutron energy spectrum of <sup>241</sup>Am-Be source obtained from our analysis is shown in figure 8.
If L = 2$`r_c`$ is taken as the distance in the ion’s path which contributes significantly in nucleation of superheated drops in calculating the thermodunamic efficeincy of nucleation, our experimental analysis produces the value close to 0.01. It may be noted here that Apfel et al. obtained this value ranging from 0.03 to 0.05.
5. DISCUSSION
The result shows that $`\eta `$ increases with temperature. At low temperature, the threshold energy for nucleation (W) is high, which indicates that a larger energy is required to cause nucleation. According to equation 3 as temperature increases, W decreases and more and more neutrons from the low energy range of the spectrum are taking part in nucleation. So $`\eta `$ increases with temperature. The sharp increase of $`\eta `$ near 25<sup>o</sup>C corresponds to detection of neutrons with energies ranging from highest available to those at the peak of the spectrum. At high temperature when all the neutrons of the spectrum contribute in nucleation, $`\eta `$ should be constant with temperature. But at about 45<sup>o</sup>C, $`\eta `$ increases again. We suspected that the sample becomes sensitive to gamma rays coming out of the Am-Be source. In a separate experiment, we indeed observed that R-12 becomes sensitive to gamma rays at about 45$`{}_{}{}^{o}C`$. Since, as mentioned before, the analysis has been performed upto a maximum temperature of 42.5$`{}_{}{}^{o}C`$, the contribution due to gamma rays is absent in the measured energy spectrum in this work.
Figure 3 shows that the d$`\eta `$/dT vs. T graph resembles the neutron energy spectrum of <sup>241</sup>Am-Be where the second peak corresponds to the gamma sensitivity of the sample. The ambient temperature of the superheated drops was converted to the energy of the neutrons following the method described in Section 4. So by using superheated drops at different temperatures, it is possible to obtain the neutron energy spectrum. This indicates important use of SDD in neutron spectrometry. The maximum uncertianty in neutron energy as could be found from figure 7 is within 5$`\%`$ in the entire region of our investigation. This method can be used to determine any other neutron energy outside the present range, only then one has to consider the ion with maximum dE/dx and the rest of the analysis is same as this. The present study also helps to select the suitable material (liquid) for a given neutron energy spectrum.
It has been observed in this experiment that the nucleation rate of superheated drops rapidly changes for samples exposed to thermal shock compared to samples whose temperature was changed slowly. The liquid appeared to be more fragile when temperature was changed rapidly. Though this is not quite unexpected in the exact physics of this phenomenon, why the liquid becomes more fragile under heat shock, requires further investigation.
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19. Mala Das,B. Roy, B. K. Chatterjee and S. C. Roy, Rad.Meas. 30, 35 (1999).
20. B. Roy, B. K. Chatterjee and S. C. Roy, Rad. Meas. 29, 173 (1998).
FIGURE CAPTIONS
Fig. 1: Variation of threshold energy (W) required for nucleation in R12 as a function of temperature (T).
Fig. 2: Observed variation of neutron detection efficiency ($`\eta `$) as a function of temperature (T) in R12.
Fig. 3: Variation of the derivative of neutron detection efficiency (d$`\eta `$/dT) as a function of temperature (T).
Fig. 4: Variation of stopping power (dE/dx) of different ions (C, Cl, F) in Freon-12 as a function of neutron energy.
Fig. 5: Variation of average stopping power $`(dE/dx)_{average}`$ over three different ions in Freon-12 as a function of neutron energy.
Fig. 6: Variation of stopping power (dE/dx) of ion in R12 as a function of temperature (T) of the sample, for different arbitrary values of $`k`$.
Fig. 7: Variation of neutron energy as a function of temperature (T) of the sample, for different $`k`$.
Fig. 8: The neutron energy spectrum of <sup>241</sup>Am-Be obtained from the experiment.
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# limit-from𝐽-holomorphic Curves And Periodic Reeb Orbits Project 19871044 Supported by NSF
## 1 Introduction and results
A contact structure on a manifold is a field of a tangent hyperplanes (contact hyperplanes) that is nondegenerate at any point. Locally such a field is defined as the field of zeros of a $`1`$form $`\lambda `$, called a contact form. The nondegeneracy condition is $`d\lambda `$ is nondegenerate on the hyperplanes on which $`\lambda `$ vanishes; equivalently, in $`(2n1)`$space:
$$\lambda (d\lambda )^{n1}0$$
The important example of contact manifold is the well-known projective cotangent bundles definded as follows:
Let $`N=T^{}M`$ be the cotangent bundle of the smooth connected compact manifold $`M`$. $`N`$ carries a canonical symplectic structure $`\omega =d\lambda `$ where $`\lambda =_{i=1}^ny_idx_i`$ is the Liouville form on $`N`$, see . Let $`P=PT^{}M`$ be the oriented projective cotangent bundle of $`M`$, i.e. $`P=_{xM}PT_x^{}M`$. It is well known that $`P`$ carries a canonical contact structure induced by the Liouville form and the projection $`\pi :T^{}MPT^{}M`$.
Let $`(\mathrm{\Sigma },\lambda )`$ be a smooth closed oriented manifold of dimension $`2n1`$ with a contact form $`\lambda `$. Associated to $`\lambda `$ there are two important structures. First of all the so-called Reed vectorfield $`X_\lambda `$ defined by
$$i_X\lambda 1,i_Xd\lambda 0$$
and secondly the contact structure $`\xi =\xi _\lambda \mathrm{\Sigma }`$ given by
$$\xi _\lambda =\mathrm{ker}(\lambda )T\mathrm{\Sigma }$$
by a result of Gray, , the contact structure is very stable. In fact, if $`(\lambda _t)_{t[0,1]}`$ is a smooth arc of contact forms inducing the arc of contact structures $`(\xi _t)_{t[0,1]}`$, there exists a smooth arc $`(\psi _t)_{t[0,1]}`$ of diffeomorphisms with $`\psi _0=Id`$, such that
$$T\mathrm{\Psi }_t(\xi _0)=\xi _t$$
(1.1)
here it is important that $`\mathrm{\Sigma }`$ is compact. From (1.1) and the fact that $`\mathrm{\Psi }_0=Id`$ it follows immediately that there exists a smooth family of maps $`[0,1]\times \mathrm{\Sigma }(0,\mathrm{}):(t,m)f_t(m)`$ such that
$$\mathrm{\Psi }_t^{}\lambda _t=f_t\lambda _0$$
(1.2)
In contrast to the contact structure the dynamics of the Reeb vectorfield changes drastically under small perturbation and in general the flows associated to $`X_t`$ and $`X_s`$ for $`ts`$ will not be conjugated, see.
Let $`M`$ be a Riemann manifold with Riemann metric, then it is well known that there exists a canonical contact structure in the unit sphere of its tangent bundle and the motion of geodesic line lifts to a geodesic flow on the unit sphere bundles. Therefore the closed orbit of geodesic flow or Reeb flow on the sphere bundle projects to a closed geodesics in the Riemann manifolds, conversely the closed geodesic orbit lifts to a closed Reeb orbit. The classical work of Ljusternik and Fet states that every simply connected Riemannian manifold has at least one closed geodesics, this with the Cartan and Hadamard’s results on non-simply closed Riemann manifold implies that any closed Riemann manifolds has a closed geodesics, i.e., the sphere bundle of a closed Riemann manifold with standard contact form carries at least one closed Reeb orbits which is a lift of closed geodesics of base manifold. Its proof depends on the classical minimax principle of Ljusternik and Schnirelman or minimalization of Hadamard and Cartan,. In sympletic geometry, Gromov introduces the global methods to proves the existences of symplectic fixed points or periodic orbits which depends on the nonlinear Fredholm alternative of $`J`$holomorphic curves in the symplectic manifolds. In this paper we use the $`J`$holomorphic curve’s method to prove
###### Theorem 1.1
Every closed contact manifold $`\mathrm{\Sigma }`$ with contact form $`\lambda `$ carries at least one closed orbit.
Theorem1.1 was conjectured by Weinstein in under the assumption $`H^1(\mathrm{\Sigma })=0`$. Note that Viterbo first proved the above result for any contact manifolds $`\mathrm{\Sigma }`$ of restricted type in $`R^{2n}`$ after Rabinowitz and Weinstein . After Viterbo’s work many results were obtained in etc by using variational method or Gromov’s nonlinear Fredholm alternative, see survey paper . Through $`J`$holomorphic curves, especially, Hofer in proved the following first striking results.
###### Corollary 1.1
(Hofer)The three dimentional sphere with any contact form carries at least one closed Reeb orbit.
###### Theorem 1.2
(Ljusternik-Fet) Every simply connected closed Riemannian manifold has at least one closed geodesics.
Therefore we get a new proof on the well-known Ljusternik-Fet Theorem without using the classical minimax principle, an alternative proof can be found in .
###### Theorem 1.3
(Cartan-Hadamard) Every non-simply closed Riemannian manifold has at least one closed geodesics.
Our method can not conclude that the geodesics is minimal.
Sketch of proofs: We work in the framework as in . In Section 2, we study the linear Cauchy-Riemann operator and sketch some basic properties. In section 3, first we construct a Lagrangian submanifold $`W`$ under the assumption that there does not exists closed Reeb orbit in $`(\mathrm{\Sigma },\lambda )`$; second, we study the space $`𝒟(V,W)`$ of contractible disks in manifold $`V`$ with boundary in Lagrangian submanifold $`W`$ and construct a Fredholm section of tangent bundle of $`𝒟(V,W)`$. In section 4, following , we prove that the Fredholm section is not proper by using a special anti-holomorphic section as in . In section 5-6, we use a geometric argument to deduce the boundary $`C^0`$estimates on $`W`$. In the final section, we use nonlinear Fredholm trick in to complete our proof.
Since the proofs in this paper are very difficult, we suggest the reader first read the Gromov’s paper, Audin and Lafondaine’s book, and Hummel’s book.
###### Note 1.1
The related problem with Weinstein conjecture(see) is Arnold chord conjecture(see) which was discussed in and finally solved in . The generalized Arnold conjecture corresponding to Theorem 1.1 was also solved in similar method of this paper. These results was reported in the Second International Conference on Nonlinear Analysis, 14-19 June 1999, Tianjin, China; First $`3\times 3`$ Canada-China Math Congress, Tsinghua University, Beijing, August 23-28,1999; Differential Geometry Seminar in Nankai Institute of Mathematics, Oct. 24-31, 2000; Symplectic Geometry Seminar In Nankai Institute of Math., Dec. 28-31, 2000; International conference on Symplectic geometry in Sichun Uni., June 24-July, 2001. Some technique part of proofs was carried in ICTP from August to October, 2001. The auther is deeply grateful to thank for the all inviters, especially to Professor Y. M. Long.
## 2 Linear Fredholm Theory
For $`100<k<\mathrm{}`$ consider the Hilbert space $`V_k`$ consisting of all maps $`uH^{k,2}(D,C\times C^n)`$, such that $`u(z)\{izR\}\times R^nC\times C^n`$ for almost all $`zD`$. $`L_{k1}`$ denotes the usual Sobolev space $`H_{k1}(D,C\times C^n)`$. We define an operator $`\overline{}:V_kL_{k1}`$ by
$$\overline{}u=u_s+iu_t$$
(2.1)
where the coordinates on $`D`$ are $`(s,t)=s+it`$, $`D=\{z||z|1\}`$. The following result is well known(see\[34, p96,Th3.3.1\]).
###### Proposition 2.1
$`\overline{}:V_kL_{k1}`$ is a surjective real linear Fredholm operator of index $`n+3`$. The kernel consists of $`(a_0+isz\overline{a}_0z^2,s_1,\mathrm{},s_n)`$, $`a_0C`$, $`s,s_1,\mathrm{},s_n`$.
Let $`(C^n,\sigma =Im(,))`$ be the standard symplectic space. We consider a real $`n`$dimensional plane $`R^nC^n`$. It is called Lagrangian if the skew-scalar product of any two vectors of $`R^n`$ equals zero. For example, the plane $`\{(p,q)|p=0\}`$ and $`\{(p,q)|q=0\}`$ are two transversal Lagrangian subspaces. The manifold of all (nonoriented) Lagrangian subspaces of $`R^{2n}`$ is called the Lagrangian-Grassmanian $`\mathrm{\Lambda }(n)`$. One can prove that the fundamental group of $`\mathrm{\Lambda }(n)`$ is free cyclic, i.e. $`\pi _1(\mathrm{\Lambda }(n))=Z`$. Next assume $`(\mathrm{\Gamma }(z))_{zD}`$ is a smooth map associating to a point $`zD`$ a Lagrangian subspace $`\mathrm{\Gamma }(z)`$ of $`C^n`$, i.e. $`(\mathrm{\Gamma }(z))_{zD}`$ defines a smooth curve $`\alpha `$ in the Lagrangian-Grassmanian manifold $`\mathrm{\Lambda }(n)`$. Since $`\pi _1(\mathrm{\Lambda }(n))=Z`$, one have $`[\alpha ]=ke`$, we call integer $`k`$ the Maslov index of the closed curve $`\alpha `$ and denote it by $`m(\mathrm{\Gamma })`$ (see\[2, p116-118\]). Note that the Maslov index of the closed curve $`\alpha `$ is just the two times of the rotation numbers(see\[2, p116-118\] or \[34, p96,Th3.3.1\]).
Now let $`z:S^1R\times R^nC\times C^n`$ be a smooth curve. Then it defines a constant loop $`\alpha `$ in Lagrangian-Grassmanian manifold $`\mathrm{\Lambda }(n+1)`$. This loop defines the Maslov index $`m(\alpha )`$ of the map $`z`$ which is easily seen to be zero.
Now Let $`(V,\omega )`$ be a symplectic manifold, $`WV`$ a closed Lagrangian submanifold. Let $`(\overline{V},\overline{\omega })=(D\times V,\omega _0+\omega )`$ and $`\overline{W}=D\times W`$. Let $`\overline{u}=(id,u):(D,D)(D\times V,D\times W)`$ be a smooth map homotopic to the map $`(id,u_0)`$, here $`u_0:(D,D)pWV`$. Then $`\overline{u}^{}TV`$ is a symplectic vector bundle on $`D`$ and $`(\overline{u}|_D)^{}T\overline{W}`$ be a Lagrangian subbundle in $`\overline{u}^{}T\overline{V}|_D`$. Since $`\overline{u}:(D,D)(\overline{V},\overline{W})`$ is homotopic to $`\overline{u}_0`$, here$`u_0(z)=(z,p)`$, i.e., there exists a homotopy $`h:[0,1]\times (D,D)(\overline{V},\overline{W})`$ such that $`h(0,z)=(z,p),h(1,z)=\overline{u}(z)`$, we can take a trivialization of the symplectic vector bundle $`h^{}T\overline{V}`$ on $`[0,1]\times (D,D)`$ as
$$\mathrm{\Phi }(h^{}T\overline{V})=[0,1]\times D\times C\times C^n$$
and
$$\mathrm{\Phi }((h|_{[0,1]\times D})^{}T\overline{W})[0,1]\times S^1\times C\times C^n$$
Let
$$\pi _2:[0,1]\times D\times C\times C^nC\times C^n$$
then
$$\stackrel{~}{h}:(s,z)[0,1]\times S^1\pi _2\mathrm{\Phi }(h|_{[0,1]\times D})^{}T\overline{W}|(s,z)\mathrm{\Lambda }(n+1).$$
###### Lemma 2.1
Let $`\overline{u}:(D,D)(\overline{V},\overline{W})`$ be a $`C^k`$map $`(k1)`$ as above. Then,
$$m(\stackrel{~}{u})=2.$$
Proof. Since $`\overline{u}`$ is homotopic to $`\overline{u}_0`$ in $`\overline{V}`$ relative to $`\overline{W}`$, by the above argument we have a homotopy $`\mathrm{\Phi }_s`$ of trivializations such that
$$\mathrm{\Phi }_s(\overline{u}^{}TV)=D\times C\times C^n$$
and
$$\mathrm{\Phi }_s((\overline{u}|_D)^{}T\overline{W})S^1\times C\times C^n$$
Moreover
$$\mathrm{\Phi }_0(\overline{u}|_D)^{}T\overline{W}=S^1\times izR\times R^n$$
So, the homotopy induces a homotopy $`\stackrel{~}{h}`$ in Lagrangian-Grassmanian manifold. Note that $`m(\stackrel{~}{h}(0,))=0`$. By the homotopy invariance of Maslov index, we know that $`m(\stackrel{~}{u}|_D)=2`$.
Consider the partial differential equation
$`\overline{}\overline{u}+A(z)\overline{u}=0onD`$ (2.2)
$`\overline{u}(z)\mathrm{\Gamma }(z)(izR\times R^n)forzD`$ (2.3)
$`\mathrm{\Gamma }(z)GL(2(n+1),R)Sp(2(n+1))`$ (2.4)
$`m(\mathrm{\Gamma })=2`$ (2.5)
For $`100<k<\mathrm{}`$ consider the Banach space $`\overline{V}_k`$ consisting of all maps $`uH^{k,2}(D,C^n)`$ such that $`u(z)\mathrm{\Gamma }(z)`$ for almost all $`zD`$. Let $`L_{k1}`$ the usual Sobolev space $`H_{k1}(D,C\times C^n)`$
###### Proposition 2.2
$`\overline{}:\overline{V}_kL_{k1}`$ is a real linear Fredholm operator of index n+3.
## 3 Nonlinear Fredholm Theory
### 3.1 Constructions of Lagrangian submanifolds
Let $`(\mathrm{\Sigma },\lambda )`$ be a contact manifolds with contact form $`\lambda `$ and $`X`$ its Reeb vector field, then $`X`$ integrates to a Reeb flow $`\eta _t`$ for $`tR^1`$. Consider the form $`d(e^a\lambda )`$ at the point $`(a,\sigma )`$ on the manifold $`(R\times \mathrm{\Sigma })`$, then one can check that $`d(e^a\lambda )`$ is a symplectic form on $`R\times \mathrm{\Sigma }`$. Moreover One can check that
$`i_X(e^a\lambda )=e^a`$ (3.1)
$`i_X(d(e^a\lambda ))=de^a`$ (3.2)
So, the symplectization of Reeb vector field $`X`$ is the Hamilton vector field of $`e^a`$ with respect to the symplectic form $`d(e^a\lambda )`$. Therefore the Reeb flow lifts to the Hamilton flow $`h_s`$ on $`R\times \mathrm{\Sigma }`$(see).
Let
$$(V^{},\omega ^{})=((R\times \mathrm{\Sigma })\times (R\times \mathrm{\Sigma }),d(e^a\lambda )d(e^b\lambda ))$$
and
$$=\{((0,\sigma ),(0,\sigma ))|(0,\sigma )R\times \mathrm{\Sigma }\}.$$
Let
$$L^{}=\times R,L_s^{}=\times \{s\}.$$
Then define
$`G^{}:L^{}V^{}`$ (3.3)
$`G^{}(l^{})=G^{}(((\sigma ,0),(\sigma ,0)),s)=((0,\sigma ),(0,\eta _s(\sigma )))`$ (3.4)
Then
$$W^{}=G^{}(L^{})=\{((0,\sigma ),(0,\eta _s(\sigma )))|(0,\sigma )R\times \mathrm{\Sigma },sR\}$$
$$W_s^{}=G^{}(L_s^{})=\{((0,\sigma ),(0,\eta _s(\sigma )))|(0,\sigma )R\times \mathrm{\Sigma }\}$$
for fixed $`sR`$.
###### Lemma 3.1
There does not exist any Reeb closed orbit in $`(\mathrm{\Sigma },\lambda )`$ if and only if $`W_s^{}W_s^{}^{}`$ is empty for $`ss^{}`$.
Proof. First if there exists a closed Reeb orbit in $`(\mathrm{\Sigma },\lambda )`$, i.e., there exists $`\sigma _0\mathrm{\Sigma }`$, $`t_0>0`$ such that $`\sigma _0=\eta _{t_0}(\sigma _0)`$, then $`((0,\sigma _0),(0,\sigma _0))W_0^{}W_{t_0}^{}`$. Second if there exists $`s_0s_0^{}`$ such that $`W_{s_0}^{}W_{s_0^{}}^{}\mathrm{}`$, i.e., there exists $`\sigma _0`$ such that
$$((0,\sigma _0),(0,\eta _{s_0}(\sigma _0))=((0,\sigma _0),(0,\eta _{s_0^{}}(\sigma _0)),$$
then $`\eta _{(s_0s_0^{})}(\sigma _0)=\sigma _0`$, i.e., $`\eta _t(\sigma _0)`$ is a closed Reeb orbit.
###### Lemma 3.2
If there does not exist any closed Reeb orbit in $`(\mathrm{\Sigma },\lambda )`$ then there exists a smooth Lagrangian injective immersion $`G^{}:W^{}V^{}`$ with $`G^{}(((0,\sigma ),(0,\sigma )),s)=((0,\sigma ),(0,\eta _s(\sigma )))`$ such that
$$G_{s_1,s_2}^{}:\times (s_1,s_2)V^{}$$
(3.5)
is a regular exact Lagrangian embedding for any finite real number $`s_1`$, $`s_2`$, here we denote by $`W^{}(s_1,s_2)=G_{s_1,s_2}^{}(\times (s_1,s_2))`$.
Proof. One check
$$G_{}^{}{}_{}{}^{}((e^a\lambda e^b\lambda ))=\lambda \eta (,)^{}\lambda =\lambda (\eta _s^{}\lambda +i_X\lambda ds)=ds$$
(3.6)
since $`\eta _s^{}\lambda =\lambda `$. This implies that $`G^{}`$ is an exact Lagrangian embedding, this proves Lemma 3.2.
Now set
$$c(s,t)=\epsilon te^{s^2}$$
(3.7)
$$\psi _0(s,t)=se^{c(s,t)}c_s=2e^{(\epsilon te^{s^2})s^2}\epsilon ts^2=\epsilon \psi _0^{}$$
(3.8)
here $`\psi _0^{}(s,t)=2ts^2e^{(\epsilon te^{s^2})s^2}`$;
$$\psi _1(s,t)=_{\mathrm{}}^s\psi _0(\tau ,t)𝑑\tau =\epsilon _{\mathrm{}}^s\psi _0^{}(\tau ,t)𝑑\tau =\epsilon \psi _1^{}$$
(3.9)
here $`\psi _1^{}=_{\mathrm{}}^s\psi _0^{}`$;
$$\psi (s,t)=\frac{\psi _1}{t}se^{c(s,t)}c_t=\epsilon \psi ^{}$$
(3.10)
here $`\psi ^{}(s,t)=\frac{\psi _1}{t}se^{\epsilon te^{s^2}}e^{s^2}`$;
$$\mathrm{\Psi }^{}=se^{c(s,t)};\stackrel{~}{l}^{}=\psi (s,t)dt.$$
(3.11)
Now we construct an isotopy of Lagrangian injective immersions as follows:
$`F^{}:\times R\times [0,1](R\times \mathrm{\Sigma })\times (R\times \mathrm{\Sigma })`$ (3.12)
$`F^{}(((0,\sigma ),(0,\sigma )),s,t)=((c(s,t),\sigma ),(c(s,t),\eta _s(\sigma )))`$ (3.13)
$`F_t^{}(((0,\sigma ),(0,\sigma )),s)=F^{}(((0,\sigma ),(0,\sigma )),s,t)`$ (3.14)
###### Lemma 3.3
If there does not exist any closed Reeb orbit in $`(\mathrm{\Sigma },\lambda )`$ then for the choice of $`c(s,t)`$ satisfying $`_0^sc(s,t)𝑑s`$ or $`_s^0c(s,t)𝑑s`$ exists and is smooth on $`(s,t)`$, $`F^{}`$ is an exact isotopy of Lagrangian embeddings. Moreover if $`c(s,0)c(s,1)`$, then $`F_0^{}(\mathrm{\Sigma }\times R)F_1^{}(\mathrm{\Sigma }\times R)=\mathrm{}`$.
Proof. Let $`F_t^{}=F^{}(,t):\times R(R\times \mathrm{\Sigma })\times (R\times \mathrm{\Sigma })`$. It is obvious that $`F_t^{}`$ is an embedding. We check that
$`F_{}^{}{}_{}{}^{}(e^a\lambda e^b\lambda )`$ $`=`$ $`e^{c(s,t)}ds`$ (3.15)
$`=`$ $`\{d(se^{c(s,t)})sde^{c(s,t)}\}`$ (3.16)
$`=`$ $`\{d(se^{c(s,t)})se^{c(s,t)}c_sdsse^{c(s,t)}c_tdt\}`$ (3.17)
$`=`$ $`\{d(se^{c(s,t)})d_s\psi _1se^{c(s,t)}c_tdt\}`$ (3.18)
$`=`$ $`\{d((se^{c(s,t)})\psi _1)+{\displaystyle \frac{\psi _1}{t}}dtse^{c(s,t)}c_tdt\}`$ (3.19)
$`=`$ $`\{d\mathrm{\Psi }^{}+{\displaystyle \frac{\psi _1}{t}}dtse^{c(s,t)}c_tdt\}`$ (3.20)
$`=`$ $`d\mathrm{\Psi }^{}\psi (s,t)dt`$ (3.21)
$`=`$ $`d\mathrm{\Psi }^{}+\stackrel{~}{l}^{}`$ (3.22)
here $`\psi _i,\psi _i^{}`$, and $`\stackrel{~}{l}^{},\mathrm{\Psi }^{}`$ as in (3.8-3.11).
Let $`(V^{},\omega ^{})`$, $`W^{}`$ as above and $`(V,\omega )=(V^{}\times C,\omega ^{}\omega _0)`$. As in \[13, p330,2.3.$`B_3^{}`$\](see also \[3, p291-292\]), we use figure eight trick invented by Gromov to construct a Lagrangian submanifold in $`V`$ through the Lagrange isotopy $`F^{}`$ in $`V^{}`$. Fix a positive $`\delta <1`$ and take a $`C^{\mathrm{}}`$-map $`\rho :S^1[0,1]`$, where the circle $`S^1`$ ia parametrized by $`\mathrm{\Theta }[1,1]`$, such that the $`\delta `$neighborhood $`I_0`$ of $`0S^1`$ goes to $`0[0,1]`$ and $`\delta `$neighbourhood $`I_1`$ of $`\pm 1S^1`$ goes $`1[0,1]`$. Let
$`\stackrel{~}{l}`$ $`=`$ $`\psi (s,\rho (\mathrm{\Theta }))\rho ^{}(\mathrm{\Theta })d\mathrm{\Theta }`$ (3.23)
$`=`$ $`\mathrm{\Phi }d\mathrm{\Theta }`$ (3.24)
be the pull-back of the form $`\stackrel{~}{l}^{}=\psi (s,t)dt`$ to $`W^{}\times S^1`$ under the map $`(w^{},\mathrm{\Theta })(w^{},\rho (\mathrm{\Theta }))`$ and assume without loss of generality $`\mathrm{\Phi }`$ vanishes on $`W^{}\times (I_0I_1)`$.
Next, consider a map $`\alpha `$ of the annulus $`S^1\times [\mathrm{\Phi }_{},\mathrm{\Phi }_+]`$ into $`R^2`$, where $`\mathrm{\Phi }_{}`$ and $`\mathrm{\Phi }_+`$ are the lower and the upper bound of the fuction $`\mathrm{\Phi }`$ correspondingly, such that
$`(i)`$ The pull-back under $`\alpha `$ of the form $`dxdy`$ on $`R^2`$ equals $`d\mathrm{\Phi }d\mathrm{\Theta }`$.
$`(ii)`$ The map $`\alpha `$ is bijective on $`I\times [\mathrm{\Phi }_{},\mathrm{\Phi }_+]`$ where $`IS^1`$ is some closed subset, such that $`II_0I_1=S^1`$; furthermore, the origin $`0R^2`$ is a unique double point of the map $`\alpha `$ on $`S^1\times 0`$, that is
$$0=\alpha (0,0)=\alpha (\pm 1,0),$$
and $`\alpha `$ is injective on $`S^1=S^1\times 0`$ minus $`\{0,\pm 1\}`$.
$`(iii)`$ The curve $`S_0^1=\alpha (S^1\times 0)R^2`$ “bounds” zero area in $`R^2`$, that is $`_{S_0^1}x𝑑y=0`$, for the $`1`$form $`xdy`$ on $`R^2`$.
###### Proposition 3.1
Let $`V^{}`$, $`W^{}`$ and $`F^{}`$ as above. Then there exists an exact Lagrangian embedding $`F:W^{}\times S^1V^{}\times R^2`$ given by $`F(w^{},\mathrm{\Theta })=(F^{}(w^{},\rho (\mathrm{\Theta })),\alpha (\mathrm{\Theta },\mathrm{\Phi }))`$.
Proof. We follow as in \[13, 2.3$`B_3^{}`$\]. Now let $`F^{}:W^{}\times S^1V^{}\times R^2`$ be given by $`(w^{},\mathrm{\Theta })(F^{}(w,\rho (\mathrm{\Theta })),\alpha (\mathrm{\Theta },\mathrm{\Phi }))`$. Then
$`(i)^{}`$ The pull-back under $`F^{}`$ of the form $`\omega =\omega ^{}+dxdy`$ equals $`d\stackrel{~}{l}^{}d\mathrm{\Phi }d\mathrm{\Theta }=0`$ on $`W^{}\times S^1`$.
$`(ii)^{}`$ The set of double points of $`F^{}`$ is $`W_0^{}W_1^{}V^{}=V^{}\times 0V^{}\times R^2`$.
$`(iii)^{}`$ If $`F^{}`$ has no double point then the Lagrangian submanifold $`W=F^{}(W^{}\times S^1)(V^{}\times R^2,\omega ^{}+dxdy)`$ is exact if and only if $`W_0^{}V^{}`$ is such.
This completes the proof of Proposition 3.1.
### 3.2 Formulation of Hilbert bundles
Let $`(\mathrm{\Sigma },\lambda )`$ be a closed $`(2n1)`$ dimensional manifold with a contact form $`\lambda `$. Let $`S\mathrm{\Sigma }=R\times \mathrm{\Sigma }`$ and put $`\xi =\mathrm{ker}(\lambda )`$. Let $`J_\lambda ^{}`$ be an almost complex structure on $`S\mathrm{\Sigma }`$ tamed by the symplectic form $`d(e^a\lambda )`$.
We define a metric $`g_\lambda `$ on $`S\mathrm{\Sigma }=R\times \mathrm{\Sigma }`$ by
$$g_\lambda =d(e^a\lambda )(,J_\lambda )$$
(3.25)
which is adapted to $`J_\lambda `$ and $`d(e^a\lambda )`$ but not complete.
In the following we denote by $`(V^{},\omega ^{})=((R\times \mathrm{\Sigma })\times (R\times \mathrm{\Sigma }),d(e^a\lambda _1e^b\lambda _2))`$ and $`(V,\omega )=(V^{}\times R^2,\omega ^{}+dxdy)`$ with the metric $`g=g^{}g_0=g_{\lambda _1}g_{\lambda _2}g_0`$ induced by $`\omega (,J)`$($`J=J^{}i=J_{\lambda _1}(J_{\lambda _2})i`$ and $`WV`$ a Lagrangian submanifold which was constructed in section 3.1.
Let $`\overline{V}=D\times V`$, then $`\pi _1:\overline{V}D`$ be a symplectic vector bundle. Let $`\overline{J}`$ be an almost complex structure on $`\overline{V}`$ such that $`\pi _1:\overline{V}D`$ is a holomorphic map and each fibre $`\overline{V}_z=\pi _1(z)`$ is a $`\overline{J}`$ complex submanifold. Let $`H^k(D)`$ be the space of $`H^k`$maps from $`D`$ to $`\overline{V}`$, here $`H^k`$ represents Sobolev derivatives up to order $`k`$. Let $`\overline{W}=D\times W`$, $`\overline{p}=\{1\}\times p`$, $`W^\pm =\{\pm i\}\times W`$ and
$$𝒟^k=\{\overline{u}H^k(D)|\overline{u}(x)\overline{W}a.eforxDand\overline{u}(1)=\overline{p},\overline{u}(\pm i)\{\pm i\}\times W\}$$
for $`k100`$.
###### Lemma 3.4
Let $`W`$ be a closed Lagrangian submanifold in $`V`$. Then, $`𝒟^k`$ is a pseudo-Hilbert manifold with the tangent bundle
$$T𝒟^k=\underset{\overline{u}𝒟^k}{}\mathrm{\Lambda }^{k1}$$
(3.26)
here
$$\mathrm{\Lambda }^{k1}=\{\overline{w}H^{k1}(\overline{u}^{}(T\overline{V})|\overline{w}(1)=0,and\overline{w}(\pm i)T\overline{W}\}$$
###### Note 3.1
Since $`W`$ is not regular we know that $`𝒟^k`$ is in general complete, however it is enough for our purpose.
Proof: See \[3, p309-310\] or follow step by step from \[20, ch1\].
Now we consider a section from $`𝒟^k`$ to $`T𝒟^k`$ follows as in \[13, p327,2.2\] or \[13, p310\], i.e., let $`\overline{}:𝒟^kT𝒟^k`$ be the Cauchy-Riemmann section
$$\overline{}\overline{u}=\frac{\overline{u}}{s}+J\frac{\overline{u}}{t}$$
(3.27)
for $`\overline{u}𝒟^k`$.
###### Theorem 3.1
The Cauchy-Riemann section $`\overline{}`$ defined in (3.27) is a Fredholm section of Index zero.
Proof. According to the definition of the Fredholm section, we need to prove that $`\overline{u}𝒟^k`$, the linearization $`D\overline{}(\overline{u})`$ of $`\overline{}`$ at $`\overline{u}`$ is a linear Fredholm operator. Note that
$$D\overline{}(\overline{u})=D\overline{}_{[\overline{u}]}$$
(3.28)
where
$$(D\overline{}_{[\overline{u}]})v=\frac{\overline{v}}{s}+J\frac{\overline{v}}{t}+A(\overline{u})\overline{v}$$
(3.29)
with
$$\overline{v}|_D(\overline{u}|_D)^{}T\overline{W}$$
here $`A(\overline{u})`$ is $`2n\times 2n`$ matrix induced by the torsion of almost complex structure, see \[13, p324,2.1\] for the computation.
Observe that the linearization $`D\overline{}(\overline{u})`$ of $`\overline{}`$ at $`\overline{u}`$ is equivalent to the following Lagrangian boundary value problem
$`{\displaystyle \frac{\overline{v}}{s}}+\overline{J}{\displaystyle \frac{\overline{v}}{t}}+A(\overline{u})\overline{v}=\overline{f},\overline{v}\mathrm{\Lambda }^k(\overline{u}^{}T\overline{V})`$ (3.30)
$`\overline{v}(t)T_{\overline{u}(t)}W,tD`$ (3.31)
One can check that (3.31) defines a linear Fredholm operator. In fact, by proposition 2.2 and Lemma 2.1, since the operator $`A(\overline{u})`$ is a compact, we know that the operator $`\overline{}`$ is a nonlinear Fredholm operator of the index zero.
###### Definition 3.1
Let $`X`$ be a Banach manifold and $`P:YX`$ the Banach vector bundle. A Fredholm section $`F:XY`$ is proper if $`F^1(0)`$ is a compact set and is called generic if $`F`$ intersects the zero section transversally, see \[13, p327-328,2.2B\].
###### Definition 3.2
$`deg(F,y)=\mathrm{}\{F^1(0)\}mod2`$ is called the Fredholm degree of a Fredholm section (see\[13, p327-328,2.2B\]).
###### Theorem 3.2
Assum that $`\overline{J}=iJ`$ on $`\overline{V}`$ and $`i`$ is complex structure on $`D`$ and $`J`$ the almost complex structure on $`V`$. Assume that $`J`$ is integrable at $`pV`$. Then the Fredholm section $`F=\overline{}_{\overline{J}}:𝒟^kT𝒟^k`$ constructed in (3.27) has degree one, i.e.,
$$deg(F,0)=1$$
Proof: We assume that $`\overline{u}:D\overline{V}`$ be a $`\overline{J}`$holomorphic disk with boundary $`\overline{u}(D)\overline{W}`$ and by the assumption that $`\overline{u}`$ is homotopic to the map $`\overline{u}_1=(id,p)`$. Since almost complex structure $`\overline{J}`$ splits and is tamed by the symplectic form $`\overline{\omega }`$, by stokes formula, we conclude the second component $`u:DV`$ is a constant map. Because $`u(1)=p`$, We know that $`F^1(0)=(id,p)`$. Next we show that the linearizatioon $`DF_{(id,p)}`$ of $`F`$ at $`(id,p)`$ is an isomorphism from $`T_{(id,p)}𝒟^k`$ to $`E`$. This is equivalent to solve the equations
$`{\displaystyle \frac{\overline{v}}{s}}+\overline{J}{\displaystyle \frac{\overline{v}}{t}}=f`$ (3.32)
$`\overline{v}|_DT_{(id,p)}\overline{W}`$ (3.33)
here $`\overline{J}=i+J(p)`$ since $`J`$ is integrable at $`p`$. By Lemma 2.1, we know that $`DF_{(id,p)}`$ is an isomorphism. Therefore $`deg(F,0)=1`$.
## 4 Anti-holomorphic sections
In this section we construct a Fredholm section which is not proper as in \[13, p329-330,2.3.B\](see also \[3, p315, 5.3\]).
Let $`(V^{},\omega ^{})=(S\mathrm{\Sigma }\times S\mathrm{\Sigma },d(e^a\lambda _1e^b\lambda _2))`$ and $`(V,\omega )=(V^{}\times C,\omega ^{}\omega _0)`$, $`W`$ as in section3 and $`J=J^{}i`$, $`g=g^{}g_0`$, $`g_0`$ the standard metric on $`C`$.
Now let $`cC`$ be a non-zero vector. We consider $`c`$ as an anti-holomorphic homomorphism $`c:TDTV^{}TC`$, i.e., $`c(\frac{}{\overline{z}})=(0,c\frac{}{z}).`$ Since the constant section $`c`$ is not a section of the Hilbert bundle in section 3 due to $`c`$ is not tangent to the Lagrangian submanifold $`W`$, we must modify it as follows:
Let $`c`$ as in section 4.1, we define
$`c_{\chi ,\delta }(z,v)=\{\begin{array}{cc}c\hfill & \text{if }|z|12\delta \text{,}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ (4.3)
Then by using the cut off function $`\phi _h(z)`$ and its convolution with section $`c_{\chi ,\delta }`$, we obtain a smooth section $`c_\delta `$ satisfying
$`c_\delta (z,v)=\{\begin{array}{cc}c\hfill & \text{if }|z|13\delta \text{,}\hfill \\ 0\hfill & \text{if }|z|1\delta \text{.}\hfill \end{array}`$ (4.6)
$`|c_\delta ||c|`$ (4.7)
for $`h`$ small enough, for the convolution theory see \[18, ch1,p16-17,Th1.3.1\]. Then one can easily check that $`\overline{c}_\delta =(0,0,c_\delta )`$ is an anti-holomorphic section tangent to $`\overline{W}`$.
Now we modify the almost conplex structure on the $`V`$. Let $`J_1`$, $`J_2`$ be the almost complex structures on $`V`$ tamed by $`\omega `$. Let $`g_i=\omega (,J_i)`$ the metrics by $`\omega `$ and $`J_i`$. We assume there exists a constant $`c_1`$ such that
$$c_1^1g_1g_2c_1g_1$$
(4.8)
Let
$`J_{\chi ,\delta }(z,v)=\{\begin{array}{cc}iJ_1\hfill & \text{if }|z|12\delta \text{,}\hfill \\ iJ_2\hfill & \text{otherwise}\hfill \end{array}`$ (4.11)
Then by using the cut off function $`\phi _h(z)`$ and its convolution with section $`J_{\chi ,\delta }`$, we obtain a smooth section $`J_\delta `$ satisfying
$`J_\delta (z,v)=\{\begin{array}{cc}iJ_1\hfill & \text{if }|z|13\delta \text{,}\hfill \\ iJ_2\hfill & \text{if }|z|1\delta \text{.}\hfill \end{array}`$ (4.14)
for $`h`$ small enough, for the convolution theory see \[18, ch1,p16-17,Th1.3.1\].
Now we get an almost conplex structure $`\overline{J}=iJ_\delta `$ on the symplectic fibration $`D\times VD`$ such that $`\pi _1:D\times VD`$ is a holomorphic fibration and $`\pi _1^1(z)`$ is an almost complex submanifold. Let $`g_\delta =\overline{\omega }(,\overline{J})`$, $`\overline{g}_i=g_0g_i`$ be the metrics on $`D\times V`$, $`g_0`$ is metric on $`D`$. We assume there exists a constants $`c_2`$ such that
$$c_2^1\overline{g}_ig_\delta c_2g_i,i=1,2.$$
(4.15)
Now we consider the equations
$`\overline{v}=(id,v)=(id,v^{},f):DD\times V^{}\times C`$ (4.16)
$`\overline{}_{J_\delta }v=c_\delta `$ (4.17)
$`\overline{}_J^{}v^{}=0,\overline{}f=c_\delta onD_{12\delta }`$ (4.18)
$`v|_D:DW`$ (4.19)
here $`v`$ homotopic to constant map $`\{p\}`$ relative to $`W`$. Note that $`WV\times B_R(0)`$ for $`2\pi R(\epsilon )^2`$, here $`R(\epsilon )0`$ as $`\epsilon 0`$ and $`\epsilon `$ as in section 3.1.
###### Lemma 4.1
Let $`\overline{v}=(id,v)`$ be the solutions of (4.19), then one has the following estimates
$`E(v)=\{{\displaystyle _D}(g^{}({\displaystyle \frac{v^{}}{x}},J^{}{\displaystyle \frac{v^{}}{x}})+g^{}({\displaystyle \frac{v^{}}{y}},J^{}{\displaystyle \frac{v^{}}{y}})`$
$`+g_0({\displaystyle \frac{f}{x}},i{\displaystyle \frac{f}{x}})+g_0({\displaystyle \frac{f}{y}},i{\displaystyle \frac{f}{y}}))d\sigma \}4\pi R(\epsilon )^2.`$ (4.20)
Proof: Since $`v(z)=(v^{}(z),f(z))`$ satisfy (4.19) and $`v(z)=(v^{}(z),f(z))V^{}\times C`$ is homotopic to constant map $`v_0:D\{p\}W`$ in $`(V,W)`$, by the Stokes formula
$$_Dv^{}(\omega ^{}\omega _0)=0$$
(4.21)
Note that the metric $`g`$ is adapted to the symplectic form $`\omega `$ and $`J`$, i.e.,
$$g=\omega (,J)$$
(4.22)
By the simple algebraic computation, we have
$$_Dv^{}\omega =\frac{1}{4}_{D^2}(|v|^2|\overline{}v|^2)=0$$
(4.23)
and
$$|v|=\frac{1}{2}(|v|^2+|\overline{}v|^2$$
(4.24)
Then
$`E(v)`$ $`=`$ $`{\displaystyle _D}|v|`$ (4.25)
$`=`$ $`{\displaystyle _D}\{{\displaystyle \frac{1}{2}}(|v|^2+|\overline{}v|^2)\})d\sigma `$
$`=`$ $`{\displaystyle _D}|c_\delta |_{\overline{g}}^2𝑑\sigma `$
By Cauchy integral formula,
$$f(z)=\frac{1}{2\pi i}_D\frac{f(\xi )}{\xi z}𝑑\xi +\frac{1}{2\pi i}_D\frac{\overline{}f(\xi )}{\xi z}𝑑\xi d\overline{\xi }$$
(4.26)
Since $`f`$ is smooth up to the boundary, we integrate the two sides on $`D_r`$ for $`r<1`$, one get
$`{\displaystyle _{D_r}}f(z)𝑑z`$ $`=`$ $`{\displaystyle _{D_r}}{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}{\displaystyle \frac{f(\xi )}{\xi z}}𝑑\xi 𝑑z+{\displaystyle _{D_r}}{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}{\displaystyle \frac{\overline{}f(\xi )}{\xi z}}𝑑\xi d\overline{\xi }`$ (4.27)
$`=`$ $`0+{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}{\displaystyle _{D_r}}{\displaystyle \frac{\overline{}f(\xi )}{\xi z}}𝑑z𝑑\xi d\overline{\xi }`$ (4.28)
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}2\pi i\overline{}f(\xi )𝑑\xi d\overline{\xi }`$ (4.29)
Let $`r1`$, we get
$`{\displaystyle _D}f(z)𝑑z={\displaystyle _D}\overline{}f(\xi )𝑑\xi d\overline{\xi }`$ (4.30)
By the equations (4.19), one get
$$\overline{}f=conD_{12\delta }$$
(4.31)
So, we have
$$2\pi i(12\delta )c=_Df(z)𝑑z_{DD_{12\delta }}\overline{}f(\xi )𝑑\xi d\overline{\xi }$$
(4.32)
So,
$`|c|`$ $``$ $`{\displaystyle \frac{1}{2\pi (12\delta )}}|{\displaystyle _D}f(z)𝑑z|+|{\displaystyle _{DD_{12\delta }}}\overline{}f(\xi )𝑑\xi d\overline{\xi }|`$ (4.33)
$``$ $`{\displaystyle \frac{1}{2\pi (12\delta )}}2\pi |diam(pr_2(W))+c_1c_2|c|(\pi \pi (12\delta )^2))`$ (4.34)
Therefore, one has
$`|c|`$ $``$ $`c(\delta )R(\epsilon )`$ (4.35)
and
$`E(v)`$ $`=`$ $`\pi {\displaystyle _D}|c_\delta |_{\overline{g}}^2`$ (4.36)
$`=`$ $`\pi c(\delta )^2R(\epsilon )^2.`$ (4.37)
This finishes the proof of Lemma.
###### Proposition 4.1
For $`|c|2c(\delta )R(\epsilon )`$, then the equations (4.19) has no solutions.
Proof. By 4.35, it is obvious.
###### Theorem 4.1
The Fredholm section $`F_1=\overline{}_{\overline{J}}+\overline{c}_\delta :𝒟^kE`$ is not proper.
Proof. By the Proposition 4.1 and Theorem 3.2, it is obvious(see also \[13, p330,$`2.3B_1`$\] or \[3, p316\]).
## 5 $`J`$holomorphic section
Recall that $`W(K,K)WV^{}\times R^2`$ as in section 3. The Riemann metric $`g`$ on $`V^{}\times R^2`$ induces a metric $`g|W`$.
Now let $`cC`$ be a non-zero vector and $`c_\delta `$ the induced anti-holomorphic section. We consider the nonlinear inhomogeneous equations (4.19) and transform it into $`\overline{J}`$holomorphic map by considering its graph as in \[13, p319,1.4.C\] or \[3, p312,Lemma5.2.3\].
Denote by $`Y^{(1)}D\times V`$ the bundle of homomorphisms $`T_s(D)T_v(V)`$. If $`D`$ and $`V`$ are given the disk and the almost Kähler manifold, then we distinguish the subbundle $`X^{(1)}Y^{(1)}`$ which consists of complex linear homomorphisms and we denote $`\overline{X}^{(1)}D\times V`$ the quotient bundle $`Y^{(1)}/X^{(1)}`$. Now, we assign to each $`C^1`$-map $`v:DV`$ the section $`\overline{}v`$ of the bundle $`\overline{X}^{(1)}`$ over the graph $`\mathrm{\Gamma }_vD\times V`$ by composing the differential of $`v`$ with the quotient homomorphism $`Y^{(1)}\overline{X}^{(1)}`$. If $`c_\delta :D\times V\overline{X}`$ is a $`H^k`$ section we write $`\overline{}v=c_\delta `$ for the equation $`\overline{}v=c_\delta |\mathrm{\Gamma }_v`$.
###### Lemma 5.1
(Gromov\[13, 1.4.$`C^{}`$\])There exists a unique almost complex structure $`J_g`$ on $`D\times V`$(which also depends on the given structures in $`D`$ and in $`V`$), such that the (germs of) $`J_\delta `$holomorphic sections $`v:DD\times V`$ are exactly and only the solutions of the equations $`\overline{}v=c_\delta `$. Furthermore, the fibres $`z\times VD\times V`$ are $`J_\delta `$holomorphic( i.e. the subbundles $`T(z\times V)T(D\times V)`$ are $`J_\delta `$complex) and the structure $`J_\delta |z\times V`$ equals the original structure on $`V=z\times V`$. Moreover $`J_\delta `$ is tamed by $`k\omega _0\omega `$ for $`k`$ large enough which is independent of $`\delta `$.
## 6 Gromov’s $`C^0`$convergence theorem
### 6.1 Analysis of Gromov’s figure eight
Since $`W^{}S\mathrm{\Sigma }\times S\mathrm{\Sigma }`$ is an exact Lagrangian submanifold and $`F_t^{}`$ is an exact Lagrangian isotopy(see section 3.1). Now we carefully check the Gromov’s construction of Lagrangian submanifold $`WV^{}\times R^2`$ from the exact Lagrangian isotopy of $`W^{}`$ in section 3.
Let $`S^1T^{}S^1`$ be a zero section and $`S^1=_{i=1}^4S_i`$ be a partition of the zero section $`S^1`$ such that $`S_1=I_0`$, $`S_3=I_1`$. Write $`S^1\{I_0I_1\}=I_2I_3`$ and $`I_0=(\delta ,\frac{5}{6}\delta ](\frac{5\delta }{6},+\frac{5\delta }{6})[\frac{5\delta }{6},\delta )=I_0^{}I_0^{}I_0^+`$, similarly $`I_1=(1\delta ,1\frac{5}{6}\delta ](1\frac{5\delta }{6},1+\frac{5\delta }{6})[1+\frac{5\delta }{6},1+\delta )=I_1^{}I_1^{}I_1^+`$. Let $`S_2=I_0^+I_2I_1^{}`$, $`S_4=I_1^+I_3I_0^+`$. Moreover, we can assume that the double points of map $`\alpha `$ in Gromov’s figure eight is contained in $`(\overline{I}_0^{}\overline{I}_1^{})\times [\mathrm{\Phi }_{},\mathrm{\Phi }_+]`$, here $`\overline{I}_0^{}=(\frac{5\delta }{12},+\frac{5\delta }{12})`$ and $`\overline{I}_1^{}=(1\frac{5\delta }{12},1+\frac{5\delta }{12})`$. Recall that $`\alpha :(S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$ is an exact symplectic immersion, i.e., $`\alpha ^{}(ydx)\mathrm{\Psi }d\mathrm{\Theta }=dh`$, $`h:T^{}S^1R`$. By the construction of figure eight, we can assume that $`\alpha _i^{}=\alpha |((S^1I_i^{})\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])`$ is an embedding for $`i=0,1`$. Let $`Y=\alpha (S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$ and $`Y_i=\alpha (S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`\alpha _i=\alpha |Y_i(S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])`$. So, $`\alpha _i`$ puts the function $`h`$ to the function $`h_{i0}=\alpha _i^1h`$ on $`Y_i`$. We extend the function $`h_{i0}`$ to whole plane $`R^2`$. In the following we take the liouville form $`\beta _{i0}=ydxdh_{i0}`$ on $`R^2`$. This does not change the symplectic form $`dxdy`$ on $`R^2`$. But we have $`\alpha _i^{}\beta =\mathrm{\Phi }d\mathrm{\Theta }`$ on $`(S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])`$ for $`i=1,2,3,4`$. Finally, note that
$`F:W^{}\times S^1V^{}\times R^2;`$ (6.1)
$`F(w^{},\mathrm{\Theta })=(F_{\rho (\mathrm{\Theta })}^{}(w^{}),\alpha (\mathrm{\Theta },\mathrm{\Phi }(w^{},\rho (\mathrm{\Theta })).`$ (6.2)
Since $`\rho (\mathrm{\Theta })=0`$ for $`\mathrm{\Theta }I_0`$ and $`\rho (\mathrm{\Theta })=1`$ for $`\mathrm{\Theta }I_1`$, we know that $`\mathrm{\Phi }(w^{},\rho (\mathrm{\Theta }))=0`$ for $`\mathrm{\Theta }I_0I_1`$. Therefore,
$`F(W^{}\times I_0)=W^{}\times \alpha (I_0);F(W^{}\times I_1)=W^{}\times \alpha (I_1).`$ (6.3)
### 6.2 Gromov’s Schwartz lemma
In our proof we need a crucial tools, i.e., Gromov’s Schwartz Lemma as in . We first consider the case without boundary.
###### Proposition 6.1
Let $`(V,J,\mu )`$ be as in section 4 and $`V_K`$ the compact part of $`V`$. There exist constants $`\epsilon _0`$ and $`C`$(depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^\alpha `$ norm of $`J`$ and $`A_0`$) such that every $`J`$holomorphic map of the unit disc to an $`\epsilon _0`$-ball of $`V`$ with center in $`V_K`$ and area less than $`A_0`$ has its derivatives up to order $`k+1+\alpha `$ on $`D_{\frac{1}{2}}(0)`$ bounded by $`C`$.
For a proof, see.
Now we consider the Gromov’s Schwartz Lemma for $`J`$holomorphic map with boundary in a closed Lagrangian submanifold as in .
###### Proposition 6.2
Let $`(V,J,\mu )`$ as above and $`LV`$ be a closed Lagrangian submanifold and $`V_K`$ one compact part of $`V`$. There exist constants $`\epsilon _0`$ and $`C`$(depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^\alpha `$ norm of $`J`$ and $`K,A_0`$) such that every $`J`$holomorphic map of the half unit disc $`D^+`$ to a $`\epsilon _0`$-ball of $`V`$ with boundary in $`L`$ and area less than $`A_0`$ has its derivatives up to order $`k+1+\alpha `$ on $`D_{\frac{1}{2}}^+(0)`$ bounded by $`C`$.
For a proof see .
Since in our case $`W`$ is a non-compact Lagrangian submanifold, Proposition 6.2 can not be used directly but the proofs of Proposition 6.1-2 is still holds in our case.
###### Lemma 6.1
Recall that $`V=V^{}\times R^2`$. Let $`(V,J,\mu )`$ as above and $`WV`$ be as above and $`V_c`$ the compact set in $`V`$. Let $`\overline{V}=D\times V`$, $`\overline{W}=D\times W`$, and $`\overline{V}_c=D\times V_c`$. Let $`Y=\alpha (S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`Y_i=\alpha (S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`\{X_j\}_{j=1}^q`$ be a Darboux covering of $`\mathrm{\Sigma }`$ and $`V_{ij}^{}=(R\times X_i)\times (R\times X_j)`$. Let $`D=S^{1+}S^1`$. There exist constant $`c_0`$ such that every $`J`$holomorphic map $`v`$ of the half unit disc $`D^+`$ to the $`D\times V_j^{}\times R^2`$ with its boundary $`v((1,1))(S^{1\pm })\times F(\times R\times S_i)\overline{W},i=1,..,4`$ has
$$area(v(D^+))c_0l^2(v(^{}D^+)).$$
(6.4)
here $`^{}D^+=D[1,1]`$ and $`l(v(^{}D^+))=length(v(^{}D^+))`$.
Proof. Let $`\overline{W}_{i\pm }=S^{1\pm }\times F(W^{}\times S_i)`$. Let $`v=(v_1,v_2):D^+\overline{V}=D\times V`$ be the $`J`$holomorphic map with $`v(D^+)\overline{W}_{i\pm }D\times W`$, then
$`area(v)`$ $`=`$ $`{\displaystyle _{D^+}}v^{}d(\alpha _0\alpha )`$ (6.5)
$`=`$ $`{\displaystyle _{D^+}}𝑑v^{}(\alpha _0\alpha )`$ (6.6)
$`=`$ $`{\displaystyle _{D^+}}v^{}(\alpha _0\alpha )`$ (6.7)
$`=`$ $`{\displaystyle _{D^+}}v_1^{}\alpha _0+{\displaystyle _{D^+}}v_2^{}\alpha `$ (6.8)
$`=`$ $`{\displaystyle _{^{}D^+[1,+1]}}v_1^{}\alpha _0+{\displaystyle _{^{}D^+[1,+1]}}v_2^{}(e^a\lambda ydxdh_{i0})`$ (6.9)
$`=`$ $`{\displaystyle _{^{}D^+[1,+1]}}v_1^{}\alpha _0+{\displaystyle _{^{}D^+}}v_2^{}(e^a\lambda ydxdh_{i0})+B_1,`$ (6.10)
here $`B_1=_{[1,+1]}v_{2}^{}{}_{}{}^{}(d\mathrm{\Psi }^{})`$. Now take a zig-zag curve $`C`$ in $`V_j^{}\times Y_i`$ connecting $`v_2(1)`$ and $`v_2(+1)`$ such that
$`{\displaystyle _C}(e^a\lambda +ydx)`$ $`=`$ $`B_1`$ (6.11)
$`length(C)`$ $``$ $`k_1length(v_2(^{}D^+))`$ (6.12)
Now take a minimal surface $`M`$ in $`V_{ij}^{}\times R^2`$ bounded by $`v_2(^{}D^+)C`$, then by the isoperimetric ineqality(see\[\[14, p283\]), we get
$`area(M)`$ $``$ $`m_1length(C+v_2(^{}D^+))^2`$ (6.13)
$``$ $`m_2length(v_2(^{}D^+))^2,`$ (6.14)
here we use the (6.12).
Since $`area(M)_M\omega `$ and $`_M\omega =_{D^+}v_2^{}\omega =area(v)`$, this proves the lemma.
###### Lemma 6.2
Let $`v`$ as in Lemma 6.1, then we have
$$area(v(D^+)c_0(dist(v(0),v(^{}D^+)))^2,$$
(6.15)
here $`c_0`$ depends only on $`\mathrm{\Sigma },J,\omega ,\mathrm{},`$etc, not on $`v`$.
Proof. By the standard argument as in \[3, p79\].
The following estimates is a crucial step in our proof.
###### Lemma 6.3
Recall that $`V=V^{}\times R^2`$. Let $`(V,J,\mu )`$ as above and $`WV`$ be as above and $`V_c`$ the compact set in $`V`$. Let $`\overline{V}=D\times V`$, $`\overline{W}=D\times W`$, and $`\overline{V}_c=D\times V_c`$. Let $`Y=\alpha (S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`Y_i=\alpha (S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`D=S^{1+}S^1`$. There exist constant $`c_0`$ such that every $`J`$holomorphic map $`v`$ of the half unit disc $`D^+`$ to the $`D\times V^{}\times R^2`$ with its boundary $`v((1,1))(S^{1\pm })\times F(\times R\times S_i)\overline{W},i=1,..,4`$ has
$$area(v(D^+))c_0l^2(v(^{}D^+)).$$
(6.16)
here $`^{}D^+=D[1,1]`$ and $`l(v(^{}D^+))=length(v(^{}D^+))`$.
Proof. We first assume that $`\epsilon `$ in section 3.1 is small enough. Let $`l_0`$ is a constant small enough. If $`length(^{}D^+)l_0`$, then Lemma 6.3 holds. If $`length(^{}D^+)l_0`$ and $`v(D^+)D\times V_{ij}^{}\times R^2`$, then Lemma6.3 reduces to Lemma6.1. If $`length(^{}D^+)l_0`$ and $`v(D^+)\overline{}D\times V_{ij}^{}\times R^2`$, then Lemma6.2 imples $`area(v)\tau _0>100\pi R(\epsilon )^2`$, this is a contradiction. Therefore we proved the lemma.
###### Proposition 6.3
Let $`(V,J,\mu )`$ and $`WV`$ be as in section 4 and $`V_K`$ the compact part of $`V`$. Let $`\overline{V}`$, $`\overline{V}_K`$ and $`\overline{W}`$ as section 5.1. There exist constants $`\epsilon _0`$ (depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^\alpha `$ norm of $`J`$) and $`C`$(depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^{k+\alpha }`$ norm of $`J`$) such that every $`J`$holomorphic map of the half unit disc $`D^+`$ to the $`D\times V^{}\times R^2`$ with its boundary $`v((1,1))(S^{1\pm })\times F(\times R\times S_i)\overline{W},i=1,..,4`$ has its derivatives up to order $`k+1+\alpha `$ on $`D_{\frac{1}{2}}^+(0)`$ bounded by $`C`$.
Proof. One uses Lemma 6.3 and Gromov’s proof on Schwartz lemma to yield proposition 6.3.
### 6.3 Removal singularity of $`J`$curves
In our proof we need another crucial tools, i.e., Gromov’s removal singularity theorem. We first consider the case without boundary.
###### Proposition 6.4
Let $`(V,J,\mu )`$ be as in section 4 and $`V_K`$ the compact part of $`V`$. If $`v:D\{0\}V_K`$ be a $`J`$holomorphic disk with bounded energy and bounded image, then $`v`$ extends to a $`J`$holomorphic map from the unit disc $`D`$ to $`V_K`$.
For a proof, see.
Now we consider the Gromov’s removal singularity theorem for $`J`$holomorphic map with boundary in a closed Lagrangian submanifold as in .
###### Proposition 6.5
Let $`(V,J,\mu )`$ as above and $`LV`$ be a closed Lagrangian submanifold and $`V_K`$ one compact part of $`V`$. If $`v:(D^+\{0\},^{\prime \prime }D^+\{0\})(V_K,L)`$ be a $`J`$holomorphic half-disk with bounded energy and bounded image, then $`v`$ extends to a $`J`$holomorphic map from the half unit disc $`(D^+,^{\prime \prime }D^+)`$ to $`(V_K,L)`$.
For a proof see .
###### Proposition 6.6
Let $`(V,J,\mu )`$ and $`WV`$ be as in section 4 and $`V_c`$ the compact set in $`V`$. Let $`\overline{V}=D\times V`$, $`\overline{W}=D\times W`$, and $`\overline{V}_c=D\times V_c`$. Then every $`J`$holomorphic map $`v`$ of the half unit disc $`D^+\{0\}`$ to the $`\overline{V}`$ with center in $`\overline{V}_c`$ and its boundary $`v((1,1)\{0\})(S^{1\pm })\times F(\times [K,K]\times S_i)\overline{W}`$ and
$$area(v(D^+\{0\}))E$$
(6.17)
extends to a $`J`$holomorphic map $`\stackrel{~}{v}:(D^+,^{\prime \prime }D)(\overline{V}_c,\overline{W})`$.
Proof. This is ordinary Gromov’s removal singularity theorem by $`K`$assumption.
### 6.4 $`C^0`$Convergence Theorem
We now recall that the well-known Gromov’s compactness theorem for cusp’s curves for the compact symplectic manifolds with closed Lagrangian submanifolds in it. For reader’s convenience, we first recall the “weak-convergence” for closed curves.
Cusp-curves. Take a system of disjoint simple closed curves $`\gamma _i`$ in a closed surface $`S`$ for $`i=1,\mathrm{},k`$, and denote by $`S^0`$ the surface obtained from $`S_{i=1}^k\gamma _i`$. Denote by $`\overline{S}`$ the space obtained from $`S`$ by shrinking every $`\gamma _i`$ to a single point and observe the obvious map $`\alpha :S^0\overline{S}`$ gluing pairs of points $`s_i^{}`$ and $`s_i^{\prime \prime }`$ in $`S^0`$, such that $`\overline{s}_i=\alpha (s_i^{})=\alpha (s_i^{\prime \prime })\overline{S}`$ are singular (or cuspidal) points in $`\overline{S}`$(see).
An almost complex structure in $`\overline{S}`$ by definition is that in $`S^0`$. A continuous map $`\beta :\overline{S}V`$ is called a (parametrized $`J`$holomorphic) cusp-curve in $`V`$ if the composed map $`\beta \alpha :S^0V`$ is holomorphic.
Weak convergence. A sequence of closed $`J`$curves $`C_jV`$ is said to weakly converge to a cusp-curve $`\overline{C}V`$ if the following four conditions are satisfied
(i) all curves $`C_j`$ are parametrized by a fixed surface $`S`$ whose almost complex structure depends on $`j`$, say $`C_j=f_j(S)`$ for some holomorphic maps
$$f_j:(S,J_j)(V,J)$$
(ii) There are disjoint simple closed curves $`\gamma _iS`$, $`i=1,\mathrm{},k`$, such that $`\overline{C}=\overline{f}(\overline{S})`$ for a map $`\overline{f}:\overline{S}V`$ which is holomorphic for some almost complex structure $`\overline{J}`$ on $`\overline{S}`$.
(iii) The structures $`J_j`$ uniformly $`C^{\mathrm{}}`$converge to $`\overline{J}`$ on compact subsets in $`S_{i=1}^k\gamma _i`$.
(iv) The maps $`f_j`$ uniformly $`C^{\mathrm{}}`$converge to $`\overline{f}`$ on compact subsets in $`S_{i=1}^k\gamma _i`$. Moreover, $`f_j`$ uniformly $`C^0`$converge on entire $`S`$ to the composed map $`S\overline{S}\stackrel{\overline{f}}{}V`$. Furthermore,
$$Area_\mu f_j(S)Area_\mu \overline{f}(\overline{S})forj\mathrm{},$$
where $`\mu `$ is a Riemannian metric in $`V`$ and where the area is counted with the geometric multiplicity(see).
Gromov’s Compactness theorem for closed curves. Let $`C_j`$ be a sequence of closed $`J`$curves of a fixed genus in a compact manifold $`(V,J,\mu ).`$ If the areas of $`C_j`$ are uniformly bounded,
$$Area_\mu A,j=1,..,$$
then some subsequence weakly converges to a cusp-curve $`\overline{C}`$ in $`V`$.
Cusp-curves with boundary. Let $`T`$ be a compact complex manifold with boundary of dimension $`1`$(i.e., it has an atlas of holomorphic charts onto open subsets of $`C`$ or of a closed half plane). Its double is a compact Riemann surface $`S`$ with a natureal anti-holomorphic involution $`\tau `$ which exchanges $`T`$ and $`ST`$ while fixing the boundary $`T`$. IF$`f:TV`$ is a continous map, holomorphic in the interior of $`T`$, it is convenient to extend $`f`$ to $`S`$ by
$$f=f\tau $$
Take a totally real submanifold $`W(V,J)`$ and consider compact holomorphic curves $`CV`$ with boundaries, $`(\overline{C},\overline{C})(V,W)`$, which are, topologically speaking, obtained by shrinking to points some (short) closed loops in $`C`$ and also some (short) segments in $`C`$ between boundary points. This is seen by looking on the double $`C_CC`$.
Gromov’s Compactness theorem for curves with boundary. Let $`V`$ be a closed Riemannian manifold, $`W`$ a totally real closed submanifold of $`V`$. Let $`C_j`$ be a sequence of $`J`$curves with boundary in $`W`$ of a fixed genus in a compact manifold $`(V,J,\mu )`$. If the areas of $`C_j`$ are uniformly bounded,
$$Area_\mu A,j=1,..,$$
then some subsequence weakly converges to a cusp-curve $`\overline{C}`$ in $`V`$.
The proofs of Gromov’s compactness theorem can found in \[al, 13\]. In our case the Lagrangian submanifold $`W`$ is not compact, Gromov’s compactness theorem can not be applied directly but its proof is still effective since the $`W`$ has the special geometry. In the following we modify Gromov’s proof to prove the $`C^0`$compactness theorem in our case.
Now we state the $`C^0`$convergence theorem in our case.
###### Theorem 6.1
Let $`(V,J,\omega ,\mu )`$ and $`W`$ as in section4. Let $`C_j`$ be a sequence of $`\overline{J}_\delta `$holomorphic section $`v_j=(id,((a_j^1,u_j^1),(a_j^2,u_j^2),f_j)):DD\times V`$ with $`v_j:DD\times W`$ and $`v_j(1)=(1,p)D\times W`$. constructed from section 4. Then the areas of $`C_j`$ are uniformly bounded,i.e.,
$$Area_\mu (C_j)A,j=1,..,$$
and some subsequence weakly converges to a cusp-section $`\overline{C}`$ in $`V`$(see).
Proof. We follow the proofs in . Write $`v_j=(id,(a_j^1,u_j^1),(a_j^2,u_j^2),f_j))`$ then $`|a_{ij}^2|a_0`$ by the ordinary Monotone inequality of minimal surface without boundary, see following Proposition 7.1. Similarly $`|f_j|R_1`$ by using the fact $`f_j(D)`$ is bounded in $`B_{R_1}(0)`$ and $`_D|f_j|4\pi R^2`$ via monotone inequality for minimal surfaces. So, we assume that $`v_j(D)V_c`$ for a compact set $`V_c`$.
1. Removal of a net.
1a. Let $`\overline{V}=D\times V`$ and $`v_j`$ be regular curves. First we study induced metrics $`\mu _j`$ in $`v_j`$. We apply the ordinary monotone inequality for minimal surfaces without boundary to small concentric balls $`B_\epsilon (A_j,\mu _j)`$ for $`0<\epsilon \epsilon _0`$ and conclude by the standard argument to the inequality
$$Area(B_\epsilon )\epsilon ^2,for\epsilon \epsilon _0;$$
Using this we easily find a interior $`\epsilon `$net $`F_j(v_j,\mu _j)`$ containing $`N`$ points for a fixed integer $`N=(\overline{V},\overline{J},\mu )`$, such that every topological annulus $`Av_jF_j`$ satisfies
$$Diam_\mu A10length_\mu A.$$
(6.18)
Furthermore, let $`A`$ be conformally equivalent to the cylinder $`S^1\times [0,l]`$ where $`S^1`$ is the circle of the unit length, and let $`S_t^1A`$ be the curve in $`A`$ corresponding to the circle $`S^1\times t`$for $`t[0,l]`$. Then obviously
$$_a^b(lengthS_t^1)^2𝑑tArea(A)C_5.$$
(6.19)
for all $`[a,b][0,l]`$. Hence, the annulus $`A_tA`$ between the curves $`S_t^1`$ and $`S_{lt}^1`$ satisfies
$$diam_\mu A_t20(\frac{C_5}{t})$$
(6.20)
for all $`t[0,l]`$.
$`1b`$. We consider the sets $`v_j((S^{1\pm })\times F(W^{}\times I_i^\pm )),i=0,1`$. By the construction of Gromov’s figure eight, there exists a finite components, denote it by
$$v_j((S^{1\pm })\times F(\times R\times I_i^\pm ))=\{\overline{\gamma }_{ij}^k\},i=0,1.$$
(6.21)
we choose one point in $`\overline{\gamma }_{ij}^k`$ as a boundary puncture point in $`v_j`$ for each $`i,k`$.
Consider the concentric $`\epsilon `$ half-disks or quadrature $`B_\epsilon (p)`$ with center $`p`$ on $`\overline{\gamma }_{ij}^k`$, then
$$Area(B_\epsilon (p))\tau _0$$
(6.22)
Since $`Area(v_j)E_0`$, there exists a uniform finite puncture points.
So, we find a boundary net $`G_jv_j`$ containing $`N_1`$ points for a fixed integer $`N_1(\overline{V},\overline{J},\mu )`$, such that every topological quadrature or half annulus $`Bv_j\{F_j,G_j\}`$ satisfies
$`^{\prime \prime }B=B\overline{W}(S^{1\pm })\times F(\times R\times S_i),i=1,2,3,4.`$ (6.23)
2. Poincare’s metrics. 2a. Now, let $`\mu _j^{}`$ be a metric of constant curvature $`1`$ in $`v_j(D)F_jG_j`$ conformally equivalent to $`\mu _j`$. Then for every $`\mu _j^{}`$ball $`B_\rho `$ in $`v_jF_jG_j`$ of radius $`\rho 0.1`$, there exists an annulus $`A`$ contained in $`v_jF_jG_j`$ such that $`B_\rho A_t`$ for $`t=0.01|log|`$(see Lemma 3.2.2in \[al, chVIII\]). This implies with $`(6.3)`$ the uniform continuity of the (inclusion) maps $`(v_jF_j,\mu _j^{})(\overline{V},\overline{\mu })`$, and hence a uniform bound on the $`r^{th}`$ order differentials for every $`r=0,1,2,\mathrm{}`$.
2b. Similarly, for every $`\mu _j^{}`$half ball $`B_\rho ^+`$ in $`v_jF_jG_j`$ of radius $`\rho 0.1`$, there exists a half annulus or quadrature $`B`$ contained in $`v_jF_jG_j`$ such that $`B_\rho ^+B`$ with
$`^{\prime \prime }B=B\overline{W}(S^{1\pm })\times F(\times R\times S_i),i=1,2,3,4.`$ (6.24)
Then, by Gromov’s Schwartz Lemma, i.e., Proposition 6.1-6.3 implies the uniform bound on the $`r^{th}`$ order differentials for every $`r=0,1,2,\mathrm{}`$.
3. Convergence of metrics. Next, by the standard (and obvious ) properties of hyperbolic surfaces there is a subsequence(see), which is still denoted by $`v_j`$, such that
$`(a)`$. There exist $`k`$ closed geodesics or geodesic arcs with boundaries in $`v_jF_j`$, say
$$\gamma _i^j(v_jF_j,\mu _j^{}),i=1,\mathrm{},k,j=1,2,\mathrm{},$$
whose $`\mu _j^{}`$length converges to zero as $`j\mathrm{}`$, where $`k`$ is a fixed integer.
$`(b)`$. There exist $`k`$ closed curves or geodesic arcs with boundaries in $`S`$ of a fixed surface, say $`\gamma _j`$ in $`S`$, and an almost complex structure $`\overline{J}`$ on the corresponding (singular) surface $`\overline{S}`$, such that the almost complex structure $`J_j`$ on $`v_jF_j`$ induced from $`(V,J)`$ $`C^{\mathrm{}}`$converge to $`\overline{J}`$ outside $`_{j=1}^k\gamma _j`$. Namely, there exist continuous maps $`g^j:v_j\overline{S}`$ which are homeomorphisms outside the geodesics $`\gamma _i^j`$, which pinch these geodesics to the corresponding singular points of $`\overline{S}`$(that are the images of $`\gamma _i`$) and which send $`F_j`$ to a fixed subset $`F`$ in the nonsingular locus of $`\overline{S}`$. Now, the convergence $`J_j\overline{J}`$ is understood as the uniform $`C^{\mathrm{}}`$convergence $`g_{}^j(J_j)\overline{J}`$ on the compact subsets in the non-singular locus $`\overline{S}^{}`$ of $`\overline{S}`$ which is identified with $`S_{i=1}^k\gamma _i`$.
4. $`C^0`$interior convergence. The limit cusp-curve $`\overline{v}:\overline{S}^{}\overline{V}`$, that is a holomorphic map which is constructed by first taking the maps
$$\overline{v}_j=(g_j)^1:S_{i=1}^k\gamma _i\overline{V}$$
Near the nodes of $`\overline{S}`$ including interior nodes and boundary nodes, by the properties of hyperbolic metric $`\mu ^{}`$ on $`\overline{S}`$, the neighbourhoods of interior nodes are corresponding to the annulis of the geodesic cycles. By the reparametrization of $`v_j`$, called $`\overline{v}_j`$ which is defined on $`S`$ and extends the maps $`\overline{v}_j:SS_jV`$(see). Now let $`\{z_i|i=1,\mathrm{},n\}`$ be the interior nodes of $`\overline{S}`$. Then the arguments in yield the $`C^0`$interir convergece near $`z_i`$.
5. $`C^0`$boundary convergence. Now it is possible that the boundary of the cusp curve $`\overline{v}`$ does not remain in $`\overline{W}`$. Write $`\overline{v}(z)=(h,((a_1,u_1),(a_2,u_2)),f)(z)`$, here $`h(z)=z`$ or $`h(z)z_i`$, $`i=1,\mathrm{},n`$, $`z_i`$ is cusp-point or bubble point. We can assume that $`\overline{p}=(1,p)\overline{v}_n`$ is a puncture boundary point. Let $`\overline{v}_1`$ be the component of $`\overline{v}`$ which through the point $`\overline{p}`$. Let $`D=\{z|z=re^{i\theta },0r1,0\theta 2\pi \}`$. We assume that $`\overline{v}_1:D\{e^{i\theta _i}\}_{i=1}^kV_c`$, here $`e^{i\theta _i}`$ is node or puncture point. Near $`e^{i\theta _i}`$, we take a small disk $`D_i`$ in $`D`$ containing only one puncture or node point $`e^{i\theta _i}`$. By the reparametrization and the convergence procedure, we can assume that $`\overline{v}_{1i}=(\overline{v}_1|D_i)`$ as a map from $`D^+\{0\}V_c`$ with $`\overline{v}_1([1,1]\{0\})S^1\times F(W^{}\times S^1)`$ and $`area(\overline{v}_{1i})a_0`$, $`a_0`$ small enough. Since $`Area(\overline{v}_{1i})a_0`$, there exist curves $`c_k`$ near $`0`$ such that $`l(\overline{v}_{1i}(c_k))\delta _1`$. By the construction of convergence, we can assume that $`l(\overline{v}_n(c_k))2\delta _1`$. If $`\overline{v}_{1i}(c_k)(S^1)\times F(\times [N_0,N_0]\times S^1)`$, we have $`\overline{v}_n(c_k)(S^1)\times F(\times [2N_0,2N_0]\times S^1)`$ for $`n`$ large enough. Now $`\overline{v}_n(c_k)`$ cuts $`\overline{v}_n(D)`$ as two parts, one part corresponds to $`\overline{v}_{1i}`$, say $`\overline{u}_n(D)`$. Then $`area(\overline{u}_n(D))=area(h_{n1})+|\mathrm{\Psi }^{}(u_{n2}(c_k^1))\mathrm{\Psi }^{}(u_{n2}(c_k^2))|`$, here $`c_k=\{c_k^1,c_k^2\}`$. Then by the proof of Lemma6.1-6.3, we know that $`\overline{u}_n(Dc_k)(S^1)\times F(\times [100N_0,100N_0]\times S^1)`$. So, $`\overline{v}_{1i}([1,1]\{0\})S^1\times F(\times [100N_0,100N_0]\times S^1)`$. By proposition 6.6, one singularity of $`\overline{v}_1`$ is deleted. We repeat this procedure, we proved that $`\overline{v}_1`$ is extended to whole $`D`$. So, the boundary node or puncture points of $`\overline{v}`$ are removed. Then by choosing the sub-sub-sequences of $`\mu _j^{}`$ and $`\overline{v}_j`$, we know that $`\overline{v}_j`$ converges to $`\overline{v}`$ in $`C^0`$ near the boundary node or puncture point. This proved the $`C^0`$boundary convergence. Since $`\overline{v}_j(1)=\overline{p}`$, $`\overline{p}\overline{v}(D)`$, $`\overline{v}(D)\overline{W}`$.
6. Convergence of area. Finally by the $`C^0`$convergence and $`area(v_j)=_Dv_j^{}\overline{\omega }`$, one easily deduces
$$area(v(S))=\underset{j\mathrm{}}{lim}(v_j(S_j)).$$
### 6.5 Bounded image of $`J`$holomorphic curves in $`W`$
###### Proposition 6.7
Let $`v`$ be the solutions of equations (4.16), then
$$d_W(p,v(D^2))=max\{d_W(p,q)|qf(D^2)\}d_0<+\mathrm{}$$
Proof. It follows directly from Gromov’s $`C^0`$convergence theorem.
## 7 Proof of Theorem 1.1
###### Proposition 7.1
If $`J`$holomorphic curves $`C\overline{V}`$ with boundary
$$CD^2\times ([0,\epsilon ]\times \mathrm{\Sigma })\times ([0,\epsilon ]\times \mathrm{\Sigma })\times R^2$$
and
$$C(D^2\times (\{3\}\times \mathrm{\Sigma })\times (R\times \mathrm{\Sigma })\times R^2)\mathrm{}$$
or
$$C(D^2\times (R\times \mathrm{\Sigma })\times (\{3\}\times \mathrm{\Sigma })\times R^2)\mathrm{}$$
Then
$$area(C)2l_0.$$
Proof. It is obvious by monotone inequality argument for minimal surfaces.
###### Note 7.1
we first observe that any $`J`$holomorphic curves with boundary in $`R^+\times \mathrm{\Sigma }`$ meet the hypersurface $`\{3\}\times \mathrm{\Sigma }`$ has energy at least $`2l_0`$, so we take $`\epsilon `$ small enough such that the Gromov’s figure eight contained in $`B_{R(\epsilon )}C`$ for $`\epsilon `$ small enough and the energy of solutions in section 4 is smaller than $`l_0`$. we specify the constant $`a_0`$, $`\epsilon `$ in section 3.1-3 such that the above conditions satisfied.
###### Theorem 7.1
There exists a non-constant $`J`$holomorphic map $`u:(D,D)(V^{}\times C,W)`$ with $`E(u)4\pi R(\epsilon )^2`$ for $`\epsilon `$ small enough such that $`4\pi R(\epsilon )^2l_0`$.
Proof. By Proposition 5.1, we know that the image $`\overline{v}(D)`$ of solutions of equations (4.19) remains a bounded or compact part of the non-compact Lagrangian submanifold $`W`$. Then, all arguments in for the case $`W`$ is closed in $`S\mathrm{\Sigma }\times S\mathrm{\Sigma }\times R^2`$ can be extended to our case, especially Gromov’s $`C^0`$converngence theorem holds. But the results in section 4 shows the solutions of equations (4.19) must denegerate to a cusp curves, i.e., we obtain a Sacks-Uhlenbeck-Gromov’s bubble, i.e., $`J`$holomorphic sphere or disk with boundary in $`W`$, the exactness of $`\omega `$ rules out the possibility of $`J`$holomorphic sphere. For the more detail, see the proof of Theorem 2.3.B in .
Proof of Theorem 1.1. If $`(\mathrm{\Sigma },\lambda )`$ has not closed Reeb orbit, then we can construct a Lagrangian submanifold $`W`$ in $`V=V^{}\times C`$, see section 3. Then as in section 4, we construct an anti-holomorphic section $`c`$ and for large vector $`cC`$ we know that the nonlinear Fredholm section or Cauchy-Riemann section has no solution, this implies that the section is non-proper, see section 4. The non-properness of the section and the Gromov’s compactness theorem in section 6 implies the existences of the cusp-curves. So, we must have the $`J`$holomorphic sphere or $`J`$holomorphic disk with bounadry in $`W`$. Since the symplectic manifold $`V`$ is an exact symplectic mainifold and $`W`$ is an exact Lagrangian submanifold in $`V`$, by Stokes formula, we know that the possibility of $`J`$holomorphic sphere or disk elimitated. So our priori assumption does not hold which implies the contact maifold $`(\mathrm{\Sigma },\lambda )`$ has at least closed Reeb orbit. This finishes the proof of Theorem 1.1.
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# Natural Chaotic Inflation in Supergravity
## I Introduction
The inflationary expansion of the early universe is the most attractive ingredient in modern cosmology. This is not only because it naturally solves the longstanding problems in cosmology, that is the horizon and flatness problems, but also because it accounts for the origin of density fluctuations as observed by the Comic Background Explorer(COBE) satellite . Among various types of inflation models proposed so far, chaotic inflation model is the most attractive since it can realize an inflationary expansion even in the presence of large quantum fluctuations at the Planck time. In fact, many authors have used the chaotic inflation model to discuss a number of interesting phenomena such as preheating , superheavy particle production , and primordial gravitational waves in the inflationary cosmology .
On the other hand, supersymmetry (SUSY) is widely discussed as the most interesting candidate for the physics beyond the standard model since it ensures the stability of the large hierarchy between the electroweak and the Planck scales against radiative corrections. This kind of stability is also very important to keep the flatness of inflaton potential at the quantum level. Therefore, it is quite natural to consider the inflation model in the framework of supergravity.
However, the above two ideas, i.e. chaotic inflation and supergravity, have not been naturally realized simultaneously. The main reason is that the minimal supergravity potential has an exponential factor, $`\mathrm{exp}(\frac{\phi _i^{}\phi _i}{M_G^2})`$, which prevents any scalar fields $`\phi _i`$ from having values larger than the gravitational scale $`M_G2.4\times 10^{18}`$GeV. However, the inflaton $`\phi `$ is supposed to have a value much larger than $`M_G`$ at the Planck time to cause the chaotic inflation. Thus, the above effect makes it very difficult to incorporate the chaotic inflation in the framework of supergravity. In fact, all of the existing models for chaotic inflation use rather specific Kähler potential, and one needs a fine tuning in the Kähler potential since there is no symmetry reason for having such specific forms of Kähler potentials. Thus, it is very important to find a natural chaotic inflation model without any fine tuning.
In this letter, we propose a natural chaotic inflation model where the form of Kähler potential is determined by a symmetry. With this Kähler potential the inflaton $`\phi `$ may have a large value $`\phi M_G`$ to begin the chaotic inflation. Our models, in fact, need two small parameters for successful inflation. However, we emphasize that the smallness of these parameters is justified by symmetries and hence the model is natural in ’t Hooft’s sense .
The existence of a natural chaotic inflation model may open a new branch of inflation-model building in supergravity, since most of the model building in supergravity has been concentrated on other types of inflation models (e.g. hybrid inflation model etc. ). Furthermore, we consider that future astrophysical observations will be able to select types of inflation models.
Our model is based on the Nambu-Goldstone-like shift symmetry of the inflaton chiral multiplet $`\mathrm{\Phi }(x,\theta )`$. Namely, we assume that the Kähler potential $`K(\mathrm{\Phi },\mathrm{\Phi }^{})`$ is invariant under the shift of $`\mathrm{\Phi }`$,
$$\mathrm{\Phi }\mathrm{\Phi }+iCM_G,$$
(1)
where $`C`$ is a dimensionless real parameter. Thus, the Kähler potential is a function of $`\mathrm{\Phi }+\mathrm{\Phi }^{}`$, $`K(\mathrm{\Phi },\mathrm{\Phi }^{})=K(\mathrm{\Phi }+\mathrm{\Phi }^{})`$. It is now clear that the supergravity effect $`e^{K(\mathrm{\Phi }+\mathrm{\Phi }^{})}`$ discussed above does not prevent the imaginary part of the scalar components of $`\mathrm{\Phi }`$ from having a larger value than $`M_G`$. We identify it with the inflaton field $`\phi `$. We also stress that the present model overcomes the so-called $`\eta `$ problem and it is an alternative to other inflation models such as D-term inflation models and running inflaton mass models . However, as long as the shift symmetry is exact, the inflaton $`\phi `$ never has a potential and hence it never causes the inflation. Therefore, we have to introduce a small breaking term of the shift symmetry in the theory. The simplest choice is to introduce a small mass term for $`\mathrm{\Phi }`$ in the superpotential,
$$W=m\mathrm{\Phi }^2.$$
(2)
Then, we have the potential,
$$V=e^K\left\{\left(\frac{^2K}{\mathrm{\Phi }\mathrm{\Phi }^{}}\right)^1D_\mathrm{\Phi }WD_\mathrm{\Phi }^{}W^{}3|W|^2\right\},$$
(3)
with
$$D_\mathrm{\Phi }W=\frac{W}{\mathrm{\Phi }}+\frac{K(\mathrm{\Phi }+\mathrm{\Phi }^{})}{\mathrm{\Phi }}W.$$
(4)
Here, $`\mathrm{\Phi }`$ denotes the scalar component of the superfield $`\mathrm{\Phi }`$ and we have set $`M_G`$ to be unity. We easily see that $`V\mathrm{}`$ as $`|\phi |\mathrm{}`$ with $`\mathrm{\Phi }+\mathrm{\Phi }^{}=0`$ and the chaotic inflation does not take place, where $`\phi =i(\mathrm{\Phi }\mathrm{\Phi }^{})/\sqrt{2}`$.
In this letter, we propose instead the following small mass term in the superpotential introducing a new chiral multiplet $`X(x,\theta )`$,
$$W=mX\mathrm{\Phi }.$$
(5)
Notice that the present model possesses $`U(1)_\mathrm{R}`$ symmetry under which
$`X(\theta )`$ $``$ $`e^{2i\alpha }X(\theta e^{i\alpha }),`$ (6)
$`\mathrm{\Phi }(\theta )`$ $``$ $`\mathrm{\Phi }(\theta e^{i\alpha }),`$ (7)
and $`Z_2`$ symmetry under which
$`X(\theta )`$ $``$ $`X(\theta e^{i\alpha }),`$ (8)
$`\mathrm{\Phi }(\theta )`$ $``$ $`\mathrm{\Phi }(\theta e^{i\alpha }).`$ (9)
The above superpotential is not invariant under the shift symmetry of $`\mathrm{\Phi }`$. However, we should stress that the present model is completely natural in ’t Hooft’s sense , since we have an enhanced symmetry (the shift symmetry) in the limit $`m0`$. That is, we consider that the small parameter $`m`$ is originated from small breaking of the shift symmetry in a more fundamental theory. We consider that as long as $`m𝒪(1)`$, the corrections from the breaking term eq.(5) to the Kähler potential are negligibly small.<sup>*</sup><sup>*</sup>*The Kähler potential may have also the induced breaking terms such as $`K|m\mathrm{\Phi }|^2+\mathrm{}`$. However, these breaking terms are negligible in the present analysis as long as $`|\phi |m^1`$. Then, we assume that the Kähler potential has the shift symmetry eq.(1) and the above $`U(1)_\mathrm{R}\times Z_2`$ symmetry neglecting the breaking effects,
$$K(\mathrm{\Phi },\mathrm{\Phi }^{},X,X^{})=K[(\mathrm{\Phi }+\mathrm{\Phi }^{})^2,XX^{}].$$
(10)
In the following analysis we take, for simplicity,
$$K=\frac{1}{2}(\mathrm{\Phi }+\mathrm{\Phi }^{})^2+XX^{}+\mathrm{}.$$
(11)
## II Dynamics of inflation
The Lagrangian density $`L(\mathrm{\Phi },X)`$ is now given by
$$L(\mathrm{\Phi },X)=_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }^{}+_\mu X^\mu X^{}V(\mathrm{\Phi },X),$$
(12)
with the potential $`V(\mathrm{\Phi },X)`$ given by
$$V(\mathrm{\Phi },X)=m^2e^K\left[|\mathrm{\Phi }|^2(1+|X|^4)+|X|^2\left\{1|\mathrm{\Phi }|^2+(\mathrm{\Phi }+\mathrm{\Phi }^{})^2(1+|\mathrm{\Phi }|^2)\right\}\right],$$
(13)
where we have neglected higher order terms in the Kähler potential eq.(11) whose effects will be discussed later. Here, $`X`$ denotes also the scalar component of the superfield $`X`$. Now, we decompose the complex scalar field $`\mathrm{\Phi }`$ into two real scalar fields as,
$$\mathrm{\Phi }=\frac{1}{\sqrt{2}}(\eta +i\phi ).$$
(14)
Then, the Lagrangian density $`L(\eta ,\phi ,X)`$ is given by
$$L(\eta ,\phi ,X)=\frac{1}{2}_\mu \eta ^\mu \eta +\frac{1}{2}_\mu \phi ^\mu \phi +_\mu X^\mu X^{}V(\eta ,\phi ,X),$$
(15)
with the potential $`V(\eta ,\phi ,X)`$ given by
$`V(\eta ,\phi ,X)=m^2\mathrm{exp}\left(\eta ^2+|X|^2\right)`$ (17)
$`\times \left[{\displaystyle \frac{1}{2}}(\eta ^2+\phi ^2)(1+|X|^4)+|X|^2\left\{1{\displaystyle \frac{1}{2}}(\eta ^2+\phi ^2)+2\eta ^2\left(1+{\displaystyle \frac{1}{2}}(\eta ^2+\phi ^2)\right)\right\}\right].`$
Note that $`\eta `$ and $`|X|`$ should be taken as $`|\eta |,|X|𝒪(1)`$ because of the presence of $`e^K`$ factor. On the other hand, $`\phi `$ can take a value much larger than $`𝒪(1)`$ since $`e^K`$ does not contain $`\phi `$. For $`\eta ,|X|𝒪(1)`$, we can rewrite the potential as
$$V(\eta ,\phi ,X)\frac{1}{2}m^2\phi ^2(1+\eta ^2)+m^2|X|^2.$$
(18)
At around the Planck time, we may have a region where
$$\dot{\phi }^2(\phi )^2V(\phi )1.(\text{initial chaotic situation})$$
(19)
Here the dot represents the time derivative. In this region the classical description of the $`\phi `$ field dynamics is feasible because of $`|\phi |𝒪(1)`$ though quantum fluctuations are $`\delta \phi 𝒪(1)`$. Then, as the universe expands, the potential energy dominates and the universe begins inflation.
Since the initial values of the inflaton $`\phi (0)`$ is determined so that $`V(\phi (0))\frac{1}{2}m^2\phi (0)^21`$, $`\phi (0)m^11`$. (Notice that one has only to demand $`\phi (0)15.0`$ in order to solve the flatness and horizon problems. ) For such large $`\phi `$ the effective mass of $`\eta `$ becomes much larger than $`m`$ and hence it quickly settles down to $`\eta =0`$. On the other hand, the $`X`$ field has a relatively light mass $`m`$ and slowly rolls down toward the origin ($`X=0`$). With $`\eta =0`$, the potential eq.(18) is written as
$$V(\phi ,X)\frac{1}{2}m^2\phi ^2+m^2|X|^2.$$
(20)
Since $`\phi 1`$ and $`|X|<1`$, the $`\phi `$ field dominates the potential and the chaotic inflation takes place. The Hubble parameter is given by
$$H\frac{m\phi }{\sqrt{3}}.$$
(21)
During the inflation both $`\phi `$ and $`X`$ satisfy the slow roll condition ($`|\frac{V^{\prime \prime }}{V}|1,\frac{1}{2}|\frac{V^{}}{V}|^21`$ where the dash represents the derivative of $`\phi `$ or $`X`$) and hence the time evolutions are described by
$`3H{\displaystyle \frac{d\phi }{dt}}`$ $``$ $`m^2\phi ,`$ (22)
$`3H{\displaystyle \frac{dX}{dt}}`$ $``$ $`m^2X.`$ (23)
Here and hereafter, we assume that $`X`$ is real and positive making use of the freedom of the phase choice. From the above equations we obtain
$$\left(\frac{X}{X(0)}\right)\left(\frac{\phi }{\phi (0)}\right),$$
(24)
where $`\phi (0)`$ and $`X(0)`$ are the initial values of $`\phi `$ and $`X`$ fields. Therefore, $`X`$ decreases faster than $`\phi `$. At the end of the inflation, i.e. $`\phi 1`$ ($`|\frac{V^{\prime \prime }}{V}|\frac{1}{2}|\frac{V^{}}{V}|^21`$), $`X`$ is given by
$$Xm,$$
(25)
where we have used $`X(0)1`$ and $`\phi (0)m^1`$. We see that the $`X`$ field becomes much smaller than $`1`$ ($`m10^5`$ as shown below). The density fluctuations produced by this chaotic inflation is estimated as
$$\frac{\delta \rho }{\rho }\frac{1}{5\sqrt{3}\pi }\frac{m}{2\sqrt{2}}(\phi ^2+X^2).$$
(26)
Since $`X\phi `$, the amplitude of the density fluctuations is determined only by the $`\phi `$ field and the normalization at the COBE scale ($`\delta \rho /\rho 2\times 10^5`$ for $`\phi _{\mathrm{COBE}}14`$ ) gives The spectral index $`n_s0.96`$ for $`\phi _{\mathrm{COBE}}14`$.
$$m10^{13}\mathrm{GeV}.$$
(27)
After the inflation ends, an inflaton field $`\phi `$ begins to oscillate and its successive decays cause reheating of the universe. In the present model the reheating takes place efficiently if we introduce the following superpotential:
$$W=\lambda XH\overline{H},$$
(28)
where $`H`$ and $`\overline{H}`$ are a pair of Higgs doublets whose $`R`$-charge are assumed to be zero and $`\lambda `$ is a constant.$`D_XW`$ is changed to $`D_XW=(m\mathrm{\Phi }+\lambda H\overline{H})(1+|X|^2)`$. $`|H|`$ and $`|\overline{H}|`$ take values $`𝒪(1)`$ due to the factor of $`e^{K(H,\overline{H})}`$ as the $`X`$ field. Therefore, the $`m\mathrm{\Phi }`$ term dominates $`D_XW`$ unless $`\lambda 𝒪(1)`$, since $`|\mathrm{\Phi }(0)|m^1`$ at the beginning of the universe, and the chaotic inflation begins. Once the inflation takes place, $`H`$ and $`\overline{H}`$ acquire masses of the order the Hubble scale and rapidly go to zero. Thus, the above superpotential eq.(28) does not affect the dynamics of the inflation. Then, we have the coupling of the inflaton $`\phi `$ to the Higgs doublets as
$$L\lambda m\phi H\overline{H},$$
(29)
which gives the reheating temperature
$$T_R10^9\mathrm{GeV}\left(\frac{\lambda }{10^5}\right)\left(\frac{m}{10^{13}\mathrm{GeV}}\right)^{1/2}.$$
(30)
In order to avoid the overproduction of gravitinos, the reheating temperature $`T_R`$ must be lower than $`10^9`$GeV , which requires the small coupling $`\lambda 10^5`$. The small coupling $`\lambda `$ is naturally understood in ’t Hooft’s sense provided that $`H\overline{H}`$ is even under the $`Z_2`$ symmetry in eq.(9).
So far we have taken the minimal Kähler potential and neglected higher order terms like $`(\mathrm{\Phi }+\mathrm{\Phi }^{})^4,|X|^4`$, and $`\mathrm{}`$. Here, we make a comment on the higher terms in the Kähler potential. Since the leading quadratic terms make the expectation values of $`\eta `$ and $`X`$ fields less than 1, the inflation dynamics is almost unchanged in the presence of the higher terms. The only relevant difference comes from the $`\zeta |X|^4`$ ( $`\zeta `$: constant) term which induces the effective mass of $`X`$ given by
$$m_X^2=2m^22\zeta m^2\phi ^22\zeta m^2\phi ^2.$$
(31)
Thus, $`\zeta `$ should be negative to ensure the positiveness of $`m_X^2`$. If $`|\zeta |1`$, the effective mass becomes larger than the Hubble parameter and the $`X`$ quickly settles down to $`X=0`$ without slow-roll.
## III Conclusion
We have shown that a chaotic inflation naturally takes place if we assume that the Kähler potential has the Nambu-Goldstone-like shift symmetry of the inflaton chiral multiplet $`\mathrm{\Phi }`$ and introduce a small breaking term of the shift symmetry in the superpotential eq.(5). Unlike other inflation models the chaotic inflation model has no initial value problem and hence it is the most attractive. However, it had been difficult to construct a natural chaotic inflation model in the framework of supergravity because the supergravity potential generally becomes very steep beyond the Planck scale. Therefore, the existence of a natural chaotic inflation model may open a new branch of inflation-model building in supergravity. Furthermore, the chaotic inflation is known to produce gravitational waves ( tensor metric perturbations ) which might be detectable in future astrophysical observations .
### Acknowledgments
M.Y. is grateful to J. Yokoyama for useful discussions. M.K. and T.Y. are supported in part by the Grant-in-Aid, Priority Area “Supersymmetry and Unified Theory of Elementary Particles”(#707). M.Y. is partially supported by the Japanese Society for the Promotion of Science.
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# Leontovich Relations in Thermal Field Theory
## I Introduction
Naively one expects that high energy particles, weakly interacting with a medium, can be treated perturbatively. The basic process (production, absorption, or scattering) should follow from lowest order perturbation theory as in vacuum. The only role of the medium is to provide the particles with which the high energy particle interacts. In other words, medium effects enter the cross section only via the distribution functions of the in-medium particles. Note, however, that in the case of bosons the cross sections can be infrared enhanced due to Bose condensation. Famous examples are the effective medium dependent masses of neutrinos interacting with the solar or terrestrial matter leading to medium induced neutrino oscillations. These masses follow directly from integrating over the electron distribution . Note that it is not even necessary that these distributions are in equilibrium. Another example is the production of dileptons with high invariant masses from a quark-gluon plasma (QGP), which is given to lowest order by the annihilation of bare quarks and anti-quarks (Born term) . Due to phase space the main contribution comes from quarks in the high energy tail of the distributions. At lower invariant masses ($`M<1`$ GeV) bare quarks are not sufficient and additional medium effects will lead to interesting structures in the dilepton production rate .
The quantities mentioned above, namely the effective neutrino mass and the dilepton production, are infrared finite to lowest order perturbation theory. In some cases, however, the lowest order contribution suffers from infrared divergences. In this case additional medium effects even for high energy, weakly interacting particles are essential. If the quantity under consideration has a logarithmic infrared singularity within naive perturbation theory, a finite result can be obtained by using hard thermal loop (HTL) resummed propagators for soft momentum transfers . However, some quantities, e.g. damping rates , which exhibit a higher infrared singularity, cannot be calculated to leading order using the HTL method. Also contributions beyond the leading logarithm might require a non-perturbative treatment .
Here we will demonstrate, considering two examples, that the leading logarithmic contributions can be calculated easily using generalized Kramers-Kronig relations. In contrast to the HTL improved perturbation theory the knowledge of the resummed propagator is required only in the high frequency limit. The method presented here is much more general as exact expressions can be derived, which allow in principle non-perturbative results, if only the full propagator in the high frequency limit is known.
In the next section we discuss the usefulness of the Leontovich relation for thermal field theoretic calculations in the case of the collisional energy loss of charged particles in a relativistic plasma. This quantity has already been considered by Kirzhnits using the Leontovich method within the language of plasma physics and recently by the author within thermal field theory . Therefore we will discuss this quantity only briefly focusing on the use of the Leontovich relation. In section 3 we show in more detail that the production rate of high energy photons from a quark QGP can be treated in a similar way. For this purpose we generalize the Leontovich relation used so far only for gauge bosons to quarks.
## II Collisional Energy Loss
The energy loss of a fast charged particle in a medium is a well studied subject . Recently the energy loss of energetic particles, such as leptons and partons, in relativistic plasmas has attracted great interest. In relativistic heavy ion collisions the energy loss of a high energy quark or gluon coming from primary hard collisions in the fireball will lead to jet quenching. Jets therefore serve as a direct probe for the fireball and may provide a signature for the quark-gluon plasma formation . In Supernovae explosions the energy loss of neutrinos, having a weak charge, in the plasma surrounding the stellar core might be an important mechanism for triggering the explosion .
The total energy loss of a particle in a medium can be decomposed into a collisional and a radiative contribution. While the first one originates from the energy transfer to the medium particles, the latter one is caused by radiation from the fast particle. Here we want to consider only the collisional component.
The collisional energy loss is defined as the energy transferred per unit length from the fast particle to the medium in a single collision. It is assumed that the fast particle loses only a small fraction of its energy in each collision. In quantum field theory the collisional energy loss is defined as
$$\frac{dE}{dx}=\frac{1}{v}𝑑\mathrm{\Gamma }\omega ,$$
(1)
where $`v`$ is the velocity of the incident particle with energy $`E`$ and $`\omega =EE^{}`$ the energy transfer to the medium. The interaction rate $`\mathrm{\Gamma }`$ is identical with the inverse mean free path. It can be calculated either from the matrix element of the process responsible for the energy loss or equivalently from the imaginary part of the self energy $`\mathrm{\Sigma }`$ of the particle with four momentum $`P=(E,𝐩)`$, and mass $`M`$, ($`p=|𝐩|`$)
$$\mathrm{\Gamma }(E)=\frac{1}{2E}[1n_F(E)]tr[(P/+M)Im\mathrm{\Sigma }(E,p)],$$
(2)
where $`n_F(E)=1/[\mathrm{exp}(E/T)+1]`$ is the Fermi distribution in the case of a fermion propagating through a plasma of temperature $`T`$. In the following we restrict ourselves to electrons or muons with high energies $`ET`$ in an electron-positron plasma.
To lowest order the interaction rate is caused by elastic scattering of the fast lepton off the thermal electrons and positrons via one-photon exchange. Due to the massless photon this rate is quadratically infrared divergent in naive perturbation theory and cannot be regulated using a HTL photon propagator containing Debye screening. The collisional energy loss, on the other hand, due to the additional factor $`\omega `$ in the integrand of (1) is only logarithmically infrared divergent within naive perturbation theory and finite within the HTL improved perturbation theory. Such a quantity can be calculated by introducing a separation scale $`eTk^{}T`$ for the momentum transfer . Restricting to the leading logarithmic approximation it is sufficient to consider the soft momentum transfer $`k<k^{}`$ only. Since the final result must be independent of $`k^{}`$, it follows from the soft contribution simply by replacing $`k^{}`$ by the maximum momentum transfer. In the soft part of the energy loss a dressed propagator, containing medium effects such as Debye screening, has to be used to regulate the infrared singularity. The exchange of a soft collective photon or plasma mode corresponds to the energy loss by polarization of the medium, also known as Fermi density effect .
The collisional energy loss, caused by the exchange of a single dressed photon, follows from a one-loop approximation for $`\mathrm{\Sigma }`$. Here we allow for the most general photon propagator, indicated by the blob in Fig.1. Then we find for ultrarelativistic electrons or muons ($`v=1`$)
$$\left(\frac{dE}{dx}\right)_{soft}=\frac{e^2}{4\pi }_0^k^{}𝑑kk_k^k𝑑\omega \omega \left[\rho _l(\omega ,k)+\left(1\frac{\omega ^2}{k^2}\right)\rho _t(\omega ,k)\right],$$
(3)
where $`\rho _{l,t}`$ are the spectral functions of the full photon propagator, defined as
$$\rho _{l,t}(\omega ,k)=\frac{1}{\pi }ImD_{l,t}(\omega ,k).$$
(4)
The full photon propagator fulfills the Kramers-Kronig relation
$$D_{l,t}(k_0,k)=_{\mathrm{}}^{\mathrm{}}𝑑\omega \frac{\rho _{l,t}(\omega ,k)}{k_0\omega +i\epsilon }.$$
(5)
At finite temperature the photon propagator has only two independent components, given in Coulomb gauge by the longitudinal and transverse propagators
$`D_l(k_0,k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2\mathrm{\Pi }_l(k_0,k)+i\epsilon }},`$ (6)
$`D_t(k_0,k)`$ $`=`$ $`{\displaystyle \frac{1}{k_0^2k^2\mathrm{\Pi }_t(k_0,k)+i\epsilon }},`$ (7)
where $`\mathrm{\Pi }_{l,t}`$ are the longitudinal and transverse components of the polarization tensor. It should be noted that the soft collisional energy loss, discussed here, follows according to (3) only from the exchange of one dressed space-like ($`\omega ^2k^2<0`$) photon from the particle to the medium. However, the medium particles may undergo further interactions. The physical process corresponding to the imaginary part of the self energy of Fig.1 can be found by using cutting rules. An example is shown in Fig.2. There is no diagram, where to or more photons are emitted from the fast particle, as in the case of the radiative energy loss.
Now we introduce the photon response function
$$R(k_0,k)=k^2D_l(k_0,k)+(k_0^2k^2)D_t(k_0,k).$$
(8)
Making the substitution $`kq=\sqrt{k^2\omega ^2}`$, i.e. introducing the magnitude of the four momentum of the exchanged photon, and using $`ImR(\omega )=ImR(\omega )`$, which follows from the general property $`\rho _{l,t}(\omega )=\rho _{l,t}(\omega )`$ , we find
$$\left(\frac{dE}{dx}\right)_{soft}=\frac{e^2}{2\pi ^2}_0^q^{}𝑑qq_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,\sqrt{q^2+\omega ^2})}{q^2+\omega ^2}$$
(9)
with $`q^{}T`$. Eq. (9) agrees with Ref., which is based on plasma physics arguments if we replace there $`Q^2`$ by $`e^2/4\pi `$.
The response function $`R`$ fulfills the following Kramers-Kronig relation
$$R(k_0,k)=\stackrel{~}{R}+\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,k)}{\omega ^2k_0^2i\epsilon },$$
(10)
which can be shown to be equivalent to (4), if we use $`\rho _{l,t}(\omega )=\rho _{l,t}(\omega )`$. Here $`\stackrel{~}{R}=lim_{k_0\mathrm{}}ReR(k_0,k)`$.
The Kramers-Kronig relation (10) can be generalized by making a Lorentz transformation from $`\omega `$ and $`𝐤`$ to $`\omega ^{}`$ and $`𝐤^{}`$ given in a system which moves with the velocity $`𝐮`$ relative to the initial system. Following the arguments in Ref., choosing $`𝐮𝐤=0`$ and $`|𝐮|=1`$, and utilizing that $`R(\omega ,𝐤)`$ depends only on $`k`$ in an isotropic and homogeneous medium we obtain the Leontovich relation
$$R(k_0,\sqrt{k^2+k_0^2})=R_{\mathrm{}}+\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,\sqrt{k^2+\omega ^2})}{\omega ^2k_0^2i\epsilon },$$
(11)
where $`R_{\mathrm{}}=lim_{k_0\mathrm{}}ReR(k_0,\sqrt{k^2+k_0^2})`$.
The $`\omega `$-integral
$$I=\frac{2}{\pi }_0^{\mathrm{}}𝑑\omega \omega \frac{ImR(\omega ,\sqrt{q^2+\omega ^2})}{q^2+\omega ^2}$$
(12)
appearing in the energy loss (9) agrees with the integral on the right hand side of the Leontovich relation, if we replace $`k_0`$ by $`iq`$ and $`\sqrt{k^2+k_0^2}`$ by 0, i.e. $`k^2=q^2`$, in (11). Therefore we can write
$$I=R(iq,0)R_{\mathrm{}}.$$
(13)
The zero momentum limit of the response function vanishes due to the fact that there is no preferred direction in the medium at vanishing momentum . Consequently, all we have to know is the response function in the high frequency limit $`R_{\mathrm{}}`$ to find the collisional energy loss in the leading logarithmic approximation. Defining
$$\omega _0^2\underset{k_0\mathrm{}}{lim}\mathrm{\Pi }_t(k_0,\sqrt{q^2+k_0^2})$$
(14)
we obtain from (8)
$$R_{\mathrm{}}=\frac{\omega _0^2}{q^2+\omega _0^2}.$$
(15)
Using the Kramers-Kronig relation for the transverse dielectric function, which is related to the transverse polarization tensor , Kirzhnits argued that $`\omega _0`$ is independent of $`q`$ and can be considered as the effective thermal mass of the transverse high frequency plasma excitations, which is given by $`\omega _0^2=e^2n1/\mathrm{\Omega }`$ in the relativistic limit. Here $`n`$ is the number density of the medium and $`\mathrm{\Omega }`$ the energy of the plasma particles.
Combining (15) with (9) and replacing $`q^{}\omega _0`$ by $`q_{max}`$, which is proportional to $`\sqrt{ET}`$ in the relativistic limit $`E\mathrm{\Omega }`$ , we end up with the final result for total collisional energy loss
$$\frac{dE}{dx}=\frac{e^2}{4\pi }\omega _0^2\mathrm{ln}\frac{q_{max}}{\omega _0}.$$
(16)
This result is “exact” in the sense that it is independent of any approximation to the full photon propagator. To logarithmic accuracy the final result just depends on the parameter $`\omega _0`$.
To proceed from here we have to make an approximation for $`\omega _0`$. Since in the high frequency limit medium effects are small, we calculate $`\omega _0`$ by lowest order perturbation theory. Adopting the gauge invariant expression for the transverse polarization tensor in the high temperature limit , we obtain $`\omega _0^2=3m_\gamma ^2/2`$, where the plasma frequency is given by $`m_\gamma =eT/3`$. As expected $`\omega _0`$ is equivalent to the high frequency mass of the transverse photon $`\omega _0^2=e^2n1/\mathrm{\Omega }`$ . Inserting the high temperature result for $`\omega _0`$ in (16) leads to an estimate for the collisional energy loss which agrees to leading logarithm with the one found in the HTL approximation
$$\frac{dE}{dx}=\frac{e^2}{4\pi }\omega _0^2\left(\mathrm{ln}\frac{\sqrt{ET}}{\omega _0}+0.120\right).$$
(17)
In contrast to the HTL method we only had to know the transverse polarization tensor in the high frequency limit. Moreover, we observe that the HTL result has already the same form as the exact result (16), which includes infinitely many higher order diagrams such as the one in Fig.2. Assuming that the exact high frequency transverse polarization tensor can be approximated by its high temperature limit we find that the complete collisional energy loss can be estimated by its lowest order HTL result (17) and that higher order diagrams can be neglected at least within the leading logarithm approximation.
In Ref. this result has been applied to the energy loss of energetic partons in a QGP. It has been shown that the radiative energy loss caused by bremsstrahlung from the fast particle, which increases linearly with the distance $`L`$ over which the parton propagates, dominates over the collisional one for $`L>1`$ fm.
Another application of (16) has been discussed in Ref. in connection with the neutrino energy loss in matter. Here the collisional energy loss provides a reliable estimate of the total energy loss since bremsstrahlung from the neutrino, i.e. emission of a $`Z`$ boson, is suppressed by the large mass of the gauge boson at least for temperatures below $`T100`$ GeV.
## III High Energy Photons
Now we want to discuss another example, namely the production of high energy real photons in the QGP, which might also serve as a promising signature for the QGP formation in relativistic heavy ion collisions . Due the weak interaction of photons with the QGP photons present as well as jet quenching a direct probe for the hot fireball. To lowest order the production rate of real photons is given by the diagrams of Fig.3 (Compton scattering, annihilation with gluon absorption). Here an intermediate bare quark appears, which can be assumed to be massless as the bare mass of up and down quarks can be neglected compared to the temperature $`T`$ of the QGP. This leads to a logarithmic infrared divergence in the production rate, which is regulated by medium effects.
The production rate corresponding to the processes of Fig.3 can be calculated by using the HTL resummed quark propagator in the case of a soft quark exchange, i.e. momentum exchange much smaller than $`T`$. This corresponds to a one-loop calculation including a HTL quark propagator, that contains an effective quark mass of the order $`gT`$, which cuts off the logarithmic singularity. For the hard momentum transfer the tree level scattering matrix elements convoluted with the parton distribution functions can be used . In this way the production rate of energetic photons to leading logarithmic order $`\alpha _s\mathrm{ln}(1/\alpha _s)`$ in the strong coupling constant has been obtained.
Here we will derive the photon rate to leading logarithm, using the Leontovich relation for the full quark propagator, which allows a more general and after all simpler derivation of the rate than applying the HTL method. For this purpose we calculate the photon production rate from the imaginary part of the polarization tensor according to
$$E\frac{dR}{d^3p}=\frac{2}{(2\pi )^3}\frac{1}{e^{E/T}1}Im\mathrm{\Pi }_\mu ^\mu (E).$$
(18)
This expression is exact to all orders in $`\alpha _s`$ and to leading order in $`\alpha `$.
Here we focus only on the soft contribution to the photon rate since the hard contribution can be calculated perturbatively from Fig.1 restricting to the logarithmic approximation. Since the energy of the produced photon is high ($`ET`$) the polarization tensor is given by Fig.4. The blob denotes the full non-perturbative quark propagator. Due to kinematics there is only one full quark propagator, since the other one has to be hard, and no vertex correction since the high energy photon resolves the vertex completely. By cutting this polarization tensor one observes that all processes are taken into account, where the soft quark interacts with the medium in all possible ways. For example it can absorb a thermal gluon as in Fig.3. But also bremsstrahlung from the thermal particles and other higher order processes are included.
Now we want to calculate the photon production rate from Fig.4 for the most general full quark propagator without assuming any approximation for it. For this purpose, we proceed similarly as in the case of the collisional energy loss for energetic charged particles in a plasma. We start from an exact expression for the imaginary part of the soft polarization tensor in the case of two massless quark flavors using $`ET`$
$$Im\mathrm{\Pi }_\mu ^\mu (E)=\frac{5e^2}{12\pi }_0^{\mathrm{}}𝑑k_k^k𝑑\omega [(k\omega )\rho _+(\omega ,k)+(k+\omega )\rho _{}(\omega ,k)]\theta (q_c^2k^2+\omega ^2),$$
(19)
where $`q_cT`$ is the separation scale between the soft and the hard contribution, $`\omega `$ and $`k`$ the energy and the magnitude of the three momentum of the soft quark, and $`\rho _\pm `$ the spectral functions of the full quark propagator $`S(K)`$ in the helicity representation
$$\rho _\pm (\omega ,k)=\frac{1}{\pi }Im\frac{1}{D_\pm (\omega ,k)}$$
(20)
with ($`K=(k_0,𝐤)`$)
$$S(K)=\frac{\gamma _0\widehat{𝐤}𝜸}{2D_+(k_0,k)}+\frac{\gamma _0+\widehat{𝐤}𝜸}{2D_{}(k_0,k)}.$$
(21)
Replacing again $`k`$ by $`q=\sqrt{k^2\omega ^2}`$ in the integral of (19) we find
$$Im\mathrm{\Pi }_\mu ^\mu (E)=\frac{5e^2}{12\pi ^2}_0^{q_c}𝑑qq_{\mathrm{}}^{\mathrm{}}𝑑\omega \frac{\omega }{\omega ^2+q^2}ImQ(\omega ,\sqrt{\omega ^2+q^2}),$$
(22)
where the quark response function $`Q`$ is given by
$$Q(k_0,k)=\frac{k}{k_0}\left(\frac{kk_0}{D_+(k_0,k)}+\frac{k+k_0}{D_{}(k_0,k)}\right).$$
(23)
This response function fulfills the same Kramers-Kronig relation as the photon response function (10)
$$Q(k_0,k)=\stackrel{~}{Q}+\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑\omega \omega \frac{ImQ(\omega ,k)}{\omega ^2k_0^2iϵ},$$
(24)
where $`\stackrel{~}{Q}=lim_{k_0\mathrm{}}ReQ(k_0,k)`$. This relation is a direct consequence of the definition of the spectral functions
$$\frac{1}{D_\pm (k_0,k)}=_{\mathrm{}}^{\mathrm{}}𝑑\omega \frac{\rho _\pm (\omega ,k)}{k_0\omega +i\epsilon }$$
(25)
using $`\rho _+(\omega ,k)=\rho _{}(\omega ,k)`$ .
The Kramers-Kronig relation (22) can be generalized again by the Lorentz transformation of section 2, from which we obtain the Leontovich relation analogously to (11)
$$Q(k_0,\sqrt{k_0^2+k^2})=Q_{\mathrm{}}+\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑\omega \omega \frac{ImQ(\omega ,\sqrt{\omega ^2+k^2})}{\omega ^2k_0^2i\epsilon },$$
(26)
where $`Q_{\mathrm{}}=lim_{k_0\mathrm{}}ReQ(k_0,\sqrt{k_0^2+k^2})`$. This relation is more restrictive than the Kramers-Kronig relation (24) and will be used to evaluate the photon production in the following from (22).
The $`\omega `$-integral in (22) agrees with the integral on the right hand side of the Leontovich relation, if we replace again $`k_0`$ by $`iq`$ and $`\sqrt{k_0^2+k^2}`$ by 0, i.e. $`k^2=q^2`$, in (26). Since $`D_+(k_0,k=0)=D_{}(k_0,k=0)`$ , $`Q(iq,0)=0`$. Hence the imaginary part of the polarization tensor containing the most general in-medium quark propagator is given by the simple expression
$$Im\mathrm{\Pi }_\mu ^\mu (E)=\frac{5e^2}{12\pi }_0^{q_c}𝑑qqQ_{\mathrm{}}.$$
(27)
Using the Leontovich relation we were able to express the soft part of the photon production rate by an integral over the real part of the response function only in the high frequency limit just below the light cone. This expression is exact as long as we do not assume any approximations for the response function. The advantage of this method is that we do not have to know the quark response function or the quark propagator over the entire energy range, but only in the high frequency limit.
Starting from the most general expression for the full quark propagator (see e.g. Ref.)
$$D_\pm (k_0,k)=(k_0\pm k)[1+a(k_0,k)]b(k_0,k)$$
(28)
and using that the scalar functions $`a`$ and $`b`$ fulfill the following inequalities in the high frequency limit, where the medium effects on the quark propagator become small, $`lim_{k_0\mathrm{}}a1`$ and $`lim_{k_0\mathrm{}}bk_0`$, we find
$$Q_{\mathrm{}}=\underset{k_0\mathrm{}}{lim}\frac{2b_{\mathrm{}}}{q^2/k_02b_{\mathrm{}}}.$$
(29)
Here $`b_{\mathrm{}}=lim_{k_0\mathrm{}}b(k_0,\sqrt{k_0^2+q^2})`$. To proceed we have to make an approximation for the full quark propagator or equivalently for the response function in order to determine $`b_{\mathrm{}}`$.
In the high frequency limit the response function should be calculable perturbatively, as it is also the case for the photon or gluon response function which is related to the dielectric function of the medium . For in the high frequency limit the dielectric function has to be close to its vacuum value and can be computed therefore perturbatively. The same argument holds for the quark response function in the high frequency limit. Note that the quark response function to lowest order, determined from the one-loop quark self energy, is infrared finite. In the high temperature limit the gauge invariant result
$$b_{\mathrm{}}=\frac{m_q^2}{k_0}Q_{\mathrm{}}=\frac{2m_q^2}{q^2+2m_q^2}$$
(30)
is found. Here $`2m_q^2=g^2T^2/3`$ is the square of the effective high frequency quark mass. Combining (27) with (30) we get
$$Im\mathrm{\Pi }_\mu ^\mu (E)=\frac{5e^2}{12\pi }m_q^2\mathrm{ln}\frac{q_c^2}{2m_q^2},$$
(31)
where we assumed $`q_cm_q`$. This result has also be found by lowest order HTL perturbation theory, where the factor $`1/2`$ under the logarithm could be derived only numerically . Using, however, the Leontovich relation, where one needs to know the response function only in the high frequency limit, this factor, related to the high frequency effective quark mass, has been obtained analytically.
Combining the soft part with the hard part, calculated perturbativeky from Fig.3 in Ref., we obtain the final result for the production rate of energetic photons to leading logarithm
$$E\frac{dR}{d^3p}=\frac{5}{18\pi ^2}\alpha \alpha _sT^2e^{E/T}\mathrm{ln}\frac{0.2317E}{\alpha _sT}.$$
(32)
Here the separation scale $`q_c`$ serving as an infrared cutoff for the hard part drops out since the hard part and the soft part have the same factors in front of the logarithm. This had to be expected by physical reasons, since the final result has to be independent of the arbitrary separation scale $`q_c`$ .
The result (32) agrees with the lowest order HTL contribution . Unfortunately, this result is of no practical relevance, as for physical values of the strong coupling constant, terms beyond the leading logarithm dominate. These contributions come from higher order diagrams, i.e. two-loop diagrams of the HTL perturbative expansion describing e.g. bremsstrahlung, which show a strong infrared sensitivity . They are not included in the soft part of the polarization tensor using the full quark propagator (27), since they come from the exchange of a hard quark . For realistic values of the coupling these contributions even dominate clearly over the lowest order ones, in particular at high photon energies $`E`$ as the two-loop contributions are proportional to $`ET`$ in contrast to the one-loop contribution which is proportional to $`T^2`$. As a matter of fact, presumably infinitely many diagrams within the HTL resummed perturbative expansion contribute to order $`\alpha _s`$ . Therefore it would be desirable to extend the arguments given here also to hard momentum transfer. Maybe a resummation of all higher order diagrams contributing to order $`\alpha _s`$ leads to a cancellation between these diagrams and a suppression of the bremsstrahlung and higher processes in the hard photon production.
## IV Conclusions
Lowest order perturbation theory at finite temperature in the weak coupling limit works only for quantities, which are infrared finite in naive perturbation theory. Examples are effective masses and the production of dileptons with high invariant masses.
Quantities of energetic particles ($`ET`$) which are logarithmically infrared divergent in naive perturbation theory, such as the collisional energy loss or the photon production rate, can be calculated within the leading logarithm approximation using the lowest order HTL improved perturbation theory. To extract the leading logarithm it is sufficient to consider soft momentum transfers described by HTL propagators. However, for realistic values of the coupling constants, in particular in the case of the strong coupling constant, contributions beyond the leading logarithm become important and can even dominate, as in the case of the photon production. For the photon production rate these contributions have been shown to be non-perturbative, i.e. infinitely many higher order diagrams containing HTL propagators and vertices contribute to the same order in the coupling constant.
Quantities, such as damping rates, which exhibit a higher degree of infrared divergence in naive perturbation theory cannot be treated within the HTL improved perturbation theory. So many properties of (high energy) particles in a medium cannot be calculated perturbatively even in the weak coupling limit.
Here we presented a method for computing within thermal field theory quantities of energetic particles, which are in naive perturbation theory logarithmically infrared divergent, by using generalized Kramers-Kronig relations. These so-called Leontovich relations follow from the usual Kramers-Kronig relations for the thermal propagators performing a Lorentz transformation. Applying these more restrictive relations one is able to express the soft part of the quantities under consideration, such as the collisional energy loss or the photon production, by simple integrals. For evaluating these integrals one needs only the self energy of the high energy particle in the high frequency limit just below the light cone. In this way an exact expression, i.e. independent of any approximation for the full propagator, for these quantities within the logarithmic approximation is obtained. Assuming perturbation theory to hold for the high frequency self energies and using the HTL result for them, the results yielded within the lowest order HTL improved perturbation theory are reproduced. In contrast to the HTL resummation technique the method presented here enables analytical calculations of the soft contributions. Also it is more general, allowing in principle non-perturbative results, if only the full self energy in the high frequency limit is known. Since the application of perturbation theory at finite temperature fails in many cases even in the weak coupling limit, it would be desirable to have non-perturbative or even exact statements. Therefore it might be worthwhile to extend the methods presented here also to hard momentum transfers, going beyond the leading logarithm.
ACKNOWLEDGMENTS
The author is grateful to G. Raffelt for drawing his attention to the paper by D.A. Kirzhnits and for helpful discussions and to the Max-Planck-Institut für Physik (Werner-Heisenberg-Institut) for their hospitality.
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# Copolymer at a selective interface and two dimensional wetting: a grand canonical approach
## I Introduction
The behaviour of polymers at interfaces is currently a subject of great activity . In this paper, we focus our interest on two problems, namely (i) the adsorption of a periodic copolymer chain on a selective fluid-fluid interface (ii) the adsorption of a homopolymer chain on a hard wall (wetting). These problems have been previously solved for one dimensional geometries, provided certain assumptions are fulfilled (self avoidance is neglected, wetting is considered in the SOS approximation,…). A quick guide to the recent litterature can be found in references and references therein. Here, we will tackle these models, with the same assumptions as above, through the use of a grand canonical method. This method was initiated in the context of the DNA helix-coil transition and was recently used in the context of RNA folding . We will recover previously known results about models (i) and (ii). Moreover, this approach enables us to obtain the loop length distribution function in the localized regime. We also find in this way some analytic expressions of finite size corrections at the critical temperature, which are in agreement with numerical transfer matrix calculations. The lay out of the paper is as follows. Section II deals with the periodic $`AB`$ copolymer chain at a selective interface. In section III, we consider the wetting of a homogeneous wall by a homopolymer chain.
## II Periodic copolymer chain at a selective interface
### A Introduction
We consider a periodic $`AB`$ copolymer chain of $`2N`$ monomers at a selective fluid-fluid interface, located at $`z_0=\frac{1}{2}`$. The upper (resp. lower) fluid is a good solvent for $`A`$-type (resp. $`B`$-type) monomers, and will herefrom be called fluid $`\alpha `$ (resp. fluid $`\beta `$). This situation is modelled through the Hamiltonian
$$_{2N}=\underset{i=1}{\overset{2N}{}}q_i\mathrm{sgn}[z(i)\frac{1}{2}]$$
(1)
where $`\mathrm{sgn}(u)`$ denotes the sign function, and the “charges” $`q_i`$ of the periodic chain are chosen as $`q_{2n+1}=q_A>0`$ and $`q_{2n}=q_B<0`$, together with $`q_Aq_B`$. In equation (1), $`z(i)`$ denotes the position of monomer $`i`$ of the chain, which we represent as an (unidimensional) random walk (RW) with the properties $`z(i)=..,2,1,0,+1,+2,..`$, and $`z(i+1)z(i)=\pm 1`$. We emphasize once more that self avoidance is neglected in this model.
The partition function then reads
$`Z_{2N}={\displaystyle \underset{(RW)}{}}e^{\beta _{2N}}`$ (2)
where $`\beta =\frac{1}{T}`$ and the sum runs over all random walks (RW) starting with an $`A`$ monomer at $`z(i=1)=1`$. The thermodynamical properties of this model can be obtained through the use of transfer matrix methods , as given in Appendix $`\mathrm{A}`$. Here, we follow references , and consider a particular chain configuration as a collection of non interacting loops. This class of problems is solved in a grand canonical approach, by defining a grand canonical partition function
$`Z(K)={\displaystyle \underset{N=1}{\overset{+\mathrm{}}{}}}K^NZ_{2N}`$ (3)
and noting that the radius of convergence $`K_{}(\beta )`$ of the series $`Z(K)`$ gives the free energy per monomer in the thermodynamic limit
$`f={\displaystyle \frac{\mathrm{log}(K_{}(\beta ))}{2\beta }}`$ (4)
Moreover, one can easily calculate in this way quantities such as the loop length distribution function, which are not accessible through the transfer matrix approach.
### B Calculation of the grand canonical partition function
The partition function $`Z_{2N}`$ for a chain of $`(2N)`$ monomers can then be decomposed according to the position $`z(2N)`$ of the end point of the chain. For simplicity, we set $`z(2N)=2n`$, and write
$`Z_{2N}={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}Z_{2N}(2n)`$ (5)
Given the expression of the Hamiltonian (1), the chain configurations can be further decomposed into loops. A preliminary remark is as follows: if $`𝒩(l,z)`$ denotes the number of random walks of $`l`$ steps going from $`z=0`$ to $`z0`$ in the presence of an absorbing barrier at $`z=0^{}`$, the image method yields
$`𝒩(2l,2p)=C_{2l}^{l+p}C_{2l}^{l+p+1}`$ (6)
$`𝒩(2l1,2p1)=C_{2l1}^{l+p1}C_{2l1}^{l+p}`$ (7)
depending on the parity of $`(l,z)`$.
Let us temporarily suppose that the chain end is in fluid $`\alpha `$ $`(2n2)`$. The first loop, located by hypothesis in fluid $`\alpha `$ contains $`(2l_11)`$ monomers, namely $`l_1`$ A monomers and $`(l_11)`$ B monomers. The second loop, located in fluid $`\beta `$, contains $`(2l_21)`$ monomers, namely $`(l_21)`$ A monomers and $`l_2`$ B monomers…. The last loop is denoted by $`l_{2p}`$. Finally, one must also specify the number $`2l`$ of monomers between the interface and the chain end. We may then write
$`Z_{2N}(2n)`$ $`={\displaystyle \underset{p=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{l_{2p}=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\delta (2N({\displaystyle \underset{i=1}{\overset{2p}{}}}(2l_i1)+2l))e^{\beta q_B+\beta (q_Aq_B)l_1}𝒩(2l_12,0)`$ (10)
$`e^{\beta q_A\beta (q_Aq_B)l_2}𝒩(2l_22,0)\mathrm{}..e^{\beta q_A\beta (q_Aq_B)l_{2p}}𝒩(2l_{2p}2,0)`$
$`e^{\beta (q_Aq_B)l}𝒩(2l1,2n1)`$
Plugging equation (10) into equations (3) and (5), we get
$$Z(K)=\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}Z(K,2n)$$
(11)
where
$`Z(K,2n)=`$ $`{\displaystyle \underset{N=1}{\overset{+\mathrm{}}{}}}K^NZ_{2N}(2n)`$ (12)
$`=`$ $`{\displaystyle \frac{1}{1e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})B(K^{1/2}e^{\beta q_0})}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}(Ke^{2\beta q_0})^l𝒩(2l1,2n1)`$ (13)
with
$`q_0={\displaystyle \frac{q_Aq_B}{2}}`$ (14)
and
$`B(y)={\displaystyle \underset{l=1}{\overset{+\mathrm{}}{}}}y^{2l1}𝒩(2l2,0)={\displaystyle \frac{1}{2y}}\left(1\sqrt{14y^2}\right)`$ (15)
We finally get, for ($`2n2`$)
$`Z(K,2n)={\displaystyle \underset{N=1}{\overset{+\mathrm{}}{}}}K^NZ_{2N}(2n)={\displaystyle \frac{1}{1e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})B(K^{1/2}e^{\beta q_0})}}\left({\displaystyle \frac{B(K^{1/2}e^{\beta q_0})}{K^{1/2}e^{\beta q_0}}}1\right)^n`$ (16)
Similarly, if the chain end is in fluid $`\beta `$ ($`2n0`$), we find
$`Z(K,2n)={\displaystyle \underset{N=1}{\overset{+\mathrm{}}{}}}K^NZ_{2N}(2n)={\displaystyle \frac{e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})B(K^{1/2}e^{\beta q_0})}{1e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})B(K^{1/2}e^{\beta q_0})}}\left({\displaystyle \frac{B(K^{1/2}e^{\beta q_0})}{K^{1/2}e^{\beta q_0}}}1\right)^{|n|}`$ (17)
### C Thermodynamics
We finally get from (16) and (17)
$`Z(K)={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}Z(K,2n)={\displaystyle \frac{r_+(K^{1/2}e^{\beta q_0})+e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})r_{}(K^{1/2}e^{\beta q_0})}{1e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})B(K^{1/2}e^{\beta q_0})}}`$ (18)
where
$`r_+(y)={\displaystyle \frac{\frac{B(y)}{y}1}{2\frac{B(y)}{y}}}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{\sqrt{14y^2}}}1\right)`$ (19)
$`r_{}(y)={\displaystyle \frac{B(y)}{2\frac{B(y)}{y}}}={\displaystyle \frac{y}{\sqrt{14y^2}}}`$ (20)
The radius of convergence $`K_{}(\beta )`$ of the grand canonical partition function $`Z(K)`$ is either given by the singularity of the numerator of equation (18), or by the singularity of the denominator. The former corresponds to the value
$`K_{deloc}(\beta )={\displaystyle \frac{1}{4}}e^{2\beta q_0}`$ (21)
The latter comes from the resummation of the loops, and is given by
$`K_{}(\beta )={\displaystyle \frac{\mathrm{sinh}(\beta q_A)\mathrm{sinh}(\beta q_B)}{\mathrm{sinh}^2\beta (q_A+q_B))}}`$ (22)
If the solution $`K_{}(\beta )`$ is smaller than $`K_{deloc}`$ then the chain is localized. This actually happens for $`\beta >\beta _c`$ (i.e. $`T<T_c`$), where $`\beta _c`$ is the solution of
$`e^{2\beta _cq_B}+e^{2\beta _cq_A}2=0`$ (23)
In particular, in the limit $`q_0=(q_Aq_B)/20`$, the critical temperature $`T_c`$ diverges as
$`T_c\underset{q_00}{}{\displaystyle \frac{q_A^2}{q_0}}`$ (24)
Note that the scaling of $`T_c`$ with $`q_0`$ differs from equation (10) of reference , since the model used in this reference allows the monomers to sit at the interface, in contrast to equation (1). The physics is nevertheless the same: a symmetric copolymer chain ($`q_A=q_B`$) is always localized at the interface.
Since the radius of convergence $`K_{}(\beta )`$ and free energy per monomer are linked by equation (4), we get
$`f(T)={\displaystyle \frac{1}{\beta }}\mathrm{ln}{\displaystyle \frac{\sqrt{\mathrm{sinh}(\beta q_A)\mathrm{sinh}(\beta q_B)}}{\mathrm{sinh}\beta (q_A+q_B)}}`$ (25)
for $`TT_c`$, which in turn yields
$`f(T)\underset{TT_c}{}f_{deloc}(T_c)C(T_cT)^2`$ (26)
where $`f_{deloc}(T_c)=T_c\mathrm{ln}2q_0`$ is the free energy of the delocalized phase (since we have chosen $`q_0>0`$, the chain is then in fluid $`\alpha `$). Equation (26) corresponds to a critical exponent $`\alpha =0`$. The jump in the specific heat per monomer is given by
$`C_{\mathrm{}}(T_c^{})=\left({\displaystyle \frac{(2e^{2\beta _cq_A}1)\mathrm{ln}(2e^{2\beta _cq_A})2\beta _cq_A}{2(e^{2\beta _cq_A}1)}}\right)^2`$ (27)
$`C_{\mathrm{}}(T_c^+)=0`$ (28)
### D Density profile
We denote the normalized density profile by $`\rho (2n)`$ and write
$`\rho (2n)=\underset{N\mathrm{}}{lim}{\displaystyle \frac{Z_{2N}(2n)}{Z_{2N}}}`$ (29)
Inverting equations (16) and (17), we have
$`Z_{2N}(2n)={\displaystyle _{\text{ circle around }0}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{Z(K,2n)}{K^{N+1}}}`$ (30)
After deformation of the contour, the dominant contribution for large $`N`$ comes from the pole at $`K_{}(\beta )`$ (see equation (22)), and reads
$`Z_{2N}(2n)`$ $`\underset{N\text{large}}{}{\displaystyle _{\text{ circle around }K_{}}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{Z(K,2n)}{K^{N+1}}}`$ (31)
Summing over $`n`$ leads to
$`Z_{2N}`$ $`\underset{N\text{large}}{}{\displaystyle _{\text{ circle around }K_{}}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{Z(K)}{K^{N+1}}}`$ (32)
The density profile $`\rho (2n)`$ is then obtained as
$`\rho (2n)={\displaystyle \frac{1}{r_+(K_{}^{1/2}e^{\beta q_0})+e^{\beta (q_A+q_B)}B(K_{}^{1/2}e^{\beta q_0})r_{}(K_{}^{1/2}e^{\beta q_0})}}\left({\displaystyle \frac{B(K_{}^{1/2}e^{\beta q_0})}{K_{}^{1/2}e^{\beta q_0}}}1\right)^n`$ (33)
for $`n1`$ and
$`\rho (2n)={\displaystyle \frac{e^{\beta (q_A+q_B)}B(K_{}^{1/2}e^{\beta q_0})B(K_{}^{1/2}e^{\beta q_0})}{r_+(K_{}^{1/2}e^{\beta q_0})+e^{\beta (q_A+q_B)}B(K_{}^{1/2}e^{\beta q_0})r_{}(K_{}^{1/2}e^{\beta q_0})}}\left({\displaystyle \frac{B(K_{}^{1/2}e^{\beta q_0})}{K_{}^{1/2}e^{\beta q_0}}}1\right)^{|n|}`$ (34)
for $`n0`$. Inserting the value of $`K_{}(\beta )`$, we get
$`\rho (2n)={\displaystyle \frac{(e^{2\beta q_B}+e^{2\beta q_A}2)(e^{2\beta q_A}+e^{2\beta q_B}2)}{(e^{2\beta (q_A+q_B)}1)(1e^{2\beta q_A})(1e^{2\beta q_B})}}\left({\displaystyle \frac{1e^{2\beta q_A}}{e^{2\beta q_B}1}}\right)^n\text{for}n1`$ (35)
$`\rho (2n)={\displaystyle \frac{(e^{2\beta q_B}+e^{2\beta q_A}2)(e^{2\beta q_A}+e^{2\beta q_B}2)}{(e^{2\beta (q_A+q_B)}1)(1e^{2\beta q_A})(1e^{2\beta q_B})}}\left({\displaystyle \frac{1e^{2\beta q_B}}{e^{2\beta q_A}1}}\right)^n\text{for}n0`$ (36)
i.e. the density around the interface decays exponentially from the interface with characteristic lengths $`\xi _\alpha `$ (resp. $`\xi _\beta `$) in fluid $`\alpha `$ (resp. fluid $`\beta `$), with
$`\xi _\alpha (T)=`$ $`{\displaystyle \frac{2}{\mathrm{ln}\left(\frac{e^{2\beta q_B}1}{1e^{2\beta q_A}}\right)}}`$ (37)
$`\xi _\beta (T)=`$ $`{\displaystyle \frac{2}{\mathrm{ln}\left(\frac{e^{2\beta q_A}1}{1e^{2\beta q_B}}\right)}}`$ (38)
Close to $`T_c`$, we get
$`\xi _\alpha (T)\underset{TT_c}{}{\displaystyle \frac{1}{T_cT}}`$ (39)
and
$`\xi _\beta (T_c)={\displaystyle \frac{T_c}{(q_A+q_B)}}`$ (40)
The order parameter of the transition can be chosen as the probability $`_\beta `$ to be in fluid $`\beta `$ and vanishes linearly at the transition:
$`_\beta ={\displaystyle \underset{n=\mathrm{}}{\overset{0}{}}}\rho (2n)\underset{TT_c}{}(T_cT)`$ (41)
### E Probability distribution of loop lengths
Contrary to the direct matrix transfer approach, the grand canonical approach allows one to obtain informations on the statistical properties of loops in each fluid. Without loss of generality, we may consider the partition function $`Z_{2N}(2n=0)`$ where the chain of length $`2N`$ is attached at the interface at both ends, and decompose it into
$`Z_{2N}(2n=0)={\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}Z_{2N}(2n=0,2l_11,2l_21)`$ (42)
where $`Z_{2N}(2n=0,2l_11,2l_21)`$ represents the partition function with the constraints that the first loop in fluid $`\alpha `$ contains exactly $`(2l_11)`$ monomers and that the first loop in fluid $`\beta `$ contains exactly $`(2l_21)`$ monomers.
We again define generating functions by
$`\widehat{Z}(K,s_1,s_2)={\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}K^Ne^{s_1(2l_11)s_2(2l_21)}Z_{2N}(2n=0,2l_11,2l_21)`$ (43)
The generalization of equation (17) to the case $`(s_1,s_2)(0,0)`$ then yields
$`\widehat{Z}(K,s_1,s_2)={\displaystyle \frac{e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0}e^{s_1})B(K^{1/2}e^{\beta q_0}e^{s_2})}{1e^{\beta (q_A+q_B)}B(K^{1/2}e^{\beta q_0})B(K^{1/2}e^{\beta q_0})}}`$ (44)
Again, by inversion and deformation of the contour in the complex plane, we get the asymptotic behavior for large $`N`$
$`\widehat{Z}_{2N}(s_1,s_2)=`$ $`{\displaystyle _{\text{ circle around }0}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{\widehat{Z}(K,s_1,s_2)}{K^{N+1}}}`$ (46)
$`\underset{N\text{large}}{}{\displaystyle _{\text{circle around }K_{}}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{\widehat{Z}(K,s_1,s_2)}{K^{N+1}}}`$
The generating functions $`\widehat{P}_\alpha `$ and $`\widehat{P}_\beta `$ of the loop length distribution functions in the two fluids are given by
$`\widehat{P}_\alpha (s){\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}e^{s(2l_11)}P_\alpha (2l_11)=\underset{N\mathrm{}}{lim}\left({\displaystyle \frac{\widehat{Z}_{2N}(s_1=s,s_2=0)}{\widehat{Z}_{2N}(s_1=0,s_2=0)}}\right)`$ (47)
$`\widehat{P}_\beta (s){\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}e^{s(2l_21)}P_\beta (2l_21)=\underset{N\mathrm{}}{lim}\left({\displaystyle \frac{\widehat{Z}_{2N}(s_1=0,s_2=s)}{\widehat{Z}_{2N}(s_1=0,s_2=0)}}\right)`$ (48)
More explicitely, we have
$`\widehat{P}_\alpha (s)={\displaystyle \frac{B(K_{}^{1/2}e^{\beta q_0}e^s)}{B(K_{}^{1/2}e^{\beta q_0})}}=e^s{\displaystyle \frac{1\sqrt{1e^{\omega (\beta )2s}}}{1\sqrt{1e^{\omega (\beta )}}}}`$ (49)
$`\widehat{P}_\beta (s)={\displaystyle \frac{B(K_{}^{1/2}e^{\beta q_0}e^s)}{B(K_{}^{1/2}e^{\beta q_0})}}=e^s{\displaystyle \frac{1\sqrt{1e^{\omega (\beta )4\beta q_02s}}}{1\sqrt{1e^{\omega (\beta )4\beta q_0}}}}`$ (50)
with
$`\omega (\beta )=2\beta q_02\mathrm{ln}2\mathrm{ln}K_{}=2\beta (f_{deloc}f)=\mathrm{ln}\left({\displaystyle \frac{(e^{2\beta (q_A+q_B)}1)^2}{4e^{2\beta q_A}(e^{2\beta q_A}1)(e^{2\beta q_B}1)}}\right)`$ (51)
We note that equations (49) and (50) exhibit several loop size scales, that we consider below.
At the critical temperature ($`T=T_c`$), the distribution of loop lengths $`(l_\alpha )`$ in fluid $`\alpha `$ is simply the loop distribution of a free random walk
$`\widehat{P}_\alpha (s)\underset{T=T_c}{=}e^s\left[1\sqrt{1e^{2s}}\right]=1\sqrt{2s}+O(s)`$ (52)
yielding an algebraic decay for large $`l_\alpha `$
$`P_\alpha (l_\alpha )\underset{l_\alpha \mathrm{}}{}{\displaystyle \frac{1}{l_\alpha ^{3/2}}}`$ (53)
A consequence of this critical distribution ($`T=T_c`$) is as follows. For a chain of length $`N`$, the number of loops scales as $`\sqrt{N}`$, the longest ($`\alpha `$) loop being of order $`N`$.
In the critical region $`TT_c`$ where $`\omega (\beta )(T_cT)^2`$, the moments of $`(l_\alpha )`$ diverge as
$`\overline{l_\alpha ^n}={\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}(2l1)^nP_\alpha (2l1)\underset{TT_c}{}(T_cT)\left({\displaystyle \frac{1}{(T_cT)^2}}\right)^n`$ (54)
i.e. for large $`l_\alpha `$ the correct scaling variable is $`\lambda =l_\alpha (T_cT)^2`$, but the normalisation of the region of finite $`\lambda `$ varies as $`(T_cT)`$. This can be understood by considering the normalisation of the region of loops bigger than a given scale $`l_c\frac{1}{(T_cT)^2}`$ for a free random walk (53)
$`{\displaystyle _{l_c}^{\mathrm{}}}𝑑l{\displaystyle \frac{1}{l^{3/2}}}{\displaystyle \frac{1}{l_c^{1/2}}}(T_cT)`$ (55)
In other words, the phase transition is driven by a small fraction of large loops.
It is of interest to note that the weight of the small $`l_\alpha `$ loops is finite, even at $`T_c`$. For instance, the probability to have $`l_\alpha =1`$ is
$`P_\alpha (l_\alpha =1)={\displaystyle \frac{e^{2\beta q_B}1}{e^{2\beta q_B}e^{2\beta q_A}}}\underset{TT_c}{}{\displaystyle \frac{1}{2}}`$ (56)
At zero temperature we of course recover the ground state
$`P_\alpha (l_\alpha )\underset{T=0}{}\delta _{l_\alpha ,1}`$ (57)
In marked contrast, the critical distribution of loop lengths in the ($`\beta `$) fluid is given by
$`\widehat{P}_\beta (s)\underset{T=T_c}{}e^s{\displaystyle \frac{1\sqrt{1e^{4\beta _cq_02s}}}{1\sqrt{1e^{4\beta _cq_0}}}}`$ (58)
leading to finite $`\overline{l_\beta ^n}`$ moments.
### F Correlation function
To further emphasize the meaning of the scale corresponding to $`\omega (\beta )`$ in equation (51), we consider the connected correlation function
$`𝒞(ji)=<\mathrm{sgn}[z(i){\displaystyle \frac{1}{2}}]\mathrm{sgn}[z(j){\displaystyle \frac{1}{2}}]><\mathrm{sgn}[z(i){\displaystyle \frac{1}{2}}]><\mathrm{sgn}[z(j){\displaystyle \frac{1}{2}}]>`$ (59)
Its generating function reads
$`\widehat{𝒞}(K)={\displaystyle \underset{j=i}{\overset{\mathrm{}}{}}}K^{ji}𝒞(ji)`$ (60)
and can be expressed, via a loop decomposition, in terms of the functions $`\widehat{P}_\alpha (s)`$ and $`\widehat{P}_\beta (s)`$ defined in (49,50), with the replacement $`(s=\mathrm{ln}K)`$. Inverting (60), we have
$`𝒞(ji)={\displaystyle _{\text{circle around }0}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{\widehat{𝒞}(K)}{K^{1+ji}}}`$ (61)
At large $`(ji)`$ separation, the behaviour of $`𝒞(ji)`$ is dominated by the singularity at $`K=e^{\frac{\omega (\beta )}{2}}`$ of $`\widehat{P}_\alpha (\mathrm{ln}K)`$, leading to
$`𝒞(ji)\underset{(ji)\mathrm{}}{}e^{(ji)\frac{\omega (\beta )}{2}}`$ (62)
This shows that $`\omega (\beta )`$ is the inverse of the correlation length along the chain.
### G Finite size properties
In this section, we again consider that both chain ends are fixed at the interface. This means in particular that one has $`Z_{2N}=Z_{2N}(0),Z(K)=Z(K,0),\mathrm{}.`$.
#### 1 Free energy
In the inversion formula of equation (30), we set $`2n=0`$ and get
$`Z_{2N}(0)={\displaystyle _{\text{ circle around }0}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{Z(K,0)}{K^{N+1}}}`$ (63)
where $`Z(K,0)`$ given in (17) presents a simple pole at $`K_{}(\beta )`$, and a cut on the real axis, namely $`[K_{deloc}(\beta ),+\mathrm{}]`$.
In the localized phase, and after deformation of the contour, the dominant contribution for large $`N`$ comes from the pole at $`K_{}(\beta )`$, and the correction is given by the leading contribution of the cuts. Simplifying as explained above the notations, we have
$`Z_{2N}`$ $`\underset{N\text{large}}{}{\displaystyle _{\text{circle around }K_{}}}{\displaystyle \frac{dK}{2i\pi }}{\displaystyle \frac{Z(K)}{K^{N+1}}}+O\left({\displaystyle \frac{1}{K_{deloc}^N}}\right)`$ (65)
$`\underset{N\text{large}}{}{\displaystyle \frac{1}{K_{}^N}}\left(\rho (0)+O\left({\displaystyle \frac{K_{}}{K_{deloc}}}\right)^N\right)`$
where $`\rho (0)`$ given in (36) represents the probability to be at $`2n=0`$.
In the delocalized phase, the finite size properties are given by the contributions of the cuts
$`Z_{2N}`$ $`={\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dx}{2i\pi }}{\displaystyle \frac{1}{K_{deloc}^N(1+x)^{N+1}}}\left[Z(K_{deloc}(1+x+io))Z(K_{deloc}(1+xio))\right]`$ (66)
To obtain the dominant behavior at large $`N`$, we will set $`v=N\mathrm{ln}(1+x)`$.
For $`T>T_c`$, using
$`Z(K_{deloc}(1+x+io))Z(K_{deloc}(1+xio))={\displaystyle \frac{2i\sqrt{x}(1+G(\beta ))}{G(\beta )^2}}\left(1+O(x)\right)`$ (67)
with
$`G(\beta )=e^{2\beta q_A}(1\sqrt{1e^{4\beta q_0}})1`$ (68)
we get the asymptotic behavior
$`Z_{2N}`$ $`\underset{N\text{large}}{}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dv}{2Ni\pi }}{\displaystyle \frac{1}{K_{deloc}^N}}e^v\left[Z(K_{deloc}(e^{\frac{v}{N}}+io))Z(K_{deloc}(e^{\frac{v}{N}}io))\right]`$ (70)
$`\underset{N\text{large}}{}{\displaystyle \frac{1}{K_{deloc}^N\sqrt{\pi }N^{3/2}}}{\displaystyle \frac{(1+G(\beta ))}{2G(\beta )^2}}`$
For $`T=T_c`$, $`G(\beta _c)`$ vanishes and we have
$`Z(K_{deloc}(1+x+io))Z(K_{deloc}(1+xio))=i{\displaystyle \frac{2}{\sqrt{x}}}\left(1+O(x)\right)`$ (71)
leading to
$`Z_{2N}`$ $`\underset{N\text{large}}{}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dv}{2Ni\pi }}{\displaystyle \frac{1}{K_{deloc}^N}}e^v\left[Z(K_{deloc}(e^{\frac{v}{N}}+io))Z(K_{deloc}(e^{\frac{v}{N}}io))\right]`$ (73)
$`\underset{N\text{large}}{}{\displaystyle \frac{1}{K_{deloc}^N\sqrt{\pi N}}}`$
This shows that the probability of presence at the interface decays with the length $`N`$ as $`1/N^{1/2}`$ at $`T=T_c`$ (73), whereas it decays as $`1/N^{3/2}`$ for $`T>T_c`$.
In the following, we will need the general scaling form which contains (70) and (73) as special cases
$`Z_{2N}\underset{N\text{large}}{}{\displaystyle \frac{1+G(\beta )}{K_{deloc}^N\sqrt{N}H(\beta )}}q\left(\sqrt{N}{\displaystyle \frac{G(\beta )}{\sqrt{H(\beta )}}}\right)`$ (74)
where the scaling function $`q(x)`$ reads
$`q(x)={\displaystyle \frac{1}{\pi }}{\displaystyle _0^+\mathrm{}}𝑑v{\displaystyle \frac{e^v\sqrt{v}}{v+x^2}}`$ (75)
and
$`H(\beta )={\displaystyle \frac{e^{4\beta q_B}(1e^{2\beta q_A})}{\sqrt{1e^{4\beta q_0}}(1+\sqrt{1e^{4\beta q_0}})^2}}`$ (76)
and where $`G(\beta )`$ has been defined in (68).
In particular, near $`T_c`$, we have for $`TT_c`$ the scaling form
$`Z_{2N}\underset{N\text{large},TT_c^+}{}{\displaystyle \frac{1}{K_{deloc}^N\sqrt{N}}}q\left({\displaystyle \frac{(TT_c)}{T_c}}\sqrt{NC_{\mathrm{}}(T_c^{})}\right)`$ (77)
leading to the finite size free-energy
$`f_{2N}={\displaystyle \frac{T\mathrm{ln}Z_{2N}}{2N}}\underset{N\text{large},TT_c^+}{}f_{deloc}(T_c)+{\displaystyle \frac{T_c}{4N}}\mathrm{ln}N{\displaystyle \frac{T_c}{2N}}\mathrm{ln}q\left({\displaystyle \frac{(TT_c)}{T_c}}\sqrt{NC_{\mathrm{}}^{}(T_c)}\right)`$ (78)
#### 2 Specific heat
In the thermodynamic limit, the specific heat $`C_{\mathrm{}}(T)=T\frac{d^2f(T)}{dT^2}`$ presents a jump at $`T_c`$ given in equation (27). To compute the finite-size specific heat $`C_{2N}(T_c)`$ at $`T=T_c`$, we use the scaling form (78)
$`C_{2N}(T_c)=`$ $`T{\displaystyle \frac{d^2f_{2N}}{dT^2}}|_{T=T_c}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{q^{\prime \prime }(0)}{q(0)}}\left({\displaystyle \frac{q^{}(0)}{q(0)}}\right)^2\right)C_{\mathrm{}}(T_c^{})+{\displaystyle \frac{q^{}(0)}{q(0)}}{\displaystyle \frac{\sqrt{C_{\mathrm{}}(T_c^{})}}{\sqrt{N}}}+..`$ (79)
$`=`$ $`{\displaystyle \frac{4\pi }{2}}C_{\mathrm{}}(T_c^{})\sqrt{\pi }{\displaystyle \frac{\sqrt{C_{\mathrm{}}(T_c^{})}}{\sqrt{N}}}+..`$ (80)
Similar relations between $`C_{2N}(T_c)`$ and $`C_{\mathrm{}}(T_c)`$ are found in phase transitions displaying a jump in the specific heat .
#### 3 Order parameter and uniform susceptibility
A possible order parameter at the transition has been mentionned in equation (41), and linked to the probability to be in fluid $`\beta `$. We consider here
$`M_{2N}={\displaystyle \underset{i=1}{\overset{2N}{}}}\text{ sgn}\left[z(i){\displaystyle \frac{1}{2}}\right]`$ (81)
and the order parameter $`m_{2N}=\frac{<M_{2N}>}{2N}`$. Note that $`m_{2N}`$ is related to the quantity $`_\beta `$ of (41) by ($`m_{2N}=12_\beta `$). Given the expression of the Hamiltonian (1), we get
$`m_{\mathrm{}}=<\text{ sgn}\left[z(i){\displaystyle \frac{1}{2}}\right]>=`$ $`_hf(T,q_Aq_A+h,q_Bq_Bh)|_{h=0}`$ (82)
$`=`$ $`{\displaystyle \frac{\mathrm{sinh}\beta (q_Aq_B)}{2\mathrm{sinh}\beta q_A\mathrm{sinh}\beta q_B}}\underset{TT_c^{}}{}1Cte(T_cT)+..`$ (83)
In the thermodynamic limit, the corresponding susceptibility reads
$`\chi _{\mathrm{}}={\displaystyle \frac{1}{2TN}}\left(<M_{2N}^2>(<M_{2N}>)^2\right)=`$ $`_h^2f(T,q_Aq_A+h,q_Bq_Bh)|_{h=0}`$ (84)
$`=`$ $`{\displaystyle \frac{1}{2T}}\left({\displaystyle \frac{1}{\mathrm{sinh}^2\beta q_A}}+{\displaystyle \frac{1}{\mathrm{sinh}^2\beta q_B}}\right)`$ (85)
leading, at criticality, to
$`\chi _{\mathrm{}}(T_c^{})={\displaystyle \frac{4\beta _c}{(1e^{2\beta _cq_A})^2}}`$ (86)
To compute the finite-size behavior of the order parameter $`m_{2N}`$ and of the susceptibility $`\chi _{2N}`$, we use again the scaling form (78), but where now $`T_c`$ is a function of $`h`$ defined by the equation for $`T_c`$ where $`q_Aq_A+h,q_Bq_Bh`$. Since the calculations exactly parallel the ones described above, we only quote the results
$`m_{2N}(T_c)=1{\displaystyle \frac{\sqrt{2\pi }}{(1e^{2\beta _cq_A})}}{\displaystyle \frac{1}{\sqrt{2N}}}\mathrm{}`$ (87)
and
$`\chi _{2N}(T_c)={\displaystyle \frac{4\pi }{2}}\chi _{\mathrm{}}(T_c^{})`$ (88)
#### 4 Numerical results
We have done numerical calculations using transfer matrix methods with both chain ends fixed at the interface. We chose $`q_A=1`$ and $`q_B=0.5`$, yielding a critical temperature $`T_c=1/(\mathrm{ln}(\frac{1+\sqrt{5}}{2}))=2.07809\mathrm{}`$ (see equation (23)). The corresponding values of the critical specific heat and susceptibility are $`C_{\mathrm{}}(T_c^{})=0.11056\mathrm{}`$ and $`\chi _{\mathrm{}}(T_c^{})=5.03932\mathrm{}`$ (see equations (27), (86)).
Results for $`C_{2N}(T)`$ and $`\chi _{2N}(T)`$ in the critical region are shown in Figures 1 and 2, for chain length up to $`2N=18000`$ and are in quantitative agreement with the above values. We get for instance a value $`C_{2N}(T_c)0.04739..`$ as compared to the value $`0.047454\mathrm{}`$ coming from equation (80), and $`\chi _{2N}(T_c)2.13376\mathrm{}`$ as compared to the value $`2.16289\mathrm{}`$ coming from equation (88). We have also tested equation (87): Figure 3 shows the behaviour of $`\sqrt{2N}\frac{(1m_{2N}(T))}{2}`$ in the critical region. At $`T_c`$, the theoretical value coming from equation (87) is 2.0279… as compared to the “experimental” value 2.0130.. of Figure 3.
## III The wetting transition
### A Calculation of the grand canonical partition function
We consider the SOS version of the 2D wetting problem near a attracting wall at $`h=0`$. Our presentation will be rather brief, and we follow the notations of and . Let us denote by $`\{h(i),i=1,2,..,N\}`$ the $`N`$ consecutive heigths with the properties $`h(i)=0,1,2\mathrm{}`$ and $`h(i+1)h(i)=+1,0,1`$. This height model can also be considered as describing the adsorption of a polymer chain onto the wall. We further assume that the first height is fixed ($`h(i=1)=1`$). The partition function of the model reads
$`Z_N={\displaystyle \underset{\{h(i)\}}{}}\mathrm{exp}\left(\beta J{\displaystyle \underset{i=1}{\overset{N1}{}}}|h(i+1)h(i)|+\beta u_0{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _{h_i,0}\right)`$ (89)
Following closely the steps of the previous section, we first express the partition function $`Z_N`$ as a function of its end-point ($`h(N)=h`$), and write
$`Z_N={\displaystyle \underset{h=0}{\overset{+\mathrm{}}{}}}Z_N(h)`$ (90)
At this stage, it is convenient to introduce some notations and relations.
Let $`(l)`$ represent the partition function for a chain starting from $`h=1`$ and arriving at $`h=0`$ for the first time after $`l`$ steps. We then denote by $`(l1,h)`$ the partition function for a chain of $`(l1)`$ steps going from $`h=1`$ to $`h1`$ in the presence of an absorbing boundary at $`h=0`$. Setting $`t=e^{\beta J}`$, the method of images yields
$`(l)=t\left(R(l1,0)R(l1,2)\right)`$ (91)
and
$`(l1,h)=R(l1,h1)R(l1,h+1)`$ (92)
where $`R(l,h)`$ denotes the partition function of the chain in the absence of the wall. We explicitly have
$`R(l,h)={\displaystyle \underset{p=|h|}{\overset{l}{}}}C_l^pC_p^{\frac{p+|h|}{2}}t^p`$ (93)
normalized to $`_{h=l}^{+l}R(l,h)=(1+2t)^l`$.
We now decompose $`Z_N(h)`$ into free and adsorbed segments as
$`Z_N(h)=`$ $`{\displaystyle \underset{p=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{l_{2p}=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\delta (N({\displaystyle \underset{i=1}{\overset{2p}{}}}l_i+l))`$ (95)
$`(l_1)(e^{\beta u_0l_2}t)(l_3)(e^{\beta u_0l_4}t)\mathrm{}(e^{\beta u_0l_{2p}}t)(l1,h)\text{for}h1`$
$`Z_N(0)=`$ $`{\displaystyle \underset{p=0}{\overset{+\mathrm{}}{}}}{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{l_{2p}=1}{\overset{\mathrm{}}{}}}\delta (N({\displaystyle \underset{i=1}{\overset{2p}{}}}l_i+l))(l_1)(e^{\beta u_0l_2}t)\mathrm{}(l_{2p1})e^{\beta u_0l_{2p}}`$ (96)
It is again convenient to consider the grand canonical partition function
$`Z(K)={\displaystyle \underset{N=1}{\overset{+\mathrm{}}{}}}K^NZ_N`$ (97)
Using
$`\widehat{}(K)={\displaystyle \underset{l=1}{\overset{+\mathrm{}}{}}}K^l(l)={\displaystyle \frac{1\sqrt{14\left(\frac{tK}{1K}\right)^2}}{2\frac{tK}{1K}}}`$ (98)
we finally have
$`Z(K,2h)`$ $`={\displaystyle \frac{1}{1\widehat{}(K)\frac{tKe^{\beta u_0}}{1Ke^{\beta u_0}}}}{\displaystyle \frac{1}{t}}\left({\displaystyle \frac{1K}{tK}}\widehat{}(K)1\right)^h`$ (99)
$`Z(K,2h+1)`$ $`={\displaystyle \frac{1}{1\widehat{}(K)\frac{tKe^{\beta u_0}}{1Ke^{\beta u_0}}}}{\displaystyle \frac{1}{t}}\widehat{}(K)\left({\displaystyle \frac{1K}{tK}}\widehat{}(K)1\right)^h`$ (100)
### B Thermodynamics
Equations (99) and (100) give
$`Z(K)={\displaystyle \underset{h=0}{\overset{+\mathrm{}}{}}}Z(K,h)={\displaystyle \frac{1}{1\widehat{}(K)\frac{tKe^{\beta u_0}}{1Ke^{\beta u_0}}}}{\displaystyle \frac{1+\widehat{}(K)}{t}}{\displaystyle \frac{1}{2\frac{1K}{tK}\widehat{}(K)}}`$ (101)
The radius of convergence $`K_{}(\beta )`$ of $`Z(K)`$ is either given by the solution of the equation $`\left(2\frac{1K}{tK}\widehat{}(K)=0\right)`$, or by the solution of the equation $`\left(1\widehat{}(K)\frac{tKe^{\beta u_0}}{1Ke^{\beta u_0}}=0\right)`$. The former corresponds to a value
$`K_{deloc}={\displaystyle \frac{1}{1+2t}}`$ (102)
and physically describes the delocalized phase. The latter is given by
$`K_{}(\beta )={\displaystyle \frac{2e^{\beta u_0}}{1+\sqrt{\frac{e^{\beta u_0}1+4t^2}{e^{\beta u_0}1}}}}`$ (103)
The chain is localized close to the wall as long as $`K_{}(\beta )<K_{deloc}`$, i.e. $`T<T_c=1/\beta _c`$ with
$`e^{\beta _cu_0}={\displaystyle \frac{1+2t}{1+t}}`$ (104)
in agreement with references .
### C Density profile
We follow equation (29) and define a normalized density profile by
$`\rho (h)=\underset{N\mathrm{}}{lim}{\displaystyle \frac{Z_N(h)}{Z_N}}`$ (105)
Inverting equation (99, we get
$`\rho (h)=\rho (0)\left[{\displaystyle \frac{1}{2t}}\left(\sqrt{{\displaystyle \frac{e^{\beta u_0}1+4t^2}{e^{\beta u_0}1}}}1\right)\right]^h`$ (106)
in the localized phase. The characteristic length reads
$`\xi (T)={\displaystyle \frac{1}{\mathrm{ln}\left[\frac{1}{2t}\left(\sqrt{\frac{e^{\beta u_0}1+4t^2}{e^{\beta u_0}1}}1\right)\right]}}`$ (107)
and diverges at the transition as
$`\xi (T)\underset{TT_c}{}{\displaystyle \frac{1}{T_cT}}`$ (108)
in agreement with equation (7.8) of reference . The fraction of adsorbed monomers vanishes as
$`\rho (h=0)\underset{TT_c}{}T_cT`$ (109)
### D Probability distributions of adsorbed and desorbed segments
Since the detailed procedure has been given in section II E, our presentation will be rather sketchy, and we only give the main results. Let us denote by A (resp. D) the adsorbed (resp. desorbed) monomers of the chain. The probability distribution of desorbed loop lengths will be denoted by $`P_D(l)`$, and its Laplace transform by $`\widehat{P_D}(s)`$. We get
$`\widehat{P_D}(s)={\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}e^{sl}P_D(l)={\displaystyle \frac{\widehat{}(K_{}e^s)}{\widehat{}(K_{})}}`$ (110)
yielding, at criticality, to
$`\widehat{P_D}(s)\underset{T=T_c}{=}\widehat{}({\displaystyle \frac{1}{1+2t}}e^s)=1\sqrt{{\displaystyle \frac{1+2t}{t}}}\sqrt{s}+O(s)`$ (111)
This corresponds to an algebraic decay of $`P_D(l)`$ for large $`(l)`$, as in equation (53). Similarly, in the region $`TT_c`$ the moments of $`P_D(l)`$ diverge as
$`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}l^nP_D(l)\underset{TT_c}{}(T_cT)\left({\displaystyle \frac{1}{(T_cT)^2}}\right)^n`$ (112)
As in the copolymer case, the correlation length along the chain scales as $`\left(\frac{1}{(T_cT)^2}\right)`$, in agreement with equation (6.22) of reference . As previously found, the small loops have a finite weight at $`T_c`$. One has for instance
$`P_D(l=1)\underset{T=T_c}{=}{\displaystyle \frac{t}{1+2t}}`$ (113)
As for the probability distribution of the adsorbed segment lengths, we find an exponential form, even at $`T_c`$.
$`P_A(l)=(1K_{}e^{\beta u_0})(K_{}e^{\beta u_0})^{l1}`$ (114)
where $`K_{}`$ is given in equation (103).
We conclude this section on wetting by pointing out that a finite size relation for the specific heat, quite similar to the one derived in equation (80), has been found in reference .
## IV Conclusion
We have presented a grand canonical approach to the localization of a single polymer chain either at a fluid-fluid interface, or at an attracting hard wall. This method, which rests on the absence of self avoidance in the models, provides detailed informations on the loop length distribution function. The extension to various disordered situations would be of great interest.
## A Transfer matrix approach
In this appendix, we will briefly present the transfer matrix calculations for the periodic $`AB`$ copolymer chain. We decompose $`Z_N`$ according to the position $`n`$ of the end point of the chain, ($`Z_N=_{n=\mathrm{}}^+\mathrm{}Z_N(n)`$), and using the Hamiltonian of equation (1), we obtain the recursion relations
$`Z_{N+1}(n)=e^{\beta q_{N+1}\mathrm{sgn}[n\frac{1}{2}]}\left(Z_N(n+1)+Z_N(n1)\right)`$ (A1)
It is convenient to introduce the Fourier transforms for the two half-spaces
$`\widehat{Z}_N^+(k)={\displaystyle \underset{n=1}{\overset{+\mathrm{}}{}}}e^{ik(n1)}Z_N(n)`$ (A2)
$`\widehat{Z}_N^{}(k)={\displaystyle \underset{n=\mathrm{}}{\overset{0}{}}}e^{ikn}Z_N(n)`$ (A3)
as well as their inverse
$`Z_N(n)={\displaystyle _0^{2\pi }}{\displaystyle \frac{dk}{2\pi }}e^{ik(n1)}\widehat{Z}_N^+(k)\text{for}n1`$ (A4)
$`Z_N(n)={\displaystyle _0^{2\pi }}{\displaystyle \frac{dk}{2\pi }}e^{ikn}\widehat{Z}_N^{}(k)\text{for}n0`$ (A5)
The recursion relations now read
$`\widehat{Z}_{N+1}^+(k)=e^{\beta q_{N+1}}\left(2\mathrm{cos}k\widehat{Z}_N^+(k)e^{ik}Z_N(1)+Z_N(0)\right)`$ (A6)
$`\widehat{Z}_{N+1}^{}(k)=e^{\beta q_{N+1}}\left(2\mathrm{cos}k\widehat{Z}_N^{}(k)e^{ik}Z_N(0)Z_N(1)\right)`$ (A7)
To study the thermodynamic limit for our problem, ($`q_{2p+1}=q_A>0`$ and $`q_{2p}=q_B<0`$), we look for a stationnary solution $`\widehat{Z}_{2N}^\pm (k)`$ with eigenvalue $`\lambda `$. We find
$`\widehat{Z}^+(k)={\displaystyle \frac{1}{\lambda e^{\beta (q_Bq_A)}4\mathrm{cos}^2k}}[e^{ik}Z(2)+e^{ik}Z(0)]`$ (A9)
$`\widehat{Z}^{}(k)={\displaystyle \frac{1}{\lambda e^{\beta (q_Bq_A)}4\mathrm{cos}^2k}}[(e^{2\beta q_A}1e^{i2k})Z(0)+e^{2\beta q_A}Z(2)]`$
Using the inversion relations (A5), we obtain the following consistency equations for $`Z(2)`$ and $`Z(0)`$
$`Z(2)=I_2^+(\lambda )Z(2)+I_0^+(\lambda )Z(0)`$ (A10)
$`Z(0)=e^{2\beta q_A}I_0^{}(\lambda )Z(2)+\left((e^{2\beta q_A}1)I_0^{}(\lambda )I_2^{}(\lambda )\right)Z(0)`$ (A11)
whith the notations
$`I_p^\pm (\lambda )={\displaystyle \frac{1}{4}}J_p\left(y={\displaystyle \frac{\lambda e^{\pm \beta (q_Bq_A)}}{4}}\right)`$ (A12)
$`J_0(y)={\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle \frac{1}{y\mathrm{cos}^2k}}={\displaystyle \frac{1}{\sqrt{y(y1)}}}`$ (A13)
$`J_2(y)={\displaystyle _\pi ^{+\pi }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle \frac{\mathrm{cos}2k}{y\mathrm{cos}^2k}}=2\sqrt{{\displaystyle \frac{y}{y1}}}2{\displaystyle \frac{1}{\sqrt{y(y1)}}}`$ (A14)
For a localized eigenstate, both $`Z(2)`$ and $`Z(0)`$ are non zero. As a consequence, the determinant of the system (A11) has to vanish, and this gives an equation for the eigenvalue $`\lambda `$ of the localized state, which reads in our case ($`q_Aq_B`$)
$`\lambda ={\displaystyle \frac{1}{K_{}(\beta )}}`$ (A15)
where $`K_{}(\beta )`$ is given in equation (22). The free energy is then given by
$`f(T)=\underset{N\mathrm{}}{lim}{\displaystyle \frac{\mathrm{ln}Z_{2N}}{\beta 2N}}={\displaystyle \frac{\mathrm{ln}\lambda }{2\beta }}`$ (A16)
in agreement with equation (4).
As a final remark on this approach, we mention how one may find the polymer density profile $`\rho (2n)`$ in the localized phase. The Fourier transforms of $`\rho (2n)`$ on the two-half spaces are indeed directly related to the localized eigenvector $`(\widehat{Z}^+(k))`$ given in (A9) by
$`{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}e^{ik(2n1)}\rho (2n)={\displaystyle \frac{\widehat{Z}^+(k)}{\widehat{Z}^+(0)+\widehat{Z}^{}(0)}}`$ (A17)
$`{\displaystyle \underset{n=\mathrm{}}{\overset{0}{}}}e^{ik2n}\rho (2n)={\displaystyle \frac{\widehat{Z}^{}(k)}{\widehat{Z}^+(0)+\widehat{Z}^{}(0)}}`$ (A18)
in agreement with the result (36). All the results of section II can be thus recovered, except for the results of the loop length distribution function of section II E.
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# Early galaxy evolution from deep wide field star counts
## 1 Introduction
This paper is part of a global analysis of star counts developed to constrain consistently scenarii of galaxy formation and evolution.The central tool of this approach is the ”Besançon” model of population synthesis. This model is gradually tuned to fit an increasing number of observational constraints while keeping compatibility with previous fits and theoretical prescriptions. In the present paper we address the problem of the halo (dark or visible) by trying to compare the properties of the spheroid population (the visible halo) with the dark matter halo, as traced by microlensing at high galactic latitudes and by the rotation curve. If the dark matter is made at least partly of stellar remnants, as shown by recent statistics of microlensing at high galactic latitudes Aubourg et al. (1993); Alcock et al. (1997, 1998), the density trend of this matter should be close to a power law with index of 2 (as expected from a flat rotation curve). It is natural to think of a similar shape for the stellar spheroid.
Constraints on the overall shape of the dark halo are poor. Cosmological simulations of halo formation generally predict that halos are flattened by about c/a $``$ 0.7 Rix (1996). But the axis ratio depends on how much the halo matter is dissipative, the more dissipative, the flatter the halo. Direct determinations of the dark matter distribution in polar ring galaxies show flattened halos with c/a $``$ 0.5 Sackett et al. (1994); Rix (1996).
Concerning the spheroid population, most previous analyses suggest rather steep density slopes with power indices between 3.0 and 3.5. However, these analyses are based on rather small samples of well identified tracers. The estimated flattening also cover a wide range between 0.6 to 1.0:
The distribution of galactic globular clusters appears to be well fitted by a power law density with index $`n3.5`$ and flattening of 1. Harris (1976); Zinn (1985). Hawkins RR Lyrae observations Hawkins (1984) showed $`n=3.1\pm 0.2`$ with a flattening of 0.9. Saha Saha (1985), using a spherically symmetric model, found $`n3`$ out to 25 kpc but then the RR Lyrae density falls off more rapidly beyond 25 kpc. Another study of RR Lyrae by Wetterer Wetterer and McGraw (1996) showed that a spherically symmetric model yields $`n3`$ whereas an ellipsoidal distribution yields $`n3.5`$. Sluis Sluis and Arnold (1998) counted blue horizontal branch (BHB) stars and RR Lyrae and found $`c/a0.5`$ and $`n=3.2\pm 0.3`$. Still from BHB star counts, Sommer-Larsen Sommer-Larsen and Christensen (1987) derived $`c/a0.8`$ and $`n3`$ up to 40 kpc, Preston Preston et al. (1991) found that $`c/a`$ increases from 0.5 to 1 up to 20 kpc with $`n=3.5`$. Soubiran Soubiran (1993) showed that $`n=3.5\pm 0.5`$ is compatible with the kinematical behavior of a star sample near the north galactic pole. K dwarf counts with HST yield $`c/a=0.8\pm 0.1`$ and $`n=3.06\pm 0.22`$ Gould et al. (1998). All of these studies were based on a few hundred objects at most.
In order to find new constraints on the spheroid density law, we undertook a photometric and astrometric sample survey in various galactic directions. We complemented these data with existing deep photometric star counts in several high and intermediate latitude fields. Most such counts contain large numbers of halo dwarfs, but they cannot be distinguished from thick disc dwarfs by their colours but at faint magnitudes. Since no large optical surveys were available at magnitude fainter than 20, we used heterogeneous data coming from various studies (often of extragalactic aim) in various photometric systems.
The population synthesis model used here permits to perform a global analysis of these heterogeneous data, since observational data can be simulated in each field with the true observational conditions (photometric system, errors and selection effects). The synthetic approach allows also to estimate the biases and expected contaminations by other populations.
In section 2 we describe the model of population synthesis and external constraints on the spheroid population. In section 3 we describe the data sets and the comparison method. In section 4 we discuss the results and their implications for the dark matter halo.
## 2 The model of population synthesis
We have used a revised version of the Besançon model of population synthesis. Previous versions were described in Bienaymé et al. Bienaymé et al. (1987b, a) and Haywood et al. Haywood et al. (1997a).
The model is based on a semi-empirical approach, where physical constraints and current knowledge of the formation and evolution scenario of the Galaxy are used as a first approximation for the population synthesis. The model involves 4 populations (disc, thick disc, halo and bulge) each deserving a specific treatment. The bulge population which is irrelevant for this spheroid analysis will be described elsewhere.
### 2.1 The disc population
A standard evolution model is used to produce the disc population, based on a set of usual parameters : an initial mass function (IMF), a star formation rate (SFR), a set of evolutionary tracks (see Haywood et al., 1997 and references therein). The disc population is assumed to evolve during 10 Gyr. A set of IMF slopes and SFR’s are tentatively assumed and tested against star counts. The tuning of disc parameters against relevant observational data was described in Haywood et al. Haywood et al. (1997b, a).
The model fixes the distribution of stars in the space of intrinsic parameters : effective temperature, gravity, absolute magnitude, mass and age. These parameters are converted into colours in various systems through stellar atmosphere models corrected to fit empirical data Lejeune et al. (1997, 1998). While some errors still remain in the resulting colours for some spectral types, the overall agreement is good in the major part of the HR diagram.
Since the Haywood et al. model was based on evolutionary tracks at solar metallicities, inverse blanketing corrections are introduced to give to the disc a metallicity distribution in agreement with Twarog Twarog (1980) age/metallicity distribution (mean and dispersion about the mean).
The model returns the present-day distribution of stars as a function of intrinsic parameters in a unit volume column centered at the sun position. Since the evolution model does not account for orbital evolution, stars are redistributed in the reference volume over the z axis. The key for redistributing stars along the z-axis is age : an empirical relation associates z velocity dispersions to ages. Then the Boltzmann equation is used to convert z velocity distributions into z density. The model is dynamically self-consistent in the sense that the potential used in the Boltzmann equation is the one generated by the total mass distribution of stellar populations. The self consistency is established iteratively. We slice the disc populations into seven isothermal populations of different ages, from 0 to 10 Gyr. Each sub-population (except the youngest one, which cannot be considered as relaxed) has its velocity dispersion imposed by the age/velocity dispersion relation. We then deduce the scale height of each sub-population using the Boltzmann equation. The overall scheme is described in Bienaymé Bienaymé et al. (1987b).
Resulting density laws are used to correct the evolution model distribution in and off the plane, then to compute the stellar densities all over the Galaxy.
### 2.2 The thick disc population
A detailed analysis of the thick disc population from photometric and astrometric star counts has been given elsewhere Ojha et al. (1994b, a, 1996); Robin et al. (1996); Ojha et al. (1999). The kinematics, metallicity, and density law were measured allowing us to constrain the origin for this population. In this series of papers, evidence was given that the majority of thick disc stars should originate from a merging event at the beginning of the life of the thin disc, after the first collapse. One or several satellite galaxies have heated the thin disc, then the gas re-collapsed and reformed a new thin disc Robin et al. (1996).
In the population synthesis process, the thick disc population is modeled as originating from a single epoch of star formation. We use Bergbush & Vandenberg Bergbush and VandenBerg (1992) oxygen enhanced evolutionary tracks. No strong constraint exists on the thick disc age until now. We assume an age of 11 Gyr, which is slightly older than the disc and younger than the halo. The initial mass function is modeled by a simple power law with a slope about $`\alpha =12`$, referring to the notation $`\varphi (m)m^\alpha `$.
The thick disc metallicity can be chosen between -0.4 and -1.5 dex in the simulations. The standard value of -0.7 dex is usually adopted, following in situ spectroscopic determination from Gilmore et al. Gilmore et al. (1995) and photometric star count determinations Robin et al. (1996); Buser et al. (1999). The low metallicity tail of the thick disc seems to represent a weak contribution to general star counts Morrison (1993b). It was neglected here. An internal metallicity dispersion among the thick disc population is allowed. The standard value for this dispersion is 0.25 dex. No evidence has been found for a significant metallicity gradient in the thick disc population (Robin et al, 1996).
The thick disc density law is assumed to be a truncated exponential : at large distances the law is exponential. At short distances it is a parabola. This formula ensures the continuity and derivability of the density law (contrarily to a true exponential) and eases the computation of the potential. The scale height of the exponential can vary between 600 and 2600 pc. The standard value, 760 pc, has been obtained from star count fitting in various directions (Robin et al, 1996). Nevertheless, it can be shown that star counts when restricted to a small number of galactic directions and a small magnitude range do not give a strong constraint on the scale height, but rather on the parameter : (local density)$`\times `$(scale height)<sup>2</sup>. At present there is no accurate determination of the thick disc density in the solar neighbourhood, independently from the scale height. But reasonable values range between 700 to 1200 pc for the scale height and 1 to 4% for the local density relative to the thin disc.
### 2.3 The spheroid
We assume a homogeneous population of spheroid stars with a short period of star formation. We thus use the Bergbush and Vandenberg Bergbush and VandenBerg (1992) oxygen enhanced models, assuming an age of roughly 14 Gyr (until more constraints on the age are available), a mean metallicity of -1.7 dex and a dispersion of 0.25 about this value. No galactocentric gradient is assumed. The IMF has to be constrained either from globular clusters (if they are representative of the spheroid population) or from deep star counts. This point is discussed in section 3.
The density of spheroid stars is modeled by a power law :
$$\rho (R,z)=\rho _0\times (R^2+\frac{z^2}{ϵ^2})^{n/2}$$
where $`\rho _0`$ is the local density, $`n`$ is the power law index and $`ϵ`$ is the flattening.
The local density can be constrained by local measurements of high velocity stars, or by remote counts of giants (spectroscopically selected) or dwarfs (photometrically selected). The local density cannot be determined independently from the other density parameters with our limited number of data sets. Thus we have used independent constraints from the literature on the local spheroid density.
### 2.4 The local spheroid density
The local stellar spheroid density, $`\rho _0`$, is bounded by observational data on halo dwarfs and giants. Figure 1 shows the luminosity function obtained by different authors. We only selected recent results obtained in good conditions from sufficiently large samples. Bahcall et al Bahcall and Casertano (1986) and Gizis & Reid Gizis and Reid (1999) derived their values from high proper motion dwarf samples. Dahn et al Dahn et al. (1995) determined accurate parallaxes for local late-type subdwarfs and deduced the local luminosity function of halo stars in the absolute visual magnitude range 9 to 14. These three results are biased by the kinematic selection. They took the bias correction into account but this correction is model dependent and introduces an unknown uncertainty into the result. We expect that the differences between the three measurements rely upon this correction. On the giant side, Morrison Morrison (1993a) used a non kinematically-biased sample of halo giants, selected from their metallicity to estimate the spheroid local density. In figure 1a we show the luminosity function from Bergbush and Vandenberg (1992) for a population of 14 Gyr with a metallicity of -1.75 and an IMF slope $`\alpha `$ of 2. If we let the local halo density vary from a factor of 0.75 to 1.25 relative to this reference model (dotted lines in figure 1a), we get a good agreement with the specified observations given their uncertainties. In the next section we allow the local density to vary within these limits.
## 3 Data sets and fitting methods
Obtaining good constraints on the spheroid density law requires a good photometric accuracy. This generally depends on using CCD detectors on large telescopes, on fields as wide as possible to cover large samples, and a large range of galactocentric distances. This can be obtained with a number of data sets at various galactic latitudes and longitudes. We have collected such data sets from the literature. Most have been made for extragalactic purposes.
### 3.1 Available data
The main data characteristics are summarized in table 1. The photometric systems are close to the Johnson-Cousins system. Spheroid star selection was based on their magnitude and colour (either B-V, V-R, or V-I, depending on the available observations), in order to avoid presence of contamination by other populations. Aiming at model independent results, the model was used essentially to select colour and magnitude ranges and fields where the contamination by thick disc stars remains negligible under any reasonable thick disc hypothesis. For this reason all data brighter than magnitude 20 at intermediate and low latitudes were excluded. A small number of disc white dwarfs is also present in the selection but the proportion is at most a few percent and has no consequence on the result.
Our survey program include at the moment two fields, one towards the north galactic pole, another at intermediate latitude (l=150,b=60). The NGP field is the deepest up to now : it is complete and free from galaxy contamination up to magnitude 24. A full description of these data sets will be given in a forthcoming paper.
The other selected data sets are the six fields of the DMS survey Hall et al. (1996); Osmer et al. (1998) observed in V and R bands at medium latitude, 4 fields from the Canada-France Redshift survey (CFRS, Lilly et al. (1995); Le Fevre et al. (1995); Hammer et al. (1995)) dedicated to galaxy counts, two fields from the Koo and Kron investigation for quasars Koo and Kron (1982); Koo et al. (1986), another field from Reid and Majewski near the north galactic pole Reid and Majewski (1993).
The absolute visual magnitude of halo stars in the selected samples ranges between 3 and 8, except our north galactic pole field which reaches M$`{}_{V}{}^{}11`$. All these fields taken together cover a large part of the (R,z) plane, as can be seen in figure 2 where the distributions in R and z of 90% of halo stars in each field of view are drawn.
### 3.2 Analysis method
Population synthesis simulations have been computed in every observed field using photometric errors as close as possible to the true observational errors, generally with photometric errors growing as a function of the magnitude and assumed to be Gaussian. Monte Carlo simulations are done in a solid angle much larger than the data in order to minimize the Poisson noise.
Then we compare the number of stars produced by the model with the observations in the selected region of the plane (magnitude, colour) and we compute the likelihood of the observed data to be a realization of the model (following the method described in Bienaymé et al, 1987a, appendix C). The likelihood has been computed for a set of models, varying the power law index between 2.0 and 3.5, the flattening between 0.3 and 1.0, the local density between 0.5 and 1.25 times the standard value as defined in section 2.4, and the spheroid IMF slope $`\alpha `$ from 1.0 to 2.2.
The confidence limits of estimated parameters are determined by the likelihood level which can be reached by pure random change of the sample : a series of simulated random samples are produced using the set of model parameters. The rms dispersion of the likelihood about the mean of this series gives an estimate of the likelihood fluctuations due to the random noise. It is then used to compute the confidence limit. Resulting errors are not strictly speaking standard errors, they give only an order of magnitude.
## 4 Results and discussion
### 4.1 Constraints on the spheroid density law
Figure 3 gives the value of the likelihood as a function of the flattening, power law index and local normalization. On the left, iso-likelihood contours are drawn for four values of the local normalization (0.5, 0.75, 1.0 and 1.25 $`\times \rho _0`$). On the right, we show the likelihood values as a function of the power law index for the best fit value of the flattening.
Comparing the results of different local normalization we conclude that the choice of the local normalization sensitively displaces the best fit power law index and flattening, but their likelihood are not similar. The best fit model is obtained either with a local density of 0.75 $`\rho _0`$, a power law index of 2.44 and an axis ratio of 0.76, or a local density of 0.5 $`\rho _0`$, a power law index of 2.24 and an axis ratio of 0.86. The values obtained with a standard local density $`\rho _0`$ are slightly worse but stay within 1 sigma confidence level. They are a power law index of 2.62 and an axis ratio of 0.70. The best fit local densities 0.5$`\rho _0`$ and 0.75$`\rho _0`$ agrees with the Bahcall & Casertano determination of the local luminosity function as seen in figure 1, but conflict with Dahn et al., which favors a local density of 1.25 $`\rho _0`$. However, in the present study the statistics is dominated by stars with absolute magnitudes in the range 3 to 8, a range poorly represented in the Dahn et al. sample. Only deeper counts could give constraints on the fainter part of the luminosity function.
It is worth having a look at the colour distributions as predicted by the best fit model compared with the observational data. Figure 4 and 5 show the colour distributions observed (dots) and predicted (heavy solid line) by the best fit model ($`ϵ`$=0.76, n=2.44, 0.75$`\rho _0`$) in the selected magnitude interval in each tested field. Superimposed we show the distribution of the spheroid population alone as predicted by the model (light solid line). We see here that some photometric systems are not closely matched by the model, as seen by slight shifts between model and data in some cases. But the way we have selected spheroid stars in the blue peak of the distribution cannot introduce a bias even in case of colour shifts.
### 4.2 Sensitivity to the IMF slope
Whatever the assumed IMF slope in the range 1-2.2, the maximum likelihood is obtained for the same density law parameters. There is a slight likelihood variation related to the choice of the IMF, but it is only due to the deepest magnitude bin towards the pole.
A separate analysis of star counts deeper than 22 towards the pole can help determine precisely the IMF slope. In this magnitude range spheroid stars with absolute magnitude 10-11 contribute substantially to star counts, while their frequency is sensitive to the IMF as can be seen on figure 1b. The analysis is slightly different from the determination of the density law. In this range of absolute magnitude the subdwarf sequence turns redwards making the colour index a good luminosity indicator. The V-I distribution is used as an additional constraint. A V-I histogram is built with a bin 0.1 magnitude wide over the range 0. to 3. The density law is adopted from the above analysis, so the free parameters are the halo and thick disc IMF slopes. Since these two populations are quite well separated in the (V,V-I) plane, the two IMF slope estimates are de-correlated. Table 2 gives the resulting slope estimates with their likelihood in the magnitude range 22-24 for the different spheroid density laws determined previously.
Spheroid models with a local density 0.75 $`\rho _0`$ and 1.0 $`\rho _0`$ give the maximum likelihood, well in agreement with the previous result. However the model with 0.50 $`\rho _0`$ is noticeably worse. Eventually, the resulting IMF slopes do not depend significantly on the assumed density laws and the likelihood is well peaked around the maximum indicating a robust determination. We conclude that the IMF of the halo, in the mass range \[0.1, 0.8\] is :
$$\varphi (m)m^{1.9}$$
while the IMF slope of the thick disc seems to be slightly smaller and similar to the disc’s Haywood (1994). These values do not account for binarity. Thus the true IMF should slightly steepen. We leave the value uncorrected until more data are available on the binary fraction among low mass spheroid stars.
This result is the first direct measurement of the mass function of field star spheroid with a good statistics, thanks to the wide field of the CCD mosaic. Several previous determinations used kinematically selected samples (see section 2.4 for references) or deep fields. But the latter were limited to narrow fields : the first attempt by Richer & Fahlman Richer and Fahlman (1992) lead to a very steep IMF slope of $`\alpha =4.5\pm 1.2`$ which had given hope for a dark matter halo of brown dwarfs. Later results have given shallower slopes but the uncertainties were not significantly decreased.
Gould et al. Gould et al. (1998) analyzed a sample of 166 stars in 53 HST WFPC fields, making difficult the de-correlation between structural parameters of the spheroid and its mass function. They found a luminosity function down by a factor two from the present one and deduced an IMF slope of $`\alpha =0.75`$ (in our notation). Their result relies upon the assumption that the spheroid has a mean metallicity of -1.0, which looks too high considering most direct measurement of its abundances. This high metallicity induces an overestimate of the luminosity at a given colour, hence of the distance, as well as an overestimate of the mass relatively to a smaller assumed metallicity.
### 4.3 Variations of density law with galactocentric position
If we independently check the results obtained in inner fields and in outer fields, we are able to search for solutions with varying power law index and flattening over the galactic radius. Contrarily to Preston et al. (1991) we find no evidence for varying power law index or flattening. However, a round spheroid is ruled out by the inner field data as well as by fields at low latitudes. Thus our results are compatible with a true power law and a constant flattening all along the tested galactic radius.
When comparing data sets from different sources in close galactic fields, discrepancies appear which are larger than expected on the basis of pure random noise. This may be due either to data incompleteness, or to systematic errors in the photometry (including mismatch of the standard photometric system), or to true inhomogeneities in the spheroid distribution. Currently available data are not sufficient to discriminate between these different causes. Homogeneous wide field surveys will be necessary to clarify these aspects. The scope of the current investigation is for this reason limited to large scale average characteristics.
### 4.4 Contamination by other populations
The blue peak at these magnitudes may be contaminated by disc white dwarfs or by thick disc main sequence stars. The former are very few compared to the density of the halo. The contamination by disc white dwarfs, as determined by the model, is at most 5% in the magnitude range 22-24.
The contamination from the thick disc has been estimated using our best fit thick disc model as adjusted on medium deep star counts. The contamination can reach about 30% in the magnitude range 18-20 but becomes negligible at magnitude larger than 20 as seen in figure 6. Hence, we do not take into account magnitudes lower than 20 in our study.
Had the thick disc contamination been underestimated, then the contribution assigned to the halo in the blue peak would be too large, resulting in a possible distortion of the density law. In order to evaluate how this would affect our conclusions, we have investigated different thick disc models which could fake the halo contribution to the blue peak. Attempts were limited to realistic thick discs roughly fitting the red peak. We have selected two extreme thick disc parameters for which the contamination to the blue peak becomes significant. A thick disc with a local density of 3.9%, a scale height of 1150 pc and an IMF slope of 1. (referred to as model B). A thick disc with a scale height of 2 kpc, a local density of 0.5% and an IMF slope of 1.75. (model C). With such thick disc models, the process of adjusting the spheroid density law parameters end up to tiny local density of about 25% of the standard value and very small power law index of the order of 1.5. Surprisingly, the fit is good on star counts up to magnitude 22, showing that a large range of parameters can reproduce a wide set of star counts. However at magnitude 22-24, model B and C are unable to reproduce the counts, as exemplified in figure 7, where at the top star counts at the pole in the magnitude range 22-24 are overestimated in the blue peak, and in the range 18-22 (bottom) the fit is still acceptable.
So, the degeneracy between thick disc models and halo parameters holds only if star counts are not deep enough. Keeping reasonable values for the thick disc parameters leads to a small contamination with no risk of underestimation of the halo density. If the thick disc contribution is higher than expected from standard models then we would overestimate the local halo density and power law index, strengthening our conclusion towards a flat spheroid with a small power law index.
## 5 Conclusions
We have shown that, with the data available up to now, the spheroid star distribution follows a power law with an index smaller than previously thought. It is moderately flattened, as already found by several investigations. The best fit power law index is found 2.44 for a flattening of 0.76. We cannot exclude a spheroid population having a power law as low as 2 and an axis ratio of 0.5 at the 2 sigma level. Assuming a high exponent of 3.5, as suggested by some globular cluster and RR Lyrae data, model predictions deviate significantly from observations in the external parts of the galaxy whatever reasonable flattening is adopted.
The IMF slope of the spheroid is found to be $`\alpha =1.9\pm 0.2`$, value which gives a local density of 1.64 10<sup>-4</sup> stars pc<sup>-3</sup> and a mass density of 4.15 10<sup>-5</sup> Mpc<sup>-3</sup> for the stellar halo, yet this value ignores possible old white dwarfs. With this slope the expected mass density of brown dwarfs in the halo makes a negligible part of the dark matter halo, as already estimated from microlensing surveys.
Recent searches for ancient halo white dwarfs have given a hope to identify the microlensing events with such objects Ibata et al. (1999, 2000); Hodgkin et al. (2000). Ibata et al. Ibata et al. (2000) conclude that old white dwarfs may constitute a significant fraction of about 10% of the dark matter halo. In the mean time, microlensing experiments have narrowed the range of the estimated halo baryonic fraction to 20 to 50% Alcock et al. (2000). These two results are well in agreement according to the uncertainties.
So, as star count data progresses in depth and extent, the picture of the spheroid star population that comes out points to characteristics quite compatible with what we know about the distribution of baryonic dark matter if it is made of stellar remnants, suggesting a common dynamical origin. The visible spheroid and its heavy counterpart of dark remnants can make a significant but not dominant part of the so-called dark matter halo.
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# Closed cosmologies with a perfect fluid and a scalar field
## I Introduction
Cosmological models with a scalar field and an exponential potential are of fundamental importance in the study of the universe. These models are motivated by the fact that they arise naturally in alternative theories of gravity and occur as the low-energy limit in supergravity theories . By using qualitative techniques, the well-known power-law inflationary solution has been shown to be an attractor for all initially expanding Bianchi models (except a subclass of the Bianchi type IX models which will recollapse) in the class of spatially homogeneous Bianchi models . More recently, cosmological models which contain both barotropic matter and a scalar field with an exponential potential have been studied , partially motivated by the fact that there exist spatially flat isotropic scaling solutions in which the energy density due to the scalar field is proportional to the energy density of the perfect fluid . In the stability of these cosmological scaling solutions within the class of spatially homogeneous cosmological models with a perfect fluid subject to the equation of state $`p=(\gamma 1)\mu `$ (where $`\gamma `$ is a constant satisfying $`0<\gamma <2`$) was studied and it was found that when $`\gamma >2/3`$, and particularly for realistic matter with $`\gamma 1`$, the scaling solutions are unstable; essentially they are unstable to curvature perturbations, although they are stable to shear perturbations. In addition, in homogeneous and isotropic spacetimes with non-zero spatial curvature were studied.
It is clearly of interest to study more general cosmological models. One class of models of particular interest are those with positive spatial curvature. These models have attracted less attention since they are more complicated mathematically. Positive-curvature Friedmann-Robertson-Walker (FRW) models , Kantowski-Sachs models and Bianchi type IX models have been studied using qualitative methods, although rigorous analyses using a new set of compact variables have not been carried out. The Bianchi type IX models are known to have very complicated dynamics, exhibiting the characteristics of chaos , and are hence beyond the scope of the present study. Recently positive-curvature FRW models and Kantowski-Sachs models with a perfect fluid and a cosmological constant have been investigated using qualitative methods and utilizing compactified variables.
The outline of this paper is as follows. We shall first comprehensively study the qualitative properties of the class of positive-curvature FRW models with a barotropic fluid and a non-interacting scalar field with an exponential potential, extending and generalising work by Turner , who used different basic variables. We shall then analyse the qualitative properties of the Kantowski-Sachs models. Positive-curvature FRW models and Kantowski-Sachs models belong to the class of spherically symmetric models, and hence the present work is a natural extension of recent work . Indeed, it turns out that understanding the dynamics of the Kantowski-Sachs models is crucial for understanding the global dynamics of general spherically symmetric similarity models .
### A Matter model
The matter content of the models is taken to be a perfect fluid and a scalar field with exponential potential. The corresponding energy-momentum tensor is
$`T_{ab}`$ $`=`$ $`(T_{\mathrm{pf}ab}+T_{\mathrm{sf}ab}),`$ (1)
$`T_{\mathrm{pf}ab}`$ $`=`$ $`\mu u_au_b+p\left(u_au_b+g_{ab}\right),`$ (2)
$`T_{\mathrm{sf}ab}`$ $`=`$ $`\varphi _{,a}\varphi _{,b}\left({\displaystyle \frac{1}{2}}\varphi _{,c}\varphi ^{,c}+V(\varphi )\right)g_{ab},`$ (3)
$`V(\varphi )`$ $`=`$ $`V_0e^{\kappa \varphi },`$ (4)
where $`\kappa `$ is a non-negative constant, and the pressure is given by $`p=(\gamma 1)\mu `$, with the equation-of-state parameter in the range $`1\gamma 2`$. The fluid energy density $`\mu `$ and the scalar field $`\varphi `$ are functions of a timelike coordinate $`t`$. A dot denotes differentiation with respect to $`t`$ and throughout units are used in which $`c=8\pi G=1`$. The matter components are assumed to be non-coupled, and thus they are separately conserved:
$`_aT_{\mathrm{pf}}^{ab}=`$ $`0`$ $`=_aT_{\mathrm{sf}}^{ab}.`$ (5)
For convenience, we define
$$X=\frac{1}{\sqrt{2}}\dot{\varphi }.$$
(6)
## II Closed Friedmann models
We start our investigation of closed cosmological models with a perfect fluid and a scalar field by looking at the closed FRW models. The line element for these models can be written
$$ds^2=dt^2+S(t)^2dr^2+S(t)^2\mathrm{sin}^2rd\mathrm{\Omega }^2.$$
(7)
The expansion of the fluid congruence is given by $`\theta =3\dot{S}/S`$, and the evolution equation for the curvature $`K9/(\theta ^2S^2)`$ is
$$\dot{K}=\frac{2K}{3\theta }(3\dot{\theta }+\theta ^2).$$
(8)
The conservation equations yield
$`\dot{\mu }`$ $`=`$ $`\gamma \theta \mu ,`$ (9)
$`\dot{X}`$ $`=`$ $`\theta X+{\displaystyle \frac{\kappa }{\sqrt{2}}}V.`$ (10)
From the field equations we obtain
$`\mu `$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[(1+K)\theta ^23X^23V\right],`$ (11)
$`\dot{\theta }`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left\{\theta ^2+{\displaystyle \frac{1}{2}}D^2+{\displaystyle \frac{3}{2}}\left[3X^23V+3(\gamma 1)\mu \right]\right\}.`$ (12)
Assuming $`\mu 0`$, the Friedmann equation, Eq. (11), shows that $`D=\sqrt{(1+K)\theta ^2}`$ is a dominant quantity. Thus, compact variables can be defined according to
$$Q_0=\frac{\theta }{D},U=\frac{\sqrt{3}X}{D},W=\frac{\sqrt{3V}}{D}.$$
(13)
Note also that the curvature is given by
$$K=\frac{1Q_0^2}{Q_0^2}.$$
(14)
The Friedmann equation becomes
$$\mathrm{\Omega }_D=\frac{3\mu }{D^2}=1U^2W^2.$$
(15)
Defining a new independent variable, $`{}_{}{}^{}=d/d\tau =\frac{3}{D}d/dt`$, the evolution equation for $`D`$
$`D^{}`$ $`=`$ $`3Q_0\left(U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)D`$ (16)
decouples. Thus, a reduced set of evolution equations is obtained:
$`Q_0^{}`$ $`=`$ $`(1Q_0^2)\left[13\left(U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right],`$ (17)
$`U^{}`$ $`=`$ $`3Q_0U\left[1+\left(U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right]+\sqrt{{\displaystyle \frac{3}{2}}}\kappa W^2,`$ (18)
$`W^{}`$ $`=`$ $`3Q_0\left(U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)W\sqrt{{\displaystyle \frac{3}{2}}}\kappa UW.`$ (19)
There is also an auxiliary evolution equation
$$\mathrm{\Omega }_D^{}=3Q_0\left[(1\mathrm{\Omega }_D)\gamma 2U^2\right]\mathrm{\Omega }_D,$$
(20)
and it is straight-forward to consider the set of variables $`(Q_0,U,\mathrm{\Omega }_D)`$, rather than $`(Q_0,U,W)`$ . Note that by setting $`\kappa =0,U=0`$, and identifying $`\mathrm{\Lambda }=V_0`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=W^2`$, the evolution equations corresponding to closed FRW models with a cosmological constant are obtained .
It is also useful to consider the deceleration parameter, given by
$$q_{\mathrm{pf}}\left(1+3\frac{u_{\mathrm{pf}}^a_a\theta _{\mathrm{pf}}}{\theta _{\mathrm{pf}}^2}\right)=\frac{1}{Q_0^2}\left[13\left(U^2+\frac{\gamma }{2}\mathrm{\Omega }_D\right)\right],$$
(21)
for $`Q_00`$. From this expression, we can see that there is an inflationary region ($`q_{\mathrm{pf}}<0`$) in the state space whenever $`\mathrm{\Omega }_D<\frac{2}{3\gamma }(13U^2)`$. However, as will be seen below, it is only when $`\kappa ^2<2`$ that there exist attractors that are inflationary. Note also that for $`Q_00`$,
$$Q_0^{}=(1Q_0^2)Q_0^2q_{\mathrm{pf}},$$
(22)
so that $`Q_0^{}<0`$ whenever $`q_{\mathrm{pf}}>0`$ in which case $`Q_0`$ is itself monotonic. When $`q_{\mathrm{pf}}<0`$, that is in the inflationary region, $`Q_0`$ need not be monotonic – for example, see the orbits close to $`{}_{+}{}^{}\mathrm{\Phi }`$ in Figs. 1 and 2.
The dynamical system Eqs. (18) is symmetric under the transformation
$$(\tau ,Q_0,U,W)(\tau ,Q_0,U,W).$$
(23)
Thus, it is sufficient to discuss the behaviour in one part of the state space, the dynamics in the other part being obtained by Eq. (23).
Furthermore, note that
$`M`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_D}{1Q_0^2}}`$ (24)
$`M^{}`$ $`=`$ $`(3\gamma 2)Q_0M`$ (25)
is a monotonic function in the regions $`Q_0<0`$ and $`Q_0>0`$ for $`\mathrm{\Omega }_D0`$. As there are no equilibrium points with $`Q_0=0`$ when $`\gamma >2/3`$, $`M`$ acts as a monotonic function in the interior of the state space. Consequently there can be no periodic or recurrent orbits in the interior state space and global results can be deduced. In addition, from the expression for the monotonic function $`M`$ we can see immediately that either $`Q_0^21`$ or $`\mathrm{\Omega }_D0`$ asymptotically.
### A Equilibrium points of the closed FRW dynamical system
A number of equilibrium points can be found for the dynamical system, Eqs. (18). In what follows, $`ϵ=\pm 1`$ denotes the sign of $`Q_0`$, whereas $`\mathrm{\Omega }_\varphi =U^2+W^2`$ is a density parameter associated with the scalar field. In Table I, the various equilibrium points are summarized. The subscripts on the labels have the following significance: The left subscript gives the sign of $`Q_0`$ and indicates whether the corresponding model is expanding (+) or contracting ($``$). The right subscript gives the sign of $`U`$; i.e., the sign of $`\dot{\varphi }`$.
The equilibrium points labeled $`{}_{\pm }{}^{}\mathrm{F}`$ correspond to the flat Friedmann solution. For these points, the scalar field vanishes ($`U=0=W`$). There is an orbit from $`{}_{+}{}^{}\mathrm{F}`$ to $`{}_{}{}^{}\mathrm{F}`$ and this orbit represents the closed FRW solution with no scalar field, starting from a Big Bang at $`{}_{+}{}^{}\mathrm{F}`$ and recollapsing to a “Big Crunch” at $`{}_{}{}^{}\mathrm{F}`$.
The K points represent exact solutions with a massless scalar field ($`W=0`$). As the fluid is negligible ($`\mathrm{\Omega }_D=0`$), these solutions are dominated by the kinetic term $`U`$. They correspond to Jacobs analogues of Kasner solutions in which $`U`$ takes on the role of a shearing mode .
There are equilibrium points with non-vanishing potential, where the scalar field dominates ($`\mathrm{\Omega }_D=0`$). These points $`{}_{\pm }{}^{}\mathrm{\Phi }`$ are only physical when $`\kappa ^2<6`$. For $`\kappa ^2=6`$, $`{}_{+}{}^{}\mathrm{\Phi }`$ coincides with $`{}_{+}{}^{}\mathrm{K}_{+}^{}`$, and $`{}_{}{}^{}\mathrm{\Phi }`$ with $`{}_{}{}^{}\mathrm{K}_{}^{}`$. For $`\kappa ^2>6`$, they are outside the physical part of the state space. The equilibrium point with $`ϵ=+1`$ (i.e., $`{}_{+}{}^{}\mathrm{\Phi }`$) is a sink, and for $`\kappa ^2<2`$ it corresponds to the power-law inflationary attractor solution (cf. the expression for $`q_{\mathrm{pf}}`$ in Table II).
There are also points $`{}_{\pm }{}^{}\mathrm{X}`$ for which the matter is unimportant, but the curvature is non-vanishing ($`Q_0^21`$), and tracks the scalar field. The corresponding solutions are called curvature scaling solutions . These solutions only exist when $`\kappa ^2<2`$. For $`\kappa ^2=2`$, $`{}_{+}{}^{}\mathrm{X}`$ coincides with $`{}_{+}{}^{}\mathrm{\Phi }`$, and $`{}_{}{}^{}\mathrm{X}`$ with $`{}_{}{}^{}\mathrm{\Phi }`$. Above this value of $`\kappa `$, these equilibrium points are outside the physical part of the state space.
When $`\kappa ^2>3\gamma `$, there is a flat matter scaling solution, for which both the fluid and the scalar field are dynamically important. The corresponding equilibrium points are denoted $`{}_{\pm }{}^{}\mathrm{FS}`$, and for $`\kappa ^2=3\gamma `$ they coincide with $`{}_{\pm }{}^{}\mathrm{\Phi }`$.
Finally, there are equilibrium points $`\mathrm{S}_\pm `$ corresponding to static solutions, analogous to the Einstein static universe. These are only physical when $`\gamma <2/3`$. For $`\gamma =2/3`$, a set of equilibrium points appears along the line $`U=0=W`$, signaling a change of stability when the points $`\mathrm{S}_\pm `$ leave the physical state space. In what follows, we will only consider equations of state for which $`1\gamma 2`$, hence we will not consider these points further.
Table II presents some physical quantities for the various equilibrium points, and Table III lists their eigenvalues. The sources and sinks of the dynamical system when $`\gamma >2/3`$ are listed in Table IV, and the global behaviour for different values of $`\kappa `$ can be summarized as follows: The state space when $`0<\kappa ^2<2/3`$ is depicted in Fig. 1, where the features in the rear part of the state space have been suppressed. Dashed and full curves represent orbits in the boundary submanifolds, while dotted curves represent orbits in the interior. Both $`{}_{+}{}^{}\mathrm{K}_{}^{}`$ and $`{}_{+}{}^{}\mathrm{K}_{+}^{}`$ are past attractors, while $`{}_{}{}^{}\mathrm{K}_{+}^{}`$ and $`{}_{+}{}^{}\mathrm{\Phi }`$ act as future attractors. There are also orbits from $`{}_{+}{}^{}\mathrm{K}_{}^{}`$, whose future attractor is $`{}_{}{}^{}\mathrm{K}_{}^{}`$ at the rear of the figure. Note that orbits future asymptotic to $`{}_{+}{}^{}\mathrm{\Phi }`$ correspond to solutions that exhibit power-law inflation ($`1<q_{\mathrm{pf}}<0`$, see Table II). Observe that the outgoing eigenvector directions from the saddle point $`{}_{+}{}^{}\mathrm{F}`$ span a separatrix surface in the interior of the state space. Similarly, $`{}_{+}{}^{}\mathrm{X}`$ is a saddle for which the ingoing eigenvector directions span another separatrix surface. These separatices confine orbits in the interior state space to specific regions. For example, there is one region where all orbits are past asymptotic to $`{}_{+}{}^{}\mathrm{K}_{+}^{}`$ and future asymptotic to $`{}_{}{}^{}\mathrm{K}_{+}^{}`$.
When $`\kappa ^2>2/3`$, the separatrix surface associated with $`{}_{+}{}^{}\mathrm{X}`$ changes structure, see Fig. 2. This is easiest seen by considering the separatrix orbits in the $`\mathrm{\Omega }_D=0`$ submanifold, see Fig. 3. When $`\kappa ^2<2/3`$, there is one separatrix orbit from $`{}_{}{}^{}\mathrm{\Phi }`$ to $`{}_{+}{}^{}\mathrm{X}`$, and one from $`{}_{}{}^{}\mathrm{X}`$ to $`{}_{+}{}^{}\mathrm{\Phi }`$. For $`\kappa ^2=2/3`$, these two orbits coalesce into a single orbit from $`{}_{}{}^{}\mathrm{X}`$ to $`{}_{+}{}^{}\mathrm{X}`$. This orbit corresponds to a special “Bouncing Universe” solution, existing only for this particular value of $`\kappa `$. When $`\kappa ^2>2/3`$, the separatrix orbit to $`{}_{+}{}^{}\mathrm{X}`$ starts at $`{}_{+}{}^{}\mathrm{K}_{}^{}`$, and the orbit from $`{}_{}{}^{}\mathrm{X}`$ goes to $`{}_{}{}^{}\mathrm{K}_{+}^{}`$. Thus, there is a bifurcation for $`\kappa ^2=2/3`$; however, we note that there is no stability change of equilibrium points involved. A similar behaviour of separatrix surfaces has been found for Bianchi type IX models .
When $`\kappa `$ increases, the equilibrium point $`{}_{+}{}^{}\mathrm{X}`$ approaches $`{}_{+}{}^{}\mathrm{\Phi }`$. For $`\kappa ^2=2`$ these two points coincide, and the state space for $`\kappa ^2>2`$ is depicted in Fig. 4. Both $`{}_{+}{}^{}\mathrm{K}_{}^{}`$ and $`{}_{+}{}^{}\mathrm{K}_{+}^{}`$ still act as past attractors, while $`{}_{}{}^{}\mathrm{K}_{+}^{}`$ is a future attractor. However, the stability of $`{}_{+}{}^{}\mathrm{\Phi }`$ has changed; it has become a saddle. There is still a separatrix surface associated with $`{}_{+}{}^{}\mathrm{F}`$. Thus, orbits having $`{}_{+}{}^{}\mathrm{K}_{+}^{}`$ as their past attractor all end at $`{}_{}{}^{}\mathrm{K}_{+}^{}`$.
For $`\kappa ^2>3\gamma `$, the equilibrium point $`{}_{+}{}^{}\mathrm{FS}`$, corresponding to the matter-scaling solution, appears from $`{}_{+}{}^{}\mathrm{\Phi }`$, see Fig. 5. The point $`{}_{+}{}^{}\mathrm{FS}`$ is a spiral sink with an out-going eigenvector direction entering the interior state space. Note that this equilibrium point thus is stable in the flat ($`Q_0=1`$) submanifold, but unstable to curvature perturbations (i.e., perturbations in the $`Q_0`$ direction). The scalar-field dominated point $`{}_{+}{}^{}\mathrm{\Phi }`$ still is a saddle, and now there is a separatrix surface spanned by the out-going eigenvector directions there. Thus, there are two separatrix surfaces, both of which are spiraling around the out-going eigenvector direction of $`{}_{+}{}^{}\mathrm{FS}`$.
When $`\kappa `$ increases, the $`{}_{+}{}^{}\mathrm{\Phi }`$ equilibrium point comes closer and closer to $`{}_{+}{}^{}\mathrm{K}_{+}^{}`$, and for $`\kappa ^2=6`$ they coincide. The state space when $`\kappa ^2>6`$ is given in Fig. 6. There is only one past and one future attractor, namely $`{}_{+}{}^{}\mathrm{K}_{}^{}`$ and $`{}_{}{}^{}\mathrm{K}_{+}^{}`$, respectively.
To summarize, when $`\kappa ^2>2`$ all solutions start from and recollape to a singularity ($`KK`$). Thus, in this case solutions can neither expand forever nor inflate. When $`\kappa ^2<2`$, there are also ever-expanding ($`\mathrm{K}\mathrm{\Phi }`$) (and ever-collapsing $`\mathrm{\Phi }\mathrm{K}`$) solutions in addition to the recollapsing solutions. Inflation occurs when $`\ddot{S}>0`$, i.e., $`3\dot{\theta }+\theta ^2>0`$, which leads to the condition
$$(3\gamma 2)3(2\gamma )U^2+3\gamma W^2>0.$$
(26)
This corresponds to a parabolic region along the ridge of the state spaces, Figs. 1, 2 and 46. The only equilibrium points within this region are $`{}_{\pm }{}^{}\mathrm{\Phi }`$ for $`\kappa ^2<2`$, corresponding to power-law inflation. Consequently, in the case $`\kappa ^2<2`$ there is a subclass of solutions that inflate.
## III Kantowski-Sachs models
We now turn our attention to the Kantowski-Sachs (KS) models. The line element can be written
$$ds^2=dt^2+D_1(t)^2dx^2+D_2(t)^2d\mathrm{\Omega }^2,$$
(27)
where
$$D_1=\mathrm{exp}\left[\beta ^0(t)2\beta ^+(t)\right],D_2=\mathrm{exp}\left[\beta ^0(t)+\beta ^+(t)\right].$$
(28)
The kinematic quantities of the fluid congruence are related to the Misner variables ($`\beta ^0`$, $`\beta ^+`$) by
$$\theta =3\dot{\beta }^0,\sigma _+=3\dot{\beta }^+,$$
(29)
and the evolution equations for the metric functions $`B_1D_1^1`$ and $`B_2D_2^1`$ become
$$\dot{B}_1=\frac{1}{3}(\theta +\sigma _+)B_1,\dot{B}_2=\frac{1}{3}(\theta +\sigma _+)B_2.$$
(30)
The conservation equations give
$`\dot{\mu }`$ $`=`$ $`\gamma \theta \mu ,`$ (31)
$`\dot{X}`$ $`=`$ $`\theta X+{\displaystyle \frac{\kappa }{\sqrt{2}}}V,`$ (32)
and the field equations yield
$`\mu `$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(\theta ^2\sigma _+^2+3B_2^23X^23V\right),`$ (33)
$`\dot{\theta }`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(\theta ^2+2\sigma _+^2+6X^23V+{\displaystyle \frac{3}{2}}(3\gamma 2)\mu \right),`$ (34)
$`\dot{\sigma }_+`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(\theta ^2\sigma _+^23\theta \sigma _+3X^23V3\mu \right).`$ (35)
The Friedmann equation, Eq. (33), together with the assumption $`\mu 0`$ shows that $`D=\sqrt{\theta ^2+3B_2^2}`$ is a dominant quantity. Consequently, compact variables are introduced according to
$$Q_0=\frac{\theta }{D},Q_+=\frac{\sigma _+}{D},U=\frac{\sqrt{3}X}{D},W=\frac{\sqrt{3V}}{D}.$$
(36)
The curvature variable $`K=3B_2^2\theta ^2=(1Q_0^2)/Q_0^2`$ shows that the flat solutions correspond to $`Q_0^2=1`$. The Friedmann equation becomes
$$\mathrm{\Omega }_D=\frac{3\mu }{D^2}=1Q_+^2U^2W^2.$$
(37)
By introducing a new independent variable, $`\tau `$, where $`{}_{}{}^{}=d/d\tau =\frac{3}{D}d/dt`$, the evolution equation for $`D`$,
$$D^{}=\left[Q_+(1Q_0^2)+3Q_0\left(Q_+^2+U^2+\frac{\gamma }{2}\mathrm{\Omega }_D\right)\right]D,$$
(38)
decouples, and a reduced set of evolution equations is obtained:
$`Q_0^{}`$ $`=`$ $`(1Q_0^2)\left[1+Q_0Q_+3\left(Q_+^2+U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right],`$ (39)
$`Q_+^{}`$ $`=`$ $`(1Q_0^2)(1Q_+^2)+3Q_0Q_+\left[1+\left(Q_+^2+U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right],`$ (40)
$`U^{}`$ $`=`$ $`U\left\{(1Q_0^2)Q_++3Q_0\left[1+\left(Q_+^2+U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right]\right\}+\sqrt{{\displaystyle \frac{3}{2}}}\kappa W^2,`$ (41)
$`W^{}`$ $`=`$ $`W\left\{(1Q_0^2)Q_++3Q_0\left(Q_+^2+U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right\}\sqrt{{\displaystyle \frac{3}{2}}}\kappa UW.`$ (42)
There is also an auxiliary evolution equation:
$`\mathrm{\Omega }_D^{}`$ $`=`$ $`\mathrm{\Omega }_D\left\{3\gamma Q_02\left[Q_+(1Q_0^2)+3Q_0\left(Q_+^2+U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right]\right\}.`$ (43)
Note that by setting $`\kappa =0,U=0`$, and identifying $`\mathrm{\Lambda }=V_0`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=W^2`$, the evolution equation corresponding to Kantowski-Sachs models with a cosmological constant are obtained . The deceleration parameter is given by
$`q_{\mathrm{pf}}`$ $`=`$ $`{\displaystyle \frac{1}{Q_0^2}}\left[13\left(Q_+^2+U^2+{\displaystyle \frac{\gamma }{2}}\mathrm{\Omega }_D\right)\right].`$ (44)
Note that the dynamical system Eqs. (40) is symmetric under the transformation
$$(\tau ,Q_0,Q_+,U,W)(\tau ,Q_0,Q_+,U,W).$$
(45)
Thus, it is sufficient to discuss the behaviour in one part of the state space, the dynamics in the other part being obtained by Eq. (45).
The function
$`M`$ $`=`$ $`Q_+^{2(3\gamma 2)}(1Q_0^2)^{3(2\gamma )}\mathrm{\Omega }_D^4,`$ (46)
$`M^{}`$ $`=`$ $`2Q_+^1\left[(3\gamma 2)(1Q_0^2)+3(2\gamma )Q_+^2\right]M`$ (47)
is monotonic in the regions $`Q_+>0`$ and $`Q_+<0`$, since $`2/3<\gamma <2`$. Noting that
$$Q_+^{}|_{Q_+=0}=(1Q_0^2)<0,$$
(48)
we conclude that the submanifold $`Q_+=0`$ is not invariant, but acts as a membrane. Thus, the existence of $`M`$ rules out any periodic or recurrent orbits in the interior of the state space and again global results are possible. From the expression for the monotonic function $`M`$ we can immediately see that asymptotically $`Q_+0`$, $`Q_0^21`$ or $`\mathrm{\Omega }_D0`$.
### A Equilibrium points of the KS dynamical system
The dynamical system, Eqs. (40), has several equilibrium points, which are displayed in Table. V. As before, $`ϵ=\pm 1`$ denotes the sign of $`Q_0`$, while $`\mathrm{\Omega }_\varphi U^2+W^2`$. Again, the left subscript gives the sign of $`Q_0`$ and indicates whether the corresponding model is expanding or contracting. The values of $`\mathrm{\Omega }_D`$, $`\mathrm{\Omega }_\varphi `$, and $`q_{\mathrm{pf}}`$ for each of the equilibrium points are given in Table VI while the eigenvalues are displayed in Table VII. Note that all of the equilibrium points correspond to exact self-similar cosmological models .
As for the closed FRW models, $`{}_{\pm }{}^{}\mathrm{F}`$ denotes the flat Friedmann solution. Note that the closed FRW solution without a scalar field does not appear as a submanifold of the Kantowski-Sachs models without a scalar field. Consequently, there is no orbit connecting $`{}_{+}{}^{}\mathrm{F}`$ with $`{}_{}{}^{}\mathrm{F}`$.
There are two sets $`{}_{\pm }{}^{}\mathrm{K}`$ of vacuum ($`\mathrm{\Omega }_D=0`$) equilibrium points, parameterized by the constant $`U_0`$, corresponding to kinetic dominated solutions. These sets are analogues of the “Kasner rings” that are present for various Bianchi models.
The flat scalar-field dominated points $`{}_{\pm }{}^{}\mathrm{\Phi }`$, already encountered for the closed FRW models, appear in the Kantowski-Sachs case as well. As for the closed FRW models, they are physical when $`\kappa ^2<6`$ and inflationary when $`\kappa ^2<2`$.
There are also equilibrium points $`{}_{\pm }{}^{}\mathrm{\Xi }`$, corresponding to curvature scaling solutions (i.e. they have $`\mathrm{\Omega }_D=0`$, $`Q_0^2<1`$) which are physical when $`\kappa ^2<2`$. They are also inflationary, but in other respects they resemble the points $`{}_{\pm }{}^{}\mathrm{X}`$ of the closed FRW models.
As for the closed FRW models, the equilibrium points $`{}_{\pm }{}^{}\mathrm{FS}`$, corresponding to the flat matter-scaling solution, enter the physical part of the state space when $`\kappa ^2>3\gamma `$.
Finally, there are also equilibrium points corresponding to the self-similar Kantowski-Sachs solution. This solution is only physical when $`\gamma <2/3`$, and so we will not consider them further.
The eigenvalues for each of the equilibrium points are given in Table VII. The sources and sinks of the dynamical system when $`\gamma >2/3`$ are listed in Table VIII (all of the other equilibrium points are saddles). Thus, there is always two segments of the equilibrium set $`{}_{+}{}^{}\mathrm{K}`$ that act as sources for orbits. Similarly there are two segments on $`{}_{}{}^{}\mathrm{K}`$ that are sinks. When $`\kappa ^2>2`$, these are the only attractors, and all solutions start from and recollape to a singularity ($`KK`$).
When $`\kappa ^2<2`$, which then implies that $`\kappa ^2<3\gamma `$ for $`\gamma >2/3`$, the equilibrium points $`{}_{\pm }{}^{}\mathrm{\Phi }`$ are attractors. Thus, for $`\kappa ^2<2`$, there are also ever-expanding ($`{}_{+}{}^{}\mathrm{K}{}_{+}{}^{}\mathrm{\Phi }`$) and ever-collapsing ($`{}_{}{}^{}\mathrm{\Phi }{}_{}{}^{}\mathrm{K}`$) solutions.
From the expression for the monotonic function $`M`$ we deduce that all orbits asymptotically have $`Q_+0`$, $`Q_0^21`$ or $`\mathrm{\Omega }_D0`$. Indeed, the existence of the monotonic function ensures that there are no periodic orbits and that generically orbits asymptote towards the local attractors (sinks and sources). Therefore, we can determine the global dynamics of the models.
To summarize, when $`\kappa ^2>2`$ all solutions start from and recollape to a singularity ($`KK`$). Thus, in this case solutions can neither isotropize nor inflate. When $`\kappa ^2<2`$, there are also ever-expanding ($`\mathrm{K}\mathrm{\Phi }`$) (and ever-collapsing $`\mathrm{\Phi }\mathrm{K}`$) solutions in addition to the recollapsing solutions. Again, the $`\mathrm{\Phi }`$ points correspond to power-law inflation when $`\kappa ^2<2`$. Consequently, in this case there is a subclass of solutions that isotropize and inflate.
The global asymptotic dynamics is similar to that in the case of positive-curvature FRW models. However, due to the presence of shear, the intermediate or transient dynamics can be quite different. In the Kantowski-Sachs case the state space is four-dimensional and so we cannot display the phase portraits graphically (as in the FRW case). However, as an illustration we shall present the phase portraits in the three-dimensional fluid vacuum and the massless scalar field invariant sets in order to compare intermediate behaviours.
### B Fluid vacuum
The fluid vacuum ($`\mathrm{\Omega }_D=0`$) is an invariant submanifold, as seen from Eq. (43). Using the Friedmann equation to eliminate $`W`$, we obtain a three-dimensional dynamical system in $`(Q_0,Q_+,U)`$:
$`Q_0^{}`$ $`=`$ $`(1Q_0^2)\left[1+Q_0Q_+3(Q_+^2+U^2)\right],`$ (49)
$`Q_+^{}`$ $`=`$ $`(1Q_0^2)(1Q_+^2)3Q_0Q_+(1Q_+^2U^2),`$ (50)
$`U^{}`$ $`=`$ $`(1Q_0^2)Q_+U+\left(\sqrt{{\displaystyle \frac{3}{2}}}\kappa 3Q_0U\right)(1Q_+^2U^2).`$ (51)
From table VI, it is immediately seen that the equilibrium points that are contained in this submanifold are $`{}_{\pm }{}^{}\mathrm{K}`$, $`{}_{\pm }{}^{}\mathrm{\Phi }`$, and $`{}_{\pm }{}^{}\mathrm{\Xi }`$. The state space is depicted in Figs. 7, 8 and 9. Note that $`Z=0`$, where $`Z`$ is defined by
$$Z=Q_+\frac{1}{2}Q_0+\frac{\sqrt{6}}{2\kappa }U,$$
(52)
is an invariant submanifold, and that both $`{}_{\pm }{}^{}\mathrm{\Phi }`$ and $`{}_{\pm }{}^{}\mathrm{\Xi }`$ are contained in this submanifold.
### C Massless case
The massless case corresponds to the invariant submanifold $`W=0`$, which leads to a three-dimensional system in ($`Q_0,Q_+,U`$):
$`Q_0^{}`$ $`=`$ $`(1Q_0^2)\left[{\displaystyle \frac{3\gamma 2}{2}}Q_0Q_++{\displaystyle \frac{3}{2}}(2\gamma )(Q_+^2+U^2)\right],`$ (53)
$`Q_+^{}`$ $`=`$ $`(1Q_0^2)(1Q_+^2){\displaystyle \frac{3}{2}}(2\gamma )Q_0Q_+\mathrm{\Omega }_D,`$ (54)
$`U^{}`$ $`=`$ $`U\left[(1Q_0^2)Q_+{\displaystyle \frac{3}{2}}(2\gamma )Q_0\mathrm{\Omega }_D\right].`$ (55)
From table V, it is immediately seen that the equilibrium points that are contained in this submanifold are $`{}_{\pm }{}^{}\mathrm{K}`$ and $`{}_{\pm }{}^{}\mathrm{F}`$. The state space is depicted in Fig. 10.
## IV Discussion
We have studied closed cosmological models with a perfect fluid satisfying a linear equation of state with $`2/3<\gamma <2`$ and a scalar field with an exponential potential. We have utilized a new set of normalised variables which lead to the compactification of state space, enabling us to apply the theory of dynamical systems to determine the qualitative properties of the models. In all cases we have been able to find monotonic functions which, together with a local analysis of the equilibrium points, enable us to determine the global properties of the models.
We first studied the closed FRW cosmological models. We found that when $`\kappa ^2>2`$, all solutions start from and recollape to a singularity ($`KK`$). In this case solutions generically do not inflate. When $`\kappa ^2<2`$, solutions can either recollapse ($`KK`$) or expand forever ($`\mathrm{K}\mathrm{\Phi }`$) towards power-law inflation solutions (or collapse forever $`\mathrm{\Phi }\mathrm{K}`$); consequently, in this case there is a subclass of solutions that inflate. A number of phase portraits were displayed.
These results generalise previous qualitative work on positive-curvature FRW models with a scalar field (only) and with a scalar field plus a barotropic perfect fluid in which compactified variables were not utilized, and rigorous analyses of perfect fluid (only) models using compactified variables , and completes and generalises more recent work using different compactified variables . We also note that positive-curvature FRW models with a perfect fluid and a positive cosmological constant have been investigated recently using qualitative methods and utilizing compactified variables .
In the case of the Kantowski-Sachs models we again found that when $`\kappa ^2>2`$ all solutions start from and recollape to a singularity ($`KK`$) and can consequently neither isotropize nor inflate. When $`\kappa ^2<2`$, there are also ever-expanding ($`\mathrm{K}\mathrm{\Phi }`$) (and ever-collapsing $`\mathrm{\Phi }\mathrm{K}`$) solutions in addition to the recollapsing solutions, where again the $`\mathrm{\Phi }`$ points correspond to the flat FRW power-law inflationary solution. Consequently, in this case there is a subclass of solutions that isotropize and inflate.
The investigation of Kantowski-Sachs models complements the study of Bianchi models and completes the analysis of spatially homogeneous models. Collins studied perfect fluid Kantowski-Sachs models qualitatively using expansion-normalised variables (for which the state space was non-compact) and showed that all models start at a Big Bang and recollapse to a final “Big Crunch” singularity. This work was generalised recently by Goliath and Ellis in which Kantowski-Sachs models with a perfect fluid and a cosmological constant were investigated using qualitative methods and utilizing the compactified variables of Uggla and Zur-Muhlen ; particular attention was focussed upon whether the models isotropize, thereby explaining the presently observed near-isotropy of the universe. More importantly, Kantowski-Sachs models with a scalar field and an exponential potential, but without barotropic matter, have been studied qualitatively , although compactified variables were not utilized.
To conclude an analysis of positive-curvature spatially homogeneous cosmological models with a perfect fluid and a scalar field with an exponential potential, Bianchi type IX models would need to be studied. However, such a study is beyond the scope of the current paper. For example, Bianchi type IX models are known to have very complicated dynamics, exhibiting the characteristics of chaos . However, partial results are known. Bianchi type IX models with a scalar field (only) have been studied qualitatively, with an emphasis on whether these models can isotropize . Scalar-field models with matter have also been studied . For example, it has been shown that the power-law inflationary solution is an attractor for all initially expanding Bianchi type IX models except for a subclass of the models which recollapse . However, compact variables have not been utilized and the analyses were not rigorous.
A more rigorous treatment of the class of Bianchi type IX models with a non-tilted perfect fluid (only) using compactified variables has been possible . Although an appropriately defined normalised Hubble variable is found to be monotonic, enabling some results to be obtained, several problems remain open. More rigorous global results are possible. For example, Bianchi type IX models with matter have been shown to obey the “closed universe recollapse” conjecture , whereby initially expanding models enter a contracting phase and recollapse to a future “Big Crunch”. In addition, Ringström has proven that a curvature invariant is unbounded in the incomplete directions of inextendible null geodesics for generic vacuum Bianchi models , and rigorously shown that the Mixmaster attractor is the past attractor of Bianchi type IX models with an orthogonal perfect fluid . A complete qualitative analysis of the special class of locally rotationally symmetric Bianchi type IX perfect fluid models, which do not exhibit oscillatory or chaotic behaviour near to the initial or final singularities, has been given in , based upon an appropriately defined set of bounded variables.
The Kantowski-Sachs models exhibit similar global properties to the positive-curvature FRW models; in particular, for $`\kappa ^2>2`$ all initially expanding models reach a maximum expansion and thereafter recollapse, whereas for $`\kappa ^2<2`$ models generically recollapse or expand forever towards a flat isotropic power-law inflationary solution. The Bianchi type IX models share these qualitative properties. However, the intermediate behaviour of the Kantowski-Sachs models can be quite different to that of the FRW models. In order to illustrate the possible intermediate dynamics of the Kantowski-Sachs models, we studied the special cases of no barotropic fluid, and a massless scalar field in Secs. III B and III C, respectively (see Figs. 710).
Finally, we remark that the dynamics of the Kantowski-Sachs models obtained here will be crucial for understanding the global dynamics of general self-similar spherically symmetric models .
## V Acknowledgements
We would like to thank Peter Turner for his help with the analysis of the positive-curvature FRW models . AC would like to acknowledge financial support from NSERC of Canada. MG would like to thank the Department of Mathematics and Statistics at Dalhousie University for hospitality while this work was carried out.
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# The gauge invariant quark correlator in QCD sum rules and lattice QCD
## 1 Introduction
Gauge invariant field correlators can serve as interesting examples for studying non-perturbative aspects of QCD in the framework of lattice QCD or QCD sum rules . In addition, they are a natural extension of the local condensates which appear in the sum rule approach, due to the use of the operator product expansion (OPE) . Phenomenological implications of nonlocal condensates have previously been discussed in the literature . The two most fundamental correlators are the gauge invariant field strength correlator and the gauge invariant quark correlator. The gluon field strength correlator has already been investigated in lattice QCD and QCD sum rules . These studies allowed to extract the correlation length of the gluon field.
On the other hand, the gauge invariant quark correlator has so far only been measured on the lattice . In this work, we therefore first present a QCD sum rule analysis of the gauge invariant quark correlator. To this end, the quark correlator is related to a heavy–light meson correlator in the heavy quark effective theory (HQET) . Thus the sum rule analysis will be related to previous investigations of the same HQET correlators in the limit of the heavy quark mass going to infinity. The correlation length of the quark field is then given as the inverse of the binding energy of the light quark.
In the second, and more innovative, part of the present work, we perform a direct comparison of the lattice data with the representations of the correlation function in the QCD sum rule approach. From the so–called phenomenological parametrisation — a single resonance plus the perturbative continuum as the simplest Ansatz — we again extract the correlation length of the quark field, in agreement with the direct sum rule determination. Fitting the lattice data with the theoretical correlator in the framework of the operator product expansion, we obtain estimates of the quark and the mixed quark–gluon condensate which are in reasonable agreement with sum rule phenomenology.
Our paper is organised as follows: In the next section we discuss the relation of the gauge invariant quark correlator and the corresponding heavy quark current correlator. In section 3 we set up the expressions needed for the sum rule analysis and in section 4 we shall present our numerical results. Section 5 compares our results with recent lattice determinations of the quark correlator. In particular, a direct comparison of the OPE with the lattice data will be performed. Finally, in section 6, we conclude with a summary and an outlook.
## 2 The gauge invariant quark correlator
The central object of our investigation is the gauge invariant two–point correlation function of quark fields,
$$𝒟_{\alpha \beta }(z)0|T\{\overline{q}_\alpha ^a(y)\varphi ^{ab}(y,x)q_\beta ^b(x)\}|0,$$
(1)
where the string operator $`\varphi ^{ab}(y,x)`$ ensures the gauge invariance of the correlator and is represented by
$$\varphi ^{ab}(y,x)[𝒫e^{igT^C_0^1𝑑\lambda z^\mu A_\mu ^C(x+\lambda z)}]^{ab},$$
(2)
with $`z=yx`$. $`T^C`$ are the generators of SU(3) in the fundamental representation and $`𝒫`$ denotes path ordering of the exponential. In general, the gauge invariant field strength correlator could be defined using an arbitrary gauge string connecting the end points $`x`$ and $`y`$, but in this work we have restricted ourselves to a straight line, for it is only in this case that the correlator (1) is related to a heavy meson correlator in HQET.
Using the path–integral formalism, we are able to derive a relation between the correlator $`𝒟_{\alpha \beta }(z)`$ and the correlator of a local, gauge invariant current composed of a light quark field $`q_\alpha ^a(x)`$ and an infinitely heavy quark field $`h_\alpha ^a(x)`$ . Analogously to HQET, the heavy quark field $`h_\alpha ^a(x)`$ is constructed from the field $`Q_\alpha ^a(x)`$ with a finite mass $`m_Q`$ in the limit $`m_Q\mathrm{}`$,
$$h_\alpha ^a(x)=\underset{m_Q\mathrm{}}{lim}\frac{1}{2}(1+\overline{)}v)_{\alpha \beta }e^{im_Qvx}Q_\beta ^a(x),$$
(3)
where $`v_\mu `$ is the four–velocity of the heavy quark. In the case of infinitely heavy quarks, the coupling to the gauge fields is given by the effective HQET action $`S_{eff}=d^4x\overline{h}iv^\mu D_\mu h`$ with $`D_\mu =_\mu igA_\mu `$ . In the free field case the heavy quark propagator is given by
$`S_{\alpha \beta }^{ab}(z)`$ $``$ $`0|T\{h_\alpha ^a(y)\overline{h}_\beta ^b(x)\}|0=\delta ^{ab}{\displaystyle \frac{1}{2}}(1+\overline{)}v)_{\alpha \beta }S(z)`$ (4)
$`=`$ $`\delta ^{ab}{\displaystyle \frac{1}{2}}(1+\overline{)}v)_{\alpha \beta }{\displaystyle \frac{1}{v^0}}\theta (z^0)\delta \left(𝐳{\displaystyle \frac{z^0}{v^0}}𝐯\right),`$
where $`v^0`$ is the zero–component of the velocity. In addition, we have the relation $`v_\mu =z_\mu /|z|`$ with $`|z|\sqrt{z^2iϵ}`$ which follows immediately from the constraints $`𝐳=z^0𝐯/v^0`$ and $`v^2=1`$. In order to obtain a solution for the interacting case, we have to find the inverse of the operator $`iv^\mu D_\mu `$. Using the following crucial relation obeyed by the string operator,
$$v^\mu _\mu ^y\varphi (y,x)=v^\mu igT^CA_\mu ^C(y)\varphi (y,x),$$
(5)
together with equation (4) one can show that the solution for the interacting case is found to be
$$0|T\{h_\alpha ^a(y)\overline{h}_\beta ^b(x)e^{iS_{eff}}\}|0=S_{\alpha \beta }(z)0|[𝒫e^{igT^C_0^1𝑑\lambda z^\mu A_\mu ^C(x+\lambda z)}]^{ab}|0.$$
(6)
The physical interpretation of this result is an infinitely heavy quark moving from point $`x`$ to point $`y`$ with a four–velocity $`v`$, acquiring a phase proportional to the path–ordered exponential. The limit $`m_Q\mathrm{}`$ is necessary in order to constrain the heavy quark to a straight line and to decouple the spin interactions.
The result (6) can be employed to relate the correlator (1) to correlators of gauge invariant local currents in HQET. To this end, we define the pseudoscalar and scalar heavy meson currents as
$$j_P(x)=\overline{q}_\alpha ^a(x)(i\gamma _5)_{\alpha \beta }h_\beta ^a(x)\text{and}j_S(x)=\overline{q}_\alpha ^a(x)h_\alpha ^a(x),$$
(7)
as well as the correlators $`\stackrel{~}{𝒟}_\mathrm{\Gamma }(z)`$ with $`\mathrm{\Gamma }=P`$ or $`S`$:
$$\stackrel{~}{𝒟}_\mathrm{\Gamma }(z)0|T\{j_\mathrm{\Gamma }(y)j_\mathrm{\Gamma }^{}(x)\}|0.$$
(8)
The two correlators could have also been defined with vector and axialvector currents, but as it was shown in ref. , the vector and axialvector correlators are equal to the pseudoscalar and scalar correlators respectively in the heavy mass limit. This result is related to the fact that in the heavy mass limit the corresponding physical states become degenerate. Integrating out the heavy quark fields in (8), we obtain the expressions
$`\stackrel{~}{𝒟}_P(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}0|\overline{q}^a(y)(1\overline{)}v)\varphi ^{ab}(y,x)q^b(x)|0S(z),`$
$`\stackrel{~}{𝒟}_S(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}0|\overline{q}^a(y)(1+\overline{)}v)\varphi ^{ab}(y,x)q^b(x)|0S(z),`$ (9)
displaying the relation of the heavy meson correlators to the gauge invariant quark correlator of equation (1).
In energy space the HQET correlator $`\stackrel{~}{𝒟}_\mathrm{\Gamma }(w)`$ obeys the spectral representation
$$\stackrel{~}{𝒟}_\mathrm{\Gamma }(w)=\underset{k}{}\frac{f_{\mathrm{\Gamma },k}^2}{(E_{\mathrm{\Gamma },k}wiϵ)}+\underset{w_0^\mathrm{\Gamma }}{\overset{\mathrm{}}{}}𝑑\lambda \frac{\rho _\mathrm{\Gamma }(\lambda )}{(\lambda wiϵ)},$$
(10)
where $`w=vq`$ and the Fourier transform of the correlators $`\stackrel{~}{𝒟}_\mathrm{\Gamma }(z)`$ in coordinate space is given by
$$\stackrel{~}{𝒟}_\mathrm{\Gamma }(w)=id^4ze^{iqz}\stackrel{~}{𝒟}_\mathrm{\Gamma }(z).$$
(11)
$`E_{\mathrm{\Gamma },k}`$ represents the energy of the bound states and $`f_{\mathrm{\Gamma },k}`$ is the coupling of the state $`k`$ with quantum numbers $`\mathrm{\Gamma }`$ to the vacuum,
$$0|j_\mathrm{\Gamma }(0)|H_{\mathrm{\Gamma },k}=f_{\mathrm{\Gamma },k}.$$
(12)
The spectral densities are defined by $`\rho _\mathrm{\Gamma }(\lambda )1/\pi \mathrm{Im}\stackrel{~}{𝒟}_\mathrm{\Gamma }(\lambda +iϵ)`$, and finally $`w_0^\mathrm{\Gamma }`$ is the threshold energy of the continuum contribution. Transforming the spectral representation (10) back to coordinate space, one obtains
$`\stackrel{~}{𝒟}_\mathrm{\Gamma }(z)`$ $`=`$ $`i{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqz}\stackrel{~}{𝒟}_\mathrm{\Gamma }(w)}S(z)𝒟_\mathrm{\Gamma }(z)`$ (13)
$`=`$ $`S(z)\left\{{\displaystyle \underset{k}{}}f_{\mathrm{\Gamma },k}^2e^{iE_{\mathrm{\Gamma },k}|z|}+{\displaystyle \underset{w_0^\mathrm{\Gamma }}{\overset{\mathrm{}}{}}}𝑑\lambda \rho _\mathrm{\Gamma }(\lambda )e^{i\lambda |z|}\right\}.`$
The factorisation of the heavy quark propagator and the relations (2) ensure a representation of the gauge invariant quark correlators $`𝒟_\mathrm{\Gamma }(z)`$ as given by the expression inside the curly brackets. Let us already remark that the correlator decays as a simple exponential in the Euclidean region. The correlation length will therefore be given by the inverse of the lowest lying bound state energy.
## 3 The sum rules
We now turn to the theoretical side of the sum rules which arises from calculating the correlator of equation (8) in the framework of the operator product expansion . The perturbative contributions for the spectral density up to next–to–leading order in the strong coupling constant have been calculated in . Up to $`𝒪(m^2)`$ in the light quark mass, they are found to be:
$`\rho _\mathrm{\Gamma }^{pt}(w)`$ $`=`$ $`w^2{\displaystyle \underset{n=0}{\overset{2}{}}}\left[d_{n0}^\mathrm{\Gamma }+a(d_{n1}^\mathrm{\Gamma }+d_{n1L}^\mathrm{\Gamma }L)\right]\left({\displaystyle \frac{m}{w}}\right)^n`$
$`=`$ $`{\displaystyle \frac{N_c}{8\pi ^2}}\{w^2[\mathrm{\hspace{0.17em}4}+C_Fa(17+\frac{4}{3}\pi ^26L)]`$ (14)
$`\pm wm[\mathrm{\hspace{0.17em}4}+C_Fa(24+\frac{4}{3}\pi ^212L)]m^2[\mathrm{\hspace{0.17em}2}+C_Fa(39L)]\}`$
where $`a\alpha _s/\pi `$ and $`L\mathrm{ln}(2w/\mu )`$. Equation (3) also defines the coefficients $`d_{ni}^\mathrm{\Gamma }`$ which will be utilised below. As above and in all the following, the upper sign corresponds to the pseudoscalar and the lower sign to the scalar current.
In order to suppress contributions in the dispersion integral coming from higher exited states and from higher dimensional operators, it is convenient to apply a Borel transformation $`\widehat{B}_T`$ with $`T`$ being the Borel variable. The Borel transformation also removes the subtraction constants which are present in the dispersion relation satisfied by the correlators. Some useful formulae for the Borel transformation are collected in the appendix. For the phenomenological side of the sum rules, equation (10), we only take the lowest lying resonance and approximate the spectral density by the perturbative expression (3), assuming quark–hadron duality for $`w>w_0^\mathrm{\Gamma }`$. We then obtain
$$\widehat{𝒟}_\mathrm{\Gamma }(T)=f_\mathrm{\Gamma }^2e^{E_\mathrm{\Gamma }/T}+\underset{w_0^\mathrm{\Gamma }}{\overset{\mathrm{}}{}}𝑑\lambda \rho _\mathrm{\Gamma }^{pt}(\lambda )e^{\lambda /T}.$$
(15)
Let us remark that equation (15) takes exactly the same form as the expression of (13) inside the curly brackets, which represents the gauge invariant quark correlator, if we identify $`1/T`$ with the Euclidean space-time coordinate.
The perturbative contribution relevant for the sum rules is the full correlator minus the corresponding continuum contribution:
$$\widehat{𝒟}_\mathrm{\Gamma }^{pt}(T)\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })=\underset{0}{\overset{\mathrm{}}{}}𝑑\lambda \rho _\mathrm{\Gamma }^{pt}(\lambda )e^{\lambda /T}\underset{w_0^\mathrm{\Gamma }}{\overset{\mathrm{}}{}}𝑑\lambda \rho _\mathrm{\Gamma }^{pt}(\lambda )e^{\lambda /T},$$
(16)
where $`\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })`$ denotes the perturbative continuum part. After the Borel transformation the correlators satisfy homogeneous renormalisation group equations. Thus we can improve the perturbative expressions by resumming the logarithmic contributions. With the help of the following general integral formula ,
$$\underset{w_0^\mathrm{\Gamma }}{\overset{\mathrm{}}{}}𝑑\lambda \lambda ^{\alpha 1}\mathrm{ln}^n\frac{2\lambda }{\mu }e^{\lambda /T}=T^\alpha \underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)\mathrm{ln}^k\frac{2T}{\mu }\left[\frac{^{nk}}{\alpha ^{nk}}\mathrm{\Gamma }(\alpha ,\frac{w_0^\mathrm{\Gamma }}{T})\right],$$
(17)
the perturbative contribution to the sum rules is found to be:
$`\widehat{𝒟}_\mathrm{\Gamma }^{pt}(T)\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })=`$ (18)
$`T^3({\displaystyle \frac{a(2T)}{a(\mu )}})^{\gamma _1/\beta _1}{\displaystyle \underset{n=0}{\overset{2}{}}}\varphi (3n,y)\left[d_{n0}^\mathrm{\Gamma }+a\left(d_{n1}^\mathrm{\Gamma }+d_{n1L}^\mathrm{\Gamma }{\displaystyle \frac{\varphi ^{}(3n,y)}{\varphi (3n,y)}}\right)\right]({\displaystyle \frac{m(2T)}{T}})^n`$
where $`\varphi (\alpha ,y)\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\alpha ,y)`$, $`\varphi ^{}(\alpha ,y)/\alpha \varphi (\alpha ,y)`$ and $`y=w_0^\mathrm{\Gamma }/T`$. Some explicit expressions for the incomplete $`\mathrm{\Gamma }`$-function $`\mathrm{\Gamma }(\alpha ,y)`$ and the functions $`\varphi (\alpha ,y)`$, $`\varphi ^{}(\alpha ,y)`$ can also be found in the appendix. In our notation $`\beta _1=(11N_c2N_f)/6`$ is the first coefficient of the QCD $`\beta `$-function. The anomalous dimension $`\gamma _1`$ of both currents is easily found by reexpanding the running coupling and mass in terms of $`a(\mu )`$ and $`m(\mu )`$. The resulting expression is
$$\gamma _1=\frac{d_{n1L}^\mathrm{\Gamma }}{d_{n0}^\mathrm{\Gamma }}n\gamma _{1m}=\frac{3}{2}C_F=\mathrm{\hspace{0.17em}2},$$
(19)
where the first coefficient of the mass anomalous dimension, $`\gamma _{1m}=3C_F/2`$, has been used. The $`\mu `$-dependence of the correlators is due to the fact that the pseudoscalar and scalar currents are not renormalisation group invariant quantities. However, the product of $`m`$ times the current, $`mj_\mathrm{\Gamma }`$, is renormalisation group invariant in the full theory. Taking into account the additional multiplicative renormalisation of the heavy quark current in HQET , one again finds $`\gamma _1=\gamma _{1m}`$.
Essential for QCD sum rule analyses are the non–perturbative contributions coming from vacuum condensates. In the operator product expansion, the correlation function is expanded in powers of $`1/w`$ corresponding to higher and higher dimensional condensates. In our case the dimension three condensate $`\overline{h}h`$ vanishes since the heavy quark mass is infinite. After Borel transformation and up to operators of dimension five, the non–perturbative contribution takes the form :
$$\widehat{𝒟}_\mathrm{\Gamma }^{np}(T)=\frac{\overline{q}q_\mu }{2}\left[\mathrm{\hspace{0.17em}1}\frac{m}{4T}+\frac{3}{2}C_Fa\right]\pm \frac{g\overline{q}\sigma Gq_\mu }{32T^2}.$$
(20)
The next condensate contribution would be of dimension six. We have checked explicitly in our numerical analysis that this contribution to the sum rule is small. Thus we shall neglect all condensate contributions for dimensions higher than five in this work. For consistency, we have also omitted the known order $`𝒪(a^2)`$ contribution to the quark condensate , because the corresponding corrections to the perturbative part have not yet been calculated. We have, however, verified that also this correction only has a minor impact on our numerical results.
In the case of the non–perturbative part, the scale dependence of the correlator is implicit in the $`\mu `$-dependence of the quark condensate. Indeed, again $`m\overline{q}q`$ is scale independent and therefore $`\overline{q}q`$ scales inversely like the quark mass yielding the same $`\mu `$-dependence as for the perturbative contribution. For the mixed quark–gluon condensate, the next–to–leading order corrections have not been calculated and thus in the numerical analysis below, we shall neglect its scale dependence.
## 4 Numerical analysis
After equating the phenomenological and the theoretical contributions to the correlation functions, we end up with the following sum rule:
$$K_\mathrm{\Gamma }(T)f_\mathrm{\Gamma }^2e^{E_\mathrm{\Gamma }/T}=\widehat{𝒟}_\mathrm{\Gamma }^{pt}(T)\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })+\widehat{𝒟}_\mathrm{\Gamma }^{np}(T).$$
(21)
The binding energy $`E_\mathrm{\Gamma }`$ can be obtained by dividing the derivative of the sum rule (21) with respect to $`1/T`$ by the original sum rule:
$`E_\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{}{(1/T)}}\mathrm{ln}K_\mathrm{\Gamma }`$
$`=`$ $`{\displaystyle \frac{\frac{}{(1/T)}\left(\widehat{𝒟}_\mathrm{\Gamma }^{pt}(T)\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })+\widehat{𝒟}_\mathrm{\Gamma }^{np}(T)\right)}{\left(\widehat{𝒟}_\mathrm{\Gamma }^{pt}(T)\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })+\widehat{𝒟}_\mathrm{\Gamma }^{np}(T)\right)}}.`$ (22)
Analogously to equation (18) an expression for the derivative of the perturbative contribution can be obtained with the help of formula (17):
$`{\displaystyle \frac{}{(1/T)}}\left(\widehat{𝒟}_\mathrm{\Gamma }^{pt}(T)\widehat{𝒟}_\mathrm{\Gamma }^{co}(T,w_0^\mathrm{\Gamma })\right)=`$ (23)
$`T^4({\displaystyle \frac{a(2T)}{a(\mu )}})^{\gamma _1/\beta _1}{\displaystyle \underset{n=0}{\overset{2}{}}}\varphi (4n,y)\left[d_{n0}^\mathrm{\Gamma }+a\left(d_{n1}^\mathrm{\Gamma }+d_{n1L}^\mathrm{\Gamma }{\displaystyle \frac{\varphi ^{}(4n,y)}{\varphi (4n,y)}}\right)\right]({\displaystyle \frac{m(2T)}{T}})^n.`$
The corresponding derivative of the non–perturbative contribution is easily calculated from equation (20). Because the renormalisation of the correlators is multiplicative, it is clear that the binding energy $`E_\mathrm{\Gamma }`$ is a physical quantity in the sense that it is scale and scheme independent. On the contrary, this is not the case for the decay constant $`f_\mathrm{\Gamma }`$, as we shall discuss in more detail below.
Let us first consider the pseudoscalar correlator for a vanishing light quark mass $`m_q=0`$. As the central values for our input parameters we use $`\overline{q}q(1\text{GeV})=(235\pm 20\text{MeV})^3`$ for the quark condensate, $`g\overline{q}\sigma Gq=M_0^2\overline{q}q`$ with $`M_0^2=(0.8\pm 0.2)\text{GeV}^2`$ for the mixed condensate and $`\mathrm{\Lambda }_{\overline{MS}}^{(3)}=340\text{MeV}`$. The latter value corresponds to three light quark flavours and $`\alpha _s(M_Z)=0.119`$. In principle, the coupling constant in the next–to–leading order term could be evaluated at any scale $`\mu `$. As our central value in the numerical analysis we have chosen $`\mu =1.2\text{GeV}`$ but we shall investigate the variation of $`\mu `$ below. In figure 1, we display the pseudoscalar binding energy $`E_P`$ as a function of $`1/T`$ for different values of the continuum threshold $`w_0^P`$. A sum rule window, were both the continuum as well as the condensate contributions are reasonably small, can be found around $`1/T=2.53.0\text{GeV}^1`$. As can be seen from figure 1, best stability is achieved for $`w_0^P0.9\text{GeV}`$. Estimating the value of $`E_P`$ in the given range $`0.7\text{GeV}<w_0^P<1.1\text{GeV}`$, we obtain:
$$E_P=\mathrm{\hspace{0.17em}450}\pm 100\text{MeV}(m_q=\mathrm{\hspace{0.17em}0}).$$
(24)
Although the next–to–leading order QCD corrections to the sum rule of equation (21) are very large, of the order of 100%, in the ratio of equation (23) they cancel to a large extent. To investigate the sensitivity of our result to higher order corrections, we next study the dependence on the renormalisation scale $`\mu `$. The dependence of $`E_P`$ on the renormalisation scale $`\mu `$ is shown in figure 2 for $`w_0^P=0.9\text{GeV}`$ and $`\mu =0.8\text{GeV}`$, $`1.2\text{GeV}`$ and $`2.4\text{GeV}`$. One should not take $`\mu `$ smaller than $`0.8\text{GeV}`$ because then also the radiative corrections to the quark condensate become unacceptably large. We observe that inspite of the huge corrections to the correlation function, the scale dependence of $`E_P`$ is relatively mild. Adding this uncertainty to the error on $`E_P`$, our central result for $`E_P`$ is found to be
$$E_P=\mathrm{\hspace{0.17em}450}\pm 150\text{MeV}(m_q=\mathrm{\hspace{0.17em}0}).$$
(25)
The additional uncertainty coming from the errors on the input parameters is small compared to the estimated uncertainty and can be neglected.
Let us now investigate the influence of a finite light quark mass $`m_q`$. In figure 3, we have plotted $`E_P`$ as a function of $`1/T`$ for $`w_0^P=0.9\text{GeV}`$, $`\mu =1.2\text{GeV}`$ and three values of the light quark mass $`m_q`$, namely $`m_q=0\text{MeV}`$, $`75\text{MeV}`$ and $`150\text{MeV}`$. The last value is in the region of the mass of the strange quark and the intermediate value is interesting for comparison with lattice QCD results. It is well known from QCD sum rules that at the mass of the strange quark, the absolute value of the corresponding quark condensate is reduced by roughly 30% . This reduction has been applied for obtaining the dashed curve in figure 3. Lacking further information, for the dotted curve corresponding to $`m_q=75\text{MeV}`$ the quark condensate has been reduced by 10%.
Qualitatively, we find that the binding energy increases with increasing light quark mass and decreasing quark condensate. For $`m_q=150\text{MeV}`$ this increase turns out to be of the order of $`50\text{MeV}`$. Such a value is only about half of the mass splitting of $`B`$ and $`B_s`$ mesons in the $`B`$-meson system as well as $`D`$ and $`D_s`$ mesons in the $`D`$-meson system which is of the order of $`100\text{MeV}`$. Further comparison with the heavy meson systems will be presented below. The instability for $`1/T{}_{}{}^{>}3.5\text{GeV}^1`$ results from the running quark mass $`m_q(2T)`$ which explodes in this region due to uncontrollably large higher order corrections.
Let us next turn to the scalar correlator. In figure 4 the scalar binding energy is shown for three different values of the continuum threshold $`w_0^S`$. We observe that generally the stability of the scalar sum rule as not as good as in the pseudoscalar case. Here, the region of best stability is around $`1/T=2.0\text{GeV}^1`$. Adding the uncertainty from the scale dependence, the scalar binding energy is found to be:
$$E_S=\mathrm{\hspace{0.17em}1.0}\pm 0.2\text{GeV}(m_q=\mathrm{\hspace{0.17em}0}).$$
(26)
Compared to the pseudoscalar correlator, the dependence of $`E_S`$ on the light quark mass $`m_q`$ is somewhat stronger. Approximately, we obtain $`E_S(m_q)E_S(0)+m_q`$ for $`m_q150\text{MeV}`$. However, in the scalar case for increasing quark mass and decreasing quark condensate, the sum rule becomes less stable and thus we refrain from making more quantitative statements.
To conclude this section, let us comment on the decay constants $`f_\mathrm{\Gamma }`$. As has been already discussed above, because of the renormalisation of the heavy meson current, $`f_\mathrm{\Gamma }`$ depends on the renormalisation scale and scheme. Therefore, it is convenient to define renormalisation-group invariant decay constants $`\widehat{f}_\mathrm{\Gamma }`$. At the next–to–leading order, $`\widehat{f}_\mathrm{\Gamma }`$ takes the form :
$`\widehat{f}_\mathrm{\Gamma }=f_\mathrm{\Gamma }\alpha _s(\mu )^{\gamma _1/2\beta _1}(\mathrm{\hspace{0.17em}1}+Ka)\text{with}K={\displaystyle \frac{5}{12}}{\displaystyle \frac{2857\pi ^2}{54\beta _1}}+{\displaystyle \frac{107}{8\beta _1^2}}.`$ (27)
As central values, from the sum rule of equation (21) we then find $`\widehat{f}_P=0.3\text{GeV}^{3/2}`$ and $`\widehat{f}_S=0.5\text{GeV}^{3/2}`$. Since the next–to–leading QCD correction to the correlators is very large, this is also true for the uncertainty on $`\widehat{f}_\mathrm{\Gamma }`$. Thus we shall not dicuss the heavy meson decay constants further.
## 5 Comparison with lattice results
In ref. the two following quark–antiquark nonlocal condensates have been determined on the lattice:
$`C_0(x)`$ $`=`$ $`{\displaystyle \underset{f=1}{\overset{4}{}}}\mathrm{Tr}[\overline{q}_\alpha ^f(0)\varphi (0,x)q_\alpha ^f(x)],`$
$`C_v(x)`$ $`=`$ $`{\displaystyle \frac{x_\mu }{|x|}}{\displaystyle \underset{f=1}{\overset{4}{}}}\mathrm{Tr}[\overline{q}_\alpha ^f(0)(\gamma _E^\mu )_{\alpha \beta }\varphi (0,x)q_\beta ^f(x)].`$ (28)
The trace in (5) is taken with respect to the colour indices and $`f`$ is a flavour index for the light quarks. $`\gamma _E^\mu `$ are Euclidean Dirac matrices, defined as: $`\gamma _{E4}=\gamma ^0`$, $`\gamma _{Ei}=i\gamma ^i`$ $`(i=1,2,3)`$. Other insertions of gamma matrices than those containing $`1`$ and $`\gamma _E^\mu `$ vanish by P and T invariance.
In order to avoid confusion with the correlators introduced in section 2 the correlators $`C_0(x)`$ and $`C_v(x)`$ are called respectively “spin–zero” nonlocal condensate and “longitudinal–vector” nonlocal condensate. They are simply related to the correlators $`\stackrel{~}{𝒟}_\mathrm{\Gamma }`$ of equation (2) if the latter are continued to Euclidean space time:
$`\stackrel{~}{𝒟}_P(x)`$ $`=`$ $`{\displaystyle \frac{1}{2N_f}}\left(C_v(x)+C_0(x)\right)S(x),`$
$`\stackrel{~}{𝒟}_S(x)`$ $`=`$ $`{\displaystyle \frac{1}{2N_f}}\left(C_v(x)C_0(x)\right)S(x),`$ (29)
where $`N_f`$ is the number of quark flavours.
The lattice computations of ref. have been performed both in the quenched approximation and in full QCD using the SU(3) Wilson action for the pure–gauge sector and four degenerate flavours of staggered fermions, so that the sum over the flavour index $`f`$ goes from 1 to 4. In full QCD the nonlocal condensates have been measured on a $`16^3\times 24`$ lattice at $`\beta =5.35`$ and two different values of the quark mass: $`am_q=0.01`$ and $`am_q=0.02`$ ($`a`$ being the lattice spacing). For the quenched case the measurements have been performed on a $`16^4`$ lattice at $`\beta =6.00`$, using valence quark masses $`am_q=0.01`$, $`0.05`$, $`0.10`$, and at $`\beta =5.91`$ with a quark mass $`am_q=0.02`$. Further details, as well as a remark about the reliability of the results obtained for the longitudinal–vector nonlocal condensate, can be found in . The scalar functions $`C_0`$ and $`C_v`$ introduced in ref. have a more complicated spectral representation than $`\stackrel{~}{𝒟}_P`$ and $`\stackrel{~}{𝒟}_S`$, since they receive contributions both from scalar and pseudoscalar intermediate states. Nevertheless, the correlator $`C_0`$ has the advantage of not receiving perturbative contributions in the chiral limit $`m_q0`$ with $`m_q`$ being the mass of the light quark:
$$\underset{x0}{lim}C_0(x)\frac{N_fN_cm_q}{\pi ^2x^2}\text{and}\underset{x0}{lim}C_v(x)\frac{2N_fN_c}{\pi ^2x^3},$$
(30)
where $`N_c`$ is the number of colours.
Since the available lattice results for $`C_v(x)`$ are not conclusive, in what follows we shall concentrate on the spin–zero nonlocal condensate $`C_0(x)`$. In ref. a best fit to the data with the following function has been performed:
$$C_0(x)=A_0\mathrm{exp}(\mu _0x)+\frac{B_0}{x^2}.$$
(31)
The perturbative–like term $`B_0/x^2`$ takes the form obtained in the leading order in perturbation theory, if the chiral limit $`m_q0`$ (see equation (30)) is approached. Using the ansatz (31), the correlation length $`\lambda _01/\mu _0`$ of the spin–zero quark–antiquark nonlocal condensate can be extracted. At the lightest quark mass $`am_q=0.01`$ in full QCD one obtains the value $`\lambda _0=0.63_{0.13}^{+0.21}`$ fm corresponding to $`\mu _0=310\pm 80\text{MeV}`$. Within errors this value agrees with the value for $`E_P`$ obtained from the sum rule analysis.
Some observations are, however, necessary at this point. One should point out that the parametrisation (31) is a sort of “hybrid” parametrisation, since it contains a “perturbative” term $`B_0/x^2`$, which should reproduce the behaviour predicted by perturbation theory at short distances, and a “non–perturbative” term $`A_0\mathrm{exp}(\mu _0x)`$, which should reproduce the predicted exponential behaviour at large distances. A simple exponential term is not dominant in the range of physical distances where lattice data are taken, i.e., $`x0.2÷0.8`$ fm. This could shed some doubt on the identification of $`\mu _0=1/\lambda _0`$ with the binding energy $`E_P`$, as determined in section 4.
Nevertheless, at the distances where lattice data are taken the operator product expansion should still be a reasonable approximation. One can therefore try to fit directly the OPE expression to the lattice data. In this case however a running coupling would lead to a Landau pole around 1 fm and one should confine the comparison to distances small compared to this scale. On the other hand most treatments of non–perturbative QCD are based on the assumption that at large distances only non–perturbative effects prevail. This can be achieved by freezing the coupling at a certain value for distances larger than a critical one. Given the fact that we have only few data points we consider here only the leading order terms in the OPE. In addition, the correlator is scheme dependent and to perform a consistent comparison between the lattice data and the OPE at the next–to–leading order, a perturbative calculation in the lattice scheme would be necessary.
At the leading order the spin–zero correlator $`C_0(x)`$ up to operators of dimension seven is given by:
$$C_0(x)/N_f=\frac{N_cm_q^2}{\pi ^2x}K_1(m_qx)\left[\mathrm{\hspace{0.17em}1}+\frac{1}{8}m_q^2x^2\right]\overline{q}q+\frac{x^2}{16}g\overline{q}\sigma Gq,$$
(32)
where in the perturbative part through the Bessel function $`K_1(z)`$ we have kept all orders in the quark mass. Fits for the quark and the mixed condensate from the different sets of lattice data of ref. are given in table 1. The presented errors just correspond to the statistical $`1\sigma `$ errors resulting from the fit if the $`\chi ^2/d.o.f.`$ is normalised to one. $`m_L`$ is the lattice mass converted to physical units and $`m_q`$ the mass appearing in the OPE of equation (32). Since we work at the leading order, the scale and scheme dependence of the quark mass are not under control.
We note that the condensates come out with the correct order of magnitude. Qualitatively also the decrease of the quark condensate with increasing quark mass is found albeit somewhat stronger than expected from phenomenology . A direct extraction of the quark condensate from the uncooled values of the spin–zero quark correlator at zero distance led to the value $`\overline{q}q(1\text{GeV})=(235\pm 15\text{MeV})^3`$ in perfect agreement with phenomenological determinations. We have again parametrised the mixed quark–gluon condensate through $`g\overline{q}\sigma Gq=M_0^2\overline{q}q`$. The dependence of $`M_0^2`$ on the quark mass is controversial in the literature . From our fits we obtain an increase of $`M_0^2`$ with increasing quark mass, supporting the findings of ref. .
The non–perturbative part of the OPE (32) can be expressed at short distances by the Gaussian $`\overline{q}q\mathrm{exp}(M_0^2x^2/16)`$ which at large distances displays an exponential falloff. Fitting the lattice data for the quenched calculation with $`am_q=0.01`$ to the perturbative part of (32) and the Gaussian non–perturbative contribution, we find $`m_q=33\pm 2\text{MeV}`$, $`\overline{q}q^{1/3}=254\pm 2\text{MeV}`$ and $`M_0^2=0.74\pm 0.02\text{GeV}^2`$. This value for $`M_0^2`$ is more compatible with phenomenological determinations . However, the result shows that higher order corrections in the OPE have some importance and it gives an indication about the systematic uncertainties.
Another way to improve the hybrid expansion (31) consists in taking the higher resonances into account in the same way as it is done in the sum rule technique as developed by SVZ , namely by including the perturbative continuum above a threshold $`w_0`$. In leading order for the spin–zero nonlocal condensate this yields:
$`C_0(x)/N_f`$ $`=`$ $`f_P^2e^{E_Px}f_S^2e^{E_Sx}`$ (33)
$`+{\displaystyle \frac{N_c}{2\pi ^2x^3}}\left((2+2w_0^Px+(w_0^Px)^2)e^{w_0^Px}(2+2w_0^Sx+(w_0^Sx)^2)e^{w_0^Sx}\right)`$
$`+{\displaystyle \frac{N_cm_q}{2\pi ^2x^2}}\left((1+w_0^Px)e^{w_0^Px}+(1+w_0^Sx)e^{w_0^Sx}\right).`$
In figure 5 we have displayed the lattice data for the quenched calculation with $`am_q=0.01`$. The dashed curve is the fit with the OPE (32) and the parameters given in the first row of table 1. The solid curve is a fit with equation (33) where for simplicity, and lacking enough data points, we have neglected the scalar resonance by setting $`f_S=0`$ and $`w_0^S=0`$. This means that in the scalar channel only the perturbative term was taken into account and no resonance was singled out. Using in addition $`m_q=38\text{MeV}`$, the fit leads to $`f_P=0.245\pm 0.002\text{GeV}^{3/2}`$,
$$E_P=\mathrm{\hspace{0.17em}411}\pm 3\text{MeV},$$
(34)
and $`w_0^P=1.0\pm 0.01\text{GeV}`$. Within the uncertainties, these results are in good agreement to the sum rule determination of the last section.
In other lattice investigations of the heavy meson systems the main interest was to determine the leptonic decay constant of the heavy light meson in HQET. In both cases it was found that in order to isolate the ground state contribution one has in some way to suppress the higher state contributions. This was done by two different smearing mechanisms. Though the authors claimed that on the lattice the mass gap $`E_P`$ is not a physical quantity and therefore divergent in the continuum limit they determined this quantity for their lattice spacing. The difference of the mass gaps $`E_SE_P`$ however is a physical quantity and must therefore have a continuum limit. In table 2 we collect the results for these quantities for the different approaches and compare to our findings presented above.
## 6 Conclusions
In this paper we have investigated the gauge invariant quark correlator. We have set up a relation between the correlator and a corresponding correlator of gauge invariant currents in the heavy quark effective theory. The relevant currents are interpolating fields of pseudoscalar and scalar heavy–light mesons. With the method of QCD sum rules and cooled lattice data it is possible to extract the mass gap between the heavy quark pole mass and the lightest bound states of each current. This mass gap is a physical quantity in the sense that it is scale and scheme independent in perturbation theory to all orders in the strong coupling constant. The resulting values have been found to be $`E_P=450\pm 150\text{MeV}`$ for the pseudoscalar state and $`E_S=1.0\pm 0.2\text{GeV}`$ for the scalar state. The corresponding correlation lengths are: $`a_P=0.44_{0.11}^{+0.22}\text{fm}`$ and $`a_s=0.20_{0.03}^{+0.05}\text{fm}`$. For the cooled lattice data of an analysis of a perturbative term plus an exponential gave the correlation length $`a_P=0.63_{0.13}^{+0.21}\text{fm}`$; the modified analysis with a single resonance plus continuum of equation (33) yields the value $`E_P=411\text{MeV}`$, corresponding to $`a_P=0.48\text{fm}`$.
These results can be compared to the spectrum of heavy pseudoscalar mesons. Let us make the very simple assumption that up to corrections of $`1/m_Q`$ where $`m_Q`$ is the heavy quark pole mass, the heavy meson mass is equal to the heavy quark mass plus binding energy:
$$E_P\frac{c}{m_Q}=M_Pm_Q.$$
(35)
Here $`M_P`$ is the mass of the pseudoscalar meson and $`c`$ is a constant. Assuming in addition that this relation is valid for both the $`B`$ and $`D`$ mesons, we can solve for $`E_P`$ and $`c`$. Using $`m_b=5.0\pm 0.2\text{GeV}`$, $`m_c=1.8\pm 0.2\text{GeV}`$ and experimental values for $`M_B`$ and $`M_D`$, as central values for $`E_P`$ and $`c`$ we obtain:
$$E_P=\mathrm{\hspace{0.17em}400}\text{MeV}\text{and}c=\mathrm{\hspace{0.17em}0.6}\text{GeV}^2.$$
(36)
The result for $`E_P`$ is in good agreement with the other determinations presented above. Nevertheless, we should remark that the latter values are very sensitive to the heavy quark pole mass and with the estimated errors on $`m_b`$ and $`m_c`$, the uncertainty on $`E_P`$ is of the order of 100%. In addition, from the value of $`c`$ we see that the assumption of small $`1/m_Q`$ corrections is not valid in the charm case. Still, we find the very qualitative agreement of our results with the spectrum of pseudoscalar mesons noteworthy.
For lattice simulations of the gauge invariant quark correlator where the higher states were suppressed by a smearing procedure the mass gaps are said to diverge in the continuum limit. Though there is no indication of such a divergence with increasing $`\beta `$ (see table 2) the very large values for $`E_P`$ found in and are definitely not physical. The difference $`E_SE_P`$ is however a physical quantity. It has been estimated in and the values listed in table 2 agree with the result from the sum rules better than expected in view of the large errors.
The lattice data for $`C_0(x)`$ have directly been compared to the operator product expansion in leading order QCD. The results for the continuum values of the quark mass, the quark and the mixed condensates are well compatible with the values determined from other sources. The results from the quenched approximation are nearer to the continuum values than those from full QCD. The known decrease of the modulus of quark condensate with the quark mass is also confirmed by the lattice calculations. The ratio of the mixed to the quark condensate $`M_0^2`$ comes out to be the same in the quenched and full QCD. There is some controversy on the dependence of $`M_0^2`$ on the quark mass . The lattice data support an increase of $`M_0^2`$ with the quark mass which was found in . For distances smaller than approximately 1 fm where the operator product expansion can still be applied a good fit to the non–perturbative part of the correlator is also given by the Gaussian $`\overline{q}q\mathrm{exp}(M_0^2x^2/16)`$ which has previously been used in the nonlocal sum rule approach .
In our opinion a direct comparison of lattice QCD simulations with QCD sum rules in the framework of the operator product expansion opens a novel route to augment our knowledge on low–lying hadronic states and non–perturbative QCD. Taking the results of this work as encouraging, we intend to further pursue this approach in the future.
###### Acknowledgments.
M. Eidemüller thanks the Landesgraduiertenförderung at the University of Heidelberg for financial support. M. Jamin would like to thank the Deutsche Forschungsgemeinschaft for support.
## Appendix A Appendix
The Borel transformation in HQET is defined as
$$\widehat{B}_T\underset{w,n\mathrm{}}{lim}\frac{(w)^{n+1}}{\mathrm{\Gamma }(n+1)}\left(\frac{d}{dw}\right)^n,T=\frac{w}{n}>0\text{fixed}.$$
(37)
Using this definition, a central formula for the Borel transformation is found to be
$$\widehat{B}_T\frac{1}{(Ewiϵ)^\alpha }=\frac{1}{\mathrm{\Gamma }(\alpha )T^{\alpha 1}}e^{E/T}.$$
(38)
Below, we have collected some analytic formulae for the incomplete Gamma function and the functions $`\varphi (\alpha ,y)`$ and $`\varphi ^{}(\alpha ,y)`$ defined as
$`\varphi (\alpha ,y)`$ $``$ $`\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\alpha ,y),`$
$`\varphi ^{}(\alpha ,y)`$ $``$ $`{\displaystyle \frac{}{\alpha }}\left(\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\alpha ,y)\right),`$ (39)
which are helpful for the numerical analysis of the sum rules.
$`\mathrm{\Gamma }(2,y)`$ $`=`$ $`e^y(1+y)`$
$`\text{}\mathrm{\Gamma }(3,y)`$ $`=`$ $`e^y(2+2y+y^2)`$
$`\text{}\mathrm{\Gamma }^{}(2,y)`$ $`=`$ $`\mathrm{\Gamma }(0,y)+e^y[1+(1+y)\mathrm{ln}y]`$
$`\text{}\mathrm{\Gamma }^{}(3,y)`$ $`=`$ $`2\mathrm{\Gamma }(0,y)+e^y[3+y+(2+2y+y^2)\mathrm{ln}y]`$
$`\text{}\varphi (n,y)`$ $`=`$ $`\mathrm{\Gamma }(n)\left(1e^y{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{y^k}{k!}}\right),n=1,2,\mathrm{}`$
$`\text{}\varphi ^{}(\alpha ,y)`$ $`=`$ $`{\displaystyle _0^y}𝑑t\mathrm{ln}te^tt^{\alpha 1}`$
$`\text{}\varphi ^{}(1,y)`$ $`=`$ $`\gamma _E\mathrm{\Gamma }(0,y)e^y\mathrm{ln}y`$
$`\text{}\varphi ^{}(2,y)`$ $`=`$ $`1\gamma _E\mathrm{\Gamma }(0,y)e^y\left(1+(1+y)\mathrm{ln}y\right)`$
$`\text{}\varphi ^{}(3,y)`$ $`=`$ $`32\gamma _E2\mathrm{\Gamma }(0,y)e^y\left(3+y+\left(2+2y+y^2\right)\mathrm{ln}y\right)`$
$`\text{}\varphi ^{}(4,y)`$ $`=`$ $`116\gamma _E6\mathrm{\Gamma }(0,y)`$ (40)
$`e^y\left(11+5y+y^2+\left(6+6y+3y^2+y^3\right)\mathrm{ln}y\right).`$
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# Ramsey–Milman phenomenon, Urysohn metric spaces, and extremely amenable groups
## 1. Introduction
The concept of amenability extends from locally compact groups to arbitrary topological groups, and an interesting observation of recent times is that under such a transition the concept ‘gains in strength’ in that a number of concrete infinite-dimensional groups of importance satisfy a reinforced version of amenability such as locally compact groups cannot possibly have.
Definitions of amenability equivalent in the locally compact case diverge already for some of the most common infinite-dimensional topological groups . Nevertheless, the following choice has become standard : call a topological group $`G`$ amenable if every continuous affine action of $`G`$ on a convex compact set has a fixed point. Equivalently, there is a left invariant mean on the space $`C_{}^b(G)`$ of all bounded right uniformly continuous functions on $`G`$. This concept is in particular given substance by the following result due to de la Harpe : a von Neumann algebra $`A`$ is injective if and only if the unitary group $`U(A)`$ equipped with the ultraweak topology is amenable. (Cf. also .) Such results suggest that namely the above definition and not, for example, the one calling for an invariant mean on all bounded continuous functions on $`G`$, is the ‘proper’ choice.
In particular, a topological group $`G`$ is amenable if it has a fixed point in every compact space it acts upon. Such topological groups are said to have the fixed point on compacta property (f.p.c.) , or else called extremely amenable, in the spirit of where the concept was applied to discrete semigroups. The condition is equivalent to the existence of a left invariant multiplicative mean on $`C_{}^b(G)`$.
At the first sight, the latter property seems to be far too restrictive to be observed en masse. In particular, according to a well-known theorem of Veech , no locally compact group has the fixed point on compacta property. (For discrete groups, this was previously noted in .) Historically the first examples of extremely amenable groups , difficult to construct, looked like genuine pathologies.
Nevertheless, in recent times it was shown that a number of well-known ‘massive’ topological groups possess the fixed point on compacta property, among them
* the unitary group $`U()`$ (and the orthogonal group $`\mathrm{O}()`$) of an infinite-dimensional Hilbert space with the strong operator topology (Gromov and Milman ),
* the group $`L_1(X,U(1))`$ of measurable maps from a non-atomic Lebesgue space to the circle group, equipped with the $`L_1`$-metric (Glasner and independently, unpublished, Furstenberg and B. Weiss),
* groups $`\mathrm{Homeo}_+(𝕀)`$ and $`\mathrm{Homeo}_+()`$ of orientation-preserving homeomorphisms with the compact-open topology (the present author ),
* groups of measure-preserving automorphisms of standard sigma-finite measure spaces with the strong topology (Giordano and the present author ).
The technique used to establish the fixed point on compacta property in the above examples has been either that of concentration of measure on high-dimensional structures (pioneered in this context by Gromov and Milman ), or else infinite Ramsey theory, as in .
In this article we isolate a new and vast class of topological groups with the fixed point on compacta property: they are groups of isometries of very regular and highly homogeneous objects, the (generalized) Urysohn metric spaces.
Universal metric spaces were introduced by Urysohn in the 20’s and investigated mostly in the separable case. In particular, there is, up to an isometry, only one complete separable Urysohn metric space, which we will denote by $`𝕌`$. For a long time Urysohn spaces remained little known outside of general topology, and the most important advances at that period were due to Katĕtov , who had made the structure of the space $`𝕌`$ more transparent, and Uspenskij , who had proved that the group of isometries $`\mathrm{Iso}(𝕌)`$ with the compact-open topology forms a universal second-countable topological group. Uspenskij’s construction was later used by Gao and Kechris to deduce, among others, the following result: every Polish topological group is the group of all isometries of a suitable separable complete metric space. Recently the Urysohn spaces were linked to wider issues in geometry and analysis, particularly by Vershik who has for example shown that the completion of the set of integers equipped with a ‘sufficiently random’ metric is almost surely isometric to $`𝕌`$. A further discussion of the space $`𝕌`$ and its links with geometry is to be found in Gromov’s book .
We shall prove that the group $`\mathrm{Iso}(𝕌)`$ has the fixed point on compacta property (Theorem 4.11), and moreover the same is true of isometry groups $`\mathrm{Iso}(U)`$ of all sufficiently homogeneous generalized (non-separable) Urysohn spaces $`U`$ (Theorem 6.6). According to a recent result by Uspenskij , every topological group is contained, as a subgroup, in the group of isometries of such a generalized Urysohn space. The two results combined imply that extreme amenability is, in a sense, ubiquitous: every topological group embeds, as a topological subgroup, into a topological group with the fixed point on compacta property (Corollary 6.7).
It is known since the work of de la Harpe that a closed subgroup of an amenable topological group need not be amenable, unlike in the locally compact case. The reported results take this observation to its extreme. The possibility of such a development was conjectured in our paper .
The proof of extreme amenability of the group $`\mathrm{Iso}(𝕌)`$ applies the technique of concentration of measure, and by way of proof we establish the following generalization of a result due to Glasner and Furstenberg–Weiss: the group of all measurable maps from a non-atomic Lebesgue measure space to an amenable locally compact group $`G`$, equipped with the topology of convergence in measure (known as the Hartman–Mycielski extension of $`G`$, ), has the fixed point on compacta property (Theorem 2.2). Another component of the proof is the following, apparently new, result (Theorem 3.2): every group of isometries of a metric space can be approximated in a certain weak sense with finite groups of isometries of suitable metric spaces. In the second-countable case the result can be interpreted as a statement on approximation of topological groups: every Polish group is the limit of a net of finite groups in the space of all closed subgroups of the group $`\mathrm{Iso}(𝕌)`$ (Corollary 4.9).
Our methods lead to a new proof of the fixed point on compacta property for the infinite orthogonal groups with the strong topology, which does not use advanced geometric tools such as Gromov’s isoperimetric inequality. (Subsection 4.5.)
In order to extend the result on extreme amenability to the groups of isometries $`\mathrm{Iso}(U)`$ of generalized Urysohn metric spaces $`U`$, we recast the fixed point on compacta property of the full isometry group of a sufficiently homogeneous metric space $`X`$ as a Ramsey-type result for the space $`X`$ itself (Theorem 5.9). As a corollary, if two metric spaces, $`X`$ and $`Y`$, are both $`\omega `$-homogeneous and have, up to isometry, the same finite metric subspaces, then the groups $`\mathrm{Iso}(X)`$ and $`\mathrm{Iso}(Y)`$ have the fixed point on compacta property (or otherwise) simultaneously (Theorem 6.5).
As another application of this technique, we show that the groups of isometries of the universal discrete metric spaces do not have the fixed point on compacta property (Theorem 6.9).
The equivalence between the fixed point on compacta property of isometry groups and Ramsey-type results for metric spaces can be exploited in the other direction as well, and thus we deduce some ‘approximate’ Ramsey-type results for both spherical and Euclidean metric spaces (Subsection 6.3).
### Acknowledgements
The investigation grew out of stimulating discussions with Mikhaïl Gromov and Vladimir Uspenskij, who have both independently conjectured extreme amenability of the isometry groups of Urysohn spaces. Moreover, Theorems 3.2 and 5.7 are rigorous incarnations of two further conjectures made by VU. My thanks go to both named mathematicians, to MG also for hospitality extended at IHÉS in September 1999. I am grateful to Alekos Kechris and to Eli Glasner for spotting badly flawed fragments in the earlier versions of the article, and to Michael Megrelishvili for a useful remark. I am thankful to Sina Greenwood for the invitation to give the opening talk at the 2000 Devonport Topology Festival, which had rekindled my work on the topic, and to Michael Cowling for his hospitality at the Department of Pure Mathematics of the University of New South Wales in April 2000, where part of the work had been done. The research was partially supported by the Royal Society of New Zealand through the Marsden Fund grant VUW703, by the Australian Research Council through the large research grant ‘Random algebraic structures,’ and by the Institut des Hautes Études Scientifiques through its visitors programme.
## 2. Concentration of measure in Hartman–Mycielski groups
### 2.1.
Our starting point is the following result, mentioned in the Introduction.
###### Theorem 2.1 (Glasner ; Furstenberg–B. Weiss, unpublished).
The group $`L_1(X,\mathrm{U}(1))`$ of all measurable maps from a nonatomic Lebesgue space to the circle rotation group, equipped with the $`L_1`$-metric, has the fixed point on compacta property.
On two occasions in this article, including the proof of one of our main theorems, we will invoke suitable modifications of the above result, and it seems appropriate to state a far-reaching generalization of Theorem 2.1, even if we shall never use its full power.
In the above form the result does not extend too far: suffice to consider the additive group of the Banach space $`L_1(X)=L_1(X,)`$, with its wealth of continuous characters. However, it is not a particular metric on the group but rather the topology it generates that matters, and the topology generated by the $`L_1`$-metric on the group $`L(X,𝕋)`$ is that of convergence in measure. (This is true of every $`L_p`$-metric, $`1p<\mathrm{}`$, on the same group.) This observation leads us to state the following generalization of Glasner–Furstenberg–Weiss theorem.
###### Theorem 2.2.
Let $`G`$ be an amenable locally compact group and let $`X`$ be a non-atomic Lebesgue measure space. Then the group $`L_0(X,G)`$ of all measurable maps from $`X`$ to $`G`$, equipped with the topology of convergence in measure, has the fixed point on compacta property (is extremely amenable).
###### Remark 2.3.
The topological groups of the form $`L_0(X,G)`$, where the subscript ‘$`0`$’ stands for the topology of convergence in measure, had apparently been first considered by Hartman and Mycielski , who had observed that $`L_0(X,G)`$ contains $`G`$ as a topological subgroup (formed by all constant functions) and is path-connected and locally path–connected. Later it was shown by Keesling that if $`G`$ is separable metrizable, then the Hartman–Mycielski extension $`L_0(X,G)`$ is homeomorphic to the separable Hilbert space. The correspondence $`GL_0(X,G)`$ determines a (covariant) functor from the category of all topological groups and continuous homomorphisms to itself, and Theorem 2.2 says that the Hartman–Mycielski functor transforms amenable locally compact groups into extremely amenable topological groups.
The following particular case (where $`G=`$ or $``$) seems to be of interest.
###### Corollary 2.4.
The \[underlying topological group of \] the topological vector space $`L_0(X)`$ of all measurable functions on a non-atomic Lebesgue measure space $`X`$, equipped with the topology of convergence in measure, has the fixed point on compacta property. ∎
###### Remark 2.5.
The above result is similar to the one from where the space $`L_0(X)`$ was equipped with the topology of convergence in a suitably chosen, the so-called pathological submeasure (a subadditive set function) on $`X`$. As a result, the abelian topological group from has an even stronger property than just extreme amenability: it admits no strongly continuous unitary representations. Notice that each of the groups of the form $`L_0(X,G)`$ from Theorem 2.2 admits a faithful strongly continuous unitary representation in the Hilbert space $`L_2(X,L_2(G))`$. (This extends an observation made in for $`G=U(1)`$.)
Our proof of Theorem 2.2 relies, similarly to that of Theorem 2.1, on the technique of concentration of measure on high-dimensional structures. However, the concept of a Lévy group becomes too narrow and has to be somewhat extended. We believe that this extension goes sufficiently far to be of interest on its own. (Though we find it useful, to replace metrics with uniform structures, this is not what our generalization is about.)
### 2.2.
If $`X=(X,𝒰_X)`$ is a uniform space, then the uniform (induced) topology on $`X`$ gives rise to a Borel structure and thus one can speak of Borel measures on $`X`$. The following is a straightforward adaptation of the by now classical concept .
###### Definition 2.6.
Let $`(\mu _\alpha )`$ be a net of probability measures on a uniform space $`(X,𝒰_X)`$. Say that the net $`(\mu _\alpha )`$ has the Lévy concentration property, or simply concentrates (in $`X`$), if whenever $`A_\alpha X`$ are Borel subsets with the property
$$\underset{\alpha }{lim\; inf}\mu _\alpha (A_\alpha )>0,$$
one has for every entourage of the diagonal $`V𝒰_X`$
$$\mu _\alpha (V[A_\alpha ])1.$$
(Here, as usual, $`V[A]=\{xXaA,(x,a)V\}`$ denotes the $`V`$-neighbourhood of $`A`$.)
###### Lemma 2.7.
Let $`f:XY`$ be a uniformly continuous map between two uniform spaces, and let $`(\mu _\alpha )`$ be a net of Borel measures on $`X`$. If $`(\mu _\alpha )`$ concentrates, then the net $`(f_{}(\mu _\alpha ))`$ of push-forward measures on $`Y`$ concentrates as well. ∎
Let $`G`$ be a group of uniform isomorphisms of a uniform space $`X`$. A compactification $`K`$ of $`X`$ is called uniform if the corresponding mapping $`i:XK`$ is uniformly continuous, and equivariant (in full, $`G`$-equivariant) if $`G`$ acts on $`K`$ by homeomorphisms in such a way that $`i`$ commutes with the action. The maximal uniform compactification of a uniform space $`X`$, known as the Samuel compactification of $`X`$ and which we denote by $`\sigma X`$, is the Gelfand space of the commutative $`C^{}`$-algebra formed by all bounded uniformly continuous complex-valued functions on $`X`$. The Samuel compactification $`\sigma X`$ is equivariant no matter what the acting group $`G`$ is, because every uniform homeomorphism $`XX`$ extends to a self-homeomorphism $`\sigma X\sigma X`$ due to universality.
It is convenient to state explicitely the following result, which is in essence folk’s knowledge in theory of topological transformation groups. (Cf. etc.)
###### Theorem 2.8.
Let $`G`$ be a group of uniform isomorphisms of a uniform space $`X`$. The following conditions are equivalent.
1. Every $`G`$-equivariant uniform compactification of $`X`$ has a fixed point.
2. For every bounded uniformly continuous function $`f`$ from $`X`$ to a finite-dimensional Euclidean space, every $`\epsilon >0`$ and every finite collection $`g_1,g_2,\mathrm{},g_nG`$, there is an $`xX`$ with $`|f(x)f(g_ix)|<\epsilon `$ for all $`i=1,2,\mathrm{},n`$.
3. For every finite cover $`\gamma `$ of $`X`$, every $`V𝒰_X`$ and every finite collection $`g_1,g_2,\mathrm{},g_nG`$, there is an $`A\gamma `$ such that
$$_{i=1}^nV[g_iA]\mathrm{}.$$
###### Proof.
(i) $``$ (ii): notice that every bounded uniformly continuous function extends over the Samuel compactification.
(ii) $``$ (iii): choose a bounded uniformly continuous pseudometric $`d`$ on $`X`$ subordinated to $`V`$ (that is, $`d(x,y)<1(x,y)V`$) and apply the condition (ii) to the function from $`X`$ to $`^{|\gamma |}`$ whose components are distance functions $`xd(x,A)`$, $`A\gamma `$, with $`\epsilon =1`$.
(iii) $``$ (i): see , Proposition 2.1 (which is, in its turn, an adaptation of an argument from Section 4 in .) ∎
###### Remark 2.9.
At this point we do not concern ourselves with a topology on the acting group $`G`$, and it may well happen that if $`G`$ is a topological group acting on the uniform space $`X`$ continuously, the extension of the action to an equivariant compactification of $`X`$ is discontinuous. As an example, consider as $`G`$ the unitary group $`U(l_2)`$ with the strong topology, and as $`X`$ the unit sphere $`𝕊^{\mathrm{}}`$ in $`l_2`$ with the metric uniformity. The action of $`U(l_2)`$ on the Samuel compactification of the sphere is continuous if $`U(l_2)`$ is equipped with the uniform operator topology, but not the strong one.
A subset $`B`$ of a uniform space $`X`$ is called uniformly open if $`B=V[A]`$ for some $`AX`$ and $`V𝒰_X`$.
###### Definition 2.10.
Let us say that two nets of probability measures, $`(\mu _\alpha )`$ and $`(\nu _\alpha )`$, on the same uniform space $`X`$ are asymptotically proximal if for every uniformly open subset $`B`$ one has
$$\underset{\alpha }{lim\; sup}|\mu _\alpha (B)\nu _\alpha (B)|<1$$
###### Remark 2.11.
Two nets as above will in particular be asymptotically proximal if $`lim\; inf_\alpha (\mu _\alpha \nu _\alpha )(X)>0`$. For instance, this is so if the restrictions of $`\mu _\alpha `$ and $`\nu _\alpha `$ coincide on some Borel subsets $`(A_\alpha )`$, whose measures are uniformly (in $`\alpha `$) bounded away from zero.
The following apparently subsumes all the previously known results of the type (concentration of measure) $``$ (existence of a fixed point) .
###### Theorem 2.12.
Let a group $`G`$ act on a uniform space $`X=(X,𝒰)`$ by uniform isomorphisms. Suppose there is a net $`(\mu _\alpha )`$ of probability measures on $`X`$ such that
* $`(\mu _\alpha )`$ concentrates in $`X`$,
* for every $`gG`$ the nets $`(\mu _\alpha )`$ and $`(g\mu _\alpha )`$ are asymptotically proximal.
Then every equivariant uniform compactification of the $`G`$-space $`X`$ has a fixed point.
###### Remark 2.13.
The second condition is a rather weak invariance-type property for a family of measures, and its advantage is being easier to verify. If we require all measures $`\mu _\alpha `$ to be eventually invariant (that is, for every $`gG`$ one has $`\mu _\alpha =g\mu _\alpha `$ for sufficiently large $`\alpha `$) and compactly-supported, then we recover the concept of a Lévy transformation group from . The above stated theorem allows for a unified approach to a number of previously known results, such as a link between amenability of unitary representations and the concentration property of unit spheres , which we will not be addressing here.
###### Proof.
Let $`\gamma `$ be a finite cover of $`X`$, let $`V𝒰_X`$, and let $`g_1,\mathrm{},g_nG`$ be arbitrary. Find an entourage of the diagonal $`W𝒰_X`$ with $`WWV`$. At least one element of $`\gamma `$, denote it by $`A`$, satisfies the property
$$\underset{\alpha }{lim\; sup}\mu _\alpha (A)|\gamma |^1.$$
By proceeding to a subnet if necessary, we may assume without loss in generality that
$$\underset{\alpha }{lim\; inf}\mu _\alpha (A)|\gamma |^1.$$
In view of the assumed concentration property of the measures $`(\mu _\alpha )`$,
$$\underset{\alpha }{lim}\mu _\alpha (W[A])=1,$$
and by force of the second assumption, one has for every $`i`$
$$\underset{\alpha }{lim\; inf}(g_i\mu _\alpha )(W[A])>0.$$
By Lemma 2.7, each of the nets of measures $`(g_i\mu _\alpha )`$, $`i=1,2,\mathrm{},n`$ concentrates, and consequently
$$\underset{\alpha }{lim}(g_i\mu _\alpha )(W[W[A]])=1.$$
Since $`WWV`$, one has
$$\underset{\alpha }{lim}\mu _\alpha (g_iV[A])=1.$$
It is therefore possible to choose an $`\alpha `$ so large that each of the numbers $`\mu _\alpha (g_1V[A])`$, $`\mathrm{}`$, $`\mu _\alpha (g_nV[A])`$ is greater than $`1\frac{1}{n}`$. It follows that the intersection of all the translates of $`V[A]`$ by elements $`g_i`$, $`i=1,2,\mathrm{},n`$ is non-empty, and application of Theorem 2.8 finishes the proof. ∎
### 2.3.
Recall that the right uniform structure of a topological group, $`𝒰_{}(G)`$, has as a basis the entourages of diagonal of the form
$$V_{}=\{(x,y)G\times Gxy^1V\},$$
where $`V`$ runs over a neighbourhood basis of $`e`$ in $`G`$. The Samuel compactification of the right uniform space $`G_{}=(G,𝒰_{}(G))`$ is a compact $`G`$-space, known as the greatest ambit of $`G`$ and denoted by $`𝒮(G)`$. (Cf. .) The greatest ambit possesses a distinguished point (the image of identity of $`G`$, which we will still denote $`e`$), whose orbit is everywhere dense in it. This object has the following universal property: whenever $`X`$ is a compact $`G`$-space and $`x_0X`$, there is a unique morphism of $`G`$-spaces from $`𝒮(G)`$ to $`X`$ taking $`e`$ to $`x_0`$. It follows that a topological group $`G`$ has the fixed point on compacta property if and only if there is a fixed point in the greatest ambit $`𝒮(G)`$.
###### Corollary 2.14.
Let $`G`$ be a topological group. Suppose there is a net $`(\mu _\alpha )`$ of probability measures on $`G`$ such that, with respect to the right uniform structure $`𝒰_{}(G)`$,
* $`(\mu _\alpha )`$ concentrates,
* for every $`gG`$ the nets $`(\mu _\alpha )`$ and $`(g\mu _\alpha )`$ are asymptotically proximal.
Then $`G`$ has the fixed point on compacta property. ∎
###### Remark 2.15.
A topological group $`G`$ is called a Lévy group if it contains a family of compact subgroups, directed by inclusion and having everywhere dense union, such that the corresponding normalized Haar measures, $`\mu _\alpha `$, concentrate in $`G_{}`$. This concept was used as means to deduce the existence of fixed points for group actions on compacta by Gromov and Milman ; see also . Lévy groups satisfy a stronger property than the the second assumption of Corollary 2.14: the measures $`\mu _\alpha `$ are eventually invariant.
### 2.4.
Let $`X=(X,𝒰_X)`$ be a uniform space. Denote by $`L(𝕀,X)`$ the collection of all Borel-measurable maps $`f:𝕀X`$ equipped with the uniform structure of convergence in measure. The standard basic entourages of diagonal are of the form
$`[V,\epsilon ]`$ $`:=`$ $`\{(f,g)L(𝕀,X)\times L(𝕀,X):`$
$`\mu \{x𝕀:(f(x),g(x))V\}<\epsilon \},`$
where $`V𝒰_X`$ and $`\epsilon >0`$. This uniformity induces a topology on $`L(𝕀,X)`$, whose standard basic neighbourhoods of a given function $`f:𝕀X`$ are
$$[V,\epsilon ,f]:=\{gL(𝕀,X):\mu \{x𝕀:(f(x),g(x))V\}<\epsilon \},$$
where $`V𝒰_X`$ and $`\epsilon >0`$. (Notice that the knowledge of topology on $`X`$ alone does not suffice: to talk of convergence in measure, it is necessary to have a uniform structure, for instance, one defined by a metric on $`X`$, or else the unique compatible uniform structure in case $`X`$ is compact.)
If $`G`$ is a Hausdorff topological group, then so is $`L(𝕀,G)`$. In this case, the standard neighbourhoods of identity are of the form
$$[V,\epsilon ]:=\{gL(𝕀,X):\mu \{x𝕀:g(x)V\}<\epsilon \},$$
where $`V`$ is a neighbourhood of $`e_G`$ in $`G`$ and $`\epsilon >0`$.
Now suppose that $`X=(X,\rho )`$ is a metric space. Let us agree on the canonical choice of the metric generating the uniformity of convergence in measure (and the corresponding topology) on $`L_0(𝕀,X)`$, as follows: if $`\lambda >0`$ is an arbitrary (but fixed) number, then set
(2.1)
$$\mathrm{me}_\lambda (f,g)=inf\{\epsilon >0\mu \{x𝕀:\rho (f(x),g(x))>\epsilon \}<\lambda \epsilon \}.$$
Such metrics for different $`\lambda >0`$ are all equivalent. (Cf. , p. 115.)
###### Definition 2.16.
An action of a topological group $`G`$ on a uniform space $`X=(X,𝒰_X)`$ by uniform isomorphisms is called bounded (or motion equicontinuous ) if for every entourage $`U𝒰_X`$ one can find a neighbourhood $`We_G`$ such that for every $`xX`$,
$$WxU[x].$$
###### Remarks 2.17.
1. Every bounded action is continuous. \[If $`gG`$, $`xX`$, and a neighbourhood $`𝒪gx`$ are arbitrary, select an $`U𝒰_X`$ with $`(UU)[x]𝒪`$ and a neighbourhood $`We_G`$ with $`WyU[y]`$ for all $`yX`$. Then $`WU[x]U^2[x]𝒪`$.\]
2. The converse is not true. \[For example, the standard action of the unitary group $`U(l_2)`$ with the strong operator topology on the unit sphere $`𝕊^{\mathrm{}}`$ equipped with the metric uniformity is continuous, but not bounded. This action becomes bounded if $`U(l_2)`$ is equipped with the uniform operator topology.\]
3. However, a continuous action of a topological group $`G`$ on a compact space $`X`$ (equipped with the unique compatible uniformity) is always bounded. This fact is well-known (and easily verified).
4. The left action of a topological group $`G`$ on the right uniform space $`G_{}`$ is bounded, but in general the same is not true of the left action of $`G`$ on the left uniform space of $`G`$.
###### Lemma 2.18.
If a topological group $`G`$ acts by isometries on a metric space $`X`$, then the topological group $`L_0(𝕀,G)`$ acts by isometries on the metric space $`L(𝕀,X)`$ equipped with the metric (2.1), where the action is defined pointwise:
$$(gf)(x):=g(x)f(x),gL(𝕀,G),fL(𝕀,X).$$
If in addition the action of $`G`$ on $`X`$ is bounded (for example $`X`$ is compact), then the action of $`L_0(𝕀,G)`$ on $`L_0(𝕀,X)`$ is continuous.
###### Proof.
The first statement is self-evident. In order to establish the second claim, it is now enough to prove that for every $`fL_0(𝕀,X)`$ the orbit map
$$L_0(𝕀,G)ggfL_0(𝕀,X)$$
is continuous. Let $`\epsilon >0`$ be any. Using the boundedness of the original action, choose a $`We_G`$ such that for all $`xX`$ and $`wW`$, $`\rho (wx,x)<\epsilon `$. The set $`g[W,\lambda \epsilon /2]`$ is a neighbourhood of $`g`$ in $`L_0(𝕀,G)`$, and if $`g_1g[W,\lambda \epsilon /2]`$ is arbitrary, then for every $`xX`$ apart from a set of measure $`\lambda \epsilon `$ one has $`\rho (g(x)f(x),g_1(x)f(x))=\rho (f(x),w(x)f(x))<\epsilon `$, where $`w(x)g(x)^1g_1(x)W`$ for every $`xX`$ apart from a set of measure $`\lambda \epsilon /2`$. This means that $`\mathrm{me}_\lambda (g_1f,gf)\epsilon `$, establishing the continuity of the orbit map. ∎
### 2.5. Proof of Theorem 2.2
Fix a parametrization of the non-atomic Lebesgue measure space $`X`$, that is, a measure space isomorphism $`𝕀X`$. The required net of measures on $`L(𝕀,G)`$ will be indexed by the set of all pairs of the form $`(n,F)`$, where $`n_+`$ and $`FG`$ is a finite subset, directed as follows: $`(n_1,F_1)(n_2,F_2)`$ iff $`n_1n_2`$ and $`F_1F_2`$. Fix a left-invariant Haar measure $`\nu `$ on $`G`$. For every $`n,F`$ as above use the Følner condition and the assumed amenability of the locally compact group $`G`$ to choose a compact subset $`K=K_{n,F}G`$ with the property
$$\frac{\nu (gK\mathrm{\Delta }K)}{\nu (K)}<\frac{1}{n}$$
for each $`gF`$. Now let $`K^n`$ denote the set of all functions in $`L_0(𝕀,G)`$ taking values in $`K`$ and constant on every interval of the form $`[i/n,(i+1)/n)`$, $`i=0,1,\mathrm{},n1`$. Topologically, $`K^n`$ can be identified with the $`n`$-th power of the compact set $`K`$. Denote by $`\nu _{n,F}`$ the product measure $`(\nu |_K)^n`$ normalized to one and viewed as a probability measure on $`L_0(𝕀,G)`$ with support $`K^n`$. It remains to verify that the net of probability measures $`(\nu _{n,F})`$ on the topological group $`L_0(𝕀,G)`$ satisfies the two assumptions of Corollary 2.14.
(i) The net of measures $`(\nu _{n,F})`$ concentrates in $`L_0(𝕀,G)`$.
The following general and powerful result, due to Talagrand (, p. 76 and Prop. 2.1.1), extends the particular case of finite spaces belonging to Schechtman . Let $`Y=(Y,\mathrm{\Sigma },\mu )`$ denote a probability space. Then the product measures $`\mu ^n`$, $`n`$, on $`Y^n`$ concentrate, as $`n\mathrm{}`$, with respect to the \[uniform structure generated by the\] normalized Hamming distance on $`Y^n`$, given by
$$\rho (f,g)=\frac{1}{n}\mathrm{}\{if_ig_i\}.$$
Moreover, the (Gaussian) bounds for the rate of concentration are independent of a particular $`Y`$, cf. loco citato. In other words, there is a family of functions $`\alpha _n:[0,1][0,\frac{1}{2}]`$ (of the form $`\alpha _n(\epsilon )=C_1\mathrm{exp}(C_2\epsilon n^2)`$), independent of $`Y`$ and $`\mu `$ and such that, whenever a measurable $`AY^n`$ has the property $`\mu ^n(A)\frac{1}{2}`$, one has for every $`\epsilon >0`$
$$\mu ^n(A_\epsilon )1\alpha _n(\epsilon ),$$
where $`A_\epsilon =\{yY^n:\rho (y,A)<\epsilon \}`$.
In view of Lemma 2.7, it is therefore enough to show that the uniform structure induced on $`K^n`$ by the Hamming-type distance $`\rho `$ is finer than the restriction of the right uniform structure $`𝒰_{}(L(X,G))`$ (which of course coincides with the unique compatible uniformity on $`K^n`$). Let $`V`$ be a neighbourhood of unity in $`G`$ and let $`\epsilon >0`$. Let $`f,gK^n`$ be arbitrary and such that $`\rho (f,g)<\epsilon `$. Then clearly
(2.2) $`\mu (\{xXf(x)g(x)^1V\})`$ $``$ $`\mu (\{xXf(x)g(x)\})`$
$`=`$ $`{\displaystyle \frac{1}{n}}|\{if_ig_i\}|`$
$`=`$ $`\rho (f,g)<\epsilon ,`$
that is, $`(f,g)[V;\epsilon ]`$, establishing the claim.
(ii) For every $`gL_0(𝕀,G)`$, the nets $`(\nu _{n,F})`$ and $`(g\nu _{n,F})`$ are asymptotically proximal.
Let $`gL_0(X,G)`$. By approximating $`g`$ with simple functions, one can assume without loss in generality that the set $`F=\{g_1,\mathrm{},g_k\}`$ of values of $`g`$ is finite, and that for sufficiently large $`n`$, the function $`g`$ is constant on each $`[i/n,(i+1)/n)`$. Since for every $`g_i`$ one has $`\nu (K_{n,F}g_iK_{n,F})>(1\frac{1}{n})\nu (K_{n,F})`$, it follows that, whenever $`nk`$,
$$\nu _{n,F}(K_{n,F}^ng_iK_{n,F}^n)>\left(1\frac{1}{n}\right)^n\frac{1}{e}.$$
To finish the proof, notice that the restrictions of the measures $`\nu _{n,F}=(\nu |_K)^n`$ and $`g\nu _{n,F}`$ to $`K_{n,F}^kg_iK_{n,F}^k`$ coincide, and use Remark 2.11. ∎
## 3. Approximation by finite groups
### 3.1.
The aim of this Section is to show that every (topological) group can be approximated, albeit in a very weak sense, by finite groups. By combining the approximation result with the extreme amenability of Hartman–Mycielski groups, we shall later deduce the fixed point on compacta property for the isometry group $`\mathrm{Iso}(𝕌)`$ of the complete separable Urysohn metric space.
We will state the approximation result in a few equivalent forms. Let us say that a metric space $`X`$ is indexed by a set $`I`$ if there is a surjection $`f_X:IX`$. We will call the pair $`(X,f_X)`$ an indexed metric space. Let us say that two metric spaces, $`X`$ and $`Y`$, indexed with the same set $`I`$ are $`\epsilon `$-isometric if for every $`i,jI`$ the distances $`d_X(f_X(i),f_X(j))`$ and $`d_Y(f_Y(i),f_Y(j))`$ differ by at most $`\epsilon `$.
###### Lemma 3.1.
If metric spaces $`X`$ and $`Y`$ indexed by a set $`I`$ are $`\epsilon `$-isometric, then $`X`$ and $`Y`$ can be isometrically embedded into a metric space $`Z`$ in such a way that for each $`iI`$, $`d_Z(f_X(i),f_Y(i))\epsilon `$.
###### Proof.
Make the set-theoretic disjoint union $`Z=XY`$ into a weighted graph, by joining a pair $`(x,y)`$ with an edge in any of the following cases:
* $`x,yX`$, with weight $`\rho _X(x,y)`$;
* $`x,yY`$, with weight $`\rho _Y(x,y)`$,
* for some $`iI`$, $`x=f_X(i)`$ and $`y=f_Y(i)`$, with weight $`\epsilon `$.
The weighted graph $`Z`$ equipped with the path metric clearly contains $`X`$ and $`Y`$ as metric subspaces and satisfies the required property. ∎
###### Theorem 3.2.
Let $`g_1,\mathrm{},g_n`$ be a finite family of isometries of a metric space $`X`$. Then for every $`\epsilon >0`$ and every finite collection $`x_1,\mathrm{},x_m`$ of elements of $`X`$ there exist a finite metric space $`\stackrel{~}{X}`$, elements $`\stackrel{~}{x}_1,\mathrm{},\stackrel{~}{x}_m`$ of $`\stackrel{~}{X}`$, and isometries $`\stackrel{~}{g}_1,\mathrm{},\stackrel{~}{g}_n`$ of $`\stackrel{~}{X}`$ such that the indexed metric spaces $`\{g_ix_j:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ and $`\{\stackrel{~}{g}_i\stackrel{~}{x}_j:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ are $`\epsilon `$-isometric.
###### Proof.
We will perform the proof in several simple steps.
1. By first rescaling the metric $`\rho _X`$ on $`X`$ and then replacing it with $`\mathrm{min}\{\rho _X,1\}`$ if necessary, we can assume without loss in generality that the values of $`\rho _X`$ are bounded by $`1`$.
2. Without loss in generality, we may also assume that $`X`$ supports the structure of an abelian group equipped with a bi-invariant metric, and $`g_i`$’s are metric-preserving group automorphisms. For instance, one can replace $`X`$ with the free abelian group $`A(X)`$ on $`X`$, and extend the metric from $`X`$ to a maximal invariant metric on $`A(X)`$ bounded by $`1`$ (the so-called Graev metric, cf. ); then every isometry of $`X`$ uniquely extends to an isometric automorphism of the metric group $`A(X)`$.
3. Let $`G`$ denote a group of isometries of $`X`$ generated by $`g_1,\mathrm{},g_n`$. The semidirect product group $`GA(X)`$ is equipped with the bi-invariant metric $`\rho `$ defined by
$$\rho ((g,a),(g^{},a^{}))=\{\begin{array}{cc}1,\hfill & \text{if }gg^{}\text{,}\hfill \\ d(a,a^{}),\hfill & \text{otherwise.}\hfill \end{array}$$
\[The bi-invariance of $`\rho `$ is established through a direct calculation using the multiplication rule in the semidirect product in question:
$$(g,a)(h,b)=(gh,a+gb).]$$
As usual, we will identify $`G`$ with a subgroup of the semidirect product under the mapping $`Gg(g,0)`$, and similarly $`A(X)`$ is identified with a normal subgroup of the semidirect product under the mapping $`A(X)x(e_G,x)`$. Under such conventions, the automorphism of $`A(X)`$ determined by each $`gG`$ is just $`g`$ itself considered as an isometric isomorphism of $`A(X)`$:
$`aA(X),gag^1`$ $``$ $`(g,0)(e,a)(g,0)^1`$
$`=`$ $`(g,ga)(g^1,0)`$
$`=`$ $`(e,ga)`$
$``$ $`ga.`$
In particular, for every $`i,j`$ one has
$$g_ix_jg_i^1=g_ix_j.$$
4. Let $`F_{m+n}`$ denote the free group on $`m+n`$ generators denoted by the symbols
$$g_1,g_2,\mathrm{},g_n,x_1,x_2,\mathrm{},x_m.$$
Denote by $`\pi :F_{m+n}GA(X)`$ the homomorphism sending each generator $`g_i`$ to the corresponding element of $`G`$ and each generator $`x_j`$ to the corresponding element of $`XA(X)`$. Pull the metric $`\rho `$ back from $`GA(X)`$ to $`F_{m+n}`$ by letting
$$\rho ^{}(x,y)=\rho (\pi (x),\pi (y)).$$
The pseudometric $`\rho ^{}`$ is bi-invariant on $`F_{m+n}`$, though need not be a metric. By force of the remark at the end of step 3, the indexed pseudometric spaces
$$\{g_ix_ji=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}(X,\rho _X)(A(X),\rho )$$
and
$$\{g_ix_jg_i^1i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}(F_{m+n},\rho ^{})$$
are isometric (and thus are both metric spaces).
5. By adding to $`\rho ^{}`$ an arbitrary bi-invariant metric on $`F_{m+n}`$ normalized so as to only slightly change the values of distances between pairs of elements $`g_ix_jg_i^1`$ (for instance, let us agree on the discrete metric taking its values in $`\{0,\epsilon /2\})`$), we can assume without loss in generality that $`\rho ^{}`$ is a bi-invariant metric on $`F_{m+n}`$, while the indexed metric spaces $`\{g_ix_ji=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ and $`\{g_ix_jg_i^1i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ are $`\epsilon /2`$-isometric.
6. Now replace the metric $`\rho ^{}`$ with the maximal among all bi-invariant metrics on $`F_{m+n}`$ that coincide with $`\rho ^{}`$ on the set $`F_{m+n}^{(3)}`$ of all words of reduced length $`3`$.
To prove the existence of such a metric, say $`d`$, denote by $``$ the family of all bi-invariant metrics on $`F_{m+n}`$ whose restriction to $`F_{m+n}^{(3)}`$ coincides with $`\rho ^{}`$. Since $`\rho ^{}`$, the family $``$ is non-empty. For any two elements $`w,vF_{m+n}`$ and an arbitrary $`\varsigma `$, the value $`\varsigma (w,v)`$ is bounded from above uniformly in $`\varsigma `$ by any sum of the form $`_k\rho ^{}(a_k,b_k)`$, where $`w=_ka_k`$ and $`v=_kb_k`$ are two representations having the same length and such that $`a_k,b_kF_{m+n}^{(3)}`$. Now set
$$d(v,w)=\underset{\varsigma }{sup}\varsigma (v,w).$$
The supremum on the r.h.s. is finite, and has all the required properties.
7. Notice that for $`w,vF_{m+n}`$
$$d(w,v)=inf\underset{k}{}\rho ^{}(a_k,b_k),$$
where the infimum is taken over all possible representations of the above sort $`w=_ka_k`$, $`v=_kb_k`$, having the same length and such that $`a_k,b_kF_{m+n}^{(3)}`$.
\[Proof: the infimum on the r.h.s. is a bi-invariant pseudometric, which is greater than or equal to $`d`$, and whose restriction to $`F_{m+n}^{(3)}`$ coincides with the restriction of $`\rho ^{}`$. We conclude: this infimum is in fact a metric, and it must coincide with $`d`$.\]
8. Denote by $`\delta `$ the smallest positive value of the distance $`\rho ^{}`$ (or, equivalently, $`d`$) between any two elements of the finite set $`F_{m+n}^{(3)}`$. It follows from (7) that for every word $`xF_{m+n}`$, the value of $`d(w,e)`$ is at least $`[l(w)/3]\delta `$, where $`l(w)`$ denotes the reduced length of $`w`$.
9. Being free, the group $`F_{m+n}`$ is residually finite, that is, admits a separating family of homomorphisms into finite groups. (Cf. e.g. , Ch. 7, exercise 5.) Therefore, for every natural $`k`$ there exists a normal subgroup $`N_k`$ such that the factor-group $`F_{m+n}/N_k`$ is finite and the only word of reduced length $`k`$ contained in $`N_k`$ is identity. Denote by $`\varpi :F_{m+n}F_{m+n}/N_k`$ the factor-homomorphism and equip $`F_{m+n}/N_k`$ with the factor-distance of $`d`$ by letting
$$d^{}(g,h)=inf\{d(w,v)\varpi (w)=g,\varpi (v)=h\}.$$
10. It is immediate that $`d^{}`$ is a bi-invariant pseudometric. Moreover, for $`k3`$ it is a metric: $`d^{}(g,h)\delta \epsilon /2`$ whenever $`gh`$, cf. step 8.
11. Now assume that $`k3[\delta ^1]+4`$. Let $`w,vF_{m+n}^{(3)}`$ and $`x,yN_k`$ be arbitrary, $`xy`$. Then $`d(wx,vy)=d(e,x^1w^1vy)1`$, because $`w^1vN_k`$ and therefore $`l(x^1w^1vy)3[\delta ^1]`$ and (8) applies. Therefore, $`d^{}(\varpi (w),\varpi (v))=d(w,v)`$ and the restriction of $`\varpi `$ to $`F_{m+n}^{(3)}`$ is an isometry.
12. Now set $`\stackrel{~}{X}=F_{m+n}/N_k`$ and $`\stackrel{~}{x}_j=\varpi (x_j)`$. For every $`i=1,\mathrm{},n`$, the inner automorphism of the finite metric group $`\stackrel{~}{X}`$ determined by $`\varpi (g_i)`$ is an isometry, because the metric is bi-invariant. Denote this isometry by $`\stackrel{~}{g}_j`$. The indexed metric spaces $`(g_ix_jg_i^1)`$ and $`(\varpi (g_i)\stackrel{~}{x}_j\varpi (g_i)^1)`$, $`i=1,\mathrm{},n,j=1,\mathrm{},m`$ are isometric by force of the concluding remark in (11). Taking into account (5), we finally conclude that the indexed metric spaces $`(\stackrel{~}{g}_j\stackrel{~}{x}_i)`$ and $`(g_jx_i)`$ are $`\epsilon `$-isometric, as required. ∎
###### Remark 3.3.
A careful analysis of the proof shows the existence of an absolute constant $`C=C(m,n,\epsilon )>0`$ such that the cardinality of the finite metric space in the statement of Theorem 3.2 does not exceed $`C`$.
### 3.2.
Let $`G`$ be a group. One can introduce a natural topology on the set of (equivalence classes) of all isometric actions of $`G`$ on metric spaces (whose size has to be bounded from above; for instance, it is natural to consider actions on all metric spaces $`X`$ of density character not exceeding the cardinality of $`G`$). This topology is similar to the Fell topology on the set of (equivalence classes) of unitary representations of a group (cf. or , p. 12), and is even closer to the topology introduced by Exel and Loring on the set of representations of a $`C^{}`$-algebra .
A neighbourhood of an action $`\tau `$ of $`G`$ by isometries on a metric space $`X`$ is determined by the following set of data: a finite subset $`F=\{g_1,\mathrm{},g_n\}G`$, an $`\epsilon >0`$, and a finite collection $`X^{}=\{x_1,\mathrm{},x_m\}X`$. Say that an action $`\varsigma `$ of $`G`$ on a metric space $`Y`$ by isometries is in $`V[F;X;\epsilon ](\tau )`$, if for some finite collection $`Y^{}=\{y_1,\mathrm{},y_n\}Y`$ the metric spaces $`\{g_jx_i\}`$ and $`\{g_jy_i\}`$, naturally indexed by $`\{1,\mathrm{},m\}\times \{1,\mathrm{},n\}`$, are $`\epsilon `$-isometric. Our Theorem 3.2 can be now reformulated as follows.
###### Corollary 3.4.
Every action of a free group $`F`$ on an arbitrary metric space by isometries is the limit of a net of actions of $`F`$ by isometries on finite metric spaces. ∎
###### Remark 3.5.
It is worth noting that approximation results of the above type are not unknown. For instance, as a corollary of a criterion by Exel and Loring and the known residual finite-dimensionality of the group $`C^{}`$-algebras of free groups , every representation of such a $`C^{}`$-algebra in a Hilbert space is approximated in the Exel–Loring topology by finite-dimensional representations.
To cast the above result as one on approximation of topological groups, we need to remind the concept of the Urysohn universal metric space.
## 4. Urysohn metric spaces and their groups of isometries
We begin this Section with a summary of some known concepts and results from theory of Urysohn metric spaces (Subsections 4.1 and 4.2), after which we state a result on approximation of Polish topological groups by finite groups (Subsection 4.3), establish the fixed point on compacta property of the group of isometries of the complete separable Urysohn space (Subsection 4.4), and finally give a new proof of the fixed point on compacta propety for the infinite orthogonal groups (Subsection 4.5).
### 4.1. Urysohn metric spaces
A metric space $`X`$ is called a (generalized) Urysohn space if it has the following property: whenever $`AX`$ is a finite metric subspace of $`X`$ and $`A^{}=A\{a\}`$ is an arbitrtary one-point metric space extension of $`A`$, the embedding $`AX`$ extends to an isometric embedding $`A^{}X`$. (Cf. and , 3.11<sub>+</sub>.)
There is only one, up to an isometry, complete separable Urysohn metric space, which we denote by $`𝕌`$. This space contains an isometric copy of every separable metric space. Moreover, if $`X`$ is a separable metric space and $`AX`$ is a finite subspace, then every isometric embedding $`A𝕌`$ extends to an isometric embedding $`X𝕌`$.
A metric space $`X`$ is called $`𝐧`$-homogeneous, where $`n`$ is a natural number, if every isometry between two subspaces of $`X`$ containing at most $`n`$ elements each extends to an isometry of $`X`$ onto itself. If $`X`$ is $`n`$-homogeneous for every natural $`n`$, then it is said to be $`\omega `$-homogeneous. The complete separable Urysohn space $`𝕌`$ is $`\omega `$-homogeneous and moreover enjoys the stronger property: every isometry between two compact subspaces of $`X`$ extends to an isometry of $`X`$ onto itself. Other well-known metric spaces having the same higher homogeneity property are the unit sphere $`𝕊`$ of the infinite-dimensional Hilbert space $``$ and the infinite-dimensional Hilbert space $``$ itself (, Ch. IV, §38).
There are some obvious modifications of the concept of Urysohn metric space. For example, one can consider only metric spaces of diameter not exceeding a given positive number $`d`$. The corresponding complete separable Urysohn space will be denoted $`𝕌_d`$. Another possibility is to consider Urysohn metric spaces in the class of metric spaces whose metrics only take values in the lattice $`\epsilon `$, $`\epsilon >0`$. The corresponding object will be denoted $`𝕌^\epsilon `$ (respectively, $`𝕌_d^\epsilon `$).
Certainly, the above are not the only classes of metric spaces for which the Urysohn-type universal objects exist. For instance, the Urysohn metric spaces for the class of spherical metric spaces of a fixed diameter in the sense of Blumenthal are spheres in spaces $`l_2(\mathrm{\Gamma })`$. The infinite-dimensional Hilbert spaces play the role of Urysohn metric spaces for the class of metric spaces embeddable into Hilbert spaces.
The following construction of the Urysohn space belongs to Katĕtov . Let us say, following , that a 1-Lipschitz real-valued function $`f`$ on a metric space $`X`$ is supported on, or else controlled by, a metric subspace $`YX`$ if for every $`xX`$
$$f(x)=inf\{\rho (x,y)+f(y):yY\}.$$
Put otherwise, $`f`$ is the largest 1-Lipschitz function on $`X`$ having the prescribed restriction to $`Y`$. For instance, every distance function $`x\rho (x,x_0)`$ from a point $`x_0`$ is controlled by a singleton, $`\{x_0\}`$.
Let $`X`$ be an arbitrary metric space. Denote by $`\mathrm{E}(X)`$ the collection of all functions $`f:X`$ controlled by some finite subset of $`X`$ (depending on the function) and having the property
(4.1)
$$|f(x)f(y)|d_X(x,y)f(x)+f(y)$$
for all $`x,yX`$. If equipped with the supremum metric, $`\mathrm{E}(X)`$ becomes a metric space of the same density character as $`X`$, containing an isometric copy of $`X`$ under the Kuratowski embedding:
$$Xx[d_x:Xy\rho (x,y)]\mathrm{E}(X).$$
Besides, the space $`\mathrm{E}(X)`$ contains all one-point metric extensions of every finite metric subspace of $`X`$.
One can form an increasing sequence of iterated extensions of the form
$$X,\mathrm{E}(X),\mathrm{E}^2=\mathrm{E}(\mathrm{E}(X)),\mathrm{},\mathrm{E}^n(X)=\mathrm{E}(\mathrm{E}^{n1}(X)),\mathrm{},$$
take the union, $`\mathrm{E}^{\mathrm{}}(X)`$, and form the metric completion of it, $`\widehat{\mathrm{E}}^{\mathrm{}}(X)`$. The latter space is a generalized Urysohn space. If the metric space $`X`$ is separable, then so is $`\widehat{\mathrm{E}}^{\mathrm{}}(X)`$, and thus it is isometric to $`𝕌`$. If $`X`$ is non-separable, then the resulting metric space $`\widehat{\mathrm{E}}^{\mathrm{}}(X)`$ need not be $`\omega `$-homogeneous.
If $`X`$ is a separable metric space with $`\mathrm{diam}(X)d`$ and throughout the above construction one replaces $`\mathrm{E}(X)`$ with the metric space $`\mathrm{E}_d(X)`$ formed by all functions $`f`$ satisfying (4.1), bounded by $`d`$, and controlled by finite subspaces in a suitably modified sense, then the resulting metric space $`\widehat{\mathrm{E}}^{\mathrm{}}(X)`$ is isometric to $`𝕌_d`$.
### 4.2. Groups of isometries
A remarkable feature of the above construction, discovered by Uspenskij, is that it enables one to keep track of groups of isometries.
Given an arbitrary metric space $`X`$, the topology of pointwise convergence and the compact-open topology on the group $`\mathrm{Iso}(X)`$ of all isometries of $`X`$ onto itself coincide and turn $`\mathrm{Iso}(X)`$ into a Hausdorff topological group. The basic neighbourhoods of identity in this topology are of the form
$$V[F;\epsilon ]=\{g\mathrm{Iso}(X):xF,d_X(g(x),x)<\epsilon \},$$
where $`FX`$ is finite and $`\epsilon >0`$. If $`X`$ is separable (and thus second-countable), then so is $`\mathrm{Iso}(X)`$.
Notice that in general the action of $`\mathrm{Iso}(X)`$ on the metric space $`X`$ is not bounded (cf. Remark 2.17.2), while the action of $`\mathrm{Iso}(X)`$ by translations on the space of bounded uniformly continuous (or Lipschitz) functions on $`X`$, equipped with the supremum norm, is not, in general, continuous.
However, the isometric action of the group $`\mathrm{Iso}(X)`$ on the metric space of all 1-Lipschitz functions on $`X`$ controlled by finite subsets happens to be continuous. Indeed, if a function $`fE(X)`$ is controlled by a finite $`YX`$, then the translation $`gf`$ does not differ from $`f`$ by more than $`\epsilon `$ at any point of $`X`$, provided $`gV[Y;\epsilon ]`$. Consequently, the canonical representation of $`\mathrm{Iso}(X)`$ in $`\mathrm{E}(X)`$ by isometries defines a topological group embedding $`\mathrm{Iso}(X)\mathrm{Iso}(\mathrm{E}(X))`$.
Iterating this process countably many times, one obtains a a continuous action of $`\mathrm{Iso}(X)`$ by isometries on $`\mathrm{E}^{\mathrm{}}(X)`$, which in its turn extends to a continuous action of $`\mathrm{Iso}(X)`$ on the metric completion $`\widehat{\mathrm{E}}^{\mathrm{}}(X)𝕌`$.
We adopt terminology suggested in and say that a metric subspace $`Y`$ is $`g`$-embedded into a metric space $`X`$ if there exists an embedding of topological groups $`e:\mathrm{Iso}(Y)\mathrm{Iso}(X)`$ with the property that for every $`h\mathrm{Iso}(Y)`$ the isometry $`e(h):XX`$ is an extension of $`h`$. The above argument establishes the following result.
###### Proposition 4.1 (Uspenskij ).
Every separable metric space $`X`$ can be $`g`$-embedded into the complete separable Urysohn metric space $`𝕌`$. ∎
Since every \[second-countable\] topological group $`G`$ embeds into the isometry group of a suitable \[separable\] metric space , we arrive at the following.
###### Theorem 4.2 (Uspenskij ).
The topological group $`\mathrm{Iso}(𝕌)`$ is the universal second-countable topological group. ∎
(Cf. also , 3.11.$`\frac{2}{3}_+`$.)
Since every isometry between two compact subspaces of $`𝕌`$ can be extended to an isometry of $`𝕌`$ onto itself, we obtain the following useful corollary of Proposition 4.1.
###### Corollary 4.3.
Each isometric embedding of a compact metric space into $`𝕌`$ is a $`g`$-embedding. ∎
The question of the existence of universal topological groups of a given uncountable weight $`\tau `$ (in fact, of any uncountable weight $`\tau `$) remains open. However, recently Uspenskij has established the following result.
###### Theorem 4.4 ().
Every topological group $`G`$ embeds, as a topological subgroup, into the group of isometries $`\mathrm{Iso}(X)`$ of a suitable $`\omega `$-homogeneous Urysohn metric space $`X`$ of the same weight as $`G`$.
The construction rather resembles the proof of Theorem 4.2, but in order to achieve $`\omega `$-homogeneity of the union space, one alternates between the Katĕtov metric extension $`\mathrm{E}()`$ and the ‘homogenization’ extension, $`H()`$, which forms the nontrivial technical core of the proof and is described in the following theorem.
###### Theorem 4.5 (Uspenskij ).
Every metric space $`X`$ $`g`$-embeds into an $`\omega `$-homogeneous metric space $`H(X)`$ of the same weight as $`X`$. ∎
### 4.3. Approximation of topological groups
Now we can state yet another reformulation of the approximation Theorem 3.2.
###### Theorem 4.6.
For every finite collection of isometries $`g_1,\mathrm{},g_n`$ of the complete separable Urysohn metric space $`𝕌`$ and every neighbourhood $`V`$ of identity in $`\mathrm{Iso}(𝕌)`$ there are isometries $`h_1,\mathrm{},h_n\mathrm{Iso}(𝕌)`$ generating a finite subgroup and such that $`h_ig_i^1V`$, $`i=1,\mathrm{},n`$.
###### Proof.
One can assume that $`V=V[X;\epsilon ]`$, where $`X=\{x_1,\mathrm{},x_m\}𝕌`$ and $`\epsilon >0`$. Using Theorem 3.2, choose a finite metric space $`\stackrel{~}{X}`$, elements $`\stackrel{~}{x}_1,\mathrm{},\stackrel{~}{x}_m`$ of $`\stackrel{~}{X}`$, and isometries $`\stackrel{~}{g}_1,\mathrm{},\stackrel{~}{g}_n`$ of $`\stackrel{~}{X}`$ such that the naturally indexed finite metric spaces $`A=\{g_i^1x_ji=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ and $`B=\{\stackrel{~}{g}_i^1\stackrel{~}{x}_ji=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ are $`\epsilon /2`$-isometric.
Using Lemma 3.1, isometrically embed $`A`$ and $`\stackrel{~}{X}`$ into a finite metric space $`Z`$ in such a way that $`d_Z(g_i^1(x_i),\stackrel{~}{g}_i^1(\stackrel{~}{x}_i))\epsilon /2`$ for all $`i,j`$. Now extend the embedding $`A𝕌`$ to an isometric embedding $`Z𝕌`$. According to Corollary 4.3, the (finite) group $`\mathrm{Iso}(\stackrel{~}{X})`$ simultaneously extends to a group of isometries of $`𝕌`$. Denote the extension of the isometry $`\stackrel{~}{g}_i`$ by $`h_i`$. One has for all $`i,j`$:
$$d(x_i,h_jg_j^1(x_i))=d(\stackrel{~}{g}_i^1(x_i),g_i^1(x_i))\epsilon /2<\epsilon ,$$
and the proof is finished. ∎
Let $`G`$ be a group and let $`X`$ be a metric space. Every action of $`G`$ on $`X`$ by isometries can be viewed as a homomorphism $`\tau :G\mathrm{Iso}(X)`$. Equip the set $`\mathrm{Hom}(G,\mathrm{Iso}(X))`$ of all such homomorphisms with the topology of pointwise convergence on $`G`$, that is, the one induced from the Tychonoff product $`\mathrm{Iso}(X)^G`$. Since, in its turn, the topological space $`\mathrm{Iso}(X)`$ is a subspace of the Tychonoff product $`X^X`$, one concludes that $`\mathrm{Hom}(G,\mathrm{Iso}(X))`$ is a topological subspace of the Tychonoff product $`X^{G\times X}`$. In this form, the identification of the collection of all actions $`\tau :G\times XX`$ with a subspace of $`X^{G\times X}`$ becomes obvious.
Call an action periodic if it factors through an action of a finite group. One can reformulate Theorem 4.6 as follows.
###### Corollary 4.7.
Let $`F`$ be a free group. The set of periodic actions of $`F`$ on the Urysohn metric space $`𝕌`$ is everywhere dense in the set of all actions. ∎
Let $`F_{\mathrm{}}`$ denote the free group of countably infinite rank. The mapping associating to an action $`\tau `$ of $`F_{\mathrm{}}`$ on $`𝕌`$ the closure of $`\tau (F_{\mathrm{}})`$ in $`\mathrm{Iso}(𝕌)`$ is a surjection from $`\mathrm{Hom}(G,\mathrm{Iso}(X))`$ onto the space $`(\mathrm{Iso}(𝕌))`$ of all closed subgroups of $`\mathrm{Iso}(𝕌)`$. Equip the latter space with the corresponding quotient topology. The topology so defined satisfies the axiom $`T_0`$.
###### Corollary 4.8.
The set of finite subgroups is everywhere dense in the topological space $`(\mathrm{Iso}(𝕌))`$. ∎
This leads to an approximation result for Polish topological groups.
###### Corollary 4.9.
Let $`G`$ be a Polish topological group. Then under every isomorphic embedding into $`\mathrm{Iso}(𝕌)`$ the group $`G`$ is the limit of a net of finite subgroups. ∎
###### Remark 4.10.
At the first sight, the above may seem to contradict the general principle (in particular espoused and explained by Vershik in ) according to which approximability of an (infinite) group $`G`$ by finite groups is essentially equivalent to amenability of $`G`$. In fact, our results are in perfect agreement with this principle in that the approximating groups come from ‘without’ the group $`G`$ and thus form an approximation not to $`G`$ itself, but to a suitable topological group extension of $`G`$, which indeed turns out to be amenable (and even extremely amenable).
### 4.4. The fixed point property of the group $`\mathrm{Iso}(𝕌)`$
Theorems 4.6 and 2.2 enable us to deduce the fixed point on compacta property for the group of isometries of the complete separable Urysohn space $`𝕌`$.
###### Theorem 4.11.
The group $`\mathrm{Iso}(𝕌)`$ of all isometries of the complete separable Urysohn space $`𝕌`$, equipped with the standard (pointwise $`=`$ compact-open) topology, is extremely amenable (has the fixed point on compacta property).
###### Proof.
Let the group $`\mathrm{Iso}(𝕌)`$ act continuously on a compact space $`K`$. We will show that every finite collection of elements of $`\mathrm{Iso}(𝕌)`$ has a common fixed point in $`K`$, from which the result follows by an obvious compactness argument. Fix an arbitrary such collection, $`g_1,\mathrm{},g_n\mathrm{Iso}(𝕌)`$.
Let $`U𝒰_K`$ be an arbitrary element of the unique compatible uniform structure on $`K`$. Without loss in generality, assume that $`U`$ is closed as a subset of $`K\times K`$ (and consequently compact). Using the boundedness of the action of $`\mathrm{Iso}(X)`$ on $`K`$, choose a finite $`X𝕌`$ and an $`\epsilon >0`$ such that, whenever $`gV=V[X;\epsilon ]`$, one has $`(g\kappa ,\kappa )V`$ for all $`\kappa K`$.
By Theorem 4.6, there are isometries $`h_1,\mathrm{},h_n\mathrm{Iso}(𝕌)`$ generating a finite subgroup $`H`$ and such that $`h_ig_i^1V`$, $`i=1,\mathrm{},n`$.
Let $`\stackrel{~}{X}`$ be a finite $`H`$-invariant subset of $`𝕌`$ containing $`X`$. The iterated Katĕtov extension $`\widehat{\mathrm{E}}^{\mathrm{}}(L_0(𝕀,\stackrel{~}{X}))`$ contains $`\stackrel{~}{X}`$ as a subspace made up of all constant functions and is isometric to $`𝕌`$, and since $`\stackrel{~}{X}`$ is finite, an isometry between the two spaces can be chosen so as to extend the canonical embedding of $`\stackrel{~}{X}`$ into $`𝕌`$. Thus we obtain a chain of $`g`$-embeddings
$$\stackrel{~}{X}L_0(𝕀,\stackrel{~}{X})𝕌.$$
The group $`L_0(𝕀,H)`$ acts on $`L_0(𝕀,\stackrel{~}{X})`$ continuously and isometrically (Lemma 2.18), and this action canonically extends to a continuous isometric action of the same group on the space $`\widehat{\mathrm{E}}^{\mathrm{}}(L_0(𝕀,\stackrel{~}{X}))𝕌`$. Thus we obtain a continuous group monomorphism $`j:L_0(𝕀,H)\mathrm{Iso}(𝕌)`$ with the property that for every $`hH`$ one has $`j(h)|_{\stackrel{~}{X}}=h|_{\stackrel{~}{X}}`$.
Composing $`j`$ with the action $`\mathrm{Iso}(𝕌)\mathrm{Homeo}(K)`$, we obtain a continuous action of $`L_0(𝕀,H)`$ on $`K`$. By force of Theorem 2.2, $`L_0(𝕀,H)`$ has a common fixed point in $`K`$, say $`\kappa `$. In particular, $`\kappa `$ is fixed under the elements $`j(h_1),\mathrm{},j(h_n)\mathrm{Iso}(𝕌)`$, where we identify elements of $`H`$ with constants in $`L_0(𝕀,H)`$.
For all $`xX`$ and $`i=1,2,\mathrm{},n`$, one has $`d(j(h_i)^1(x),g_i^1(x))=d(h_i^1(x),g_i^1(x))<\epsilon `$ for all $`i`$ and $`xX`$, implying that $`j(h_i)g_i^1V`$ for $`i=1,2,\mathrm{},n`$. Consequently and by the choice of $`V=V[X;\epsilon ]`$,
$$(g_i\kappa ,\kappa )(g_i\kappa ,j(h_i)\kappa )(g_i\kappa ,(j(h_i)g_i^1)(g_i\kappa ))U$$
for all $`i`$. Denote by $`F_U`$ the (non-empty) set of all points $`xK`$ with the property $`(g_ix,x)U`$ for all $`i`$. Since $`U`$ is closed, so is $`F_UK`$. If $`U_1U_2`$, then $`F_{U_1}F_{U_2}`$. It means that $`\{F_U\}`$ is a centred system of closed subsets of the compact space $`K`$ and therefore has a common point, which is clearly fixed under $`g_1,\mathrm{},g_n`$, as required. ∎
###### Remark 4.12.
The same argument verbatim also establishes the fixed point on compacta property of the topological group $`\mathrm{Iso}(𝕌_d)`$ of isometries of the complete separable universal Urysohn space of finite diameter $`d`$.
### 4.5. A new proof of the fixed point on compacta property of the infinite orthogonal group
The above proof can be easily modified so as to result in a new proof of extreme amenability of the orthogonal group $`\mathrm{O}()`$ of an infinite-dimensional Hilbert space with the strong operator topology. This proof does not rely on such advanced tools from geometry as Gromov’s isoperimetric inequality for groups $`\mathrm{SO}(n)`$.
The following belongs to folklore.
###### Lemma 4.13.
Let $`X`$ be a metric subspace of the unit sphere $`𝕊`$ of a real Hilbert space $``$. Suppose a topological group $`G`$ acts on $`X`$ continuously by isometries. Then the action of $`G`$ extends to a strongly continuous action of $`G`$ by isometries on the sphere $`𝕊`$ (that is, to a strongly continuous orthogonal representation of $`G`$ in $``$). Put otherwise, every metric subspace of the unit sphere $`𝕊`$ of a real Hilbert space is $`g`$-embedded into $`𝕊`$. If the linear span of $`X`$ is dense in $``$, the extension is unique.
###### Proof.
Since for every $`x,yX`$ the value of the inner product is uniquely determined by the Euclidean distance between the elements,
$$(x,y)=1\frac{1}{2}\rho _X(x,y)^2,$$
there is only one way to turn the linear span $`\mathrm{lin}(X)`$ into a pre-Hilbert space so as to induce the given metric on $`X`$. The corresponding completion $`𝒦=\widehat{\mathrm{lin}}(X)`$ is isometrically isomorphic to the closed linear span of $`X`$ in $``$, that is, $`=𝒦𝒦^{}`$. As another consequence of the same observation, every isometry of $`X`$ lifts to a unique orthogonal transformation of $`𝒦`$. The resulting homomorphism $`\pi :G\mathrm{O}(𝒦)`$ is continuous if the latter group is equipped with the topology of simple convergence on $`X`$ or, which is the same, on $`\mathrm{lin}(X)`$. On the groups of isometries of metric spaces the topology of simple convergence on an everywhere dense subset coincides with the topology of simple convergence on the entire space. Consequently, the extended orthogonal representation $`\pi `$ of $`G`$ in $`𝒦`$ is strongly continuous. It remains to extend $`\pi `$ to a representation $`\left(\begin{array}{cc}\pi & 0\\ 0& \mathrm{Id}_𝒦^{}\end{array}\right)`$ of $`G`$ in $``$. The uniqueness statement is obvious. ∎
Here is an outline of the alternative proof of extreme amenability of $`\mathrm{O}()_s`$. We will be only considering the separable case $`=l_2`$; the extension to non-separable case is straightforward.
Every finite collection $`g_1,g_2,\mathrm{},g_n`$ of elements of $`\mathrm{O}(l_2)`$, viewed as isometries of the unit sphere $`𝕊`$, can be approximated (in the strong operator topology) by a collection of elements $`g_1^{},g_2^{},\mathrm{},g_n^{}`$ of a finite-dimensional orthogonal subgroup in the following sense: for a given natural $`m`$ and an $`\epsilon >0`$, one has $`g_i(e_j)g_i^{}(e_j)<\epsilon `$ for all $`i=1,2,\mathrm{},n`$, $`j=1,2,\mathrm{},m`$, where $`g_j^{}\mathrm{O}(N)`$, $`e_j`$ denote the standard basic vectors, and the rank $`N`$ is sufficiently large.
According to Lemma 2.18, the topological group $`L_0(𝕀,\mathrm{O}(N))`$ acts continuously by isometries on the metric space $`L_2(𝕀,𝕊^N)`$, equipped with the $`l_2`$-metric. (The topology induced on $`L(𝕀,𝕊^N)`$ by $`l_2`$-metric is still that of convergence in measure, because $`𝕊^N`$ is compact.) The metric space $`L_2(𝕀,𝕊^N)`$ is spherical of diameter one and thus can be embedded into $`𝕊`$ as a metric superspace of $`𝕊^N`$. Using Lemma 4.13, we obtain a chain of continuous monomorphisms of topological groups
$$\mathrm{O}(N)<L_0(𝕀,\mathrm{O}(N))<\mathrm{Iso}(L_2(𝕀,𝕊^N))<\mathrm{O}(l_2).$$
According to Theorem 2.2, the second topological group on the left is extremely amenable. It follows that the orthogonal operators $`g_1^{},g_2^{},\mathrm{},g_n^{}`$ have a common fixed point in every compact space upon which $`\mathrm{O}(l_2)`$ acts continuously. Now the proof is accomplished in the same way as in Theorem 4.11.
## 5. Ramsey-type theorems for metric spaces vs f.p.c. property
### 5.1. Ramsey–Dvoretzky–Milman property
In order to extend the result about fixed point on compacta property of the isometry group $`\mathrm{Iso}(𝕌)`$ beyond the separable case, we will obtain a new characterization of extremely amenable groups of isometries in terms of a Ramsey-type property of the metric spaces $`X`$.
The following is an adaptation from , Sect. 9.3.
###### Definition 5.1.
Let $`G`$ be a group of uniform isomorphisms of a uniform space $`X`$. We will say that the pair $`(G,X)`$ has the Ramsey–Dvoretzky–Milman property if for every bounded uniformly continuous function $`f`$ from $`X`$ to a finite-dimensional Euclidean space, every $`\epsilon >0`$, and every compact $`KX`$, the function $`f`$ is $`\epsilon `$-constant on a suitable translate of $`K`$, that is, there is a $`gG`$ such that
$$\mathrm{Osc}(fgK)<\epsilon .$$
Equivalently, ‘compact’ can be replaced with ‘finite.’
We defer two master examples (Ex. 5.6 and 5.8) in order to precede them by a few simple preliminary results. The following is established by pulling back the function $`f`$ from $`Y`$ to $`X`$.
###### Lemma 5.2.
Let $`G`$ be a group, acting by uniform isomorphisms on the uniform spaces $`X`$ and $`Y`$, and let $`f:XY`$ be an equivariant uniformly continuous map with everywhere dense range. If the pair $`(G,X)`$ has the Ramsey–Dvoretzky–Milman property, then so does $`(G,Y)`$. ∎
Denote by $`𝒰_X^{}`$ the totally bounded replica of the uniform structure $`𝒰_X`$ on $`X`$, that is, the coarsest uniform structure preserving the uniform continuity of every bounded uniformly continuous function on $`X`$. Basic entourages of the diagonal for $`𝒰_X^{}`$ are of the form
$$\{(g,h)X\times X:|f(x)f(y)|<\epsilon \},$$
where $`f:X^N`$ is bounded uniformly continuous, $`N`$.
The following reformulation of the R–D–M property is immediate.
###### Proposition 5.3.
A pair $`(G,X)`$ has the Ramsey–Dvoretzky–Milman property if and only if for every compact (equivalently: finite) $`KX`$ and every entourage $`V𝒰_X^{}`$ there is a $`gG`$ with $`gK`$ being $`V`$-small: $`gK\times gKV`$. ∎
###### Proposition 5.4.
Let $`X=(X,𝒰_X)`$ be a uniform space. A basis of entourages for the totally bounded replica $`𝒰_X^{}`$ of $`𝒰_X`$ is given by all finite covers of the form $`\{V[A]:A\gamma \}`$, where $`\gamma `$ is an arbitrary finite cover of $`X`$ and $`V𝒰_X`$.
###### Proof.
The claim consists of two parts: first, that all sets of the form
$$_{A\gamma }V[A]\times V[A],\text{ }\gamma \text{ finite, }V𝒰_X$$
are elements of $`𝒰_X^{}`$, and second, that each enourage from $`𝒰_X^{}`$ contains a set of the above type.
(1) Given $`\gamma `$, $`V`$, and $`A`$ as above, choose a bounded uniformly continuous pseudometric $`d`$ on $`X`$ such that $`(d(x,y)<1)((x,y)V)`$, and introduce a bounded uniformly continuous function $`f`$ from $`X`$ to the Euclidean space $`^{|\gamma |}`$ with each component $`f_A`$, $`A\gamma `$, defined by
$$Xxf_A(x):=d(x,A).$$
The set $`\{(x,y)X^2:|f(x)f(y)|<1\}`$ is an element of $`𝒰_X^{}`$ and a subset of $`_{A\gamma }V[A]\times V[A]`$.
(2) Let $`W𝒰_X^{}`$ be arbitrary. Choose a bounded uniformly continuous function $`f:X^N`$ and an $`\epsilon >0`$ such that $`\{(x,y)X^2:|f(x)f(y)|<\epsilon \}W`$. Partition the image $`f(X)`$ into finitely many pieces of diameter $`\epsilon /2`$ each and let $`\gamma `$ be the family of preimages of those pieces under $`f`$. Define $`V=\{(x,y)X^2:|f(x)f(y)|<\epsilon /2\}𝒰_X^{}𝒰_X`$. Clearly, $`_{A\gamma }V[A]\times V[A]W`$. ∎
As an immediate corollary, one obtains the following.
###### Proposition 5.5.
A pair $`(G,X)`$ has the Ramsey–Dvoretzky–Milman property if and only if for every compact (equivalently: finite) $`KX`$, every finite cover $`\gamma `$ of $`X`$, and every entourage $`V𝒰_X`$, there is a $`gG`$ such that $`gK`$ is contained in the $`V`$-neighbourhood of some $`A\gamma `$. ∎
Here is the first major example.
###### Example 5.6.
Let $`\mathrm{\Gamma }`$ be an infinite set, and let $`n`$ be a natural number. Choose as $`G`$ the group $`S_\mathrm{\Gamma }^f`$ of all finite permutations of $`\mathrm{\Gamma }`$, and as $`X`$ the set $`\mathrm{\Gamma }^{(n)}`$ of all $`n`$-subsets of $`\mathrm{\Gamma }`$, equipped with the finest (discrete) uniformity. Using Proposition 5.5, one can easily see that the pair $`(\mathrm{\Gamma }^{(n)},S_\mathrm{\Gamma }^f)`$ has the Ramsey–Dvoretzky–Milman property, which statement is indeed equivalent to the finite Ramsey theorem.
Recall that the basic entourages for the left uniform structure $`𝒰_{}(G)`$ on a topological group $`G`$ are of the form
$$V_{}=\{(g,h)G\times G:g^1hV\},$$
where $`V`$ is a neighbourhood of identity in $`G`$. If $`d`$ is a left invariant continuous pseudometric on $`G`$ and $`\epsilon >0`$, then the set $`V[d;\epsilon ]=\{(x,y)X^2:d(x,y)<\epsilon \}`$ is an element of $`𝒰_{}(G)`$. Since for every neighbourhood of identity $`V`$ there is a bounded left invariant continuous pseudometric $`d`$ on $`G`$ with $`(d(x,e_G)<1)(xV)`$ and consequently $`V_{}V[d;1]`$, it follows that the left uniform structure on a topological group is determined by left invariant bounded continuous pseudometrics.
If $`d`$ is a left invariant continuous pseudometric on $`G`$, then $`H_d=\{xG:d(x,e_G)=0\}`$ forms a closed subgroup of $`G`$, and the pseudometric $`d`$ induces a continuous left-invariant metric $`\widehat{d}`$ on the factor-space $`G/H_d`$ by the formula $`\widehat{d}(xH,yH):=d(x,y)`$. The canonical factor-map $`\pi :G(G/H_d,\widehat{d})`$ is uniformly continuous. Notice that in general both the topology and the uniform structure induced by $`\widehat{d}`$ are coarser than the factor-topology and the left uniform structure on $`G/H_d`$. We will denote the $`G`$-space $`G/H_d`$ equipped with the left invariant metric $`\widehat{d}`$ by $`G/d`$, which is consistent with the notation sometimes used in set-theoretic topology: in our situation, $`G/d`$ is the metric space canonically associated to the pseudometric space $`(G,d)`$.
The following result (which grew out of V.V. Uspenskij’s conjecture) reveals the link between the Ramsey–Dvoretzky–Milman property and the existence of fixed points.
###### Theorem 5.7.
For a topological group $`G`$, the following are equivalent.
1. $`G`$ has the fixed point on compacta property.
2. The pair $`(G,G_{})`$ has the Ramsey–Dvoretzky–Milman property.
3. For every left-invariant continuous pseudometric $`d`$ on $`G`$, the pair $`(G,G/d)`$ has the Ramsey–Dvoretzky–Milman property.
4. Whenever $`G`$ acts continuously and transitively by isometries on a metric space $`X`$, the pair $`(G,X)`$ has the Ramsey–Dvoretzky–Milman property.
5. For some family $`D`$ of bounded continuous left invariant pseudometrics $`d`$, generating the topology of $`G`$, each pair $`(G,G/d)`$ has the Ramsey–Dvoretzky–Milman property.
###### Proof.
(i) $``$ (ii): according to Theorem 2.8, the fixed point on compacta property of a topological group $`G`$ is equivalent to the following: for every bounded right uniformly continuous function $`f`$ on $`G`$ taking values in a finite-dimensional Euclidean space, every finite collection of elements $`g_1,g_2,\mathrm{},g_nG`$, and every $`\epsilon >0`$, there is an $`xG`$ such that
$$|f(x)f(g_ix)|<\epsilon \text{ for all }i=1,2,\mathrm{},n.$$
The mirror image of the above statement applies to left uniformly continuous functions and calls for the existence of an $`xG`$ with the property $`|f(x)f(xg_i)|<\epsilon `$ for all $`i`$. This amounts to the Ramsey–Dvoretzky–Milman property for the pair $`(G,G_{})`$ relative to the left action (with $`K=\{e_G,g_1,g_2,\mathrm{},g_n)`$).
(ii) $``$ (iii): as the canonical map $`GG/d`$ is uniformly continuous and $`G`$-equivariant, Lemma 5.2 applies.
(iii) $``$ (iv): Let $`d_X`$ denote the invariant metric on $`X`$. Fix an arbitrary point $`x_0X`$. The formula $`d(g,h):=d_X(gx_0,hx_0)`$ defines a left invariant continuous pseudometric on $`G`$, and the map $`Gggx_0X`$ factors through to a $`G`$-equivariant isometric isomorphism between $`G/d`$ and $`X`$.
(iv) $``$ (v): Trivial, as $`G`$ acts on each space $`G/d`$ continuously and transitively by isometries.
(v) $``$ (ii): Suppose we are given a finite subset $`FG`$, a finite cover $`\gamma `$ of $`G`$, and a basic element $`V_{}`$ of the left uniformity $`G_{}`$, where $`V`$ is a neighbourhood of identity in $`G`$. Choose a bounded left invariant continuous pseudometric $`dD`$ with the property $`(d(x,e_G)<1)(xV)`$. The sets $`\pi (A)`$, $`A\gamma `$, where $`\pi :GG/d`$ is the factor-map, form a finite cover of $`G/d`$, and by assumption there is a $`gG`$ such that $`g\pi (F)`$ is entirely contained in the $`1`$-neighbourhood of some $`\pi (A)`$, $`A\gamma `$. Denote, as before, $`H_d=\{gG:d(g,e_G)=0\}`$. The value of $`d`$ is independent on the choice of representatives in left cosets: $`d(xh_1,yh_2)=d(x,y)`$ for all $`x,yG`$, $`h_1,h_2H_d`$. Let $`fF`$ be any. Since $`\widehat{d}(g\pi (f),\pi (a))<1`$ for some $`aA`$, one has $`d(gf,a)<1`$, that is, $`gF`$ is contained in $`V_{}[A]=AV`$, and the Ramsey–Dvoretzky–Milman property of $`(G,G_{})`$ is thus verified. ∎
###### Example 5.8.
The second major example is given by the pair consisting of the full unitary group $`U()`$ of an infinite-dimensional Hilbert space $``$ and the unit sphere $`𝕊_{}`$ equipped with the Euclidean distance. The Ramsey–Dvoretzky–Milman property of this pair follows from Th. 5.7 and the extreme amenability of $`U()_s`$ (cf. Subsection 4.5). In fact, a direct proof of this property does not require the extreme amenability of the unitary group, and such was the original proof by Milman (who then used the R–D–M property to give a new proof of Dvoretzky theorem on almost spherical sections of convex bodies), cf. also , Sect. 9.3.
For sufficiently homogeneous spaces and their full groups of isometries Theorem 5.7 assumes a combinatorial form of a Ramsey-type result for metric spaces somewhat in the spirit of or , but in an ‘approximate’ implementation. We proceed to examine this connection now.
### 5.2. Ramsey-type properties of metric spaces
Let $`X`$ be a metric space, and let $`F`$ be a finite metric subspace of $`X`$. The stabilizer of $`F`$,
$$\mathrm{S}t_F=\{g\mathrm{Iso}(X):gx=x\text{ for each }xF\},$$
is a closed subgroup of $`\mathrm{Iso}(X)`$. Denote by $`X^F`$ the family of all isometric embeddings of $`F`$ into $`X`$, equipped with the natural action of $`\mathrm{Iso}(X)`$ on the left:
$$X^FjgjX^F.$$
The supremum metric on $`X^F`$, given by
$$d_{sup}(i,j)=\mathrm{max}\{d(i(x),j(x)):xF\},$$
is $`\mathrm{Iso}(X)`$-invariant. Denote by $`d_F`$ the pull-back of the metric $`d_{sup}`$ to $`\mathrm{Iso}(X)`$:
$$d_F(g,h):=d_{sup}(gi_F,hi_F),$$
where $`i_F:FX`$ is the canonical embedding. Left-invariant pseudometrics of the form $`d_F`$, where $`F`$ runs over all finite subspaces of $`X`$, generate the usual topology of pointwise convergence on $`\mathrm{Iso}(X)`$.
If $`X`$ is $`|F|`$-homogeneous, the following establishes an isomorphism of $`\mathrm{Iso}(X)`$-sets:
(5.1)
$$G/\mathrm{S}t_Fg\mathrm{S}t_F[g|_F:FgF]X^F.$$
In the combinatorial spirit, we will refer to \[finite\] partitions of a metric space $`X`$ as colourings of $`X`$ \[using finitely many colours\]. A subset $`YX`$ is monochromatic if $`YX`$ for some $`A\gamma `$, and monochromatic up to an $`\epsilon >0`$ if $`Y`$ is contained in the $`\epsilon `$-neighbourhood of some $`A\gamma `$.
A direct application of Theorem 5.7 now results in the following.
###### Theorem 5.9.
Let $`X`$ be an $`\omega `$-homogeneous metric space. The following conditions are equivalent.
1. The full group of isometries $`\mathrm{Iso}(X)`$ with the pointwise topology is extremely amenable.
2. Let $`FX`$ be a finite metric space, and let $`X^F`$ be coloured using finitely many colours. Then for every finite metric subspace $`GX`$ and every $`\epsilon >0`$ there is an isometric copy of $`G`$, $`G^{}X`$, such that all isometric embeddings $`FX`$ that factor through $`G^{}`$ are monochromatic up to $`\epsilon `$.
###### Remark 5.10.
There is a natural surjection from $`X^F`$ onto the collection $`X^{(F)}`$ of all subspaces of $`X`$ isometric to $`F`$, as the latter space is obtained from the former one by factoring out the group of distance-preserving permutations of $`F`$:
$$X^{(F)}X^F/\mathrm{Iso}(F).$$
In particular, if the metric space $`F`$ is rigid (for example, if no two distances between different pairs of points are the same), then the spaces $`X^{(F)}`$ and $`X^F`$ can be identified. In general, however, the distinction between the two spaces has to be maintained, and as we shall see (Theorem 6.9), some groups of isometries of $`\omega `$-homogeneous metric spaces fail to have the fixed point on compacta property namely due to the fact that the two spaces $`X^{(F)}`$ and $`X^F`$ are different.
###### Remark 5.11.
Theorem 5.9 provides at one’s disposal a rather versatile tool. The main application in this article will be to establish the extreme amenability of groups of isometries of $`\omega `$-homogeneous generalized Urysohn spaces. The result shall also be used to demonstrate that some groups of isometries are not extremely amenable. And finally, one can turn Theorem 5.9 around in order to deduce Ramsey-type results for metric spaces from the known results on extreme amenability of various topological groups established by other means. The next Section contains examples of applications of each sort.
## 6. Applications
### 6.1. Extreme amenability of the groups $`\mathrm{Iso}(U)`$.
We want to formalize the content of the condition (ii) of Theorem 5.9, as follows.
###### Definition 6.1.
Let $`F`$ and $`G`$ be finite metric spaces, let $`m_+`$, and let $`\epsilon >0`$. Denote by $`R(F,G,m,\epsilon )`$ the following property of a metric space $`X`$:
$`XR(F,G,m,\epsilon )`$ $``$ for every colouring of the set $`X^F`$ of all isometric embeddings of $`F`$ into $`X`$ with $`m`$ colours, there is an isometric embedding $`j:GX`$ such that all embeddings of $`F`$ into $`X`$ that factor through $`j`$ are monochromatic up to $`\epsilon `$.
Say that a metric space $`X`$ has property $`R`$ if $`XR(F,G,m,\epsilon )`$ for all finite metric spaces $`F,G`$ embeddable into $`X`$, for all $`m`$, and all $`\epsilon >0`$.
###### Remark 6.2.
Now Theorem 5.9 can be reformulated as follows: an $`\omega `$-transitive metric space $`X`$ has property $`R`$ if and only if the topological group $`\mathrm{Iso}(X)`$ is extremely amenable.
###### Proposition 6.3.
Let $`F`$ and $`G`$ be finite metric spaces, let $`X`$ be a metric space containing a copy of $`F`$, let $`m`$ be a natural number, and let $`\epsilon >0`$. The following are equivalent.
1. $`XR(F,G,m,\epsilon )`$.
2. There is a finite subspace $`ZX`$ containing a copy of $`F`$ such that $`ZR(F,G,m,\epsilon )`$.
###### Proof.
(i) $``$ (ii): assume $`\neg `$(ii), that is, no finite subspace $`Z`$ of $`X`$ containing a copy of $`F`$ is in $`R(F,G,m,\epsilon )`$. Denote by $`𝒵`$ the collection of all finite metric subspaces $`ZX`$ with $`Z^F\mathrm{}`$. By assumption, $`𝒵\mathrm{}`$.
Then for every $`Z𝒵`$ the set $`Z^F`$ admits a colouring with $`m`$ colours, which we will view as a function $`f_Z:Z^F\{1,2,\mathrm{},m\}`$, in such a way that the following holds:
($``$) for every isometric embedding $`i:GZ`$ and every colour $`k=1,2,\mathrm{},m`$ there is an isometric embedding $`j_k:FG`$ such that the $`\epsilon `$-neighbourhood of $`ij_k`$ in $`Z^F`$ contains no elements of colour $`k`$.
The system $`𝒵`$ is directed by inclusion, and the collection of intervals $`[K,\mathrm{})=\{Z𝒵:KZ\}`$, where $`KX`$ is finite, is a filter on $`𝒵`$, which we will denote by $``$. Since $`X`$ can be assumed infinite (otherwise there is nothing to prove), $``$ extends to a free ultrafilter $`\mathrm{\Lambda }`$ on $`𝒵`$. For every $`jX^F`$, one has $`[\{j(F)\},\mathrm{})\mathrm{\Lambda }`$, and therefore exactly one of the sets $`\{Z𝒵:f_Z(j)=i\}`$, $`1im`$ is in $`\mathrm{\Lambda }`$. Consequently, the function
$$f(j)=\underset{\mathrm{\Lambda }}{lim}f_Z(j)$$
determines a colouring of $`X^F`$ with $`m`$ colours.
Now let $`\iota :GX`$ be an arbitrary isometric embedding, and let $`k\{1,2,\mathrm{},m\}`$ be a colour. For every $`Z[\iota (G),\mathrm{})`$ choose, using ($``$), an isometric embedding $`j_{Z,k}:FG`$ with no element in the $`\epsilon `$-neighbourhood of $`\iota j_{Z,k}`$, formed in $`Z^F`$, being of $`f_Z`$-colour $`k`$. For every $`xF`$ define $`j_k(x)=lim_\mathrm{\Lambda }j_{Z,k}(x)G`$. (The metric space $`G`$ is finite.) This $`j_k`$ is an isometric embedding of $`F`$ into $`G`$ with the property that the $`\epsilon `$-neighbourhood of $`\iota j_k`$ formed in all of $`X^F`$ contains no elements of colour $`k`$. Thus, $`\neg `$(i) is established.
(ii) $``$ (i): evident. ∎
###### Corollary 6.4.
Let $`X`$ and $`Y`$ be two metric spaces having, up to isometry, the same finite metric subspaces. If $`X`$ has property $`R`$, then so does $`Y`$. ∎
###### Theorem 6.5.
Let $`X`$ and $`Y`$ be two $`\omega `$-homogeneous metric spaces, having, up to isometry, the same finite metric subspaces. Then the topological group $`\mathrm{Iso}(X)`$ has the fixed point on compacta property if and only if the topological group $`\mathrm{Iso}(Y)`$ does.
###### Proof.
Combine Theorem 5.9 and Corollary 6.4. ∎
We can finally deduce from Theorem 6.5 and Theorem 4.11 the following result, which is the raison d’être of the article.
###### Theorem 6.6.
Let $`U`$ be a generalized Urysohn metric space. If $`U`$ is $`\omega `$-homogeneous, then the group $`\mathrm{Iso}(U)`$ has the fixed point on compacta property. ∎
Modulo Uspenskij’s Theorem 4.4, the above Theorem implies the following.
###### Corollary 6.7.
Every topological group embeds, as a topological subgroup, into an extremely amenable topological group, that is, a topological group with the fixed point on compacta property. ∎
Even the following appears to be a new result.
###### Corollary 6.8.
Every topological group embeds, as a topological subgroup, into an amenable topological group. ∎
### 6.2. Groups of isometries of discrete Urysohn spaces
Here we will demonstrate how Theorem 5.9 can be used to show the absence of the fixed point on compacta property in the case where the $`\omega `$-homogeneous metric space in question fails the ‘strong’ version of Ramsey-type property.
###### Theorem 6.9.
The group of isometries of the discrete Urysohn metric space $`𝕌^\epsilon `$ does not have the fixed point on compacta property.
###### Proof.
Denote by $`\{a,b\}`$ the two-element metric space with $`d(a,b)=\epsilon `$. Partition the set $`(𝕌^\epsilon )^{\{a,b\}}`$ of all isometric embeddings of $`\{a,b\}`$ into $`𝕌^\epsilon `$ into two disjoint subsets $`A,B`$ in such a way that whenever an injection $`i:\{a,b\}(𝕌^\epsilon )`$ is in $`A`$, the ‘flip’ injection $`i\sigma _2`$ is in $`B`$, and vice versa. Since the space $`𝕌^\epsilon `$ is $`\epsilon `$-discrete, the $`\epsilon `$-neighbourhood of a subset $`X`$ is $`X`$ itself, and ‘monochromatic up to $`\epsilon `$’ means in this context simply ‘monochromatic.’ One concludes that, with respect to the colouring $`\{A,B\}`$, no pair of injections of the form $`Y=\{i,i\sigma _2\}`$ is monochromatic up to $`\epsilon `$, and thus the metric space $`(𝕌^\epsilon )^{\{a,b\}}`$, upon which the group $`\mathrm{Iso}(𝕌^\epsilon )`$ acts transitively and continuously by isometries, fails the Ramsey–Dviretzky–Milman property. ∎
###### Remark 6.10.
The same result holds for discrete Urysohn spaces of bounded diameter, $`𝕌_d^\epsilon `$. In particular, letting $`\epsilon =1=d`$, we obtain a result proved by the present author in , Th. 6.5: the group of permutations $`S_{\mathrm{}}`$ of an infinite set, equipped with the pointwise topology, is not extremely amenable. (This result seems to answer in the negative an old question by Furstenberg discussed in .)
Notice also that the groups of isometries of infinite, $`\omega `$-homogeneous metric spaces need not be extremely amenable.
The countable metric space $`𝕌_1^{}`$, equipped with the $`\{0,1\}`$-valued metric, actually satisfies a ‘weaker’ version of the Ramsey result, namely the one for finite subspaces, rather than for their injections, and this result is the well-known Finite Ramsey Theorem. (Cf. Ex. 5.6.) However, as we have just seen, the group fails the ‘stronger’ version for embeddings of finite spaces! The latter circumstance destroys the extreme amenability of $`S_{\mathrm{}}`$.
Finally notice that the topological group $`S_{\mathrm{}}`$ is amenable, because it is approximated from within by an increasing chain of finite groups of permutations whose union is everywhere dense.
### 6.3. Deducing Ramsey-type theorems for metric spaces
By force of Theorem 5.9, the immediate corollary — and in fact an equivalent form — of the fixed point on compacta property of the group $`\mathrm{Iso}(𝕌)`$ (Theorem 4.11) is the following Ramsey-type result.
###### Corollary 6.11.
Let $`F`$ be a finite metric space, and let all isometric embeddings of $`F`$ into $`𝕌`$ be coloured using finitely many colours. Then for every finite metric space $`G`$ and every $`\epsilon >0`$ there is an isometric copy $`G^{}𝕌`$ of $`G`$ such that all isometric embeddings of $`F`$ into $`𝕌`$ that factor through $`G`$ are monochromatic up to $`\epsilon `$. ∎
By restricting ourselves to considering only $`\mathrm{Iso}(F)`$-invariant collections of embeddings of $`F`$ into $`𝕌`$, we arrive at the following.
###### Corollary 6.12.
Let $`F`$ be a finite metric space. Let all subspaces of the Urysohn space $`𝕌`$ isometric to $`F`$ be coloured using finitely many colours. Then for every finite metric space $`G`$ and every $`\epsilon >0`$ there is a subspace $`G^{}𝕌`$ isometric to $`G`$ whose subspaces isometric to $`F`$ are monochromatic up to $`\epsilon `$. ∎
Applications to spherical spaces are probably more interesting. (Cf. comments in at the bottom of p. 460). The unit sphere of the infinite-dimensional Hilbert space $``$ is an $`\omega `$-homogeneous metric space, and the orthogonal group of $``$ with the strong operator topology (that is, the topology of simple convergence on the sphere) is extremely amenable . As a corollary, we obtain Ramsey-type results for the Hilbert sphere.
###### Corollary 6.13.
Let $`F`$ be a finite metric subspace of the unit sphere $`𝕊^{\mathrm{}}`$ in an infinite-dimensional Hilbert space. Let all isometric embeddings of $`F`$ into $`𝕊^{\mathrm{}}`$ be coloured using finitely many colours. Then for every finite metric subspace $`G`$ of the sphere and every $`\epsilon >0`$ there is an isometric copy $`G^{}𝕊^{\mathrm{}}`$ of $`G`$ such that all isometric embeddings of $`F`$ into $`G^{}`$ are monochromatic up to $`\epsilon `$. ∎
###### Corollary 6.14.
Let $`F`$ be a finite metric subspace of the unit sphere $`𝕊^{\mathrm{}}`$ in an infinite-dimensional Hilbert space. Let all subspaces of $`𝕊^{\mathrm{}}`$ isometric to $`F`$ be coloured using finitely many colours. Then for every finite subspace $`Y`$ of the sphere and every $`\epsilon >0`$ there is a subspace $`Y^{}𝕊^{\mathrm{}}`$ isometric to $`Y`$ whose subspaces isometric to $`F`$ are monochromatic up to $`\epsilon `$. ∎
To establish similar corollaries for metric subspaces of the infinite-dimensional Hilbert space, we need the following result. Notice that amenability of the group $`\mathrm{Iso}()`$ of affine isometries of a Hilbert space $``$ was noted in , p. 47.
###### Theorem 6.15.
The group $`\mathrm{Iso}()`$ of affine isometries of a Hilbert space $``$ of infinite dimension is extremely amenable.
###### Proof.
The topological group $`\mathrm{Iso}()`$ is isomorphic to the semidirect product $`\mathrm{O}()`$ of the full orthogonal group $`\mathrm{O}()`$ equipped with the strong operator topology and the additive group of the Hilbert space $``$ with the usual norm topology, formed with respect to the natural action of $`\mathrm{O}()`$ on $``$ by rotations. (Cf. .) Suppose $`\mathrm{Iso}()`$ acts continuously on a compact space $`K`$. Since the group $`\mathrm{O}()`$ (identified with a subgroup of $`\mathrm{Iso}()`$) is extremely amenable (; cf. also Subsection 4.5), it has a fixed point $`\kappa K`$. The mapping $`xx\kappa K`$, where $``$ is viewed as a closed normal subgroup of $`\mathrm{Iso}()`$, is $`\mathrm{Iso}()`$-equivariant, continuous, and has everywhere dense image in $`K`$, and thus $`K`$ is an equivariant $`\mathrm{Iso}()`$-compactification of the homogeneous factor-space $`\mathrm{Iso}()/\mathrm{O}()`$.
Let $`\phi :K^N`$ be an arbitrary continuous function, $`N`$. Its pull-back, $`f(x)=:\phi (x\kappa )`$, to $``$ is right uniformly continuous. (A standard result in abstract topological dynamics.) If $`\epsilon >0`$ is arbitrary, then for some neighbourhood $`V=V[F;\delta ]`$ of identity in $`\mathrm{Iso}()`$ one has $`|f(g(0))f(h(0))|<\epsilon `$ whenever $`gh^1V`$. Without loss in generality and slightly perturbing the points of $`F`$ if necessary, one can assume that elements of $`F`$ are affinely independent. Let $`x,y`$ be two arbitrary elements with the property $`xz=yz`$ for each $`zF`$. Find an isometric copy of $`F`$, say $`F^{}`$, such that $`F^{}\{0\}`$ is isometric to $`F\{x\}`$ (or, equivalently, to $`F\{y\}`$). There is an isometry $`g`$ of $``$ taking $`F^{}\{0\}`$ to $`F\{x\}`$, and an isometry $`h`$ taking $`F^{}\{0\}`$ to $`F\{y\}`$. In particular, $`gh^1|_F=\mathrm{Id}_FV`$, and consequently $`|f(x)f(y)|<\epsilon `$. Thus, the function $`f`$ is $`\epsilon `$-constant on every affine sphere of codimension $`|F|`$ having the form $`\{x:xz=r_z,zF\}_{zF}𝕊_{r_z}(z)`$. Another way to say it is that, up to $`\epsilon `$, the function $`f(x)`$ only depends on the collection of distances $`\{xz:zF\}`$.
Now let $`g_1,\mathrm{},g_n\mathrm{Iso}()`$ be an arbitrary collection of isometries. By slightly perturbing them if necessary, one can assume without loss in generality that all the vectors $`z`$ and $`g_i^1(z)`$, $`zF`$, $`i=1,2,\mathrm{},n`$, are affinely independent. Because of infinite-dimensionality of $``$, every element $`x`$ of some affine subspace of $``$ of finite codimension has the property that for every $`i=1,2,\mathrm{},n`$ and each $`zF`$, one has $`xg_i^1(z)=xz`$. Fix any such $`x`$. Then the values of $`f`$ at the points $`x,g_1(x),g_2(x)`$, $`\mathrm{},g_n(x)`$ differ by less than $`\epsilon `$. Now we can apply Theorem 2.8 to conclude that $`K`$ has a fixed point for $`\mathrm{Iso}()`$. ∎
###### Corollary 6.16.
Let $`F`$ be a finite metric subspace of the infinite-dimensional Hilbert space $``$. Let all isometric embeddings of $`F`$ into $``$ be coloured using finitely many colours. Then for every finite collection $`Y`$ of such embeddings and every $`\epsilon >0`$ there is a collection of embeddings $`Y^{}`$ congruent to $`Y`$ and monochromatic up to $`\epsilon `$. ∎
###### Corollary 6.17.
Let $`F`$ be a finite metric subspace of an infinite-dimensional Hilbert space $``$. If all subspaces of $``$ isometric to $`F`$ are coloured using finitely many colours, then for every finite subspace $`G`$ of $``$ and every $`\epsilon >0`$ there is an isometric copy $`G^{}`$ of $`G`$ in $``$ such that all subspaces of $`G^{}`$ isometric to $`F`$ are monochromatic up to $`\epsilon `$. ∎
## 7. Concluding remarks
In this article we have investigated some relationships inside the following triangle:
$$\begin{array}{c}\text{extreme amenability}\\ \\ \text{concentration}\text{Ramsey}\end{array}$$
Deeper explorations of the Ramsey–Milman phenomenon in topological transformation groups require discovering situations in which a ‘phase transition’ between concentration and dissipation occurs in families of topological groups / dynamical systems. (Cf. .) It could be, for example, that a solution to Glasner’s problem on the existence of a minimally almost periodic group topology on the integers without the fixed point on compacta property lies namely in this direction.
In connection with the Banach–Mazur problem (cf. ), it could be worth investigating the fixed point on compacta property for the groups of isometries of separable Banach spaces admitting a transitive norm.
Finally, we do not know if the results of Section 6 can be put in direct connection with the Euclidean Ramsey theory .
CORRIGENDUM
As kindly pointed out to me by C. Ward Henson, the proof of one of the main technical results (Theorem 3.2) in the paper is flawed. To quote from his message: “Suppose $`a,b`$ are elements of $`F=F_{m+n}^{(3)}`$ that have the minimum distance $`\delta `$ from each other in the $`\rho ^{}`$ metric, and let w be any word in $`F`$. Since the metric d is bi-invariant, the conjugate $`v=wab^1w^1`$ of $`ab^1`$ has $`d`$-distance $`\delta `$ from the identity. But it seems clear that the reduced length of $`v`$ could be made arbitrarily large by choosing w correctly. This contradicts what you claim in (8).”
Fortunately, the result is not particularly deep, and here is a corrected proof of the statement.
As in , we say that a metric space $`X`$ is indexed by a set $`I`$ if there is a surjection $`f_X:IX`$. We will call the pair $`(X,f_X)`$ an indexed metric space. Two metric spaces, $`X`$ and $`Y`$, indexed with the same set $`I`$ are $`\epsilon `$-isometric if for every $`i,jI`$ the distances $`d_X(f_X(i),f_X(j))`$ and $`d_Y(f_Y(i),f_Y(j))`$ differ by at most $`\epsilon `$.
Here is the result in question.
Theorem 3.2. Let $`g_1,\mathrm{},g_m`$ be a finite family of isometries of a metric space $`X`$. Then for every $`\epsilon >0`$ and every finite collection $`x_1,\mathrm{},x_n`$ of elements of $`X`$ there exist a finite metric space $`\stackrel{~}{X}`$, elements $`\stackrel{~}{x}_1,\mathrm{},\stackrel{~}{x}_n`$ of $`\stackrel{~}{X}`$, and isometries $`\stackrel{~}{g}_1,\mathrm{},\stackrel{~}{g}_m`$ of $`\stackrel{~}{X}`$ such that the indexed metric spaces $`\{g_jx_i:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ and $`\{\stackrel{~}{g}_j\stackrel{~}{x}_i:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ are $`\epsilon `$-isometric.
###### Proof.
Without loss in generality, we can assume that $`X`$ is separable (in fact, even countable). Such an $`X`$ can be $`g`$-embedded into the Urysohn space $`𝕌`$ (see , or else Prop. 4.1 in ), and therefore we can further assume that $`X=𝕌`$.
Choose any element $`\xi 𝕌`$ and isometries $`g_{m+1},g_{m+2},\mathrm{},g_{m+n}`$ of $`𝕌`$ with the property $`g_{m+i}(\xi )=x_i`$, $`i=1,2,\mathrm{},n`$.
Denote by $`F_{m+n}`$ the free non-abelian group on generators $`g_1,\mathrm{},g_m`$, $`g_{m+1},\mathrm{},g_{m+n}`$. The group $`F_{m+n}`$ acts on $`𝕌`$ by isometries.
The formula
$$d(g,h):=d_𝕌(g(\xi ),h(\xi )),g,hF_{m+n},$$
where $`d_𝕌`$ denotes the metric on the Urysohn space, defines a left-invariant pseudometric $`d`$ on the group $`F_{m+n}`$:
$`d(xg,xh)`$ $`=`$ $`d_𝕌(xg(\xi ),xh(\xi ))`$
$`=`$ $`d_𝕌(x(g(\xi )),x(h(\xi )))`$
$`=`$ $`d_𝕌(g(\xi ),h(\xi ))`$
$`=`$ $`d(g,h).`$
The indexed metric subspace $`\{g_jx_i:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ of $`𝕌`$ is isometric to the metric subspace $`\{g_jg_{m+i}:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ of $`(F_{m+n},d)`$. Indeed,
$`d_𝕌(g_jx_i,g_kx_l)`$ $`=`$ $`d_𝕌(g_jg_{m+i}(\xi ),g_kg_{m+l}(\xi ))`$
$`=`$ $`d(g_jg_{m+j},g_kg_{m+l}).`$
Notice also that the latter subspace is contained in the set $`F_{m+n}^{(2)}`$ of all words having reduced length $`2`$. (The reduced length will always mean that with regard to the generators $`g_1,\mathrm{},g_m`$, $`g_{m+1},\mathrm{},g_{m+n}`$.)
By adding to $`d`$ a left-invariant metric on $`F_{m+n}`$ taking sufficiently small values on pairs of elements of $`F_{m+n}^{(2)}`$, we can assume without loss in generality that $`d`$ is a left-invariant metric on $`F_{m+n}`$. (For instance, add a metric whose only non-zero value is $`\epsilon /3`$.)
Form the Cayley graph $`\mathrm{\Gamma }`$ of $`F_{m+n}`$ with regard to the set of generators $`Y=F_{m+n}^{(4)}`$. Vertices of $`\mathrm{\Gamma }`$ are elements of $`F_{m+n}`$, and $`x,yF_{m+n}`$ are adjacent if and only if $`x^1yF_{m+n}^{(4)}`$. This graph is connected. Now make $`\mathrm{\Gamma }`$ into a weighted graph by assigning to an edge $`(a,b)`$, $`a^1bF_{m+n}^{(4)}`$ the value $`d(a,b)d(a^1b,e)`$.
Denote by $`\rho `$ the path metric of the weighted graph $`\mathrm{\Gamma }`$. Its value for $`x,yF_{m+n}`$ is given by
(7.1)
$$\rho (x,y)=inf\underset{i=0}{\overset{N1}{}}d(a_i,a_{i+1}),$$
where the infimum is taken over all natural $`N`$ and all finite sequences $`x=a_0,a_1,\mathrm{},a_{N1},a_N=b`$, with the property $`a_i^1a_{i+1}F_{m+n}^{(4)}`$ for all $`i`$.
It is easily seen that $`\rho `$ is a left-invariant metric on the group $`F_{m+n}`$.
Generally, $`\rho d`$, but restrictions of $`\rho `$ and $`d`$ to $`F_{m+n}^{(2)}`$ coincide.
If one denotes by $`\delta >0`$ the minimal value of $`d(a,b)`$ as $`a,bF_{m+n}^{(2)}`$ and $`ab`$, then $`\rho (x,y)\delta d_w(x,y)`$, where $`d_w`$ denotes the word metric with respect to the set of generators $`Y=F_{m+n}^{(4)}`$. In particular, if an $`xF_{m+n}`$ has reduced length $`l=l(x)`$, then $`d_w(x)l/4`$ and accordingly $`\rho (x,e)\delta l/4`$. (As a consequence, the infimum in Eq. (7.1) is always achieved.)
Let $`\mathrm{\Delta }`$ denote the maximal value of the metric $`\rho `$ between pairs of elements of $`F_{m+n}^{(2)}`$. Choose a natural number $`N`$ so large that $`\delta (N4)/4\mathrm{\Delta }`$, for instance, set $`N=4\mathrm{\Delta }/\delta +4`$.
Every free group is residually finite, that is, admits a separating family of homomorphisms into finite groups. (Cf. e.g. , Ch. 7, exercise 5.) Using this fact, choose a normal subgroup $`HF_{m+n}`$ of finite index so that $`HF_{m+n}^{(N)}=\{e\}`$.
The formula
$`\overline{\rho }(xH,yH)`$ $`:=`$ $`\underset{h_1,h_2H}{inf}\rho (xh_1,yh_2)`$
$``$ $`\underset{h_1,h_2H}{inf}\rho (h_1x,h_2y)`$
$``$ $`\underset{hH}{inf}\rho (hx,y)`$
defines a left-invariant pseudometric on the finite factor-group $`F_{m+n}/H`$. The triangle inequality follows from the fact that, for all $`h^{}H`$,
$`\overline{\rho }(xH,yH)`$ $`=`$ $`\underset{hH}{inf}\rho (hx,y)`$
$``$ $`\underset{hH}{inf}[\rho (hx,h^{}z)+\rho (h^{}z,y)]`$
$`=`$ $`\underset{hH}{inf}\rho (hx,h^{}z)+\rho (h^{}z,y)`$
$`=`$ $`\underset{hH}{inf}\rho (h^1hx,z)+\rho (h^{}z,y)`$
$`=`$ $`\overline{\rho }(xH,zH)+\rho (h^{}z,y),`$
and the infimum of the r.h.s. taken over all $`h^{}H`$ equals $`\overline{\rho }(xH,zH)+\overline{\rho }(zH,yH)`$. Left-invariance of $`\overline{\rho }`$ is obvious.
Let $`x,yF_{m+n}^{(2)}`$. Closely approximate the infimum in Eq. (Proof.) by some value $`\rho (xh_1,yh_2)`$ with $`h_1,h_2H`$, then
$$\rho (xh_1,yh_2)=\rho (y^1xh_1x^1yy^1x,h_2)=\rho (y^1x,h_3),$$
where $`h_3=y^1xh_1^1x^1yh_2H`$.
The value $`\rho (y^1x,h_3)`$, $`h_3H`$, cannot get smaller than $`d(y^1x,e)=d(x,y)`$. Indeed, unless $`h_3=e`$ (in which case $`\rho (y^1x,h_3)=\rho (x,y)=d(x,y)`$), one has $`l(h_3)N`$ and so the word distance from $`y^1x`$ to $`h_3`$ is at least $`N4`$, and $`\rho (y^1x,h)\delta (N4)/4\mathrm{\Delta }d(x,y)`$.
We conclude: the restriction of the factor-homomorphism
$$\pi :F_{m+n}xxHF_{m+n}/H$$
to $`F_{m+n}^{(2)}`$ is an isometry.
One can now perturb the pseudometric on $`F_{m+n}/H`$ by adding to it a left-invariant metric taking very small values (e.g. taking the only non-zero value $`\epsilon /3`$) so as to replace $`\overline{\rho }`$ with a left-invariant metric, $`\stackrel{~}{\rho }`$.
Take now $`\stackrel{~}{X}=(F_{m+n}/H,\stackrel{~}{\rho })`$, $`\stackrel{~}{x}_i=\pi (g_{m+i})\stackrel{~}{X}`$, $`i=1,2,\mathrm{},n`$, and let $`\stackrel{~}{g}_j`$ be left translates made by the elements $`\pi (g_j)`$, $`j=1,2,\mathrm{},m`$, in the finite group $`F_{m+n}/H`$. The indexed metric space $`\{g_jg_{m+i}:i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ of $`(F_{m+n},d)`$ is $`\epsilon `$-isometric to the metric subspace $`\{\pi (g_j)\pi (g_{m+i}):i=1,2,\mathrm{},n,j=1,2,\mathrm{},m\}`$ of $`(F_{m+n}/H,\stackrel{~}{\rho })`$. Consequently, the conclusion of the Theorem is verified. ∎
Remark. Prof. Henson has also pointed out to me that in the particular case of path metric spaces associated to a graph the above result (Theorem 3.2) follows from earlier results by Hrushovski .
Acknowledgements. The author is most grateful to Prof. C. Ward Henson for his comments, as well as to all other participants of the Research Among Peers (RAP) seminar series “Polish group actions and extremely amenable groups” otganized by Prof. Slawomir Solecki ind the Department of Mathematics, University of Illinois at Urbana-Champaign in Spring 2003.
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# Calculation of Exclusive Cross Sections with the Lorentz Integral Transform Method
## I Introduction
The study of inclusive and exclusive cross sections in inelastic reactions is an essential tool in understanding the underlying dynamics of a particle system. For systems with more than two constituents a major problem in the calculation of such reactions consists in the exact knowledge of the final state wave function in the continuum. Its calculation is far more difficult than the corresponding bound state calculation. As a matter of fact, today a computation of an intermediate energy continuum wave function is out of reach for a system with more than three particles. An exact calculation, however, can be carried out in an alternative way by using proper integral transforms. They allow one to take into account final state interactions (FSI) rigorously without using final state wave functions explicitly. In fact in recent years the Lorentz integral transform (LIT) method has been successfully applied to various inclusive breakup cross sections in few-body physics. The LIT is expressed in terms of square integrable functions which are obtained from inhomogeneous differential equations. The differential equations can be solved with similar methods as a bound state problem. After having calculated the transform one obtains the cross section by inversion of the transform. A recent overview of the results obtained with the LIT method is given in .
The LIT method has not yet been used for the calculation of exclusive cross sections. The only exception is a calculation of the <sup>4</sup>He spectral function , but in this case the method proceeds along the same lines as for an inclusive reaction. However, it is in principle possible to apply the integral transform method to general exclusive reactions as shown by Efros for the Stieltjes transform . In the present paper we illustrate the details of the calculation for the $`d(e^{},ep)`$ reaction with the LIT. It will also allow us to check the precision of the obtained results by comparing them to those of a conventional calculation with an explicit $`np`$ final state wave function. Such a reliability check is necessary, since the use of integral transforms is not always unproblematic for the calculation of cross sections. In fact one may encounter problems in the inversion of the transform. For example in Ref. it was shown for the case of an inclusive reaction that the Stieltjes transform is not very appropriated, since it samples contributions over a large energy range making the inversion extremely difficult. Later the LIT was proposed and it was shown for the test case of the longitudinal inelastic deuteron form factor in electron scattering that inclusive reactions can be safely calculated with this method. Whether the LIT is as appropriated also for exclusive reactions cannot be said a priori, since such a calculation is more complicated. The aim of the present work is to investigate this question.
The paper is organized as follows. In Sec. II we describe how cross sections are calculated with the LIT. Details of the calculation for the specific reaction under consideration ($`d(e,e^{}p)n`$) are given in Sec. III. The results are illustrated in Sec. IV and a conclusion is drawn in Sec. V.
## II The LIT method
The starting point of the LIT method is the calculation of an integral transform with a Lorentz kernel
$$L(\sigma )=𝑑\omega \frac{F(\omega )}{(w\sigma _R)^2+\sigma _I^2}.$$
(1)
The function $`F`$ depends on the internal excitation energy $`\omega =E_fE_0`$ of a given particle system and contains information about the transition of the system from the ground state $`|\mathrm{\Psi }_0`$, with energy $`E_0`$, to the final state $`|\mathrm{\Psi }_f`$, with energy $`E_f`$, induced by an external probe. In case of an inclusive reaction $`F(\omega )`$ denotes the response function
$$F(\omega )=𝑑\mathrm{\Psi }_f|\mathrm{\Psi }_f|\widehat{O}|\mathrm{\Psi }_0|^2\delta (E_fE_0\omega ),$$
(2)
where $`\widehat{O}`$ is a transition operator which characterizes the specific process under consideration.
The key point of the LIT method is an evaluation of $`L(\sigma )`$ without explicit knowledge of $`F(\omega )`$. In a second step the function $`F`$ is obtained from the inversion of the transform. The great advantage of the method lies in the fact that a calculation of the generally very complicated final state wave function $`|\mathrm{\Psi }_f`$ can be avoided as will be discussed below. On the contrary a conventional calculation of $`F(\omega )`$ can only be carried out with the explicit knowledge of $`|\mathrm{\Psi }_f`$.
Using completeness one can show that in order to obtain $`L(\sigma )`$ one has to solve the following differential equation
$$(HE_0\sigma ^{})|\stackrel{~}{\mathrm{\Psi }}_1(\sigma )=\widehat{O}|\mathrm{\Psi }_0$$
(3)
with
$$\sigma =\sigma _R+i\sigma _I\sigma _R,\sigma _I>0,$$
(4)
where $`H=T+V`$ is the Hamiltonian of the system under consideration. Note that the corresponding homogeneous equation has only the trivial solution, since $`H`$ has a real eigenvalue spectrum.
The norm of the solution $`\stackrel{~}{\mathrm{\Psi }}_1(\sigma )`$ determines the LIT directly:
$$L(\sigma )=\stackrel{~}{\mathrm{\Psi }}_1(\sigma )|\stackrel{~}{\mathrm{\Psi }}_1(\sigma ).$$
(5)
Different from a Schrödinger equation at positive energies, one has for the solution of Eq. (3) a very simple boundary condition. Due to the localized source at the right hand side (rhs) of Eq. (3) $`\stackrel{~}{\mathrm{\Psi }}_1(\sigma )`$ vanishes at large distances similar to a bound state wave function. Therefore one can apply the same techniques as for the calculation of a bound state wave function.
The response function $`F(\omega )`$ serves only for the determination of inclusive cross sections. For an exclusive process one needs a more detailed information about the transition of the system. In fact one has to be able to calculate transition matrix elements of the form
$$T_{fi}(E_f)=\mathrm{\Psi }_f|\widehat{O}|\mathrm{\Psi }_0.$$
(6)
How such a calculation can be carried out with an integral transform was shown for the case of the Stieltjes transform . For the LIT the calculation proceeds analogously as outlined in the following (see also ).
For simplicity we will consider an exclusive reaction leading to a final state with two fragments, but the method can be applied also for channels with more than two particles. Besides the correct final state wave function $`|\mathrm{\Psi }_f`$ we also introduce the corresponding plane wave
$$\mathrm{\Phi }^{\mathrm{PW}}(\stackrel{}{r})=𝒜\mathrm{\Psi }_1\mathrm{\Psi }_2\frac{\mathrm{exp}(i\stackrel{}{k}\stackrel{}{r})}{(2\pi )^{3/2}},$$
(7)
where $`𝒜`$ is a proper antisymmetrizer and $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ are the internal wave functions of the two fragments, while $`\stackrel{}{k}`$ and $`\stackrel{}{r}`$ are the usual relative coordinates for momentum and position of the two-body system formed by the two fragments.
In order to calculate the $`T`$ matrix element one starts from the Lippmann-Schwinger equation for the final state
$$\mathrm{\Psi }_f|=\mathrm{\Phi }^{\mathrm{PW}}|+\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}\frac{1}{E_f+iϵH}$$
(8)
with
$$\widehat{V}=\underset{i,j}{}\widehat{V}_{ij},i_1,j_2,$$
(9)
where $`_1`$ and $`_2`$ contain all particles of the first and second fragment, respectively. Inserting the above expression in Eq. (6) one obtains a sum of two pieces, a Born term
$$T_{fi}^{\mathrm{Born}}(E_f)=\mathrm{\Phi }^{\mathrm{PW}}|\widehat{O}|\mathrm{\Psi }_0,$$
(10)
and a term depending on FSI
$$T_{fi}^{\mathrm{FSI}}(E_f)=\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}\frac{1}{E_f+iϵH}\widehat{O}|\mathrm{\Psi }_0.$$
(11)
The calculation of the Born term is rather simple. The main difficulty of the calculation is the determination of the matrix element depending on FSI.
Using the completeness of the eigenstates $`|\mathrm{\Psi }(E)`$ of $`H`$ one can rewrite $`T_{fi}^{\mathrm{FSI}}`$ as follows
$`T_{fi}^{\mathrm{FSI}}(E_f)`$ $`=`$ $`{\displaystyle 𝑑E\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\mathrm{\Psi }(E)\mathrm{\Psi }(E)|\frac{1}{E_f+iϵH}\widehat{O}|\mathrm{\Psi }_0}`$ (12)
$`=`$ $`{\displaystyle 𝑑EF_{fi}(E)\frac{1}{E_f+iϵE}},`$ (13)
with
$$F_{fi}(E)=\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\mathrm{\Psi }(E)\mathrm{\Psi }(E)|\widehat{O}|\mathrm{\Psi }_0.$$
(14)
One obtains the following formal solution for the FSI term
$$T_{fi}^{\mathrm{FSI}}(E_f)=i\pi F_{fi}(E_f)+𝒫_{E_0}^{\mathrm{}}𝑑E\frac{F_{if}(E)}{E_fE}.$$
(15)
For a calculation of $`T_{fi}`$ one needs to know $`F_{fi}`$ for any given energy. A direct calculation of $`F_{fi}`$ is of course in general far too difficult, since one has to determine final state wave functions $`\mathrm{\Psi }`$ for the whole eigenvalue spectrum of $`H`$. On the other hand an indirect calculation via the LIT is possible. To this end one performs a Lorentz integral transform of $`F_{fi}`$, i.e.
$$L(\sigma )=_{E_0}^{\mathrm{}}𝑑E\frac{F_{fi}(E)}{(E\sigma )(E\sigma ^{})},\sigma _R>E_0.$$
(16)
Inserting the definition of $`F_{fi}`$ one finds
$`L(\sigma )`$ $`=`$ $`{\displaystyle _{E_0}^{\mathrm{}}}𝑑E\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}{\displaystyle \frac{1}{H\sigma }}|\mathrm{\Psi }(E)\mathrm{\Psi }(E)|{\displaystyle \frac{1}{H\sigma ^{}}}\widehat{O}|\mathrm{\Psi }_0`$ (17)
$`=`$ $`\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}{\displaystyle \frac{1}{H\sigma }}{\displaystyle \frac{1}{H\sigma ^{}}}\widehat{O}|\mathrm{\Psi }_0`$ (18)
$`=`$ $`\stackrel{~}{\mathrm{\Psi }}_2(\sigma )|\stackrel{~}{\mathrm{\Psi }}_1(\sigma ),`$ (19)
with
$`|\stackrel{~}{\mathrm{\Psi }}_1(\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{H\sigma _R+i\sigma _I}}\widehat{O}|\mathrm{\Psi }_0,`$ (20)
$`|\stackrel{~}{\mathrm{\Psi }}_2(\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{H\sigma _R+i\sigma _I}}\widehat{V}|\mathrm{\Phi }^{\mathrm{PW}}.`$ (21)
It is evident that (20) leads essentially to the same differential equation as for the inclusive process:
$$(H\sigma _R+i\sigma _I)|\stackrel{~}{\mathrm{\Psi }}_1(\sigma )=\widehat{O}|\mathrm{\Psi }_0.$$
(22)
From Eq. (21) one obtains
$$(H\sigma _R+i\sigma _I)|\stackrel{~}{\mathrm{\Psi }}_2(\sigma )=\widehat{V}|\mathrm{\Phi }^{\mathrm{PW}},$$
(23)
which is similar to Eq. (22), but for a different source term on the rhs. For a finite range potential one has also in this case a vanishing source term for large distances. This guarantees also for $`\stackrel{~}{\mathrm{\Psi }}_2`$ an asymptotic boundary condition similar to a ground state problem. If $`\widehat{V}`$ contains also the Coulomb potential one cannot proceed exactly in the same way as shown here. In this case one has to start from a modified Lippmann-Schwinger equation, where Coulomb wave function are taken into account .
As shown by Efros the integral transform method can be extended to exclusive processes with more than two fragments. In principle one obtains equations similar to Eqs. (19,22,23). However, one cannot guarantee, as in the two fragment case, that the potential $`\widehat{V}`$ vanishes asymptotically, since two of the fragments could remain close to each other. Therefore it is necessary to choose a different solution to the problem. In fact one can rewrite $`L(\sigma )`$ as follows
$$L(\sigma )=\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\stackrel{~}{\stackrel{~}{\mathrm{\Psi }}}_1,$$
(24)
where $`|\stackrel{~}{\stackrel{~}{\mathrm{\Psi }}}_1`$ is obtained from the solution of the following differential equation
$$(H\sigma _Ri\sigma _I)|\stackrel{~}{\stackrel{~}{\mathrm{\Psi }}}_1(\sigma )=|\stackrel{~}{\mathrm{\Psi }}_1(\sigma ).$$
(25)
Since $`|\stackrel{~}{\mathrm{\Psi }}_1`$ vanishes at large distances, one can again use bound state methods for the solution of this differential equation.
There is another possibility to calculate the LIT . Starting from the following identity
$$\frac{1}{(H\sigma )(H\sigma ^{})}=\frac{1}{2i\sigma _I}\left(\frac{1}{H\sigma }\frac{1}{H\sigma ^{}}\right)$$
(26)
and defining
$$|\stackrel{~}{\mathrm{\Psi }}_1^{}(\sigma )=\frac{1}{H\sigma _Ri\sigma _I}\widehat{O}|\mathrm{\Psi }_0,$$
(27)
one gets for the transform
$$L(\sigma )=\frac{1}{2i\sigma _I}\left(\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\stackrel{~}{\mathrm{\Psi }}_1^{}(\sigma )\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\stackrel{~}{\mathrm{\Psi }}_1(\sigma )\right).$$
(28)
It is seen that one has to solve only one type of differential equation. In fact for the solution of $`|\stackrel{~}{\mathrm{\Psi }}_1^{}(\sigma )`$ it is sufficient to solve Eq. (22) a second time replacing $`\sigma ^{}`$ by $`\sigma `$.
Which of the two approaches for a nonvanishing source term is more appropriate depends also on the possibility to obtain a precise numerical solution. Small errors in the solution of Eq. (22) could lead in both cases to larger errors for the calculation of the LIT. Therefore we will also consider in Sec. IV these additional possibilities to calculate the LIT for our realistic test case of the electromagnetic deuteron breakup.
## III exclusive deuteron breakup
The exclusive deuteron breakup $`d(e,e^{}p)n`$ is governed by the four structure functions $`f_L`$, $`f_T`$, $`f_{LT}`$, and $`f_{TT}`$. In the following we will only consider the longitudinal $`f_L`$. We use the same notation as in Ref. , where the structure functions are calculated in the final $`np`$ c.m. system with $`f_L=f_L(E_{np},|\stackrel{}{q}_{c.m.}|,\theta )`$. In this system $`\stackrel{}{q}_{c.m.}`$ denotes the momentum transfer, $`E_{np}`$ the relative $`np`$ energy, and the relative $`np`$ momentum $`\stackrel{}{k}`$ has the angle $`\theta `$ with respect to $`\widehat{q}_{c.m.}`$. The structure functions are expressed in terms of the transition matrix $`T_{\mathrm{S}m\mu m_d}`$ for the process $`e+de^{}+np`$. The quantum numbers S and $`m`$ denote the spin and spin projection of the outgoing $`np`$ pair with respect to $`\widehat{k}`$, $`\mu `$ characterizes the transition operator $`\widehat{O}`$, and $`m_d`$ is the projection of the deuteron spin with respect to $`\widehat{q}_{c.m.}`$. For the longitudinal structure function, i.e. $`\mu =0`$, one has the following transition operator
$$\widehat{O}=\underset{j}{}G_{E,j}(q_\mu ^2)\mathrm{exp}(i\stackrel{}{q}_{c.m.}\stackrel{}{r}_j),$$
(29)
where $`G_{E,j}`$ denotes the electric form factor of the j-th nucleon with $`q_\mu ^2`$ being the four-momentum transfer squared. Here we use the electric dipole form factor for the proton, while the neutron electric form factor is set to zero. It is evident that the above operator does not affect the spin, i.e. $`S=1`$, and one gets
$$f_L(E_{np},q_{c.m.},\theta )=\underset{mm_d}{}T_{1m0m_d}T_{1m0m_d}^{},$$
(30)
where
$$T_{1m0m_d}=C1m|\widehat{O}|m_d$$
(31)
with
$$C=(2\pi )^{\frac{3}{2}}\sqrt{\frac{kM\alpha }{4\pi }},$$
(32)
where $`M`$ denotes the nucleon mass and $`\alpha `$ is the fine structure constant. Like $`f_L`$ the transition matrix depends on $`E_{np}`$, $`q_{c.m.}`$, and $`\theta `$. In the following we suppress the dependences on $`q_{c.m.}`$ and $`\theta `$, while the dependence on $`E_{np}`$ is made explicit in most cases.
The aim of the present work is a test of the LIT method for the exclusive deuteron breakup. Since there is no need to obtain results with a realistic potential we choose the semi-realistic TN potential for this test. It is a central potential model, which is different for the various spin-isospin (ST) channels, i. e. $`V(r)=_{\mathrm{ST}}V^{\mathrm{ST}}(r)`$. Since we have $`S=1`$, only $`V^{10}`$ and $`V^{11}`$ have to be considered here. In addition $`V^{11}(r)=0`$ for the TN potential, hence FSI effects appear only in the channel ST=10.
Due to the absence of the tensor force in the TN potential one obtains a reduced complexity of the equations to be solved. Nevertheless, as calculations with this potential model show (see e.g. ) results are sufficiently realistic in order to serve as test case.
As discussed in Sec. II there are two contributions for any $`T`$-matrix element, a Born and an FSI term. Considering an $`s`$-wave deuteron with radial wave function $`u(r)`$ one obtains for the Born term
$$T_{1m0m_d}=()^{m_d}\delta _{mm_d}G_{E,p}\sqrt{3kM}Y_{00}(\widehat{k}_{})_0^{\mathrm{}}𝑑rru(r)j_0(k_{}r),$$
(33)
where $`j_0`$ denotes the spherical Bessel function of order 0 and
$$\stackrel{}{k}_{}=\stackrel{}{k}\frac{\stackrel{}{q}_{c.m.}}{2}.$$
(34)
For the second piece, $`T^{\mathrm{FSI}}`$, it is necessary to perform a multipole decomposition. It is convenient to introduce the following expansions with projections $`M`$ with respect to $`\widehat{q}_{c.m.}`$
$`\stackrel{~}{\mathrm{\Psi }}_{1,M}`$ $`=`$ $`{\displaystyle \underset{j,L}{}}i^L[Y^{[L]}(\widehat{r})\times \chi ^{[1]}(\stackrel{}{\sigma }_1,\stackrel{}{\sigma }_2)]_M^{[j]}\sqrt{2L+1}C_{0MM}^{L1J}r^1\stackrel{~}{\psi }_{Lj}^{(1)}(r),`$ (35)
$`\stackrel{~}{\mathrm{\Psi }}_{2,M}`$ $`=`$ $`{\displaystyle \underset{j,l,m_l}{}}i^l[Y^{[l]}(\widehat{r})\times \chi ^{[1]}(\stackrel{}{\sigma }_1,\stackrel{}{\sigma }_2)]_M^{[j]}C_{m_lmM}^{l1j}Y_{lm_l}^{}(\widehat{k})r^1\stackrel{~}{\psi }_{lj}^{(2)}(r),`$ (36)
where $`\chi ^{[1]}(\stackrel{}{\sigma }_1,\stackrel{}{\sigma }_2)`$ denotes the spin wave function for a two-nucleon system with $`S=1`$. For the rhs of the differential equations (22) and (23) we perform similar expansions as for $`\stackrel{~}{\mathrm{\Psi }}_1`$ and $`\stackrel{~}{\mathrm{\Psi }}_2`$ leading to
$`\widehat{O}|m_d`$ $`=`$ $`{\displaystyle \underset{j,L}{}}i^L[Y^{[L]}(\widehat{r})\times \chi ^{[1]}(\stackrel{}{\sigma }_1,\stackrel{}{\sigma }_2)]_{m_d}^{[j]}\sqrt{2L+1}C_{0m_dm_d}^{L1j}r^1f_L(r),`$ (37)
$`f_L(r)`$ $`=`$ $`j_L({\displaystyle \frac{q_{c.m.}r}{2}})u(r)`$ (38)
and
$`\widehat{V}|\mathrm{\Phi }_M^{\mathrm{PW}}`$ $`=`$ $`{\displaystyle \underset{j,l,m_l}{}}i^l[Y^{[l]}(\widehat{r})\times \chi ^{[1]}(\stackrel{}{\sigma }_1,\stackrel{}{\sigma }_2)]_M^jC_{m_lmM}^{l1jM}Y_{lm_l}^{}(\widehat{k})r^1g_{lj}(r),`$ (39)
$`g_{lj}(r)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}rj_l(kr)V_{jl}(r),`$ (40)
where $`V_{jl}`$ is the potential for the NN partial wave $`{}_{}{}^{3}l_{j}^{}`$.
For $`\stackrel{~}{\mathrm{\Psi }}_2`$ the above multipole decompositions lead to the following coupled differential equations in real and imaginary parts
$`\left\{{\displaystyle \frac{\mathrm{}^2}{M}}\left({\displaystyle \frac{d^2}{dr^2}}{\displaystyle \frac{l(l+1)}{r^2}}\right)+V_{lj}(r)\sigma _R\right\}\mathrm{}[\stackrel{~}{\psi }_{lj}^{(2)}(\sigma ,r)]\sigma _I\mathrm{}[\stackrel{~}{\psi }_{lj}^{(2)}(\sigma ,r)]`$ $`=`$ $`g_{lj}(r)`$ (41)
$`\left\{{\displaystyle \frac{\mathrm{}^2}{M}}\left({\displaystyle \frac{d^2}{dr^2}}{\displaystyle \frac{l(l+1)}{r^2}}\right)+V_{lj}(r)\sigma _R\right\}\mathrm{}[\stackrel{~}{\psi }_{lj}^{(2)}(\sigma ,r)]+\sigma _I\mathrm{}[\stackrel{~}{\psi }_{lj}^{(2)}(\sigma ,r)]`$ $`=`$ $`0.`$ (42)
As mentioned above, because of our potential model we have to consider only the channel with ST=10, thus $`V_{lj}(r)`$ can be replaced by $`V^{10}(r)`$. In addition the Pauli principle has to be fulfilled, i.e. S+T+$`l`$ has to be odd. Therefore the differential equation has only to be solved for $`l`$ even. For channels with l odd one had to consider $`V^{11}`$, but as already mentioned $`V^{11}`$ is zero in our potential model. Note that there is no explicit dependence on $`j`$ in the coupled differential equation, thus one has $`\stackrel{~}{\psi }_{lj}^{(2)}=\stackrel{~}{\psi }_{lj^{}}^{(2)}`$.
For $`\stackrel{~}{\mathrm{\Psi }}_1`$ one finds very similar equations with the only difference that one has to replace $`g_{lj}(r)`$ by $`f_L(r)`$ on the rhs. Also here we have $`\stackrel{~}{\psi }_{Lj}^{(1)}=\stackrel{~}{\psi }_{Lj^{}}^{(1)}`$.
We solve the differential equation by adding an additional homogeneous equation determining the source terms on the rhs. In this way we obtain a coupled homogeneous differential equation system. The numerical solution leads to very precise results as shown in Ref. .
With the solutions for $`\stackrel{~}{\psi }_{Lj}^{(1)}`$ and $`\stackrel{~}{\psi }_{lj}^{(2)}`$ one obtains for the scalar product (19)
$$\stackrel{~}{\mathrm{\Psi }}_2|\stackrel{~}{\mathrm{\Psi }}_1=\underset{mm_dlm_l}{}\sqrt{2l+1}Y_{lm_l}(\widehat{k})\underset{j}{}C_{m_lmm_d}^{l1j}C_{0m_dm_d}^{l1j}\stackrel{~}{R}_{lj}(\sigma )$$
(43)
with
$$\stackrel{~}{R}_{lj}(\sigma )=_0^{\mathrm{}}𝑑r(\stackrel{~}{\psi }_{lj}^{(2)}(\sigma ,r))^{}\stackrel{~}{\psi }_{Lj}^{(1)}(\sigma ,r)\delta _{Ll}.$$
(44)
To calculate the FSI contribution to the $`T`$-matrix elements one has to invert the LIT
$$\stackrel{~}{R}_{lj}(\sigma )=_{E_0}^{\mathrm{}}𝑑E\frac{R_{lj}(E)}{(E\sigma _R)^2+\sigma _I^2}$$
(45)
in order to obtain the function $`R_{lj}(E)`$. The transform can be inverted using the following ansatz
$$R_{lj}(E)=\underset{n=1}{\overset{N}{}}c_{n,lj}\chi _{n,lj}(E,\beta ),$$
(46)
where $`\chi _{n,lj}`$ are given functions with nonlinear parameters $`\beta `$. Substituting this expansion into the rhs of Eq. (45) one obtains
$$\stackrel{~}{R}_{lj}(\sigma )=\underset{n=1}{\overset{N}{}}c_{n,lj}\stackrel{~}{\chi }_{n,lj}(\sigma ,\beta ),$$
(47)
where the $`\stackrel{~}{\chi }_{n,lj}`$ are the Lorentz integral transforms of the basis functions. The parameters $`c_{n,lj}`$ and $`\beta `$ are determined by fitting the calculated transform $`\stackrel{~}{R}_{lj}(\sigma )`$ to the above expansion at many $`\sigma _R`$ points for a fixed $`\sigma _I`$. The number of functions $`N`$ plays the role of a regularization parameter and is chosen within a stability region, i.e. where the obtained results are stable for a certain range of $`N`$ (see also ). Here we use the following set of basis functions
$$\chi _{n,lj}(E,\beta )=E^{l+\frac{1}{2}}exp(\frac{\beta E}{n}).$$
(48)
For the parametrization of the elastic monopole transition for $`l=0`$ we include an additional function in the set $`\chi `$:
$$\chi _{0,0j}(E,\beta )=\delta (EE_0).$$
(49)
Thus the sum in Eqs. (46, 47) starts in this case with $`n=0`$ instead of $`n=1`$.
Once the inversion is carried out one can make use of Eq. (15). For a specific set of $`m`$ and $`m_d`$ one finds the following FSI contribution to the $`T`$-matrix
$$T_{1m0m_d}^{\mathrm{FSI}}(E_{np})=\underset{lm_l}{}Y_{lm_l}(\widehat{k})\sqrt{\frac{2l+1}{4\pi }}t_{1m0m_d}^{\mathrm{FSI},l}(E_{np})$$
(50)
with
$$t_{1m0m_d}^{\mathrm{FSI},l}(E_{np})=C\sqrt{4\pi }G_{E,p}(q_\mu ^2)\underset{j}{}C_{m_lmm_d}^{l1j}C_{0m_dm_d}^{l1j}\left(i\pi R_{lj}(E_{np})+𝒫_{E_0}^{\mathrm{}}𝑑E\frac{R_{lj}(E)}{E_{np}E}\right).$$
(51)
For our case without tensor force the following simple relations hold
$`t_{1m0m_d}^{\mathrm{FSI},l}`$ $`=`$ $`\delta _{m,m_d}t_{1m0m_d}^{\mathrm{FSI},l}`$ (52)
$`t_{1101}^{\mathrm{FSI},l}`$ $`=`$ $`t_{1000}^{\mathrm{FSI},l}=t_{1101}^{\mathrm{FSI},l}.`$ (53)
It is worth mentioning that as a byproduct of the calculation one obtains also the NN phase shifts from our calculation, i.e. without having solved the Schrödinger equation for the scattering state. The ratio of imaginary and real parts of a given transition matrix element is equal to $`tan(\delta )`$ (see e.g., Ref. ). For our simple potential model one obtains
$$\delta _l=atan\left(\frac{\mathrm{}(t^{\mathrm{FSI},l})}{\mathrm{}(t^{\mathrm{FSI},l})+t^{\mathrm{Born},l}}\right),$$
(54)
where $`t^{\mathrm{Born},l}`$ is analogously defined as $`t^{\mathrm{FSI},l}`$, and easily evaluated from a multipole decomposition of the Born term in Eq. (33). If one is only interested in the phase shifts themselves, one can perform a simpler calculation neglecting the excitation operator $`\widehat{O}`$ (see Ref. ).
## IV discussion of results
We test the LIT method for exclusive reactions choosing for the electromagnetic deuteron breakup three different kinematics with rather strong FSI effects: (i) in the tail region beyond the quasi-elastic peak at moderate momentum transfer ($`E_{np}=120`$ MeV, $`q_{c.m.}^2=5`$ fm<sup>-2</sup>), (ii) on the photon line ($`E_\gamma =70`$ MeV), and (iii) close to the deuteron breakup threshold ($`E_{np}=1`$ MeV, $`q_{c.m.}^2=2`$ fm<sup>-2</sup>). Note for kinematics (ii) that there is no longitudinal contribution to the exclusive $`(e,e^{}p)`$ cross section, but that nevertheless the structure function $`f_L`$ does not vanish. In fact in applying Siegert’s theorem the longitudinal matrix elements are commonly used in photodisintegration and lead to the dominant contribution for the electric transitions.
In order to have a more detailed comparison between the LIT results and the results of a conventional calculation we do not simply discuss the final result for $`f_L`$, but rather study directly the FSI effect on the various multipole transitions. This allows us to make a much more precise comparison between the two calculations. Because of Eqs. (52,53) it is sufficient to consider $`t_{1101}^{\mathrm{FSI},l}`$ in the following.
Before turning to the above mentioned three kinematical cases we first illustrate results for the transform in a more general way choosing a constant $`q_{c.m.}^2`$ of 5 fm<sup>-2</sup> and various energies $`E_{np}`$. In Fig. 1 we show $`\stackrel{~}{R}_{ll}(\sigma _R,\sigma _I=20\mathrm{MeV})`$ for $`l=0`$, 2. Its inversion, $`R_{ll}`$, gives a contributes to $`t^{\mathrm{FSI},l}`$ (see Eq. (51)). For the $`l=0`$ transforms one has an interesting structure at small $`\sigma _R`$. It originates from the rather strong monopole transition strength close to the deuteron breakup threshold. It is interesting to see that the peak becomes more and more pronounced for the case that also $`E_{np}`$ moves closer to the threshold region. In addition there is a second rather sizable contribution in the low $`\sigma _R`$ range. It arises from the elastic monopole contribution. Therefore one has to pay attention in the inversion of $`\stackrel{~}{R}_{00}`$. One has to check whether $`\sigma _I`$ is small enough to resolve with sufficient precision the elastic contribution from the threshold contribution. For $`l=2`$ one has a rather different picture. One finds a peak in the quasi-elastic region. Note that for the considered momentum of $`q_{c.m.}^2`$ of 5 fm<sup>-2</sup> the quasi-elastic peak is situated at about $`E_{np}=50`$ MeV. To find such a quasi-elastic peak for $`\stackrel{~}{R}_{22}`$ is a bit surprising, since one does not expect there strong FSI effects. On the other hand FSI should be small not on an absolute scale but compared with the corresponding Born term. Furthermore, the real part of the FSI contribution is difficult to estimate from Fig. 1 because a principle value integral has to be calculated in this case (see Eq. 51).
In Fig. 2 we show $`R_{ll}(E)`$ of kinematics (i) for $`l=0`$, 2, 4 and $`\sigma _I=5`$ and 20 MeV. It is obtained from the inversion of the corresponding $`\stackrel{~}{R}_{ll}(\sigma _R,\sigma _I)`$ (see Eq. (45)). One sees that six basis functions are not sufficient for the inversion, but for a higher $`N`$ one obtains a very nice stability of the inversion. Comparing the results with different $`\sigma _I`$, one finds a small difference for $`l=0`$ in the threshold region. The differences arise because the monopole contribution has a peak at the very threshold which has to be separated from the elastic contribution at $`E_0=2.225`$ MeV. From the inversion we obtain an elastic contribution of about 1.1 fm$`^{\frac{3}{2}}`$ which is rather sizable compared to the inelastic part with a peak height of 0.045 fm$`^{\frac{3}{2}}`$. Therefore is not surprising that the higher resolution with $`\sigma _I=5`$ MeV leads to a somewhat different result. However, because of the rather high $`E_{np}`$ of 120 MeV, the difference at the threshold is rather unimportant. This is confirmed by the results for $`t^{\mathrm{FSI}}`$, which are shown in Fig. 3 as function of the number of basis functions used for the inversion. In fact the agreement among the results with $`\sigma _I=5`$ and 20 MeV is very good. It is seen that one obtains for $`N10`$ for all considered multipolarities $`l`$ and for both $`\sigma _I`$ values very similar and stable results.
Also shown in Fig. 3 is the $`t^{\mathrm{FSI}}`$ of a conventional calculation. These results are very similar to the LIT results with relative differences of less than 1%. Only for the real part of the $`l=4`$ transition the difference is a little bit larger. On the other hand one has also to consider that this matrix element is very small. In fact its size is only $`2.5`$% of the corresponding $`t^{\mathrm{Born}}`$ matrix element. Thus the relative difference for the total matrix element is of the order of $`10^4.`$ For such a small FSI effect a part of the differences could also be due to a not completely exact result of the conventional calculation. Different from the LIT method $`t^{\mathrm{FSI}}`$ is not calculated directly, but taken indirectly from the difference of $`t^{\mathrm{total}}`$ and $`t^{\mathrm{Born}}`$; here $`t^{\mathrm{total}}`$ corresponds to the transition with the correct $`np`$ final state wave function in presence of the potential. The FSI effect is much more sizable for the two other transitions. Taking also here the ratio of $`t^{\mathrm{FSI}}/t^{Born}`$ for the real parts, one finds for $`l=0`$ about $`40`$% and for $`l=2`$ about $`30`$%.
Results for kinematics (ii) are shown in Figs. 4 and 5. For the $`R_{ll}(E)`$ of Fig. 4 one finds again nice stabilities of the inversion for a larger number of basis functions. Comparing the $`R_{ll}`$ with the two different $`\sigma _I`$ one has also here differences for the monopole transition and in addition for $`l=4`$. The monopole is of course not relevant for this kinematics on the photon line, since there is no corresponding electric monopole. On the other hand it is interesting to see whether one is able to separate the strongly peaked threshold strength from the dominant elastic contribution. Due to the lower momentum transfer one obtains an even larger elastic $`R_{ll}`$ than for kinematics (i), namely a value of about 14 fm$`^{\frac{3}{2}}`$ . The $`t^{\mathrm{FSI}}`$ results are shown in Fig. 5. The real parts turn out to be very stable as function of number of inversion basis functions. They are also very similar for both $`\sigma _I`$. Here we have the following relative FSI effects comparing with the Born term: $`100`$% ($`l=0)`$, $`25`$% ($`l=2)`$, $`1.5`$% ($`l=4)`$.
In comparison to the conventional calculation one finds in Fig. 5 for all the real parts only very small differences of less than 1%. For the imaginary part of $`t^{\mathrm{FSI}}`$ the picture is a bit different. The $`l=2`$ results are very stable and agree with extremely high precision to the results of the conventional calculation. Also the $`l=4`$ results are stable, but they are a few percent larger than found in the conventional calculation. However, one should note that the matrix element is very small and hence the difference of a few percent is not relevant. In fact comparing with the above mentioned size of the real part of the total matrix element the difference between both calculations is of the order of $`10^4`$. The imaginary part of the $`l=0`$ transition shows a bit less stability with the number of inversion basis function reflecting also the above mentioned problems for the separation of the elastic contribution. On the other hand one obtains reliable results for the highest $`N`$’s.
For the third kinematics we illustrate the results in Figs. 6 and 7. Here we consider only $`l=0`$ and $`l=2`$ transitions, since FSI effects do not play any role for higher transitions at threshold. In fact the FSI contribution is already very small for $`l=2`$. The inversion results in Fig. 6 are again very stable, except for $`l=0`$ with $`\sigma _I=20`$ MeV. Of course, again it is the problem associated with the elastic contribution ($`R_{00}(E_0)`$ is about 15 fm$`^{\frac{3}{2}}`$). Figure 7 shows that one obtains very good results for the real part of $`t^{\mathrm{FSI}}`$ with $`\sigma _I=5`$ MeV, while there is somewhat less stability for the inversion results with $`\sigma _I=20`$ MeV. The comparison with the conventional calculation is also here satisfactory. There are only deviations of about 1%. Again we list the relative FSI effect comparing with the Born term: -230% ($`l=0)`$, +2.5% ($`l=2)`$. The imaginary parts are somewhat more problematic. For $`l=0`$ one has the already mentioned problem with the elastic contribution combined with the fact that one needs $`R_{00}(E)`$ close to threshold ($`E_{np}=1`$ MeV), but with $`\sigma _I=5`$ MeV one obtains a sufficiently good result as seen from the comparison to the result of the conventional calculation. Though the relative differences to the conventional calculation are rather large for the imaginary part of the $`l=2`$ transition, its value is in principle correct, since it is more or less identical to 0. Note that it is about 200 times smaller than the already very small real FSI part of the $`l=2`$ transition.
We do not show results for the angular distribution of $`f_L`$. However, from the discussion above it should be clear that the two different calculations lead for $`f_L(\theta )`$ to relative differences of considerably less than 1% for kinematics (i) and (ii) and of about 1% for kinematics (iii).
As mentioned in Sec. II one has to use a somewhat different method for the calculation of the LIT for an exclusive reaction with more than two fragments in the final state. Two other possibilities are discussed at the end of Sec. II. In both cases one has to solve different differential equations, e.g., (25) instead of Eq. (23). However, these new methods appear to be numerically more problematic. Small errors in $`\stackrel{~}{\mathrm{\Psi }}_1`$, the solution of the differential equation (22), might lead to a much larger error for the solution of Eq. (25), where $`\stackrel{~}{\mathrm{\Psi }}_1`$ serves as source term on the rhs. Also for the determination of the LIT via Eq. (28) it is important how precise $`\stackrel{~}{\mathrm{\Psi }}_1`$ and $`\stackrel{~}{\mathrm{\Psi }}_1^{}`$ are calculated, since one has to determine the difference $`\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\stackrel{~}{\mathrm{\Psi }}_1^{}(\sigma )\mathrm{\Phi }^{\mathrm{PW}}|\widehat{V}|\stackrel{~}{\mathrm{\Psi }}_1(\sigma )`$. We are able to study this question for the $`d(e,e^{}p)n`$ reaction, since we can also use these alternative ways of evaluating the LIT. As a matter of fact both alternative methods lead in our case essentially to the same results with relative differences smaller than 0.01%. In Fig. 8 we show for a few selected cases these new LIT results compared to those obtained with Eq. (19). On finds relative deviations of the order of 1%. There are larger differences for the kinematics with $`E_{np}=1`$ and 120 MeV beyond a $`\sigma _R`$ of 150 MeV, but they are rather unimportant, since both $`\stackrel{~}{R}_{00}`$ are very small there. In fact $`\stackrel{~}{R}_{00}`$ crosses zero at about 205 and 190 MeV for $`E_{np}=1`$ and 120 MeV, respectively. Altogether one can say that one does not encounter greater numerical problems in evaluating the LIT with these alternative ways. Therefore also a calculation of an exclusive reaction to a three-body channel should lead to rather reliable results with the LIT method.
## V conclusion
We have calculated the longitudinal response of the exclusive $`d(e,e^{}p)`$ reaction with the method of the Lorentz integral transform. This method allows one to include the complete FSI, however, without explicit use of final state wave functions. It is the first time that the LIT method is applied to an exclusive reaction. In the past only inclusive processes have been studied with the LIT. The great success of the method raised the question whether it can also be successfully used in exclusive reactions. The results in this work show that one obtains a very precise determination of the various transition matrix elements. Differences to the conventional calculation are generally below 1 %. Only in the case of a very small FSI effect on the transition strength, i.e. a 10<sup>-4</sup> effect compared to the corresponding Born term, one can also obtain somewhat higher differences of a few percent. However, in this case differences could, as discussed in Sec. IV, at least partly be due to a small inexactness in the conventional calculation. There is only one exception, where one can expect a somewhat larger size of the error of the LIT result. This is the case for a transition matrix element in a region with transition strength from two (or more) rather narrow lying peaks. We had chosen such a situation with our kinematics (iii), where we have a strong elastic contribution at about $`E=2.2`$ MeV and another strong peak right above breakup threshold. In such a situation one should try to improve the resolution of the transform $`L(\sigma _R,\sigma _I)`$ by choosing a smaller $`\sigma _I`$. In fact our results improve significantly from $`\sigma _I=20`$ MeV to $`\sigma _I=5`$ MeV.
In case of an exclusive reaction with more than two fragments in the final state one cannot proceed exactly in the same way as for the breakup in two fragments. In this case one has to use other ways for the determination of the LIT. We could show that also these alternative methods lead to rather precise results. Therefore, in general, we may conclude that the LIT method leads to reliable results not only for inclusive, but also for exclusive reactions.
## Acknowledgment
We thank V. D. Efros and G. Orlandini for helpful discussions.
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# Extremely Red Objects from the NICMOS/HST Parallel Imaging Survey
## 1. Introduction
Deep near-IR imaging surveys have revealed a population of extremely red objects (“EROs”; Elston, Rieke & Rieke 1988; McCarthy, Persson & West 1992; Graham & Dey 1996; Hu & Ridgeway 1994; Soifer et al. 1994; Dey, Spinrad & Dickinson 1995; Thompson et al. 1999). The nature of the extremely red population remains unclear. As it is defined largely by a single color, primarily $`\mathrm{R}\mathrm{K}`$, there is no certainty that it represents a uniform class of object, and it may contain contributions from galaxies, cool stars or substellar objects and active nuclei. The precise definition of an ERO varies among the different surveys and depends on the particular bandpasses employed. Most samples were defined by $`\mathrm{R}\mathrm{K}\mathrm{}>56`$ or $`\mathrm{I}\mathrm{K}\mathrm{}>45`$. The present work is based on a somewhat different color system, that defined by the NICMOS F160W bandpass and conventional Kron-Cousins R magnitudes. The NICMOS F160W band is similar to the Johnson H band filter and the $`\mathrm{H}\mathrm{F160W}`$ color term is negligable for a flat spectral energy distribution (in $`f_\nu `$ units) (M. Rieke, 1999, private communication). Among the resolved objects there are reasonable expectations that there should exist stellar systems at redshifts such that the K-correction applied to an old or intermediate age population will produce very red optical to near-IR colors. Alternatively even fairly modest extinction, when observed in the same redshift range, can produce very steep spectral energy distributions in the rest-frame near-UV and there are local examples of such objects among the dusty starburst population. The central issue regarding the nature of the resolved EROs is to understand to what degree these two classes of objects contribute to the overall population.
The earliest interpretation of the colors of EROs were centered around the old stellar population hypothesis (e.g. McCarthy, Persson, & West 1991; Hu & Ridgeway 1994) and redshifts of 1 - 2 were inferred from their multi-band photometry. There are clear examples which support this interpretation, such as weak radio source LBDS 53W091 at $`z=1.55`$ (Dunlop et al. 1996) and near-IR selected object CL0939+4713B at z$`=1.58`$ (Soifer et al. 1999), and a concentration of EROs at $`z=1.3`$ (Liu et al. 2000). Graham and Dey (1996) argued that the spectral energy distribution of HR10 (Hu & Ridgeway 1994) is well matched by a dusty star-forming galaxy at $`z=1.5`$. The detection of a strong sub-mm continuum from HR10 (Cimatti et al. 1998; Dey et al. 1999) provided conclusive evidence that some of EROs, if not all, are dust enshrouded starburst galaxies with inferred star formation rates of $`5002000\mathrm{h}_{50}^2\mathrm{M}_{}\mathrm{yr}^1`$ at moderate redshifts ($`z12`$). Recent deep near-IR follow-up observations of the sub-mm sources detected with the SCUBA (Smail et al. 1998) have suggested that two faint sub-mm sources may also be EROs with $`\mathrm{I}\mathrm{K}>6`$ (Smail et al. 1999). Liu et al. have measured redshifts for several other EROs and most of these are also in the $`0.8<z<1.5`$ range and some have moderately strong emission lines (Liu et al. 2000).
Two important open issues concern the surface density of EROs and the relative contribution of different classes of objects as a function of both color and apparent magnitude. We continue to lack statistically large samples of EROs. The most recent systematic large survey, covering an area of 154 square arcminutes by Thompson et al. (1999), has yielded six objects with $`K19.0`$ and $`RK>6`$. To quantify the fraction of different classes of EROs as a function of colors and magnitude, larger samples are required. The ERO surface density inferred from the Thompson et al. survey is 0.04$`\pm 0.016`$ arcmin<sup>-2</sup> for $`K19.0`$ and $`RK>6`$. Depending on the K-band magnitude limit, $`\mathrm{R}\mathrm{K}`$ color and possibly environment, the reported surface densities of EROs range from $`0.010.7`$ arcmin<sup>-2</sup>, derived from several serendipitous surveys over small areas (Hu & Ridgway 1994; Cowie et al. 1994; Beckwith et al. 1998; Thompson et al. 1999). There have been suggestions that EROs tend to cluster, particularly, in regions around high redshift AGNs compared with blank fields. The large luminosities and (admittedly uncertain) space densities imply that these objects represent a singificant consitutent of the overall galaxy population and that their contribution to the overall rate of star formation is non-negligable (e.g. Liu et al. 2000).
In this paper, we present a sample of EROs discovered using NICMOS on HST while operating in the parallel mode. Our combined NICMOS/optical survey covers only 16 square arcminutes but it provides high spatial resolution and better signal-to-noise in the near-IR than most ground-based surveys.
## 2. Observations and Reductions
### 2.1. Near Infrared Images
The images under discussion here were obtained with camera 3 on NICMOS (Thompson et al. 1998) using the F160W ($`\lambda _c=1.6\mu `$m, $`\delta \lambda =0.4\mu `$m) and F110W band-passes ($`\lambda _c=1.1\mu `$m, $`\delta \lambda =0.8\mu `$m). The data were obtained in the parallel observing mode during the period from October 1997 to November 1998, excluding the dedicated NICMOS camera 3 observing campaigns. This allowed us to observe fields which are essentially randomly distributed over the sky. The sensitivity varies from field-to-field due to different integration times and background levels. The NICMOS internal pupil adjustment mirror was set near the end of its travel, providing the best possible focus for camera 3. The PSF in the images is slightly non-gaussian, but well characterized by FWHM$`=0.25^{^{\prime \prime }}`$. We obtained four images per orbit, two each with the F110W and F160W filters. The field offset mirror (FOM) was used to dither between two fixed positions $`1.8^{^{\prime \prime }}`$ apart in a direction aligned with one axis of the detector. In addition there were small inter-orbit dither moves executed for some of the pointings. The projected size of a camera 3 pixel is $`0.204^{\prime \prime }`$, giving a $`51{}_{}{}^{\prime \prime }\times 51^{^{\prime \prime }}`$ field of view for each image. A small area is lost in the construction of the final mosaic image. The NICMOS camera 3 parallel program covered approximately 200 square arcminutes at high galactic lattitude in its 14 months of operation. In the present work, we are limited by the amount and the depth of the visible-light data that we were able to collect for the NICMOS parallel fields. As described below, most of the NICMOS F160W images easily reach F160W magnitudes fainter than 21 (Vega magnitude) with high signal-to-noise ratios (eg. 20$`\sigma `$ at H = 21.0). Thus, to detect EROs at these depths requires 3$`\sigma `$ R magnitude limits of 26 in our follow-up images. As the NICMOS camera 3 field of view is only 0.722 square arcminutes, optical followup is quite inefficient.
We used McLeod’s (1997) NicRed v1.7 package to linearize and remove the cosmic rays from the MultiAccum images. Median images were derived from more than 50 pointings and these were used to remove the dark and sky signals. Even with the optimal dark subtraction, there remain considerable frame-to-frame variations in the quality of the final images. The individual linearized, dark corrected, flatfielded and cosmic ray cleaned images were shifted, masked and combined to produce final mosaic images. Before shifting, each image was $`2\times 2`$ block-replicated and integer ($`0.1^{^{\prime \prime }}`$) offsets were applied. In this way we avoided smoothing or interpolation of the data. The MultiAccum process is not 100% efficient in rejecting cosmic ray events and so we applied a $`3\sigma `$ rejection when assembling the final mosaics. More details regarding the NICMOS imaging reduction can be found in Yan et al. (1998).
In Table 1, we list $`3\sigma `$ surface brightness limits in the F160W filter for the fields where we have obtained optical data. The 50% completeness depth in the in the shallowest field is approximately m(F160W)$`=22`$. Thus our near-IR catalog is 100% complete at the ERO selection limit (F160W $`21.5`$ as described below). We emphasize that our ERO search is entirely limited by the depth of the optical images.
### 2.2. Optical Observations
Multicolor CCD observations of several of the NICMOS parallel fields were made at Cerro-Tololo, Las Campanas, WIYN, Palomar and the W. M. Keck observatories. The date of the observations, duration of the R images, the FWHM of the seeing and the 3$`\sigma `$ limiting surface brightness of the R images are listed in Table 2. The pixel scales range from $`0.2^{^{\prime \prime }}`$ for the LCO, WIYN, and Keck cameras to $`0.4^{^{\prime \prime }}`$ for the CTIO BTC camera system. Several of the fields were observed on more than one occasions and from different sites to obtain photometric calibration. The seeing was better than $`1^{^{\prime \prime }}`$ for most of the fields. The data were reduced and calibrated using standard methods. The resulting images were interpolated and rotated to match the NICMOS fields and true color images were constructed from the F160W, F110W, and R images (or F160W, R and V where possible). This allowed for easy and efficient identification of objects with red colors. The photometry was performed on the original uninterpolated images.
## 3. Results
### 3.1. ERO Detection and Photometry
All of the ERO identifications were made visually from comparisons between the aligned near-IR and optical images. The small size of each NICMOS image ($`51^{\prime \prime }\times 51^{\prime \prime }`$) as well as its sensitive depth made visual selection effecient. All of the EROs are well detected in the near-IR bands and were either undetected or marginally detected in the R-band images. ($`\mathrm{R}\mathrm{F160W}`$) colors were measured for all of the ERO candidates using the original uninterpolated images. The colors are subject to significant uncertainties due to the low signal levels of the red objects in the R-band images. The mismatch between the point spread functions of NICMOS and ground-based CCD images adds a significant complication to the color measurement. We adopted the following method to estimate the ($`\mathrm{R}\mathrm{F160W}`$) colors. We smoothed our F160W images to the same FWHM as the corresponding optical images. The F160W and R magnitudes were then measured in identical apertures. The aperture diameter is $`2.5\times `$ the measured FWHM in the R-band images. We used the Source Extractor software (Bertin & Arnouts 1996) to measure the magnitudes.
In Table 3 we also list the magnitudes of EROs in a variety of filters. These are isophotal magnitudes measured out to 1$`\sigma `$ isophotal radius. As our sample of EROs are well detected at F160W, isophotal magnitudes are a good measure of total magnitudes. For the EROs which are undetected in the optical images, their 3$`\sigma `$ magnitude limits within a 2.5FWHM diameter aperture are used.
### 3.2. Surface Density
In a total of 22 NICMOS fields, covering 16 square arcminutes, we detect 15 objects with $`\mathrm{R}\mathrm{F160W}>5`$. Table 3 lists their equatorial coordinates, near-IR and optical magnitudes, and optical-IR colors. For some fields, we also have K-band images taken with NIRC (Matthews & Soifer, 1994) on the Keck telescope (Teplitz et al. 2000). The ERO sample is selected with $`\mathrm{R}\mathrm{F160W}>5`$ without any limits set on H magnitude. Our color selection system is similar to $`\mathrm{R}\mathrm{H}`$ used in ground-based obsevations, and so direct comparisons are possible.
Figure 1 shows the R and F160W band images of each ERO. All of our EROs are resolved in the NIC3 images and appear to be extended, except ERO 1940-6915A. This object has m(F160W) = 20.6. and appears to be unresolved in our NIC3 image. The reamining objects exhibit both elliptical and indeterminate morphologies. The detailed analyses of the ERO luminosity profiles using both the NICMOS camera 2 and camera 3 data will be presented in Yan & McCarthy (2000).
In a few fields where we have very deep optical images taken with LRIS (Oke et al. 1995) at the Keck 10 m telescope, we were able to find EROs with H magnitude as faint as 21.4. The surface density of EROs depend on the magnitude limit of the sample. Including the two apparent clusters of EROs and the point source like object ERO 1940-6915A, we estimate surface densities of 0.19$`\pm 0.11`$ arcmin.<sup>-2</sup>, 0.19$`\pm 0.11`$ arcmin.<sup>-2</sup>, 0.50$`\pm 0.18`$ arcmin.<sup>-2</sup>, and 0.94$`\pm 0.24`$ for EROs with $`\mathrm{R}\mathrm{F160W}>5`$ and $`\mathrm{F160W}<19.5`$, $`\mathrm{F160W}<20.0`$, $`\mathrm{F160W}<20.5`$, and $`<21.5`$ respectively. If we exclude the two ERO clusters and the possible stellar object, the inferred surface density is significantly lower, 0.13$`\pm 0.08`$ arcmin.<sup>-2</sup>, 0.13$`\pm 0.08`$ arcmin.<sup>-2</sup>, 0.25$`\pm 0.13`$ arcmin.<sup>-2</sup> and 0.38$`\pm 0.15`$ arcmin.<sup>-2</sup> for EROs brighter than 19.5, 20, 20.5 and 21.5 respectively. Here we did not include 2344$``$1525C and 2344$``$1525F whose R $``$ H colors are slightly less than 5. Since our areal coverage is very small, the statistical errors in these estimates are large, particularly at bright magnitudes. If EROs have F160W - K colors of $`0.51`$ as suggested by our observations, our surface density estimate is consistent within 3$`\sigma `$ with that inferred from the ground-based ERO survey by Thompson et al. (1999) of 0.04$`\pm 0.016`$ arcmin.<sup>-2</sup> for R $``$ K $`>`$ 6 and K $`<`$ 19.
### 3.3. Infrared Colors
In Figure 2, we present F110W - F160W colors for all of the galaxies detected in 19 NICMOS fields, where we have both F160W and F110W images as well as ground-based optical follow-up observations (see Table 1). The solid dots in the figure represent the EROs listed in Table 3. Overlaying on top of the data are color-magnitude tracks for a passively evolving L elliptical galaxy formed at $`z=10`$ in a single burst lasting 1 Gyr and a star-forming galaxy with a constant star formation rate of 1 M/yr respectively. Here we adopt H$`{}_{0}{}^{}=70\mathrm{kms}^1\mathrm{Mpc}^1`$ and q$`{}_{0}{}^{}=0.1`$. The open diamonds indicate the places with redshifts of 0.5, 1, 2, and 3. The error bars on F110W - F160W colors and F160W magnitudes are $`\pm 1\sigma `$, calculated from Sextractor (Bertin & Arnouts 1996). The F110W $``$ F160W colors of our R - F160W selected EROs appear to be redder than the average color of field galaxies (Figure 2). They also appear to be redder than the $`z0.9`$ cluster elliptical population studied by Stanford et al. (1998) (J $``$ H $`1.1`$). There are some objects which have red J $``$ H colors, but R $``$ H colors that do not meet our ERO definition. These objects may contain small amounts of current star formation. If the objects in our ERO sample have luminosities near L, the median redshift of our sample is $`1`$, and the total range sampled is roughly $`0.6<z<1.5`$.
Among the 22 NICMOS pointings for which we have deep optical photometry, we found two fields that each contains more than 3 objects meeting our ERO definition. Plate 1 and 2 show the VR F160W composite true color pictures of the ERO groups in 2344$``$1524 and 1631+3001. In the 2344$``$1524 field, we also obtained K band images and the selected ERO candidates have $`\mathrm{R}\mathrm{K}\mathrm{}>6`$. The 1631$`+`$3001 field has many red objects, including several very faint objects with $`\mathrm{F}(160\mathrm{W})>21.5`$, which are not selected in our ERO sample. Plates 1 and 2 clearly indicate two potential groups or clusters of EROs. Follow-up spectroscopy would be important for determining the cluster redshifts and the physical nature of these EROs. Our detections of two clusters of EROs over a small area of $`16`$ sqaure arcminutes suggest that EROs tend to be strongly clustered. Previous studies by McCarthy et al. (1992) and Thompson et al (1999) have also noted the ERO clustering phenomenon from small statistical samples. Obviously, the problem of ERO clustering will be an important goal for future deep infrared surveys covering large area.
## 4. Acknowledgments
We thank the staff of the Space Telescope Science Institute for their efforts in making this parallel program possible. In particular we thank Duccio Machetto, Peg Stanley, Doug van Orsow, and the staff of the PRESTO division. We also thank John Mackenty and members of the STScI NICMOS group for crafting the exposure sequences. Some of the data presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation. This research was supported, in part, by grants from the Space Telescope Science Institute, GO-7499.01-96A, AR-07972.01-96A and PO423101. HIT acknowledges funding by the Space Telescope Imaging Spectrograph Instrument Definition Team through the National Optical Astronomy Observatories and by the NASA Goddard Space Flight Center.
Plate 1. — This is a VR F160W composite color image of 2343$``$1524 field. The field shown here has a size of 50$`{}_{}{}^{^{\prime \prime }}\times `$50$`^{^{\prime \prime }}`$. The color image indicates a cluster of EROs.
Plate 2. — A VR F160W composite color picture of 1631+3001 field shows a cluster of EROs. The field size is the same as in Plate 1.
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# Pertinent Dirac structure for QCD sum rules of meson-baryon coupling constants
## I INTRODUCTION
In QCD sum rule approaches , the two-point correlation function with a pion
$`i{\displaystyle d^4xe^{iqx}0|\mathrm{T}[J_N(x)\overline{J}_N(0)]|\pi (p)}`$ (1)
is often used to calculate the $`\pi NN`$ coupling by facilitating a general external field method developed in Ref. . This correlation function contains three distinct Dirac structures (1) $`i\gamma _5`$ (PS), (2) $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ (T), and (3) $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ (PV), each of which can in principle be used to calculate the coupling. Currently, there is an issue of the Dirac structure dependence of the sum rule results . In calculating the coupling, one can construct either the PS sum rules beyond the chiral limit or the T sum rules . Both sum rules yield the $`\pi NN`$ coupling close to its empirical value. On the other hand, the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ sum rules contain large contributions from the continuum, which therefore do not provide reliable results.
The PS and T sum rules have been extended to calculate the meson-baryon couplings $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$ and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ by considering the two-point correlation function with a meson,
$`i{\displaystyle d^4xe^{iqx}0|\mathrm{T}[J_B(x)\overline{J}_B(0)]|(p)}.`$ (2)
Calculation of the couplings from this correlation function is somewhat limited due to the ignorance of meson wave functions when heavier mesons are involved. In the SU(3) limit however, this correlation function can be used to determine the so-called $`F/D`$ ratio unambiguously because in this limit the OPE can be exactly classified according to SU(3) relations for the couplings . The $`F/D`$ ratio is an important input in making realistic potential models for hyperon-baryon interactions as well as in analyzing the hyperon semileptonic data.
At present, there is a clear Dirac structure dependence in the calculation of the $`F/D`$ ratio using Eq. (2). In particular, we have reported from the PS sum rules $`F/D0.2`$ while from T sum rules $`F/D0.78`$ . Thus, even though the two sum rules with different Dirac structures were successful in reproducing the empirical $`\pi NN`$ coupling, their prediction for the $`F/D`$ ratio is quite different.
To resolve this issue, additional criteria to choose a proper Dirac structure are needed for reliable predictions on the $`F/D`$ ratio as well as the meson-baryon couplings. For this purpose, we first note that in Ref. the Ioffe current or its SU(3) rotated version has been used to construct sum rules Eq. (2). The Ioffe current however is a specific choice for the nucleon current among infinitely many possibilities. The Ioffe current is often used for the nucleon because it gets large contributions from the chiral breaking parameter $`\overline{q}q`$. In addition, direct instantons are believed to play less roles in this current.
Nevertheless, it may be useful to study the dependence of the sum rule results on general baryon currents. Depending on the currents, it is expected that the overlap $`\lambda _B`$ between the physical baryon state and the current may be altered but ideally the physical parameters such as meson-baryon couplings remain unchanged. Indeed, from the correlation function Eq. (2), what will actually be determined is the overlap strength multiplied by the coupling of concern. In the SU(3) symmetric limit, all the strengths depend only on the currents. They are determined from the corresponding baryon mass sum rules and all the baryon masses are the same in the SU(3) limit. Thus, in this limit, the dependence on the currents should be driven by the common overlap strength, which in return provides the coupling independent of the currents. This ideal aspect will be pursued in this work as a criterion for choosing a proper Dirac structure.
An alternative way is to calculate baryon axial charges and convert them into meson-baryon couplings using the Goldberger-Treiman relation. Ref. considered the nucleon correlation function in external axial vector field and constructed a sum rule for $`g_A1`$ using one specific Dirac structure. Recently, a new approach was proposed in Ref. where the axial vector correlation function in a one-nucleon state is considered. Both obtained an excellent agreement for $`g_A`$ of the nucleon.
This paper is organized as follows. In Section II, we construct meson-baryon coupling sum rules using general baryon currents. A brief discussion on the OPE based on chirality is given in Section III. We then briefly check in Section IV whether the discussion on the continuum threshold is still valid when the general baryon currents are used in the sum rules. In Section V, the dependence of the OPE on the baryon currents is studied. We study in the SU(3) limit whether or not the dependence on the currents are mostly contained in the overlap $`\lambda _B`$. This constraint gives us a new criterion to choose an appropriate Dirac structure. In Section VI, we calculate the couplings in the SU(3) limit from the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure. The $`F/D`$ ratio is identified in terms of the OPE. Conclusions are given in Section VII.
## II CONSTRUCTION of the QCD SUM RULES
We use the two-point correlation function with a meson,
$`i{\displaystyle d^4xe^{iqx}0|\mathrm{T}[J_B(x)\overline{J}_B(0)]|(p)}`$ (3)
where $`J_B`$ is the baryon current of concern and $`p`$ is the momentum of meson $``$. Meson states $`\pi `$ and $`\eta `$, and baryon currents for the proton, $`\mathrm{\Xi }`$ and $`\mathrm{\Sigma }`$ will be considered in this work.
The proton current is constructed from two $`u`$-quarks and one $`d`$-quark by assuming that all three quarks are in the s-wave state. In the construction of the current, one up and one down quark are combined into an isoscalar diquark. The other up quark is attached to the diquark so that quantum numbers of the proton are carried by the attached up quark. In this method, there are two possible combinations for the current. The general proton current is a linear combination of the two possibilities mediated by a real parameter $`t`$,
$`J_p(x;t)`$ $`=`$ $`2ϵ_{abc}[(u_a^T(x)Cd_b(x))\gamma _5u_c(x)+t(u_a^T(x)C\gamma _5d_b(x))u_c(x)].`$ (4)
Here, $`a,b,c`$ are color indices, $`T`$ denotes the transpose with respect to the Dirac indices, and $`C`$ the charge conjugation. The choice $`t=1`$ is called the Ioffe current . The currents for $`\mathrm{\Xi }`$ and $`\mathrm{\Sigma }`$ are obtained from the proton current via SU(3) rotations ,
$`J_\mathrm{\Xi }(x;t)`$ $`=`$ $`2ϵ_{abc}[(s_a^T(x)Cu_b(x))\gamma _5s_c(x)+t(s_a^T(x)C\gamma _5u_b(x))s_c(x)],`$ (5)
$`J_\mathrm{\Sigma }(x;t)`$ $`=`$ $`2ϵ_{abc}[(u_a^T(x)Cs_b(x))\gamma _5u_c(x)+t(u_a^T(x)C\gamma _5s_b(x))u_c(x)].`$ (6)
When going beyond the soft-meson limit, one can consider three distinct Dirac structures in correlation function in constructing sum rules: $`i\gamma _5`$ (PS), $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ (T) and $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ (PV). For the $`i\gamma _5`$ structure, the sum rules are constructed at the order $`p^2=m_\pi ^2`$ . At this order, the terms linear in quark mass $`(m_q)`$ in the OPE should be included because $`m_q`$ is the same chiral order with $`m_\pi ^2`$ via the Gell-Mann–Oakes–Renner relation,
$`2m_q\overline{q}q=m_\pi ^2f_\pi ^2.`$ (7)
On the other hand, for the T and PV structures, we construct the sum rules at the order $`𝒪(p)`$. At this order, the $`m_q`$ terms should not be included in the OPE. Technical details on the OPE calculation can be found in Refs. .
In constructing the phenomenological side, we first define $`\lambda _B(t)`$, the coupling strength between the baryon current $`J_B(x;t)`$ and the physical baryon field $`\psi _B(x)`$. Using the pseudoscalar type interaction between the meson and baryons $`g_B\overline{\psi }_Bi\gamma _5\psi _B`$, we obtain the phenomenological side of the correlation function:
$`i\gamma _5\text{structure}`$ $`\text{at the order}𝒪(p^2)i\gamma _5p^2{\displaystyle \frac{g_B\lambda _B^2(t)}{(q^2m_B^2)^2}}+\mathrm{},`$ (8)
$`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu \text{structure}`$ $`\text{at the order}𝒪(p)\gamma _5\sigma _{\mu \nu }q^\mu p^\nu {\displaystyle \frac{g_B\lambda _B^2(t)}{(q^2m_B^2)^2}}+\mathrm{},`$ (9)
$`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}\text{structure}`$ $`\text{at the order}𝒪(p)i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}{\displaystyle \frac{g_B\lambda _B^2(t)m_B}{(q^2m_B^2)^2}}+\mathrm{}.`$ (14)
The ellipsis denotes contributions from higher resonances as well as a single pole associated with transitions from the ground state to higher resonances. The continuum contributions come from transitions among higher resonances, whose spectral densities are modeled with a step function starting at the threshold $`S_0`$. Matching the OPE side with the phenomenological side and taking Borel transformation Note, we use a single dispersion relation as advocated in Ref. ., we get the sum rules of the form
$`g_B\lambda _B^2(t)\left[1+A_B(t)M^2\right]`$ $`=`$ $`e^{m_B^2/M^2}F_\text{B}^{\text{OPE}}(M^2;t)f_\text{B}^{\text{OPE}}(M^2;t)`$ (15)
where the single pole term in the phenomenological side has been denoted by $`A_B`$. Expressions for the OPE $`F_\text{B}^{\text{OPE}}(M^2;t)`$ are given in the Appendix A.
## III Chirality consideration
The OPEs given in the Appendix A have an interesting feature to discuss when $`t=1`$. Specifically, in the $`i\gamma _5`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\mu `$ sum rules, Wilson coefficients of chiral-odd operators $`\overline{q}q`$, $`f_{3\pi }`$, $`\overline{q}q\frac{\alpha _s}{\pi }𝒢^2`$, and $`m_0^2\overline{q}q`$ are all zero when $`t=1`$. Also in the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ sum rules, contributions from the chiral-even operators $`\overline{q}q^2`$ and $`m_0^2\overline{q}q^2`$ are zero. To understand this feature, it is useful to decompose the correlator according to chirality of the current,
$`J_B\overline{J}_B=J_B^R\overline{J}_B^R+J_B^R\overline{J}_B^L+J_B^L\overline{J}_B^R+J_B^L\overline{J}_B^L.`$ (16)
$`J_B^L(J_B^R)`$ denotes the left-handed (right-handed) component of the current $`J_B`$. On the other hand, Eq. (3) can be written
$`i{\displaystyle d^4xe^{iqx}0|\mathrm{T}[J_B(x)\overline{J}_B(0)]|(p)}`$ $`=`$ $`i\gamma _5\mathrm{\Pi }_{\mathrm{ps}}+i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}\mathrm{\Pi }_{\mathrm{pv}}`$ (19)
$`+`$ $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu \mathrm{\Pi }_\mathrm{T}.`$ (20)
Thus, it is easy to see that the $`i\gamma _5`$ and $`\gamma _5\sigma _{\mu \nu }`$ structures have nonzero contributions only from the chiral mixing term $`J^R\overline{J}^L+J^L\overline{J}^R`$, while the chirality conserving term $`J^R\overline{J}^R+J^L\overline{J}^L`$ contributes only to the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structure.
Now let us classify QCD operators contributing to each Dirac structure. To do that, we suppress for simplicity the color indices and write baryon current as
$$J(q^TCq)\gamma _5q+t(q^TC\gamma _5q)q.$$
(21)
Here $`q=u,d,s`$. When $`t=1`$, it is straightforward to show that
$`J_R`$ $``$ $`2(q_R^TCq_R)q_R,`$ (22)
$`J_L`$ $``$ $`2(q_L^TCq_L)q_L.`$ (23)
Thus, at this specific $`t`$, chirality of all quarks are the same as that of the baryon.
In the $`i\gamma _5`$ or $`\gamma _5\sigma _{\mu \nu }`$ sum rules, we need to consider the products $`J^R\overline{J}^L`$ and $`J^L\overline{J}^R`$. In making such products using Eqs. (22) (23), all three quark propagators should break the chirality when they move from the coordinate $`0`$ to $`x`$. Hence, it is easy to see that, among chiral-odd operators, terms such as
$$m_q\overline{q}q^2,\overline{q}q^3,m_q^2\overline{q}q,\mathrm{}$$
(24)
can contribute to the $`i\gamma _5`$ or $`\gamma _5\sigma _{\mu \nu }`$ correlator, while other chiral-odd operators such as $`\overline{q}q`$, $`f_{3\pi }`$, $`\overline{q}q\frac{\alpha _s}{\pi }𝒢^2`$, $`m_0^2\overline{q}q`$ cannot.
On the other hand, in the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structure, the product $`J^L\overline{J}^L`$ or $`J^R\overline{J}^R`$ contributes to the sum rule. Among chiral-even operators, an operator such as $`\overline{q}q^2`$ cannot be formed in the product $`J^L\overline{J}^L`$ or $`J^R\overline{J}^R`$ simply because two quarks with the same chirality cannot be combined into the quark-antiquark pair. Similarly, $`m_0^2\overline{q}q^2`$ can not be formed. This explains the disappearance of such terms in the OPE when $`t=1`$.
## IV Criterion I : sensitivity to the continuum threshold
We now analyze sum rules of the three different Dirac structures with the general baryon currents, Eqs (4) and (6). As pointed out in Refs. , sum rule results from the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structure are sensitive to the continuum threshold $`S_0`$ and therefore this structure is not reliable. On the other hand, $`i\gamma _5`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structures are insensitive to $`S_0`$. The chirality consideration suggested in Ref. implies that in the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ sum rules the large slope and the strong sensitivity to $`S_0`$ of the Borel curves can be explained if higher resonances with different parities add up. With this scenario, the higher resonances contributions cancel each other in the $`i\gamma _5`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules therefore explaining the weak sensitivity to $`S_0`$ and the small slope of the Borel curves.
Since only the Ioffe current is used in the analysis of Refs. , let us briefly check if this scenario still works when the general baryon currents are used. As the scenario does not rely on the specific form for the current, what has been claimed in Refs. must be valid even with the general baryon currents. To see this, we plot the RHS of Eq. (15) for the $`\pi NN`$ coupling from $`i\gamma _5`$, $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ and $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structures in Figs. 123, respectively. To show the dependence on $`t`$, we plot the curves for $`t=1.5,1.5`$ as well as $`t=1.0`$ (the Ioffe current). In these plots, we use the standard QCD parameters,
$`\overline{q}q`$ $`=`$ $`(0.23\mathrm{GeV})^3;{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2=(0.33\mathrm{GeV})^4,`$ (25)
$`\delta ^2`$ $`=`$ $`0.2\mathrm{GeV}^2;m_0^2=0.8\mathrm{GeV}^2.`$ (26)
For each $`t`$, the thick lines are for the continuum threshold $`S_0=2.07\mathrm{GeV}^2`$ corresponding to the Roper resonance, while the thin lines for $`S_0=2.57\mathrm{GeV}^2`$. The trend observed here is the same for the other couplings.
In Fig. 3, we observe that the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structure is sensitive to the continuum threshold even when the general current is used. The difference by changing the continuum threshold is $`15\%`$ at $`M^2=1\mathrm{GeV}^2`$. Note also that the slope is relatively large in this case. Since the coupling is determined from the intersection of the best fitting curve with the vertical line at $`M^2=0`$ (see Eq. (15)), the $`15\%`$ change at $`M^2=1\mathrm{GeV}^2`$, when it combined with the large slope, produces huge change in the extracted coupling. In contrast, from Figs. 1 and 2, the $`i\gamma _5`$ and $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structures are insensitive to $`S_0`$. Also the slopes of the curves are small. This observation is practically independent of the parameter $`t`$. At $`M^2=1\mathrm{GeV}^2`$, the difference is only $`23\%`$ level. Thus, the analysis in Refs. is still valid and the sum rule results from $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structure should be discarded under this consideration.
## V Criterion II : the dependence of the OPE on baryon currents
Using the sum rules derived in Section II, we discuss the dependence of the OPE on the baryon current (i.e. the dependence on $`t`$). For a given $`t`$, we linearly fit the RHS of Eq. (15)
$`g_B\lambda _B^2(t)\left[1+A_B(t)M^2\right]=f_\text{B}^{\text{OPE}}(M^2;t),`$
and determine $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$. Because $`f_\text{B}^{\text{OPE}}`$ is quadratic in $`t`$, $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ is also quadratic. Ideally, the physical parameter $`g_B`$ should be independent of $`t`$ if the sum rules are reliable. In other words, $`t`$ is just a parameter for the current. By changing $`t`$, only the coupling strength $`\lambda _B^2(t)`$ is expected to be affected, but not the physical parameter. This is a constraint to be satisfied when the sum rules are “good.”
To proceed, we take the SU(3) symmetric limit. Then, the strength $`\lambda _B(t)`$ should be independent of the baryons,
$`\lambda _N(t)=\lambda _\mathrm{\Xi }(t)=\lambda _\mathrm{\Sigma }(t),`$ (27)
as the baryon mass sum rules are the same in the limit. Furthermore, we have
$`\overline{s}s=\overline{q}q`$ $`;m_\eta =m_\pi ,`$ (28)
$`f_\eta =f_\pi `$ $`;f_{3\eta }=f_{3\pi },`$ (29)
$`m_N=m_\mathrm{\Xi }=m_\mathrm{\Sigma }`$ $`;m_s=m_q.`$ (30)
This SU(3) limit is particularly interesting when we select a suitable Dirac structure. Suppose we plot $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ in terms of $`t`$. If the sum rules are “good”, all $`g_B`$ should be just constants, independent of $`t`$. The functional behavior is driven only by the strength $`\lambda _B^2(t)`$. The baryon mass sum rules in the SU(3) limit constrain that all $`\lambda _B^2(t)`$ are the same irrespective of the baryons. Therefore, “good” sum rules must give $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ which are proportional to each other.
Our constraint should be satisfied when the OPE are exact. But in practice, the full OPE terms $`f_{\mathrm{full}}^{\mathrm{OPE}}`$ are separated into two groups,
$`f_{\mathrm{full}}^{\mathrm{OPE}}=f_{\mathrm{calc}}^{\mathrm{OPE}}+f_{\mathrm{rest}}^{\mathrm{OPE}},`$ (31)
where $`f_{\mathrm{calc}}^{\mathrm{OPE}}`$ denotes the calculable OPE, and $`f_{\mathrm{rest}}^{\mathrm{OPE}}`$ denotes the rest of the full OPE. In this notation, the reliability of sum rule simply means
$`f_{\mathrm{calc}}^{\mathrm{OPE}}f_{\mathrm{rest}}^{\mathrm{OPE}}.`$ (32)
The sum rules are “unreliable” if
$`f_{\mathrm{calc}}^{\mathrm{OPE}}f_{\mathrm{rest}}^{\mathrm{OPE}}.`$ (33)
In the former case, we expect that $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ derived from $`f_{\text{B;}\text{ calc}}^{\text{OPE}}`$ is almost the same as those from $`f_{\text{B;}\text{ full}}^{\text{OPE}}`$ in most region of $`t`$. On the other hand, in the latter case, $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ derived from $`f_{\text{B;}\text{ calc}}^{\text{OPE}}`$ may be quite different from those obtained from $`f_{\text{B;}\text{ full}}^{\text{OPE}}`$, and our ideal constraint may not be satisfied in most $`t`$.
Therefore, the ideal constraint can be used as a new criterion for choosing reliable sum rules. In order to apply this constraint to our sum rules, we again use the standard QCD parameters Eq. (26) and linearly fit $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ at each $`t`$. In the fitting, the continuum threshold is set to $`S_0=2.07\mathrm{GeV}^2`$, corresponding to the Roper resonance, and the Borel window is taken $`0.65M^21.24\mathrm{GeV}^2`$ as in Refs. . In this Borel window, (1) the Borel curve for each coupling is almost linear (see Figs. 12), (2) the contribution from the highest dimensional OPE term is typically $`515\%`$ in the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules and $`20\%`$ level in the $`i\gamma _5`$ sum rules, and (3) the continuum contribution is less than $`20\%`$ in both structures. It should be noted that because all the couplings are related under SU(3) rotations, we need to take a common Borel window .
Fig. 4 shows $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ as a function of $`t`$ for the $`i\gamma _5`$ sum rules. The $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ cases are shown in Fig. 5. Interesting features in the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ cases are that (1) all the curves are zero when $`t=1`$ and almost zero at $`t0.5`$, (2) each extremum of the curves coincides around $`t0.3`$. Under the chirality consideration given in Section III, we can easily understand why $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ is zero when $`t=1`$. From the figure, though not exact, one observes that the curves can be almost overlapped when multiplied by appropriate constants. For example, let us compare the $`\pi NN`$ and $`\eta NN`$ curves. When they are positive, the $`\pi NN`$ curve lies above the $`\eta NN`$ curve. When they are negative, the situation is reversed. This behavior of the $`\eta NN`$ curve can be reproduced by multiplying an appropriate constant to the $`\pi NN`$ curve. Of course, this claim can not be made when $`t0.5`$ because one curve becomes zero while the other does not. Therefore, except around $`t0.5`$, the Borel curves satisfy the ideal constraint in most region of $`t`$. Such a trend can not be observed from the $`i\gamma _5`$ sum rules. (see Fig. 4.) Therefore, we claim that the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules are more appropriate.
To support our claim that the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules are more suitable than those from the $`i\gamma _5`$ structure, one more check to do is to see the $`t`$-dependence of $`\lambda _B^2(t)`$ from baryon mass sum rules. In Fig. 6, $`\lambda _B^2(t)`$ in the SU(3) limit is plotted using chiral-odd nucleon mass sum rule Definition of the function $`E_n(xS_0/M^2)`$ is given in the Appendix A. The Wilson coefficient of the dimension 7 OPE is different from Ref. . When $`t=1`$ (the Ioffe current), however, our Wilson coefficient reduces to that of Ref. . Nevertheless, the dimension 7 condensate contributes to the sum rule only slightly. Thus, this discrepancy is marginal.:
$`m_B\lambda _B^2(t)e^{m_B^2/M^2}`$ $`=`$ $`{\displaystyle \frac{4}{(2\pi )^4}}[{\displaystyle \frac{\pi ^2}{4}}(52t+7t^2)\overline{q}qM^4E_1(x)`$ (36)
$`+{\displaystyle \frac{3}{4}}\pi ^2(1t^2)m_0^2\overline{q}qM^2E_0(x)`$
$`+{\displaystyle \frac{\pi ^4}{24}}(7+2t+5t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2],`$
where $`m_q`$ order terms are neglected and the SU(3) relations are used: $`m_N=m_\mathrm{\Xi }=m_\mathrm{\Sigma }m_B`$; $`\lambda _N=\lambda _\mathrm{\Xi }=\lambda _\mathrm{\Sigma }\lambda _B`$; $`\overline{u}u=\overline{d}d=\overline{s}s\overline{q}q`$. Comparing with Fig. 5, we confirm that the $`t`$-dependence of $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ from the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules can be reproduced from the $`t`$-dependence of $`\lambda _B^2(t)`$.
In the region $`0.5\begin{array}{c}<\hfill \\ \hfill \end{array}t\begin{array}{c}<\hfill \\ \hfill \end{array}1`$ in Fig. 6, $`\lambda _B^2(t)`$ is negative, thus not physical. In this region, the sum rules should definitely fail and a reliable prediction for a physical parameter may not be possible. At $`t0.5`$ or $`t1`$, of course, the OPE is almost zero suggesting that there are cancellations among OPE terms, i.e. the correlation function can not be well saturated by the calculated OPE. Therefore, the optimal current should be chosen away from these points.
## VI The $`F/D`$ ratio from the pseudotensor sum rules
In this section, we analyze the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules to determine the $`F/D`$ ratio. In particular, we investigate the $`t`$-dependence of the ratio using the general interpolating fields for the baryons. As already mentioned, mesons and baryons are classified according to SU(3) symmetry, which provides simple relations for the meson-baryon couplings in terms of the two parameters
$$g_{\pi N}\mathrm{and}\alpha =\frac{F}{F+D}.$$
(37)
That is,
$`g_{\eta N}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}(4\alpha 1)g_{\pi N};g_{\pi \mathrm{\Xi }}=(2\alpha 1)g_{\pi N},`$ (38)
$`g_{\eta \mathrm{\Xi }}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}(1+2\alpha )g_{\pi N};g_{\pi \mathrm{\Sigma }}=2\alpha g_{\pi N},`$ (39)
$`g_{\eta \mathrm{\Sigma }}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}}}(1\alpha )g_{\pi N}.`$ (40)
To see how these relations are reflected in the OPE of the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules (see the Appendix A for the OPE), we take the SU(3) symmetric limit to organize them in terms of two terms $`𝒪_1`$ and $`𝒪_2`$ defined as
$`𝒪_1e^{m_B^2/M^2}`$ $``$ $`{\displaystyle \frac{1}{96\pi ^2f_\pi }}(2+4t2t^2)\overline{q}qM^4E_0(x){\displaystyle \frac{f_\pi }{3}}(1+t^2)\overline{q}qM^2`$ (43)
$`{\displaystyle \frac{1}{54}}f_\pi \delta ^27(1+t^2)\overline{q}q+{\displaystyle \frac{1}{7212f_\pi }}(1+2tt^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\pi }{72}}7(1+t^2)m_0^2\overline{q}q,`$
$`𝒪_2e^{m_B^2/M^2}`$ $``$ $`{\displaystyle \frac{1}{96\pi ^2f_\pi }}12(1t^2)\overline{q}qM^4E_0(x)+{\displaystyle \frac{2f_\pi }{3}}(tt^2)\overline{q}qM^2`$ (46)
$`{\displaystyle \frac{1}{27}}f_\pi \delta ^2(313t+10t^2)\overline{q}q+{\displaystyle \frac{1}{48f_\pi }}(1t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\pi }{36}}(13t+2t^2)m_0^2\overline{q}q.`$
Specifically, we have
$`g_{\pi N}\lambda _N^2(1+A_{\pi N}M^2)=`$ $`𝒪_1`$ $`+𝒪_2,`$ (47)
$`\sqrt{3}g_{\eta N}\lambda _N^2(1+A_{\eta N}M^2)=`$ $`𝒪_1`$ $`+𝒪_2,`$ (48)
$`g_{\pi \mathrm{\Xi }}\lambda _N^2(1+A_{\pi \mathrm{\Xi }}M^2)=`$ $`𝒪_1`$ $`,`$ (49)
$`\sqrt{3}g_{\eta \mathrm{\Xi }}\lambda _N^2(1+A_{\eta \mathrm{\Xi }}M^2)=`$ $`𝒪_1`$ $`2𝒪_2,`$ (50)
$`g_{\pi \mathrm{\Sigma }}\lambda _N^2(1+A_{\pi \mathrm{\Sigma }}M^2)=`$ $`𝒪_2,`$ (51)
$`\sqrt{3}g_{\eta \mathrm{\Sigma }}\lambda _N^2(1+A_{\eta \mathrm{\Sigma }}M^2)=`$ $`2𝒪_1`$ $`+𝒪_2.`$ (52)
Note that another SU(3) relation $`\lambda _N=\lambda _\mathrm{\Xi }=\lambda _\mathrm{\Sigma }`$ has been used in writing these equations. Neglecting the unknown single pole term $`A_B`$, we identify the $`F/D`$ ratio in terms of the OPE,
$`2\alpha {\displaystyle \frac{𝒪_2}{𝒪_1+𝒪_2}}F/D{\displaystyle \frac{𝒪_2}{2𝒪_1+𝒪_2}}.`$ (53)
This is an obvious consequence of using the baryon currents constructed according to the SU(3) symmetry. Hence, it provides the consistency of our sum rules with the SU(3) relations for the couplings.
To determine the $`F/D`$ ratio, however, the unknown single pole term $`A_B`$ should be taken into account. For that purpose, we linearly fit the RHS of Eq. (52) and determine $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ for a given $`t`$. Once two of $`\left[g_B\lambda _B^2(t)\right]_{\text{fitted}}`$ are determined, their ratio can be converted to yield the $`F/D`$ ratio according to Eq. (40).
In Fig. 7, the $`F/D`$ ratio is plotted as a function of $`\mathrm{cos}\theta `$. Here, to investigate the whole range of $`\mathrm{}t+\mathrm{}`$, we introduce a new parameter $`\theta `$ defined as
$`\mathrm{tan}\theta =t.`$ (54)
Thus, the range $`0t+\mathrm{}`$ corresponds to $`0\theta \pi /2`$ while the range $`\mathrm{}t0`$ spans $`\pi /2\theta \pi `$. In Fig. 7, circles are obtained from the Borel window $`0.65M^21.24\mathrm{GeV}^2`$ with the continuum threshold $`S_0=2.07\mathrm{GeV}^2`$. To see the sensitivity to this choice, we also calculate the ratio using (1) $`0.65M^21.24\mathrm{GeV}^2`$, $`S_0=2.57\mathrm{GeV}^2`$ (triangles), (2) $`0.90M^21.50\mathrm{GeV}^2`$, $`S_0=2.07\mathrm{GeV}^2`$ (squares). We see that the $`F/D`$ ratio is insensitive to the continuum threshold, agreeing with the discussion in Section IV. Also, the calculated $`F/D`$ ratio is relatively insensitive to the choice of the Borel window. The peak around $`t0.5(\mathrm{cos}\theta 0.9)`$ can be understood from Fig. 5. Most curves are zero around this $`t`$ but not simultaneously. The $`F/D`$ ratio is basically obtained by taking a ratio of any two curves but the ratio of the two curves around $`t0.5(\mathrm{cos}\theta 0.9)`$ is not well-behaved. On the other hand, at $`t=1(\mathrm{cos}\theta =1/\sqrt{2})`$, the $`F/D`$ ratio does not diverge because all curves for the couplings in Fig. 5 go to zero linearly in $`(t1)`$.
The strong sensitivity of the $`F/D`$ ratio to $`t`$ within the region $`0.5\begin{array}{c}<\hfill \\ \hfill \end{array}t\begin{array}{c}<\hfill \\ \hfill \end{array}1(\mathrm{cos}\theta \begin{array}{c}<\hfill \\ \hfill \end{array}0.9,\text{or}0.7\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{cos}\theta )`$ is unrealistic because first of all, absolute total value of the OPE in each coupling is very small in this region. The convergence of the OPE may not be sufficient enough. Secondly, the strength $`\lambda _N^2`$ as can be seen from Fig. 6 is negative, thus not physical. Therefore, a reasonable value for the $`F/D`$ ratio should be obtained away from this region. We moderately take the realistic region as (1) $`t\begin{array}{c}<\hfill \\ \hfill \end{array}0.8(0.78\begin{array}{c}<\hfill \\ \hfill \end{array}\mathrm{cos}\theta )`$ and (2) $`1.3\begin{array}{c}<\hfill \\ \hfill \end{array}t(\mathrm{cos}\theta \begin{array}{c}<\hfill \\ \hfill \end{array}0.61)`$. The former constraint gives us the maximum value of $`F/D0.84`$, and the latter constraint gives us the minimum value of $`F/D0.63`$. Therefore, we conclude $`F/D0.60.8`$. This range includes the value from the SU(6) quark model ($`F/D=2/3`$), and is slightly higher than that extracted from semi-leptonic decay rates of hyperons ($`F/D0.57`$. It is often argued that the choice of $`t=1`$ (the Ioffe current) is optimal because the instanton effect and the continuum contribution is small, and the chiral breaking effects are maximized. If we choose $`t1`$, our estimate becomes $`F/D0.760.81`$, that is somewhat larger than the SU(6) value.
As a comparison, let us briefly consider the $`i\gamma _5`$ structure case. In this case too, we can classify the OPE of the Appendix A 1 according to Eq. (40) and identify the terms responsible for the $`F/D`$ ratio. By taking similar steps as T sum rules, we determine the $`F/D`$ ratio. Fig. 8 shows the $`F/D`$ ratio as a function of $`\mathrm{cos}\theta `$. Compared with Fig. 7, the $`F/D`$ ratio is very sensitive to $`t`$. As discussed in Section V, $`f_{\mathrm{rest}}^{\mathrm{OPE}}`$ may cause this huge $`t`$-dependence. Another possibility is due to the large contribution from direct instanton in the pseudoscalar channel. The direct instanton effect is believed to cause large OZI breaking in $`\eta `$ and $`\eta ^{}`$. To confirm it, it will be necessary to include the direct instanton effect in this pseudoscalar channel. Nevertheless, the correlation function Eq. (3) is often used in literature to calculate various couplings and our study suggests that one has to be careful in choosing a Dirac structure in that correlation function.
## VII Conclusions
In this work, we calculated the correlation function Eq. (3) for the vertices, $`\pi NN`$, $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$, and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$, using QCD sum rules. In the construction of sum rules, we used general baryon currents with no derivative instead of the Ioffe current, which enables us to discuss the dependence of sum rule results on currents. We proposed a new criterion to choose a pertinent Dirac structure by studying the dependence of the correlation function on the baryon currents. Specifically, it is imposed that a physical parameter is ideally independent of a chosen current. In checking this constraint, the SU(3) symmetric limit is quite useful as it provides simple relations among the couplings. It is found that the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure satisfies the ideal constraint relatively well, which moderately restricts the $`F/D`$ ratio within the range, $`F/D0.60.8`$. However, the $`i\gamma _5`$ sum rules beyond the chiral limit do not satisfy the constraint, which provides a large window for the value of the $`F/D`$ ratio depending on currents.
In the present study, we considered only the SU(3) limit of the meson-baryon couplings. In fact, the OPE for the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure given in the Appendix A 2 contain effects of SU(3) breaking partially as $`m_Nm_\mathrm{\Xi }m_\mathrm{\Sigma }`$, $`\lambda _N\lambda _\mathrm{\Xi }\lambda _\mathrm{\Sigma }`$, $`\overline{q}q\overline{s}s`$ and $`f_\pi f_\eta `$. If we include these differences, obtained coupling constants break the SU(3) symmetry accordingly. We, however, do not quantify this because other sources of SU(3) breaking are expected. Especially, the large strange quark mass $`(m_s)`$ may cause non-negligible SU(3) breaking effects. So far, the OPE for the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure is truncated to $`𝒪(p)`$ so that it is consistent with the chiral expansion, while effects of $`m_s`$ can only be included at $`𝒪(p^2)`$. In order to quantify SU(3) breaking effects on the meson-baryon couplings, it will be necessary to include $`𝒪(p^2)`$ contribution. The present formulation may give a solid starting point for such analyses in future.
###### Acknowledgements.
This work was supported in part by the Grant-in-Aid for scientific research (C) (2) 11640261 of the Ministry of Education, Science, Sports and Culture of Japan. H. Kim was supported by the Brain Korea 21 project. We would like to thank Prof. S.H. Lee for useful discussions.
## A Coupling sum rules from the PS, T, PV structure
Coupling sum rules for $`\pi NN`$, $`\eta NN`$, $`\pi \mathrm{\Xi }\mathrm{\Xi }`$, $`\eta \mathrm{\Xi }\mathrm{\Xi }`$, $`\pi \mathrm{\Sigma }\mathrm{\Sigma }`$, and $`\eta \mathrm{\Sigma }\mathrm{\Sigma }`$ are presented here. For the $`\eta `$ couplings, $`\eta \eta ^{}`$ mixing is not introduced because our analysis in this paper is within SU(3). In the OPE side, the quark-gluon mixed condensate is parametrized as $`\overline{q_i}g_s\sigma 𝒢q_im_0^2\overline{q_i}q_i`$ where $`q_i=u,d,s`$-quark. Also, we take the isospin symmetric limit, $`\overline{u}u=\overline{d}d\overline{q}q`$ and $`m_u=m_dm_q`$. The continuum contribution is denoted by the factor, $`E_n(xS_0/M^2)=1(1+x+\mathrm{}+x^n/n!)e^x`$ where $`S_0`$ is the continuum threshold.
### 1 Coupling sum rules from the $`i\gamma _5`$ structure
Here we present the $`i\gamma _5`$ sum rules up to dimension 8 constructed at the order $`p^2=m_\pi ^2`$. $`A_B^{\text{PS}}`$ denotes the unknown single-pole term coming from transitions between the ground state baryon and higher resonance states.
$`g_{\pi N}m_\pi ^2\lambda _N^2(1+A_{\pi N}^{\text{PS}}M^2)e^{m_N^2/M^2}`$ (A1)
$`=`$ $`{\displaystyle \frac{m_\pi ^2}{48\pi ^2f_\pi }}(52t+7t^2)\overline{q}qM^4E_0(x)+{\displaystyle \frac{3f_{3\pi }m_\pi ^2}{16\sqrt{2}\pi ^2}}(1+2tt^2)M^4E_0(x)`$ (A4)
$`{\displaystyle \frac{1}{2f_\pi }}(52tt^2)m_q\overline{q}q^2M^2{\displaystyle \frac{m_\pi ^2}{288f_\pi }}(7+2t+5t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`{\displaystyle \frac{1}{12f_\pi }}(76t7t^2)m_qm_0^2\overline{q}q^2`$
$`\sqrt{3}g_{\eta N}m_\eta ^2\lambda _N^2(1+A_{\eta N}^{\text{PS}}M^2)e^{m_N^2/M^2}`$ (A5)
$`=`$ $`{\displaystyle \frac{m_\eta ^2}{48\pi ^2f_\eta }}(7+2t+5t^2)\overline{q}qM^4E_0(x)+{\displaystyle \frac{3f_{3\eta }m_\eta ^2}{16\sqrt{2}\pi ^2}}(12t+t^2)M^4E_0(x)`$ (A8)
$`{\displaystyle \frac{1}{2f_\eta }}(714t3t^2)m_q\overline{q}q^2M^2{\displaystyle \frac{m_\eta ^2}{288f_\eta }}(52t+7t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`{\displaystyle \frac{1}{12f_\eta }}(52t5t^2)m_qm_0^2\overline{q}q^2`$
$`g_{\pi \mathrm{\Xi }}m_\pi ^2\lambda _\mathrm{\Xi }^2(1+A_{\pi \mathrm{\Xi }}^{\text{PS}}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ (A9)
$`=`$ $`{\displaystyle \frac{m_\pi ^2}{48\pi ^2f_\pi }}(1+2tt^2)\overline{q}qM^4E_0(x)+{\displaystyle \frac{3f_{3\pi }m_\pi ^2}{16\sqrt{2}\pi ^2}}(12t+t^2)M^4E_0(x)`$ (A12)
$`{\displaystyle \frac{1}{2f_\pi }}(16tt^2)m_s\overline{q}q\overline{s}sM^2{\displaystyle \frac{m_\pi ^2}{288f_\pi }}(12t+t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`{\displaystyle \frac{1}{12f_\pi }}(1+2t+t^2)m_sm_0^2\overline{q}q\overline{s}s`$
$`\sqrt{3}g_{\eta \mathrm{\Xi }}m_\eta ^2\lambda _\mathrm{\Xi }^2(1+A_{\eta \mathrm{\Xi }}^{\text{PS}}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ (A13)
$`=`$ $`{\displaystyle \frac{m_\eta ^2}{48\pi ^2f_\eta }}[(1+2tt^2)\overline{q}q+12(1t^2)\overline{s}s]M^4E_0(x)+{\displaystyle \frac{3f_{3\eta }m_\eta ^2}{16\sqrt{2}\pi ^2}}(12t+t^2)M^4E_0(x)`$ (A17)
$`{\displaystyle \frac{1}{2f_\eta }}[(1+6t+t^2)m_s\overline{q}q+2(3+4t+t^2)(m_q\overline{q}q+m_s\overline{s}s)]\overline{s}sM^2`$
$`{\displaystyle \frac{m_\eta ^2}{288f_\eta }}[(12t+t^2)\overline{q}q+12(1t^2)\overline{s}s]{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`{\displaystyle \frac{1}{12f_\eta }}[(1+2t+t^2)m_s\overline{q}q+2(3+2t+3t^2)(m_q\overline{s}s+m_s\overline{q}q)]m_0^2\overline{s}s`$
$`g_{\pi \mathrm{\Sigma }}m_\pi ^2\lambda _\mathrm{\Sigma }^2(1+A_{\pi \mathrm{\Sigma }}^{\text{PS}}M^2)e^{m_\mathrm{\Sigma }^2/M^2}`$ (A18)
$`=`$ $`{\displaystyle \frac{m_\pi ^2}{48\pi ^2f_\pi }}(6+6t^2)\overline{q}qM^4E_0(x){\displaystyle \frac{1}{2f_\pi }}(34tt^2)(m_q\overline{q}q+m_s\overline{s}s)\overline{q}qM^2`$ (A20)
$`{\displaystyle \frac{m_\pi ^2}{288f_\pi }}(6+6t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2{\displaystyle \frac{1}{12f_\pi }}(32t3t^2)(m_q\overline{s}s+m_s\overline{q}q)m_0^2\overline{q}q`$
$`\sqrt{3}g_{\eta \mathrm{\Sigma }}m_\eta ^2\lambda _\mathrm{\Sigma }^2(1+A_{\eta \mathrm{\Sigma }}^{\text{PS}}M^2)e^{m_\mathrm{\Sigma }^2/M^2}`$ (A21)
$`=`$ $`{\displaystyle \frac{m_\eta ^2}{48\pi ^2f_\eta }}[6(1t^2)\overline{q}q+2(12t+t^2)\overline{s}s]M^4E_0(x)+{\displaystyle \frac{3f_{3\eta }m_\eta ^2}{16\sqrt{2}\pi ^2}}(2+4t2t^2)M^4E_0(x)`$ (A25)
$`{\displaystyle \frac{1}{2f_\eta }}[2(1+6t+t^2)m_q\overline{s}s(3+4t+t^2)(m_q\overline{q}q+m_s\overline{s}s)]\overline{q}qM^2`$
$`{\displaystyle \frac{m_\eta ^2}{288f_\eta }}[6(1t^2)\overline{q}q+2(1+2tt^2)\overline{s}s]{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`{\displaystyle \frac{1}{12f_\eta }}[2(1+2t+t^2)m_q\overline{s}s(3+2t+3t^2)(m_q\overline{s}s+m_s\overline{q}q)]m_0^2\overline{q}q`$
### 2 Coupling sum rules from the $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ structure
The $`\gamma _5\sigma _{\mu \nu }q^\mu p^\nu `$ sum rules up to dimension 7 are the following. Again, $`A_B^\text{T}`$ denotes the unknown single-pole term contribution.
$`g_{\pi N}\lambda _N^2(1+A_{\pi N}^\text{T}M^2)e^{m_N^2/M^2}`$ (A26)
$`=`$ $`{\displaystyle \frac{1}{96\pi ^2f_\pi }}(10+4t14t^2)\overline{q}qM^4E_0(x){\displaystyle \frac{f_\pi }{3}}(12t+3t^2)\overline{q}qM^2`$ (A29)
$`{\displaystyle \frac{1}{54}}f_\pi \delta ^2(126t+27t^2)\overline{q}q+{\displaystyle \frac{1}{7212f_\pi }}(17+2t19t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\pi }{72}}(56t+11t^2)m_0^2\overline{q}q`$
$`\sqrt{3}g_{\eta N}\lambda _N^2(1+A_{\eta N}^\text{T}M^2)e^{m_N^2/M^2}`$ (A30)
$`=`$ $`{\displaystyle \frac{1}{96\pi ^2f_\eta }}(144t10t^2)\overline{q}qM^4E_0(x){\displaystyle \frac{f_\eta }{3}}(12t+t^2)\overline{q}qM^2`$ (A33)
$`{\displaystyle \frac{1}{54}}f_\eta \delta ^2(1326t+13t^2)\overline{q}q+{\displaystyle \frac{1}{7212f_\eta }}(192t17t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\eta }{72}}(96t3t^2)m_0^2\overline{q}q`$
$`g_{\pi \mathrm{\Xi }}\lambda _\mathrm{\Xi }^2(1+A_{\pi \mathrm{\Xi }}^\text{T}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ (A34)
$`=`$ $`{\displaystyle \frac{1}{96\pi ^2f_\pi }}(24t+2t^2)\overline{q}qM^4E_0(x){\displaystyle \frac{f_\pi }{3}}(1t^2)\overline{s}sM^2`$ (A37)
$`{\displaystyle \frac{1}{54}}f_\pi \delta ^2(77t^2)\overline{s}s+{\displaystyle \frac{1}{7212f_\pi }}(12t+t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\pi }{72}}(77t^2)m_0^2\overline{s}s`$
$`\sqrt{3}g_{\eta \mathrm{\Xi }}\lambda _\mathrm{\Xi }^2(1+A_{\eta \mathrm{\Xi }}^\text{T}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ (A38)
$`=`$ $`{\displaystyle \frac{1}{96\pi ^2f_\eta }}\left[(24t+2t^2)\overline{q}q+24(1+t^2)\overline{s}s\right]M^4E_0(x)`$ (A43)
$`{\displaystyle \frac{f_\eta }{3}}\left[(2+4t2t^2)\overline{q}q+3(1t^2)\overline{s}s\right]M^2`$
$`{\displaystyle \frac{1}{54}}f_\eta \delta ^2\left[26(1+2tt^2)\overline{q}q+21(1t^2)\overline{s}s\right]`$
$`+{\displaystyle \frac{1}{7212f_\eta }}\left[(12t+t^2)\overline{q}q+36(1+t^2)\overline{s}s\right]{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\eta }{72}}\left[6(1+2tt^2)\overline{q}q+9(1t^2)\overline{s}s\right]m_0^2`$
$`g_{\pi \mathrm{\Sigma }}\lambda _\mathrm{\Sigma }^2(1+A_{\pi \mathrm{\Sigma }}^\text{T}M^2)e^{m_\mathrm{\Sigma }^2/M^2}`$ (A44)
$`=`$ $`{\displaystyle \frac{1}{96\pi ^2f_\pi }}12(1t^2)\overline{q}qM^4E_0(x){\displaystyle \frac{f_\pi }{3}}\left[(1+t^2)\overline{q}q+(12t+t^2)\overline{s}s\right]M^2`$ (A47)
$`{\displaystyle \frac{1}{54}}f_\pi \delta ^2\left[7(1+t^2)\overline{q}q+13(12t+t^2)\overline{s}s\right]+{\displaystyle \frac{1}{7212f_\pi }}18(1t^2)\overline{q}q{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\pi }{72}}\left[(1+t^2)\overline{q}q+3(12t+t^2)\overline{s}s\right]m_0^2`$
$`\sqrt{3}g_{\eta \mathrm{\Sigma }}\lambda _\mathrm{\Sigma }^2(1+A_{\eta \mathrm{\Sigma }}^\text{T}M^2)e^{m_\mathrm{\Sigma }^2/M^2}`$ (A48)
$`=`$ $`{\displaystyle \frac{1}{96\pi ^2f_\eta }}\left[12(1t^2)\overline{q}q+4(1+2tt^2)\overline{s}s\right]M^4E_0(x)`$ (A53)
$`{\displaystyle \frac{f_\eta }{3}}\left[3(1+t^2)\overline{q}q+(12t+t^2)\overline{s}s\right]M^2`$
$`{\displaystyle \frac{1}{54}}f_\eta \delta ^2\left[21(1+t^2)\overline{q}q+13(12t+t^2)\overline{s}s\right]`$
$`+{\displaystyle \frac{1}{7212f_\eta }}\left[18(1t^2)\overline{q}q+2(1+2tt^2)\overline{s}s\right]{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2`$
$`+{\displaystyle \frac{f_\eta }{72}}\left[15(1+t^2)\overline{q}q+3(12t+t^2)\overline{s}s\right]m_0^2`$
### 3 Coupling sum rules from the $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ structure
The $`i\gamma _5\begin{array}{c}\text{/}\\ p\hfill \end{array}`$ sum rules up to dimension 7 are presented here.
$`g_{\pi N}m_N\lambda _N^2(1+A_{\pi N}^{\text{PV}}M^2)e^{m_N^2/M^2}`$ (A54)
$`=`$ $`{\displaystyle \frac{f_\pi }{24\pi ^2}}(10+8t+10t^2)M^6E_1(x){\displaystyle \frac{f_\pi \delta ^2}{48\pi ^2}}(2016t20t^2)M^4E_0(x)`$ (A57)
$`+{\displaystyle \frac{f_\pi }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(5+4t+5t^2)M^2+{\displaystyle \frac{1}{18f_\pi }}\overline{q}q^2(12t+3t^2)M^2`$
$`+{\displaystyle \frac{f_\pi \delta ^2}{3618}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(42t4t^2)+{\displaystyle \frac{1}{432f_\pi }}m_0^2\overline{q}q^2(56t+11t^2)`$
$`\sqrt{3}g_{\eta N}m_N\lambda _N^2(1+A_{\eta N}^{\text{PV}}M^2)e^{m_N^2/M^2}`$ (A58)
$`=`$ $`{\displaystyle \frac{f_\eta }{24\pi ^2}}(8+8t+8t^2)M^6E_1(x){\displaystyle \frac{f_\eta \delta ^2}{48\pi ^2}}(224t22t^2)M^4E_0(x)`$ (A61)
$`+{\displaystyle \frac{f_\eta }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(5+2t+5t^2)M^2+{\displaystyle \frac{1}{18f_\eta }}\overline{q}q^2(12t+t^2)M^2`$
$`+{\displaystyle \frac{f_\eta \delta ^2}{3618}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(104t10t^2)+{\displaystyle \frac{1}{432f_\eta }}m_0^2\overline{q}q^2(96t3t^2)`$
$`g_{\pi \mathrm{\Xi }}m_\mathrm{\Xi }\lambda _\mathrm{\Xi }^2(1+A_{\pi \mathrm{\Xi }}^{\text{PV}}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ (A62)
$`=`$ $`{\displaystyle \frac{f_\pi }{24\pi ^2}}(1t^2)M^6E_1(x){\displaystyle \frac{f_\pi \delta ^2}{48\pi ^2}}(1+6tt^2)M^4E_0(x)`$ (A65)
$`+{\displaystyle \frac{f_\pi }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(t)M^2+{\displaystyle \frac{1}{18f_\pi }}\overline{q}q\overline{s}s(1t^2)M^2`$
$`+{\displaystyle \frac{f_\pi \delta ^2}{3618}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(3t3t^2)+{\displaystyle \frac{1}{432f_\pi }}m_0^2\overline{q}q\overline{s}s(77t^2)`$
$`\sqrt{3}g_{\eta \mathrm{\Xi }}m_\mathrm{\Xi }\lambda _\mathrm{\Xi }^2(1+A_{\eta \mathrm{\Xi }}^{\text{PV}}M^2)e^{m_\mathrm{\Xi }^2/M^2}`$ (A66)
$`=`$ $`{\displaystyle \frac{f_\eta }{24\pi ^2}}(1916t19t^2)M^6E_1(x){\displaystyle \frac{f_\eta \delta ^2}{48\pi ^2}}(41+26t+41t^2)M^4E_0(x)`$ (A69)
$`+{\displaystyle \frac{f_\eta }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(107t10t^2)M^2+{\displaystyle \frac{1}{18f_\eta }}\left[(33t^2)\overline{q}q+(2+4t2t^2)\overline{s}s\right]\overline{s}sM^2`$
$`+{\displaystyle \frac{f_\eta \delta ^2}{3618}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(11+5t+11t^2)+{\displaystyle \frac{1}{432f_\eta }}\left[9(1t^2)\overline{q}q+6(1+2tt^2)\overline{s}s\right]m_0^2\overline{s}s`$
$`g_{\pi \mathrm{\Sigma }}m_\mathrm{\Sigma }\lambda _\mathrm{\Sigma }^2(1+A_{\pi \mathrm{\Sigma }}^{\text{PV}}M^2)e^{m_\mathrm{\Sigma }^2/M^2}`$ (A70)
$`=`$ $`{\displaystyle \frac{f_\pi }{24\pi ^2}}(9+8t+9t^2)M^6E_1(x){\displaystyle \frac{f_\pi \delta ^2}{48\pi ^2}}(2110t21t^2)M^4E_0(x)`$ (A73)
$`+{\displaystyle \frac{f_\pi }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(5+3t+5t^2)M^2+{\displaystyle \frac{1}{18f_\pi }}\left[(12t+t^2)\overline{q}q+(1+t^2)\overline{s}s\right]\overline{q}qM^2`$
$`+{\displaystyle \frac{f_\pi \delta ^2}{3618}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(73t7t^2)+{\displaystyle \frac{1}{432f_\pi }}\left[3(1+2tt^2)\overline{q}q(1t^2)\overline{s}s\right]m_0^2\overline{q}q`$
$`\sqrt{3}g_{\eta \mathrm{\Sigma }}m_\mathrm{\Sigma }\lambda _\mathrm{\Sigma }^2(1+A_{\eta \mathrm{\Sigma }}^{\text{PV}}M^2)e^{m_\mathrm{\Sigma }^2/M^2}`$ (A74)
$`=`$ $`{\displaystyle \frac{f_\eta }{24\pi ^2}}(11+8t+11t^2)M^6E_1(x){\displaystyle \frac{f_\eta \delta ^2}{48\pi ^2}}(1922t19t^2)M^4E_0(x)`$ (A77)
$`+{\displaystyle \frac{f_\eta }{72}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(5+5t+5t^2)M^2+{\displaystyle \frac{1}{18f_\eta }}\left[(12t+t^2)\overline{q}q3(1t^2)\overline{s}s\right]\overline{q}qM^2`$
$`+{\displaystyle \frac{f_\eta \delta ^2}{3618}}{\displaystyle \frac{\alpha _s}{\pi }}𝒢^2(1tt^2)+{\displaystyle \frac{1}{432f_\eta }}\left[3(1+2tt^2)\overline{q}q15(1t^2)\overline{s}s\right]m_0^2\overline{q}q`$
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# The tachyon potential in Witten’s superstring field theory
## 1 Introduction
It has been conjectured by Sen that at the stationary point of the tachyon potential for the D-brane-anti-D-brane pair or for the non-BPS D-brane of superstring theories, the negative energy density precisely cancels the brane tensions .
For the D-brane of bosonic string theory, this conjecture has been verified starting from Witten’s open string field theory and in the supersymmetric case starting from Berkovits’ superstring field theory .<sup>1</sup><sup>1</sup>1We do not agree with the results of .
In this paper we study the tachyon potential in the NS sector of Witten’s superstring field theory . Soon after this theory was proposed, it became clear that it suffered some problems . Infinities were found to arise in the calculation of tree-level scattering amplitudes. These find their origin in picture-changing operators inserted at the same point. Similar problems arise in the proof of gauge invariance of the action. Modifications of Witten’s action have been proposed in order to solve these issues , but these seem to suffer from other difficulties . In the light of these problems, the present calculation should be seen as a further study of Witten’s string field theory proposal rather than as a test of Sen’s conjecture, which has been extensively verified on a quantitative level.
We use the level truncation method of Kostelecky and Samuel retaining fields up to level 2 and terms up to level 4 in the action. We find that the tachyon potential in this theory does not have the behaviour conjectured by Sen, in fact it doesn’t even have a minimum. This is already the case for the pure tachyon contribution (in contrast to the pure tachyon contribution is Berkovits’ theory ), and this behaviour is found to persist when higher modes are included.
## 2 Witten’s superstring field theory in CFT language
In this section we review Wittens string field theory for open superstrings formulated in the conformal field theory (CFT) approach . We restrict our attention to fields in the NS sector of the theory.
In the CFT approach, a general NS-sector string field $`A`$ is represented by a CFT vertex operator of ghost number $`1`$ and picture number $`1`$. The superghosts are bosonized as $`\beta =\xi e^\varphi `$ and $`\gamma =\eta e^\varphi `$ and the string field $`A`$ is restricted to live in the “small” Hilbert space of the bosonized system, i.e. it does not include the $`\xi `$ zero mode. The string field should also satisfy a suitable reality condition in order for the action to take real values.
Witten proposed a Chern-Simons type action
$$S=AQ_BA\frac{2}{3}A^3,$$
(1)
where $`Q_B`$ is the BRST charge
$$Q_B=𝑑zj_B(z)=𝑑z\left\{c(T_m+T_{\xi \eta }+T_\varphi )+ccb+\eta e^\varphi G_m\eta \eta e^{2\varphi }b\right\},$$
and the double brackets should be interpreted as a CFT correlator<sup>2</sup><sup>2</sup>2The elementary CFT correlator is normalized as $`cc^2c(z_1)e^{2\varphi }(z_2)=2`$ with extra picture-changing insertions:
$`\mathrm{\Phi }_1\mathrm{\Phi }_2\mathrm{}\mathrm{\Phi }_{n1}\mathrm{\Phi }_n=`$ (2)
$`=Y(0)f_1^n\mathrm{\Phi }_1(0)X(0)f_2^n\mathrm{\Phi }_2(0)X(0)\mathrm{}X(0)f_{n1}^n\mathrm{\Phi }_{n1}(0)X(0)f_n^n\mathrm{\Phi }_n(0).`$
The symbols $`f_i^n`$ denote conformal transformations mapping the unit circle to wedge-formed pieces of the complex-plane:
$$f_k^n(z)=e^{2\pi i(k1)/n}\left(\frac{1+iz}{1iz}\right)^{2/n}\text{for}n1.$$
The operators $`X`$ and $`Y`$ are the picture-changing and inverse picture-changing operators, respectively:
$`X`$ $`=`$ $`\xi c+e^\varphi G_m\eta be^{2\varphi }\left(\eta be^{2\varphi }\right),`$
$`Y`$ $`=`$ $`\xi ce^{2\varphi }.`$
They are each others inverse in the following sense:
$$\underset{ϵ0}{lim}Y(z+ϵ)X(z)=1.$$
The action is formally invariant under the gauge transformations
$$\delta A=Q_Bϵ+AϵϵA,$$
where $`ϵ`$ is a NS string field of ghost number 0 in the $`1`$ picture. The proof of the gauge invariance relies on the following properties of the correlator (2)
$`Q_B(\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_n)`$ $`=`$ $`0,`$
$`\mathrm{}Q_B^2(\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_n)\mathrm{}`$ $`=`$ $`0,`$
$`\mathrm{}Q_B(A\mathrm{\Phi })\mathrm{}`$ $`=`$ $`\mathrm{}(Q_BA\mathrm{\Phi }AQ_B\mathrm{\Phi })\mathrm{},`$
$`\mathrm{}Q_B(ϵ\mathrm{\Phi })\mathrm{}`$ $`=`$ $`\mathrm{}(Q_Bϵ\mathrm{\Phi }+ϵQ_B\mathrm{\Phi })\mathrm{},`$
$`\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_{n1}\mathrm{\Phi }_n`$ $`=`$ $`\mathrm{\Phi }_n\mathrm{\Phi }_1\mathrm{}\mathrm{\Phi }_{n1}.`$ (3)
$`A`$ and $`ϵ`$ represent the string field and the gauge parameter respectively and the $`\mathrm{\Phi }_i`$ represent arbitrary operators in the GSO(+) sector.
These properties suffice to prove the gauge invariance of the action. Such a “proof” should be taken with a grain of salt however since it involves manipulating correlators with four string fields. From the definition (2), one sees that such correlators contain 2 insertions of the picture-changing operator at the origin and are hence divergent.<sup>3</sup><sup>3</sup>3A similar problem is encountered in the proof of associativity of Witten’s star product of string fields. These and other problems of the action (1) were discussed in .
## 3 Witten’s string field theory on a non-BPS D-brane
In order to describe NS excitations on a non-BPS D-brane, one should extend the theory to include also states with odd world-sheet fermion number. This can be accomplished by tensoring the fields entering in the action and gauge transformation with suitable internal Chan-Paton (CP) factors as in . The fields with CP factors attached are denoted with hats. We take the following assignments for the internal CP factors:
$`\widehat{A}`$ $`=`$ $`A_+I+A_{}\sigma _1,`$
$`\widehat{ϵ}`$ $`=`$ $`ϵ_+\sigma _3+ϵ_{}i\sigma _2,`$
$`\widehat{Q}_B`$ $`=`$ $`Q_B\sigma _3,`$
$`\widehat{Y}`$ $`=`$ $`YI,`$
$`\widehat{X}`$ $`=`$ $`X\sigma _3,`$ (4)
where $`I`$ is the $`2\times 2`$ unit matrix and $`\sigma _i`$ are the Pauli matrices. The action takes the form
$$S=\widehat{A}\widehat{Q}_B\widehat{A}\frac{2}{3}\widehat{A}^3,$$
where the double brackets should now be interpreted as
$`\widehat{\mathrm{\Phi }}_1\widehat{\mathrm{\Phi }}_2\mathrm{}\widehat{\mathrm{\Phi }}_{n1}\widehat{\mathrm{\Phi }}_n`$ $`=`$ $`{\displaystyle \frac{1}{2}}Tr\widehat{Y}(0)f_1^n\widehat{\mathrm{\Phi }}_1(0)\widehat{X}(0)f_2^n\widehat{\mathrm{\Phi }}_2(0)\times `$
$`\times \widehat{X}(0)\mathrm{}\widehat{X}(0)f_{n1}^n\widehat{\mathrm{\Phi }}_{n1}(0)\widehat{X}(0)f_n^n\widehat{\mathrm{\Phi }}_n(0).`$
The trace runs over the internal CP indices. The gauge transformations are
$$\delta \widehat{A}=\widehat{Q}_B\widehat{ϵ}+\widehat{A}\widehat{ϵ}\widehat{ϵ}\widehat{A}.$$
(5)
The CP factor assignments (4) were chosen such that gauge transformations preserve the CP structure of the string field:
$$\delta \widehat{A}=\delta A_+I+\delta A_{}\sigma _1$$
and such that the algebraic structure (3) is preserved:
$`\widehat{Q}_B(\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_n)`$ $`=`$ $`0,`$
$`\mathrm{}\widehat{Q}_{B}^{}{}_{}{}^{2}(\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_n)\mathrm{}`$ $`=`$ $`0,`$
$`\mathrm{}\widehat{Q}_B(\widehat{A}\widehat{\mathrm{\Phi }})\mathrm{}`$ $`=`$ $`\mathrm{}(\widehat{Q}_B\widehat{A}\widehat{\mathrm{\Phi }}\widehat{A}\widehat{Q}_B\widehat{\mathrm{\Phi }})\mathrm{},`$
$`\mathrm{}\widehat{Q}_B(\widehat{ϵ}\widehat{\mathrm{\Phi }})\mathrm{}`$ $`=`$ $`\mathrm{}(\widehat{Q}_B\widehat{ϵ}\widehat{\mathrm{\Phi }}+\widehat{ϵ}\widehat{Q}_B\widehat{\mathrm{\Phi }})\mathrm{},`$
$`\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_{n1}\widehat{\mathrm{\Phi }}_n`$ $`=`$ $`\widehat{\mathrm{\Phi }}_n\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_{n1}.`$
With these properties (proven in appendix B), the (formal) proof of gauge invariance goes through as in the GSO(+) sector.
## 4 The fields up to level 2
In the calculation of the tachyon potential, we can restrict the string field to lie in a subspace $`_1`$ formed by acting only with modes of the stress-energy tensor, the supercurrent and the ghost fields $`b,c,\eta ,\xi ,\varphi `$ since the other excitations can be consistently put to zero . We fix the gauge freedom (5) by imposing the Feynman-Siegel gauge $`b_0\widehat{A}=0`$ on fields with non-zero conformal weight. Taking all this together we get the following list of contributing fields up to level 2. The level of a field is just the conformal weight shifted by 1/2, in this way the tachyon is a level 0 field. We use the notation $`|q`$ for the state corresponding with the operator $`e^{q\varphi }`$.
## 5 The tachyon potential
We have calculated the tachyon potential involving fields up to level 2, including only the terms up to level 4 (the level of a term in the potential is defined to be the sum of the levels of the fields entering into it). We have performed the actual computation in the following way: the conformal tranformations of the fields were calculated by hand and the computation of all the correlation functions between these transformed fields was done with the help of Mathematica. We have written a program to compute the necessary CFT correlation functions.
We denote (subscripts refer to the level)
$`\widehat{\mathrm{\Phi }}`$ $`=`$ $`t\widehat{T}+ir\widehat{R}+is\widehat{S}+a\widehat{A}+e\widehat{E}+f\widehat{F}+k\widehat{K}+l\widehat{L}+m\widehat{M}+n\widehat{N}+p\widehat{P},`$
$`S(\widehat{\mathrm{\Phi }})`$ $`=`$ $`S_0+S_{1/2}+S_1+S_{3/2}+S_2+S_{5/2}+S_3+S_{7/2}+S_4.`$
The factors of $`i`$ are included because of the reality condition on the string field. We have obtained the following result:
$`S_0`$ $`=`$ $`{\displaystyle \frac{t^2}{2}},S_{1/2}={\displaystyle \frac{9}{2}}rt^2,S_1=2r^2,`$
$`S_{3/2}`$ $`=`$ $`8rst10at^2+2et^220ft^2,`$
$`S_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}s^2,`$
$`S_{5/2}`$ $`=`$ $`{\displaystyle \frac{16}{3\sqrt{3}}}er^2+{\displaystyle \frac{10}{3}}rs^2{\displaystyle \frac{23}{3}}krt{\displaystyle \frac{25}{2}}mrt+2nrtprt`$
$`{\displaystyle \frac{64}{3}}ast{\displaystyle \frac{80}{3}}fst,`$
$`S_3`$ $`=`$ $`4ae10f^2,`$
$`S_{7/2}`$ $`=`$ $`{\displaystyle \frac{64}{3\sqrt{3}}}aer+{\displaystyle \frac{32}{9\sqrt{3}}}e^2r{\displaystyle \frac{320}{9\sqrt{3}}}efr+{\displaystyle \frac{320}{9\sqrt{3}}}f^2r`$
$`\mathrm{\hspace{0.17em}8}krs+{\displaystyle \frac{160}{27}}lrs{\displaystyle \frac{100}{9}}mrs+{\displaystyle \frac{80}{9}}nrs+{\displaystyle \frac{152}{27}}prs+`$
$`+{\displaystyle \frac{460}{27}}akt{\displaystyle \frac{28}{27}}ekt+{\displaystyle \frac{280}{27}}fkt+{\displaystyle \frac{128}{27}}alt+{\displaystyle \frac{64}{27}}elt+{\displaystyle \frac{250}{9}}amt`$
$`{\displaystyle \frac{50}{9}}emt+{\displaystyle \frac{220}{3}}fmt{\displaystyle \frac{104}{9}}ant+{\displaystyle \frac{40}{9}}ent{\displaystyle \frac{880}{27}}fnt+`$
$`+{\displaystyle \frac{20}{9}}apt{\displaystyle \frac{4}{9}}ept{\displaystyle \frac{200}{27}}fpt,`$
$`S_4`$ $`=`$ $`6k^2+6kl3l^2{\displaystyle \frac{45}{4}}m^2+6n^2+{\displaystyle \frac{3}{2}}p^2.`$
Taking a look at the potential, we notice immediately that the $`Z_2`$ twist symmetry of Berkovits’ theory is absent in Witten’s theory. Indeed, if this symmetry were present the fields $`r`$ and $`t`$ would carry twist charge $`1`$ and $`+1`$, respectively and the $`rt^2`$ term would be constrained to vanish. The origin of this difference between the two theories can be traced to the different structure of their lagrangians: in Berkovits’ action, all operators are inserted on the boundary of the disc, while in Witten’s action there is always an insertion of the picture-changing operator in the origin. The proof of twist invariance of the former theory breaks down for Witten’s theory since it makes use of transformations that do not leave the origin invariant and hence move the picture-changing operator away from the origin.
At level 1 and 2 the fields $`r`$ and $`s`$ can be integrated out exactly to give the following effective potentials:
$`V_0(t)`$ $`=`$ $`{\displaystyle \frac{t^2}{2}},V_1(t)={\displaystyle \frac{t^2}{2}}{\displaystyle \frac{81t^4}{32}},`$
$`V_2(t)`$ $`=`$ $`{\displaystyle \frac{t^2}{2}}+{\displaystyle \frac{81t^4}{2(464t^2)^2}}{\displaystyle \frac{648t^6}{(464t^2)^2}}{\displaystyle \frac{81t^4}{4(464t^2)}}.`$
We see that the inclusion of these higher modes does not alter the fact that the potential has no minimum, in fact, its slope becomes even steeper and a singularity is encountered at level 2, see figure 1. The singularity in the tachyon potential encountered here might be contrasted with the singularities in the tachyon potential of open bosonic string theory found in . In the case at hand, the potential diverges at the singular point. Moreover, it doesn’t have the interpretation as a point where different branches of the effective potential come together since, at level 2, the equations for the fields that are integrated out are linear.
Integrating out the fields numerically for the higher levels, one finds that the more fields are included, the steeper the slope of the potential becomes. This behaviour was anticipated in the conclusions of .
## 6 Conclusions
In this letter we calculated the tachyon potential in the NS sector of Witten’s superstring field theory. We used a level truncation method and kept terms up to level four in the potential. We found that the potential does not exhibit the expected minimum, and that the situation seems to become worse the more levels are included. This result, in addition to other problems encountered within this formulation, seems to point towards Berkovits’ string field theory proposal as a more viable candidate for the description of off-shell superstring interactions.
###### Acknowledgments.
This work was supported in part by the European Commission TMR project ERBFMRXCT96-0045. We would like to thank B. Zwiebach for discussions. P.J.D.S. is aspirant FWO-Vlaanderen.
## Appendix A The conformal transformations of the fields
We now list the conformal transformations of the fields used in the calculation of the tachyon potential. To shorten the notation we denote $`w=f(z)`$.
$`fT(z)`$ $`=`$ $`(f^{}(z))^{1/2}T(w),`$
$`fR(z)`$ $`=`$ $`R(w),`$
$`fS(z)`$ $`=`$ $`(f^{}(z))^{1/2}S(w){\displaystyle \frac{1}{2}}{\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}ce^\varphi (w),`$
$`fA(z)`$ $`=`$ $`f^{}(z)A(w){\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}cc\xi e^{2\varphi }(w),`$
$`fE(z)`$ $`=`$ $`f^{}(z)E(w),`$
$`fF(z)`$ $`=`$ $`f^{}(z)F(w),`$
$`fK(z)`$ $`=`$ $`(f^{}(z))^{3/2}K(w)+2{\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}c\left(e^\varphi \right)(w)+`$
$`+\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2\right](f^{}(z))^{1/2}ce^\varphi (w),`$
$`fL(z)`$ $`=`$ $`(f^{}(z))^{3/2}L(w)+{\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}c\varphi e^\varphi (w)+`$
$`+\left[{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2{\displaystyle \frac{2}{3}}{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}\right](f^{}(z))^{1/2}ce^\varphi (w),`$
$`fM(z)`$ $`=`$ $`(f^{}(z))^{3/2}M(w)+{\displaystyle \frac{15}{12}}\left[{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}{\displaystyle \frac{3}{2}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2\right](f^{}(z))^{1/2}ce^\varphi (w),`$
$`fN(z)`$ $`=`$ $`(f^{}(z))^{3/2}N(w){\displaystyle \frac{f^{\prime \prime }(z)}{f^{}(z)}}(f^{}(z))^{1/2}ce^\varphi (w)+`$
$`+\left[2\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}\right](f^{}(z))^{1/2}ce^\varphi (w),`$
$`fP(z)`$ $`=`$ $`(f^{}(z))^{3/2}P(w)+`$
$`+\left[{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{f^{\prime \prime }}{f^{}}}\right)^2{\displaystyle \frac{1}{6}}{\displaystyle \frac{f^{\prime \prime \prime }}{f^{}}}\right](f^{}(z))^{1/2}ce^\varphi (w).`$
## Appendix B Cyclicity property of string amplitudes
The proof of the cyclicity is based on appendix A of . First we give some preliminary remarks.
It is easy to see that all the fields in the GSO($`+`$) sector are tensored with either $`I`$ or $`\sigma _3`$, and all the fields in the GSO($``$) sector with either $`\sigma _1`$ or $`i\sigma _2`$. The fields in the GSO($`+`$) sector have integer conformal weight, and the fields in GSO($``$) sector have half-integer conformal weight. We compute
$`\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_n`$ $`=`$ $`Tr\widehat{Y}f_1^n\widehat{\mathrm{\Phi }}_1\widehat{X}\mathrm{}\widehat{X}f_{n1}^n\widehat{\mathrm{\Phi }}_{n1}\widehat{X}f_n^n\widehat{\mathrm{\Phi }}_n`$
$`=`$ $`Tr\widehat{Y}f_2^n\widehat{\mathrm{\Phi }}_1\widehat{X}\mathrm{}\widehat{X}f_n^n\widehat{\mathrm{\Phi }}_{n1}\widehat{X}Rf_1^n\widehat{\mathrm{\Phi }}_n,`$
where $`R`$ is the rotation over an angle of $`2\pi `$. Next we use $`Rf_1^n\widehat{\mathrm{\Phi }}_n=\pm f_1^n\widehat{\mathrm{\Phi }}_n`$, with a plus sign if the field has integer weight, and a minus sign if the field has half-integer weight. We also use the cyclicity of the trace to move the field $`\widehat{\mathrm{\Phi }}_n`$ in front. Although we are manipulating grassmann objects, we do not get an additional minus sign, due to the fact that the total amplitude is grassmann odd.
$`\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_n`$ $`=`$ $`\pm Tr\widehat{Y}f_1^n\widehat{\mathrm{\Phi }}_nf_2^n\widehat{\mathrm{\Phi }}_1\widehat{X}\mathrm{}\widehat{X}f_n^n\widehat{\mathrm{\Phi }}_{n1}\widehat{X}`$ (6)
$`=`$ $`\pm Tr\widehat{Y}\widehat{X}f_1^n\widehat{\mathrm{\Phi }}_nf_2^n\widehat{\mathrm{\Phi }}_1\widehat{X}\mathrm{}\widehat{X}f_n^n\widehat{\mathrm{\Phi }}_{n1}`$
$`=`$ $`Tr\widehat{Y}f_1^n\widehat{\mathrm{\Phi }}_n\widehat{X}f_2^n\widehat{\mathrm{\Phi }}_1\widehat{X}\mathrm{}\widehat{X}f_n^n\widehat{\mathrm{\Phi }}_{n1}`$
$`=`$ $`\widehat{\mathrm{\Phi }}_n\widehat{\mathrm{\Phi }}_1\mathrm{}\widehat{\mathrm{\Phi }}_{n1}.`$
In the next to last line the picture changing operator commutes with the GSO($`+`$) fields and anticommutes with the GSO($``$) fields, cancelling the minus sign in front of the amplitude.
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# Bounds of Ideal Class Numbers of Real Quadratic Function Fields Project Supported by the NNSFC (No.19771052)
Abstract. Theory of continued fractions of functions is used to give lower bound for class numbers $`h(D)`$ of general real quadratic function fields $`K=k(\sqrt{D})`$ over $`k=𝐅_q(T)`$. For five series of real quadratic function fields $`K`$, the bounds of $`h(D)`$ are given more explicitly, e.g., if $`D=F^2+c,`$ then $`h(D)\text{deg}F/\text{deg}P;`$ if $`D=(SG)^2+cS,`$ then $`h(D)\text{deg}S/\text{deg}P;`$ if $`D=(A^m+a)^2+A,`$ then $`h(D)\text{deg}A/\text{deg}P,`$ where $`P`$ is irreducible polynomial splitting in $`K,c𝐅_q`$ . In addition, three types of quadratic function fields $`K`$ are found to have ideal class numbers bigger than one.
Keywords: quadratic function field, ideal class number, continued fraction of function
MR (2000) Subject Classification: 11R58; 11R29; 14H05.
China Library Classification: O156.2
I. Introduction and Main Results
Suppose that $`k=𝐅_q(T)`$ is the rational function field in indeterminate (variable) $`T`$ over $`𝐅_q`$, the finite field with $`q`$ elements ($`q`$ is a power of odd prime number). Let $`R=𝐅_q[T]`$ be the polynomial ring of $`T`$ over $`𝐅_q`$. Any finite algebraic extension $`K`$ of $`k`$ is said to be an (algebraic) function field. The integral closure of $`R`$ in $`K`$ is said to be the ring (domain) of integers of $`K`$, and is denoted by $`𝒪_K`$, which is a Dedekind domain. The fractional ideals of $`𝒪_K`$ form a multiplication group $`_K`$. Let $`𝒫_K`$ denotes the principal ideals in $`_K`$. Then the quotient group $`H(K)=_K/𝒫_K`$ is said to be the ideal class group of $`K`$. And $`h(K)=\mathrm{\#}H(K)`$ (the order of $`H(K)`$) is said to be the ideal class number of $`K`$.
A quadratic extension of $`k`$ could be expressed as $`K=k(\sqrt{D})`$, where $`DR`$ is a polynomial which is not a square (we could also assume $`D`$ is square-free). If in addition $`D`$ is monic with even degree, then $`K=k(\sqrt{D})`$ is said to be a real quadratic function field. For real quadratic number fields $`K`$, Mollin in 1987 obtained a lower bound of ideal class numbers $`h(D)=h(K)`$ by evaluating the fundamental unit (see ). Feng and Hu in obtained a similar results for function fields; they also gave an explicit bound of $`h(K)`$ for $`K=k(\sqrt{F^2+c})`$, where $`c𝐅_q^\times `$. There are also other works using continued fractions to study quadratic number fields(see \[3-6\]). We here give a theorem on lower bound of $`h(K)`$ for general real quadratic function fields $`K`$, and obtain explicit lower bound of $`h(K)`$ for six types of $`K`$ including the fields $`K=k(\sqrt{F^2+c})`$.
E. Artin in began to use continued fractions to study quadratic function fields. In we re-developed the theory of continued fractions of algebraic functions in an elementary and practicable way and studied some properties of them, which will be used here.
Suppose that $`D`$ is a monic square-free polynomial with degree $`2d.`$ By we know $`\sqrt{D}`$ has an expansion of (simple) continued fraction: $`\sqrt{D}=[a_0,a_1,\mathrm{}].`$ Then $`\alpha _i`$ $`=[a_i,a_{i+1},\mathrm{}]`$ is said to be the $`i`$th complete quotient which could be expressed as
$$\alpha _i=(\sqrt{D}+P_i)/Q_i(P_i,Q_iR),$$
and $`Q_i`$ is said to be the $`i`$th complete denominator. The fraction $`p_i/q_i=[a_0,a_1,\mathrm{},a_i]`$ is named the $`i`$th convergent. There is a positive integer $`\mathrm{}`$ such that $`a_{n+\mathrm{}}=a_n`$ for any $`1n𝐙`$ (The minimal $`\mathrm{}`$ having this property is called the $`\underset{¯}{period}`$ of the continued fraction). So the continued fraction could be written as
$`\sqrt{D}=[a_0,\overline{a_1,\mathrm{},a_{\mathrm{}}}],`$
where the underline part denotes a period, and we have also $`a_\mathrm{}i=a_i`$ for $`0<i<\mathrm{}.`$ Further more, there is a positive $`v𝐙`$ such that $`a_{n+v}=ca_n\mathrm{or}c^1a_n`$ for any $`1n𝐙`$, where $`c𝐅_q^\times `$. The minimal $`v`$ having this property is called the $`\underset{¯}{quasiperiod}`$ ($`v`$ is also the minimal integer$`(>1)`$ such that $`Q_v𝐅_q^\times )`$). We have $`v=\mathrm{}/2\mathrm{or}\mathrm{}`$ (see ).
Theorem 1. Suppose that $`K=k(\sqrt{D})`$ is a real quadratic function field, $`\text{deg}D=2d`$, and $`PR`$ is an irreducible polynomial splitting in $`K`$. Then the ideal class number $`h(D)`$ of $`K`$ has a factor $`h_1`$ satisfying
$`h_1=\text{deg}Q_i/\text{deg}P(1iv),\text{or}h_1d/\text{deg}P.`$
In particular, we have
$`h(D)min_{_{0<i<v}}\{\text{deg}Q_i/\text{deg}P,d/\text{deg}P\}.`$
Theorem 2. Suppose that $`K=k(\sqrt{D})`$ is a real quadratic function field, $`DR`$ is square-free with $`\text{deg}D=2d`$, and $`PR`$ is an irreducible polynomial splitting in $`K`$. Then the ideal class number $`h(D)`$ of $`K`$ has the following lower bound:
(1) If $`D=F^2+c,`$ then $`h(D)\text{deg}F/\text{deg}P;`$
(2) If $`D=(SG)^2+cS,`$ then $`h(D)\text{deg}S/\text{deg}P.`$
where $`c𝐅_q^\times `$, $`GR`$ with $`\text{deg}G1.`$
Theorem 3. Let $`DR`$ be square-free polynomials as the followings, where $`aR𝐅_q,`$ $`A=2a+1`$ is monic, $`m`$ is any positive integer. Assume $`P`$ is an irreducible polynomial in $`R`$ splitting in $`K=k(\sqrt{D})`$. Then the ideal class number $`h(D)`$ of $`K`$ has bound as the following:
(1) If $`D=(A^m+a)^2+A,`$ then $`h(D)\text{deg}A/\text{deg}P;`$
(2) If $`D=(A^ma)^2+A,`$ then $`h(D)\text{deg}A/\text{deg}P;`$
(3) If $`D=(A^m+a+1)^2A,`$ then $`h(D)\text{deg}A/\text{deg}P;`$
(4) If $`D=(A^ma1)^2A,`$ then $`h(D)\text{deg}A/\text{deg}P.`$
Now it is easy to find fields $`K=k(\sqrt{D})`$ with class numbers $`h(D)>1`$. In the following Corollaries 1-3, we assume $`PR`$ is irreducible, $`c𝐅_q^\times ,`$ and $`(F/P)`$ is the quadratic-residue symbol (i.e., $`(F/P)=1`$ or $`1`$ according to $`F`$ is a quadratic residue modulo $`P`$ or not).
Corollary 1. Let $`D=(PG)^2+c,`$ $`\text{deg}G2,`$ $`(c/P)=1,`$ then
$`h(D)\text{deg}(GP)/\text{deg}P>1.`$
Corollary 2. Let $`D=(SHP)^2+cS,`$ $`\text{deg}(S)>\text{deg}(P),`$ $`(cS/P)=1,`$ then
$`h(D)\text{deg}(S)/\text{deg}P>1`$.
Corollary 3. Let $`D`$ be as in Theorem 3 with $`A=SP,`$ $`\text{deg}S2,`$ then
$`h(D)\text{deg}(SP)/\text{deg}P>1.`$
As an example of Corollary 2, we have
$`h((T(T^m+1))^2+(T^m+1))m.`$
(i.e., we take $`P=T,`$ $`S=(T^m+1),`$ $`H=1,`$ $`c=1`$)
II. Lemmas and Proofs of Theorems
First, consider the expansion and property of continued fractions of $`\sqrt{D},`$ where $`DR`$ is a square-free monic polynomial with even degree $`\text{deg}(D)=2d.`$ By we know there exist uniquely determined $`f,rR`$ such that $`D=f^2+r,`$ and $`f`$ is monic, $`\text{deg}f=d,`$ $`\text{deg}r<d.`$
The following process produces the expansion of simple continued fraction $`\sqrt{D}=[a_0,a_1,\mathrm{}]:`$
1. Denote $`D=f^2+r`$ as above. Put $`a_0=f`$, then $`\sqrt{D}=a_0+\sqrt{D}a_0,`$ thus $`\alpha _1=1/(\sqrt{D}a_0)=(\sqrt{D}+a_0)/(Da_0^2)=(\sqrt{D}+P_1)/Q_1`$, where $`P_1=a_0`$, $`Q_1=Da_0^2`$.
2. Now $`\alpha _1=(f+P_1+\sqrt{D}f)/Q_1.`$ Assume $`f+P_1=a_1Q_1+r_1,\text{deg}r_1<\text{deg}Q_1.`$ Then $`\alpha _1=a_1+(\sqrt{D}(fr_1))/Q_1.`$ Thus $`\alpha _2`$ $`=Q_1/(\sqrt{D}(fr_1))`$ $`=(\sqrt{D}+P_2)/Q_2,`$ where $`P_2=(fr_1),`$ $`Q_2=(D(fr_1)^2)/Q_1=(DP_2^2)/Q_1.`$ We see $`P_2R.`$ Since $`P_2=fr_1=a_1Q_1P_1,`$ so $`DP_2^2DP_1^20(\text{mod}Q_1),`$ thus $`Q_2R.`$
Proceed continually, we could obtain the simple continued fraction of $`\sqrt{D}`$ (see ).
Lemma 1<sup></sup>. The Diophantine equation $`X^2DY^2=G`$ has a primary solution if and only if $`G=(1)^iQ_i`$ for some $`0i\mathrm{},`$ where $`Q_i`$ is the $`i`$th complete denominator of the continued fraction of $`\sqrt{D},`$ $`DR`$ is a monic square-free polynomial with even degree, $`GR`$ and $`\text{deg}G<\frac{1}{2}\text{deg}D`$. (A solution $`(X,Y)`$ is primary if $`(X,Y)=1,X,YR`$ ).
Proof of Theorem 1 . Assume $`(P)=\mathrm{}\overline{\mathrm{}},`$ where $`\mathrm{}\overline{\mathrm{}}`$ are prime ideals of $`K`$. Let $`h=h(D)`$ be the ideal class number of $`K`$, then $`\mathrm{}^h`$ is a principal ideal. Suppose that $`mh`$ is the minimal positive integer such that $`\mathrm{}^m`$ is principal, then $`m`$ is a factor of $`h`$. Since $`\{1,\sqrt{D}\}`$ is an integral basis for $`K`$, so we may assume $`\mathrm{}^m=(U+V\sqrt{D})`$ with $`U,VR.`$ Taking norm on both sides, we obtain an equation of ideals of $`k`$ : $`(U^2DV^2)=(P^m).`$ So $`U^2DV^2=cP^m`$ $`(c𝐅_q^\times )`$ since the unit group of $`k`$ (or $`R`$) is just $`𝐅_q^\times .`$
We assert that $`U`$ and $`V`$ must be relatively prime; otherwise, if $`(U,V)=CR`$ is not a constant, put $`U_1=U/C,V_1=V/C,`$ then $`\mathrm{}^m=(U+V\sqrt{D})=(C)(U_1+V_1\sqrt{D})`$, by the uniqueness of factorization of ideals, we must have $`(U_1+V_1\sqrt{D})=\mathrm{}^n`$ for some $`n<m,`$ which contradicts to the minimal assumption of $`m`$.
Thus we know $`(U,V)`$ is a primary solution of $`X^2DY^2=cP^m`$. First assume $`\text{deg}P^m<d,`$ then by Lemma 1 we know that $`cP^m=(1)^iQ_i,`$ $`\text{deg}P^m=\text{deg}Q_i`$ for some $`i`$ with $`\mathrm{\hspace{0.33em}0}i\mathrm{}.`$ Thus we have $`m=\text{deg}Q_i/\text{deg}P`$ for some $`\mathrm{\hspace{0.33em}0}iv`$ by the definition of quasi-period $`v.`$ Secondly, assume $`\text{deg}P^md`$, then we have directly $`md/\text{deg}P.`$
Proof of Theorem 2. (1) It is easy to get the expansion of simple continued fraction $`\sqrt{D}=\sqrt{F^2+c}=[F,\overline{2F/c,2F}],`$ and obtain the set of complete denominators: $`(Q_0,Q_1,Q_2)=(1,c,1).`$ Thus by Theorem 1 we know that $`h(D)`$ has a positive factor $`h_1`$ satisfy $`h_1d/\text{deg}P`$, so $`h(D)\text{deg}F/\text{deg}P`$.
(2) Expand $`\sqrt{D}=\sqrt{(SG)^2+cS}`$ as simple continued fraction: $`\sqrt{D}=[SG,\overline{2G/c,2SG}].`$ Its period is $`\mathrm{}=2`$, complete denominators are just $`(Q_0,Q_1,Q_2)=(1,cS,1).`$ Thus by Theorem 1 we know $`h(D)\text{deg}S/\text{deg}P`$ (Note that $`d=\text{deg}D\text{deg}S`$ now).
Proof of Theorem 3. (1) The polynomial $`D=(A^m+a)^2+A`$ in the theorem has good property which enables us to expand $`\sqrt{D}`$ as a simple continued fraction $`\sqrt{D}=[a_0,\overline{a_1,\mathrm{},a_{\mathrm{}}}]`$. By Theorem 1, we need only to know a quasi-period of the expansion, i.e., $`[a_0,\mathrm{},a_v].`$ It turns out that this quasi-period is quite long and demonstrates rules in three sections, so we will write down it in three sections and list $`a_n`$, $`P_n,Q_n(0nv)`$. We need to distinguish four cases $`m=4t2,4t1,4t,4t+1`$ .
The first section ($`n=0,1`$) and the second section ($`2n4t+1`$): ($`1jt`$ )
$`n`$ $`0`$ $`1`$ $`\mathrm{}`$ $`4j2`$ $`4j1`$ $`4j`$ $`4j+1`$ $`\mathrm{}`$ $`P_n`$ $`0`$ $`A^m+a`$ $`\mathrm{}`$ $`A^m+a+1`$ $`A^ma1`$ $`A^m+a+1`$ $`A^ma1`$ $`\mathrm{}`$ $`Q_n`$ $`1`$ $`A`$ $`\mathrm{}`$ $`2A^{m2j+1}`$ $`A^{2j}`$ $`2A^{m2j}`$ $`A^{2j+1}`$ $`\mathrm{}`$ $`a_n`$ $`A^m+a`$ $`2A^{m1}+1`$ $`\mathrm{}`$ $`A^{2j1}`$ $`2A^{m2j}`$ $`A^{2j}`$ $`2A^{m2j1}`$ $`\mathrm{}`$
The third section is given in four cases according to $`m`$:
(i) For $`m=4t2:`$
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^m+a+1`$ $`A^m+a`$ $`Q_n`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^m1/2`$ $`4A^m4a`$
(ii) For $`m=4t1`$ :
$`n`$ $`4t+2`$ $`4t+3`$ $`4t+4`$ $`4t+5`$ $`P_n`$ $`A^m+a+1`$ $`A^ma1`$ $`A^m+a+1`$ $`A^m+a`$ $`Q_n`$ $`2A^{m2t1}`$ $`A^{2t+2}`$ $`A`$ $`1`$ $`a_n`$ $`A^{m2}`$ $`2A`$ $`2A^{m1}+1`$ $`2A^m2a`$
(iii) For $`m=4t`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^m+a+1`$ $`A^m+a`$ $`Q_n`$ $`2A^{21}`$ $`1/2`$ $`a_n`$ $`A^{m1}1/2`$ $`4A^m4a`$
(iv) For $`m=4t+1`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`4t+4`$ $`4t+5`$ $`P_n`$ $`A^m+a+1`$ $`A^ma1`$ $`A^m+a+1`$ $`A^m+a`$ $`Q_n`$ $`2A^{m2t1}`$ $`A^{2t+2}`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^{m2}`$ $`2A`$ $`2A^{m1}+1`$ $`2A^m2a`$
(2) Similarly to (1), a quasi-period of the simple continued fraction of $`\sqrt{D}`$ is given here.
The first and second sections are combined given as the following: ($`1jt`$ )
$`n`$ $`0`$ $`1`$ $`\mathrm{}`$ $`4j2`$ $`4j1`$ $`4j`$ $`4j+1`$ $`\mathrm{}`$ $`P_n`$ $`0`$ $`A^ma`$ $`\mathrm{}`$ $`A^ma1`$ $`A^m+a+1`$ $`A^ma1`$ $`A^m+a+1`$ $`\mathrm{}`$ $`Q_n`$ $`1`$ $`A`$ $`\mathrm{}`$ $`2A^{m2j+1}`$ $`A^{2j}`$ $`2A^{m2j}`$ $`A^{2j+1}`$ $`\mathrm{}`$ $`a_n`$ $`A^ma`$ $`2A^{m1}1`$ $`\mathrm{}`$ $`A^{2j1}`$ $`2A^{m2j}`$ $`A^{2j}`$ $`2A^{m2j1}`$ $`\mathrm{}`$
The third section is given in four cases:
(i) For $`m=4t2`$ :
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^ma1`$ $`A^ma`$ $`Q_n`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^{m1}1`$ $`4A^m4a`$
(ii) For $`m=4t1`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`4t+4`$ $`4t+5`$ $`P_n`$ $`A^ma1`$ $`A^m+a+1`$ $`A^ma1`$ $`A^ma`$ $`Q_n`$ $`2A^{m2t1}`$ $`A^{2t}`$ $`2A^{m2t}`$ $`1/2`$ $`a_n`$ $`A^{m2}`$ $`2A`$ $`A^{m1}+1`$ $`4A^m+4a`$
(iii) For $`m=4t`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^ma1`$ $`A^ma`$ $`Q_n`$ $`2A`$ $`1/2`$ $`a_n`$ $`2A^{m1}1`$ $`4A^m4a`$
(iv) For $`m=4t+1`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`4t+4`$ $`4t+5`$ $`P_n`$ $`A^ma1`$ $`A^m+a+1`$ $`A^ma1`$ $`A^ma`$ $`Q_n`$ $`2A^{m2t1}`$ $`A^{2t+2}`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^{m2}`$ $`2A`$ $`A^{m1}+1/2`$ $`4A^m+4a`$ .
(3) A quasi-period of the simple continued fraction of $`\sqrt{D}`$ is given here in three sections similarly as in (1). The first and second sections of it are combined given as the following ($`1jt`$ ):
$`n`$ $`0`$ $`1`$ $`\mathrm{}`$ $`4j2`$ $`4j1`$ $`4j`$ $`4j+1`$ $`\mathrm{}`$ $`P_n`$ $`0`$ $`A^m+a+1`$ $`\mathrm{}`$ $`A^m+a`$ $`A^ma`$ $`A^m+a`$ $`A^ma`$ $`\mathrm{}`$ $`Q_n`$ $`1`$ $`A`$ $`\mathrm{}`$ $`2A^{m2j+1}`$ $`A^{2j}`$ $`2A^{m2j}`$ $`A^{2j+1}`$ $`\mathrm{}`$ $`a_n`$ $`A^m+a+1`$ $`2A^{m1}1`$ $`\mathrm{}`$ $`A^{2j1}`$ $`2A^{m2j}`$ $`A^{2j}`$ $`2A^{m2j1}`$ $`\mathrm{}`$
The third section is given in two cases.
(i) For $`m=4t+1`$ or $`m=4t+3`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^m+a`$ $`A^m+a+1`$ $`Q_n`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^{m1}1/2`$ $`4A^m+4a+1`$
(ii) For $`m=4t`$ or $`m=4t+2`$:
$`n`$ $`4t+1`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^ma`$ $`A^m+a`$ $`A^m+a+1`$ $`Q_n`$ $`A^{2t+1}`$ $`2A`$ $`1/2`$ $`a_n`$ $`2A`$ $`A^{m1}1/2`$ $`4A^m4a4`$
(4) We give a quasi-period of the simple continued fraction of $`\sqrt{D}`$ similarly as in (1). The first and second section are combined given as the following: ( $`1jt`$ )
$`n`$ $`0`$ $`1`$ $`\mathrm{}`$ $`4j2`$ $`4j1`$ $`4j`$ $`4j+1`$ $`\mathrm{}`$ $`P_n`$ $`0`$ $`A^ma1`$ $`\mathrm{}`$ $`A^ma`$ $`A^m+a`$ $`A^ma`$ $`A^m+a`$ $`\mathrm{}`$ $`Q_n`$ $`1`$ $`A`$ $`\mathrm{}`$ $`2A^{m2j+1}`$ $`A^{2j}`$ $`2A^{m2j}`$ $`A^{m2j1}`$ $`\mathrm{}`$ $`a_n`$ $`A^ma1`$ $`2A^m+1`$ $`\mathrm{}`$ $`A^{2j1}`$ $`2A^{m2j}`$ $`A^{2j}`$ $`2A^{m2j1}`$ $`\mathrm{}`$
The third section is given in two cases.
(i) For $`m=4t+1`$ or $`m=4t+3`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^ma`$ $`A^ma1`$ $`Q_n`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^{m1}1/2`$ $`4A^m+4a+4`$
(ii) For $`m=4t`$ or $`m=4t+2`$:
$`n`$ $`4t+2`$ $`4t+3`$ $`P_n`$ $`A^ma`$ $`A^ma1`$ $`Q_n`$ $`2A`$ $`1/2`$ $`a_n`$ $`A^{m1}1/2`$ $`4A^m+4a+4`$ .
Consider the above simple continued fractions of $`\sqrt{D}`$ for the four types of $`D`$, and check the complete denominators $`\{Q_n\}(0<n<v)`$ in a quasi-period, we find that the complete denominator having the minimal degree is $`\pm 2A`$ in all the cases. Since deg$`A<\frac{1}{2}`$deg $`D`$, by Theorem 1 we know $`h(D)\text{deg}A/\text{deg}P.`$ This proves Theorem 3.
References
\] R.A.Mollin, Lower bounds for class numbers of real quadratic and biquadratic fields, Proc. Amer. Math. Soc. 101 (1987) 439-444.
\] FENG Keqin, HU Weiqun, On real quadratic function fields of Chowla type with ideal class number one, Proc. Amer. Math. Soc., 127 (1999), 1301-1307.
\] JI Guangheng, LU Hongwen, Proof of class number formula by machine, Sci. in China, (A)28 (1998), 193-200.
\] S. Louboutin, Continued fraction and real quadratic fields, J. Number Theory, 30 (1998), 167-176.
\] LU Hongwen, Gauss’ conjectures on the quadratic number fields, Shanghai Sci. Tech. Pub., 1991.
\] ZHANG Xianke, L. C. Washington, Ideal class-groups and there subgroups of real quadratic fields, Sci. in China, (A)27 (1997), 522-528.
\] E. Artin, Quadratische Körper im Gebiete der höheren Kongruenzen I, II, Math. Z., 19 (1924), 153-206, 207-246.
\] WANG Kunpeng, ZHANG Xianke, The continued fractions connected with quadratic function fields (to appear).
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# Toward scalable quantum computation with cavity QED systems
## I Introduction
In the last years, numerous physical systems have been proposed as possible candidates for the implementation of a quantum computer. The desirable conditions which have to be satisfied are a reliable and easy way to prepare and detect the quantum states of the qubits, the possibility to engineer highly entangled states, the scalability to large number of qubits and a very low decoherence rate . Up to now, experimental implementations have involved linear ion traps , liquid state NMR , and cavity QED systems . In the ion trap case, only the controlled-not (C-NOT) gate between two internal states and the vibrational level of a single ion has been realized , and quantum gates involving two or more ions have not yet realized experimentally. A promising step in this direction is the very recent generation of an entangled state of four ions, even if only with $`57\%`$ fidelity . The status of liquid-state NMR quantum computing is still debated , but the fact that the signal strength becomes exponentially small with the number of qubits makes this proposal certainly not scalable to more than about ten qubits. This explains why the research of new physical implementations of a quantum computer is so active (see and references therein). Here, elaborating on the suggestions of Ref. , we propose to use the Fock’s states $`|0`$ and $`|1`$ of a high-$`Q`$ cavity mode as the two logical states of a qubit. A quantum register of $`N`$ qubits is therefore a collection of $`N`$ identical cavities in which the state of an appropriately chosen cavity mode is within the space spanned by the vacuum and the one photon state. The register transformations are achieved sending off-resonant two-level atoms through the cavities and making them mutually interacting by means of suitable classical fields. With this respect, the present proposal is similar to that of Refs. ; the important difference is that, in these papers, the logical qubits are represented by two circular Rydberg levels of the atoms. In our proposal, the role of atoms and cavity modes are exchanged. In this way, the present scheme becomes scalable in principle. In practice, its scalability can be limited by the spontaneous emission from the Rydberg levels or by other technical limitations, but the present proposal has the advantage that the needed technology is essentially already available to realize some proof-of-principle demonstrations of quantum computation with few qubits. In fact, in the present paper we shall specialize to the case of microwave cavities, for which a high level of quantum state control and engineering has been already experimentally shown . This is the reason why in our explicit calculations we shall consider microwave cavities operating in low-order modes with angular frequency $`\omega `$ in the $`10100`$ GHz range, and Rydberg atoms for which high values of the coupling constant (of the order of $`10^5`$ $`s^1`$) are possible. It is clear however that, in principle, the method can be applied to optical cavities too, in which one can have a miniaturization of the scheme and therefore a faster operation.
Some preliminary results on the possibilities offered by the present cavity QED scheme have been shown in , where some implementations of the C-NOT gate between two cavity modes have been presented. In order to give a clear and exhaustive description, here we shall review the results of , which will be extended and generalized in the present paper.
The outline of the paper is as follows. In Sec. II we review the basic properties of the considered cavity QED system. In Sec. III we show how to implement the universal C-NOT gate between two cavities, while in Sec. IV we shall discuss a different scheme for the implementation of universal two-qubit gates, using an arrangement based on that adopted in the experiment of Ref. . In Sec. V we show how single qubit operations can be realized, while Sec. VI is devoted to the implementation of useful many-qubit universal quantum gates as the Toffoli gate and the encoding and decoding network for quantum error correction schemes. Section VII is for concluding remarks, while the appendix shows the explicit implementation of the Deutsch problem .
## II The system
The interaction of a two-level atom quasi-resonant with a high-Q cavity mode is well described by the time dependent Hamiltonian :
$$(t)=\frac{\mathrm{}\omega }{2}\left[b^{}b+bb^{}\right]+\frac{\mathrm{}\omega _{eg}}{2}\left[|ee||gg|\right]+\mathrm{}\mathrm{\Omega }(t)\left[|eg|b+|ge|b^{}\right],$$
(1)
in which $`b`$ and $`\omega `$ are the annihilation operator and the angular frequency of the cavity mode respectively; $`|e`$ and $`|g`$ are the excited and lower circular Rydberg state and $`\mathrm{}\omega _{eg}`$ is their energy difference. Finally $`\mathrm{\Omega }(t)`$ is the atom-field interaction Rabi frequency which is time dependent because of the atomic motion through the cavity. In particular for a Fabry-Pérot-type cavity, with Gaussian transverse beam profile, we can assume the following continuous variation
$$\mathrm{\Omega }(t)=\mathrm{\Omega }_0e^{(\frac{t}{\tau })^2},$$
(2)
where $`2\tau `$ is the atomic transit time, which depends of course on the inverse of the atomic velocity. For $`|t|\tau `$, i.e. when the atom is outside the cavity, the energy eigenvectors of the system are $`|g|n|g,n`$ and $`|e,n`$, with $`|n`$ the generic Fock state of the cavity mode. Apart for the ground state $`|g,0`$ which remains unchanged, in the presence of the time-dependent interaction, these terms are coupled by photon emission or absorption, and the istantaneous energy eigenstates at fixed time $`t`$ are the dressed states
$$|𝒱_\pm ^{(n)}(t)=\frac{(\delta /2\pm \sqrt{(\delta /2)^2+\mathrm{\Omega }^2(t)(n+1)})|e,n+\mathrm{\Omega }(t)\sqrt{n+1}|g,n+1}{\sqrt{\delta ^2/2+2\mathrm{\Omega }^2(t)(n+1)\pm \delta \sqrt{(\delta /2)^2+\mathrm{\Omega }^2(t)(n+1)}}}$$
(3)
with eigenvalues
$$E_\pm ^{(n)}(t)=\mathrm{}\omega (n+1)\pm \mathrm{}\sqrt{(\delta /2)^2+\mathrm{\Omega }^2(t)(n+1)},$$
(4)
where $`\delta =\omega _{eg}\omega `$ is the atom-cavity detuning. Fig. 1 qualitatively shows the time dependence of the dressed levels of Eq. (4) in the case $`\delta >0`$. Now, if the atom velocity is slow enough, and the system for $`t\tau `$ is prepared in a generic energy eigenstate, then in its time evolution it will adiabatically follow this eigenstate, with negligible transitions toward other states . The exact adiabatic condition can be obtained writing the Schrödinger equation in the basis of the vectors $`|𝒱_\pm ^{(n)}(t)`$, and then neglecting terms coupling the dressed states. The resulting condition is
$$\frac{\dot{\mathrm{\Omega }}(t)\delta \sqrt{n}}{4\left[(\delta /2)^2+\mathrm{\Omega }^2(t)n\right]^{3/2}}1,$$
(5)
which, in the limit $`\mathrm{\Omega }(t)\sqrt{n}/\delta 1`$ becomes equal to that given in Ref. . The general adiabatic condition (5) shows in particular that adiabaticity can be obtained even when $`\mathrm{\Omega }(t)\sqrt{n}/\delta 1`$, provided that $`\dot{\mathrm{\Omega }}(t)`$ is sufficiently small. In the following we shall always work in this adiabatic regime.
## III The C-NOT gate
Domokos et al. have shown in Ref. that, using induced transitions between the dressed states, it is possible to implement a C-NOT gate in which a cavity containing at most one photon is the control qubit and the atom is the target qubit. This idea is the starting point for the implementation of the C-NOT between two cavities we propose here. Ref. considers an atom entering adiabatically the cavity so that the joint atom-cavity state is
$$c_1|g,0+c_2|g,1+c_3|e,0+c_4|e,1.$$
(6)
When the atom is just inside the cavity, a classical field $`S`$ of frequency $`\omega _S`$ equal to the energy difference between the dressed states $`|𝒱_+^{(1)}(t=0)`$ (originating from $`|e,1`$) and $`|𝒱_{}^{(0)}(t=0)`$ (originating from $`|g,1`$) is switched on for a time interval $`2\tau _S`$, so that the following driving Hamiltonian is added to $`(t)`$ of Eq. (1)
$$_S(t)=\mathrm{}\mathrm{\Xi }_0\mathrm{cos}(\omega _St+\phi _S)e^{(\frac{t}{\tau _S})^2}\left[|eg|+|ge|\right],$$
(7)
where $`\phi _S`$ is the phase of the classical field $`S`$ and $`\mathrm{\Xi }_0`$ is the coupling costant, which depends on the dipole moment of the transition $`eg`$ and on the intensity of $`S`$. Appropriately choosing the value of $`\tau _S`$, it is now possible to selectively couple $`S`$ with these dressed states, leaving the other components of the vector state essentially unperturbed. Moreover, with a suitable choice of the intensity $`S`$, it is possible to apply a Rabi $`\pi `$ pulse between the two states. In this way, when the atom exits the cavity, the resulting state vector in the interaction picture, apart for some phase terms, is given by Eq. (6), but with $`|e,1`$ and $`|g,1`$ exchanged. In this way one has realized a C-NOT gate in which, when the cavity (the control qubit) has one photon, the atom undergoes a NOT operation, while when the cavity contains no photons, the atomic state remains unchanged. We shall refer to this gate as the C-NOT(cavity $``$ atom).
In a similar manner, we can build also a C-NOT gate in which the roles of atom and cavity are exchanged. Let us assume in fact to tune the frequency $`\omega _S`$ to the transition between the dressed state $`|𝒱_{}^{(0)}(t=0)`$ and the state $`|g,0`$, and apply again a $`\pi `$ pulse inside the cavity as before. Now, when the atom leaves the cavity, the terms $`|g,0`$ and $`|g,1`$ in the vector state of the system are mutually exchanged with respect to the initial condition Eq. (6). The $`|e,0`$ and $`|e,1`$ components are instead not affected by the interaction with the classical source $`S`$. This means having realized a C-NOT gate in which, when the atom is in the ground state, the cavity states $`|0`$ and $`|1`$ flip, while nothing happens to the cavity state for the atom in the excited state. In analogy with the previous case, we refer to this new gate as C-NOT(atom $``$ cavity).
It is important to note that, differently from the C-NOT(cavity $``$ atom) case, in the C-NOT(atom $``$ cavity) gate the Rabi transition between the original states (i.e. $`|g,0`$ and $`|g,1`$) of the dressed states involved is forbidden by selection rules. Nevertheless, this coupling is realizable when the atom is in the cavity because the vector $`|𝒱_{}^{(0)}(t=0)`$ has also a $`|e,0`$ component. However, since this dependence is mediated by a coefficient which decreases with $`\mathrm{\Omega }_0/\delta `$ (see Eq. (3)), we have to choose a not too small value of this parameter in order to have a significative coupling constant. In particular, in our calculations we have chosen $`\mathrm{\Omega }_0/\delta 10^1`$. With such values for $`\delta `$, it is also possible to have a sufficient frequency separation between the transitions we are interested in and all the other ones. Actually, it is sufficient to set the duration of the classical pulse $`2\tau _S`$ of the order of $`20\mu s`$ to discriminate all the parasitic transitions and optimal results for the resulting quantum operation are achieved for $`\mathrm{\Omega }_0=420`$ kHz, $`\delta =0.18\mathrm{\Omega }_0`$, $`\tau 100\mu s`$ and with $`\mathrm{\Xi }_0=240`$ kHz, $`\tau _S=14\mu s`$ for the C-NOT(cavity $``$ atom), $`\mathrm{\Xi }_0=141.5`$ kHz, $`\tau _S=19\mu s`$ for the C-NOT(atom $``$ cavity). In Fig. 2 we show the time evolution of the dressed states populations for the C-NOT(atom $``$ cavity) for the above choice of parameter values.
For quantum information processing one needs to control not only the level populations, but also the relative phases. In general, during the adiabatic evolution, different dynamical phases for the different components of the vector state are generated. However, it is always possible to correct these phases by an appropriate choice of the field phase $`\phi _S`$ and eventually acting outside the cavity on the atom with suitable Stark electric fields.
We now have all the elements to realize the C-NOT gate between two distinct but identical cavities, $`A`$ and $`B`$, with the first one acting as the control qubit and the second one as the target qubit. The apparatus is sketched in Fig. 3 and it is essentially a physical realization of the logical network shown in Fig. 4. Suppose that the initial states of the two cavities are respectively $`|\varphi _A=\alpha _A|0_A+\beta _A|1_A`$ and $`|\psi _B=\alpha _B|0_B+\beta _B|1_B`$. i) A first atom, $`a_1`$, prepared in the ground state $`|g`$, enters cavity $`A`$, where it undergoes the C-NOT(cavity $``$ atom) transformation realized with the classical field source $`S_A`$, and described above. ii) Then $`a_1`$ leaves $`A`$ and enters cavity $`B`$: here the classical field $`S_B`$ is switched on in order to obtain a C-NOT(atom $``$ cavity) transformation. In the interaction picture and neglecting all the parasitic but controllable phase terms, the state of the total system at this stage is then:
$$\alpha _A|0_A\overline{|\psi }_B|g+\beta _A|1_A|\psi _B|e,$$
(8)
where $`\overline{|\psi }_B`$ is the NOT-conjugate vector of $`|\psi _B`$, that is, $`\beta _B|0_B+\alpha _B|1_B`$. iii) and iii) The atom enters again $`A`$, where it undergoes the C-NOT(cavity $``$ atom) transformation, so that the state of the system becomes
$$\alpha _A|0_A\overline{|\psi }_B|g+\beta _A|1_A|\psi _B|g=\left\{\alpha _A|0_A\overline{|\psi }_B+\beta _A|1_A|\psi _B\right\}|g.$$
(9)
In terms of the notations given in the previous section we shall refer to this gate as the C-NOT-INV($`AB`$) gate, in order to underline that $`B`$ undergoes to a NOT transformation when $`A`$ is in the $`|0_A`$ state, while nothing happens when $`A`$ is in the $`|1_A`$ state. This fact is illustrated in Fig. 4, where in the equivalent gate (II) there is a NOT transformation acting on $`B`$.
The practical realization of step iii), i.e. the return of the atom in the first cavity, is actually more complicated than what it looks from Figs. 3 and 4. The inversion of the motion of atom $`a_1`$, could be realized in principle with an atomic fountain configuration. However this implies having free-fall velocities, which are too slow in order to have the necessary interaction times within the cavities. For this reason we propose to transfer the quantum information from this atom onto a second one of the same type, but travelling in the opposite direction. With this respect, the scheme adopts the “quantum memory” scheme experimentally verified in Ref. . This quantum information transfer is implemented introducing a third cavity, the auxiliary cavity $`M`$ of Fig. 3, which, differently from $`A`$ and $`B`$, is resonant with the $`eg`$ transition. If $`M`$ is prepared in the vacuum state $`|0_M`$, and the transit time of $`a_1`$ is appropriately chosen, then the atomic state component $`|e`$ releases one photon in $`M`$ through a resonant $`\pi `$ Rabi oscillation. After that, the state of the total system (the three cavities and $`a_1`$), using the same notations of Eq. (8), will be
$$\alpha _A|0_A\overline{|\psi }_B|g|0_M+\beta _A|1_A|\psi _B|g|1_M.$$
(10)
Notice that the entanglement of $`a_1`$ with $`A`$ and $`B`$ is now transferred to the auxiliary cavity $`M`$: the state of the atom $`a_1`$ is factorized and it can be neglected from now on. At this stage, a second atom $`a_2`$ is prepared in the ground state $`|g`$ and injected into the apparatus with the same absolute value of velocity of $`a_1`$, but with the opposite direction. Entering $`M`$, it absorbes the photon left by the first atom through a similar $`\pi `$ Rabi oscillation, and the entanglement with the cavities $`A`$ and $`B`$ is transferred from $`M`$ to $`a_2`$:
$$\alpha _A|0_A\overline{|\psi }_B|g|0_M+\beta _A|1_A|\psi _B|e|0_M.$$
(11)
At this stage also the state of the cavity $`M`$ is factorized and therefore the vector state (11) is quantum logically equivalent to that of Eq. (8). In practice, the apparatus described here is essentially an atomic mirror, which permits us to reflect back the atomic state. Finally let us observe that $`a_2`$ has to cross cavity $`B`$ without interacting with it, before it could reach the cavity $`A`$. This result is achieved simply by switching off the classical source $`S_B`$; the adiabatic regime and the off-resonance condition prevents that, apart for the dynamical phase factors, the state could change during the transit of $`a_2`$ within $`B`$. The action of the C-NOT gate has been explicitely described for factorized state only just to simplify the presentation. It is clear that all the steps can be repeated for a generic entangled state of the two cavities.
Assuming the optimal values for the system parameters written above, we have solved numerically the time evolution of the total system. We describe the resulting effective C-NOT gate in the form of a matrix written in the basis of the Fock states of the two cavities, $`\{|0,0;|0,1;|1,0;|1,1\}`$ where $`|n,m=|n_A|m_B`$: The matrix has been “cleaned up” of the spurious phase factors which may appear during the evolution, and which, using the phase of the classical field $`S_B`$ and also appropriate Stark shift electrical fields, can always be suitably adjusted. Within a $`0.1\%`$ error, the optimized C-NOT matrix has the form
$$\left(\begin{array}{ccccccc}0& & e^{i\lambda }& & 0& & 0\\ & & & & & & \\ e^{i\lambda }& & 0& & 0& & 0\\ & & & & & & \\ 0& & 0& & e^{i\lambda }& & 0\\ & & & & & & \\ 0& & 0& & 0& & e^{i\lambda }\end{array}\right),$$
(12)
where the non trivial phase $`\lambda =0.07`$. The overall transformation takes place in a time of the order of one millisecond, which has to be compared with the typical decoherence timescales, that is, the atomic radiative lifetimes and the cavity relaxation times. For circular Rydberg atoms with $`n50`$, the atomic radiative lifetime is of the order of $`30ms`$ and therefore it does not represent a serious problem. The cavity damping times currently realized for microwaves have instead the same order of magnitude (some millisecond). However relaxation times of the order of $`10ms`$ will be hopefully achieved in the near future, and in this case, one would have a perfectly working C-NOT gate between two cavities. It is clear therefore that, for the present implementation of quantum information processing, the main source of decoherence in the microwave domain is just the cavity leakage. If optical cavities are instead considered, also atomic spontaneous emission may represent an important source of decoherence.
The matrix of Eq. (12) is not a pure C-NOT-INV gate, even if it is still an universal two-qubit gate ; in particular it can be transformed into a standard C-NOT gate by adding a single qubit operation on $`B`$, similar to those we shall present in the following section. Moreover, it is also possible to implement the C-NOT($`AB`$) gate (i.e. the one in which the vector component $`|1_A`$ causes the NOT transformation on $`B`$, so that the role of the target and control qubit is exchanged), simply preparing the first atom entering the apparatus in $`|e`$ rather than in $`|g`$, and then proceeding with the same identical steps of the C-NOT-INV case.
The above scheme assumes the possibility of injecting in the apparatus atomic pulses with exactly one atom: this is not experimentally achieved up to now. However, as shown in Ref. , a control of the atom number could be achieved by a modification of the Rydberg atoms preparation technique. The idea is to use an “atom counter” before the circular Rydberg state preparation, so that the preparation of the state $`|e`$ (which can have an efficiency near to $`100\%`$) is applied only when one is sure to have exactly one atom. The atom counter is realized by driving a strong transition and measuring the fluorescence, whose intensity wiil be proportional to the number of atoms. When the beam section contains zero, two or more atoms, it is discarded: the system waits then for a time of the order of a few microseconds to twenty microseconds (depending upon the atomic velocity and the precise length of the atomic beam section) until a fresh section of the beam comes in the laser beam driving the fluorescence. In this way, instead of preparing a random number at a given time, one thus prepares with a high probability a single Rydberg atom after a random delay. The average delay is minimal when the probability to have exactly one atom is maximized. With a Poissonian statistics, the optimal mean number of atoms is $`1`$. The average random delay could be of the order of $`25\mu s`$ in realistic experimental condition. This is short enough at the scale of the cavity field lifetime to play no major role in the proposed scheme.
## IV Quantum phase gate
In the preceding section we have seen how to implement a universal two-qubit gate, the C-NOT gate, between two cavity modes, using induced transitions between dressed states. However, it is possible to realize another universal two-qubit gate, the quantum phase gate (QPG) , between the two cavities, slightly elaborating on the quantum phase gate operating on qubits carried by the Rydberg atom and the two lowest Fock states of a cavity mode, recently demonstrated experimentally . Of course, since both the C-NOT and the QPG are universal quantum gates, it is always possible to implement one of them, by simply supplementing the other one with appropriate one-qubit operations. However, the interesting experimental result of Ref. suggests an alternative physical implementation of quantum logic operations between cavity qubits, which does not involve induced transitions between dressed states, and extends the scheme of to a directly scalable model.
The QPG transformation reads
$$|a,b\mathrm{exp}\left(i\varphi \delta _{a,1}\delta _{b,1}\right)|a,b,$$
(13)
where $`|a`$ and $`b`$ describe the basis states $`|0`$ and $`|1`$ of two generic qubits. This means that the QPG leaves the initial state unchanged unless when both qubits are in state $`|1`$. In Ref. , the QPG of Eq. (13) did not involve levels $`g`$ and $`e`$, but $`i`$ and $`g`$, where $`i`$ is a lower circular Rydberg level, which is uncoupled with the high-Q cavity. In this way, the gate of Eq. (13) in the case $`\varphi =\pi `$ can be realized by setting the atomic transition $`ge`$ perfectly at resonance with the relevant cavity mode (by appropriately Stark-shifting the atomic levels inside the cavity), and by selecting the atomic velocity so that the atom undergoes a complete $`2\pi `$ Rabi pulse when crossing the cavity. In fact, at resonance, such a pulse transforms the state $`|g,1`$ into $`e^{i\pi }|g,1`$, while nothing happens if the atoms is in $`i`$ or the cavity is in the vacuum state. In , the possibility to tune the phase $`\varphi `$ over a large range, by slightly detuning the cavity mode from the $`ge`$ transition, has been shown, but we shall not consider this possibility here.
We now show that this atom-cavity QPG can be used to realize a QPG between two cavity modes, by considering an arrangement very similar to that of the C-NOT gate shown in Fig. 3. The cavities $`A`$ and $`B`$ are again the two qubits, while $`M`$ is again the auxiliary cavity needed to “reflect” the atom and disentangle it. The two classical sources inside the cavities $`S_A`$ and $`S_B`$ are no more needed, while we consider the possibility to apply Stark shift electric fields inside the cavities, in order to tune the $`ge`$ transition in and out of resonance from the cavity mode. The scheme of the QPG implementation is shown in Fig. 5 and involves only two atom crossings, as in the C-NOT gate of the preceding section, and three $`\pi /2`$ pulses between the $`i`$ and $`g`$ levels (the Hadamard gates $`H`$ of Fig. 5), which can be realized with resonant classical microwave sources applied between the high-Q cavities.
Let us assume a generic state of the two cavity qubits
$$|\psi =a_0|00+a_1|01+a_2|10+a_3|11$$
(14)
and that a first atom, initially prepared in state $`i`$, is subject to a $`\pi /2`$ pulse, so to enter cavity $`A`$ in state $`\left(|i+|g\right)/\sqrt{2}`$. The cavity mode is perfectly resonant with the $`ge`$ transition and the atom velocity is selected so that the atom undergoes a $`2\pi `$ Rabi pulse if it is in state $`g`$ and the cavity contains one photon (the QPG of Ref. ). The resulting state at the exit of cavity $`A`$ is
$$\frac{|i|\psi }{\sqrt{2}}+\frac{|g}{\sqrt{2}}\left(a_0|00+a_1|01a_2|10a_3|11\right).$$
(15)
Then the atom undergoes another resonant $`\pi /2`$ pulse on the $`ig`$ transition and the state of the system becomes
$$|i\left(a_0|00+a_1|01\right)+|g\left(a_2|10+a_3|11\right).$$
(16)
Then the atom crosses cavity $`B`$, where it is subjected again to the atom-cavity QPG as in cavity $`A`$, so that the state of the system becomes
$$|i\left(a_0|00+a_1|01\right)+|g\left(a_2|10a_3|11\right).$$
(17)
At this point, as in the C-NOT case of the preceding section, in order to realize the final transformation and disentangle the atom from the cavities, one has to “reflect” it. This is again achieved with the “atomic mirror” scheme, i.e., using the auxiliary cavity $`M`$ of Fig. 5, that acts as a quantum memory and is able to transfer the entanglement from the first atom to a second atom, which is crossing the apparatus with the same absolute velocity but in the opposite direction (see Eqs. (10) and (11)).
The second atom is then subjected to a $`\pi /2`$ pulse before entering the cavities and the state of Eq. (17) becomes
$$\frac{|i}{\sqrt{2}}\left(a_0|00+a_1|01+a_2|10a_3|11\right)+\frac{|g}{\sqrt{2}}\left(a_0|00+a_1|01a_2|10+a_3|11\right).$$
(18)
The last step is the QPG between the atom and cavity $`A`$, i.e., the atom has to cross cavity $`B`$ undisturbed (this is achieved by strongly detuning the $`ge`$ transition with a Stark shift field) and then has to undergo another full Rabi cycle in cavity $`A`$. The final state is
$$\frac{|i+|g}{\sqrt{2}}\left(a_0|00+a_1|01+a_2|10a_3|11\right),$$
(19)
which is desired result, corresponding to a QPG between cavities $`A`$ and $`B`$ with conditional phase shift $`\varphi =\pi `$, and with a disentangled atom.
## V One qubit operations
One qubit operations are straightforward to implement on qubits represented by two internal atomic states because it amounts to apply suitable Rabi pulses. This task is less trivial for bosonic degrees of freedom as our cavity modes, because the two lowest Fock states for example are coupled to the more excited ones. The most practical solution is to implement one-qubit operations on the two lowest Fock states sending again atoms through the cavity. To be more specific, one has to send an atom prepared in the ground state $`|g`$ through the cavity, with the classical field $`S`$ tuned at the frequency corresponding to the transition between the states $`|𝒱_{}^{(0)}(t=0)`$ and $`|g,0`$. If one sets the time duration and the intensity of the classical source $`S`$ as in the case of the C-NOT(cavity $``$ atom), i.e. such to realize a $`\pi `$ pulse between the selected levels, one implements a “not-phase” gate, which, in the canonical basis $`\{|0,|1\}`$, is described by the following matrix
$$N(\theta )=\left(\begin{array}{ccc}0& & e^{i\theta }\\ & & \\ e^{i\theta }& & 0\end{array}\right),$$
(20)
where $`\theta `$ depends linearly on the phase $`\phi _S`$ of the classical field of Eq. (7) and it is therefore easily controllable (this scheme is simply a part of the C-NOT-INV($`AB`$) gate presented before). If otherwise the atom inside the cavity undergoes a $`\pi /2`$ instead of a $`\pi `$ pulse, one realizes the “Hadamard-phase” gate
$$H(\theta ^{})=\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}1& & e^{i\theta ^{}}\\ & & \\ e^{i\theta ^{}}& & 1\end{array}\right),$$
(21)
where also $`\theta ^{}`$ depends linearly on the classical field phase $`\phi _S`$ and is therefore controllable. $`N(\theta )`$ and $`H(\theta ^{})`$ can be used to build the more general one-qubit operation and therefore, together with C-NOT-INV($`AB`$), form an universal set of gates.
Note that the not-phase gate can be used also for an alternative realization of the C-NOT-INV($`AB`$) gate between two cavities. In the scheme described in the preceding section, one needs the auxiliary cavity M and the second atom crossing in the opposite direction in order to disentangle the first atom from the cavities. One could simplify this last stage (step iii) and iv of the preceding section) by applying an exact $`\pi /2`$ pulse when the atom has just left the second cavity and the state of the system is that of Eq. (8). The total state becomes
$`\alpha _A|0_A\overline{|\psi }_B{\displaystyle \frac{|g+|e}{\sqrt{2}}}+\beta _A|1_A|\psi _B{\displaystyle \frac{|g|e}{\sqrt{2}}}=`$ (23)
$`\left\{\alpha _A|0_A\overline{|\psi }_B+\beta _A|1_A|\psi _B\right\}|g/\sqrt{2}+\left\{\alpha _A|0_A\overline{|\psi }_B\beta _A|1_A|\psi _B\right\}|e/\sqrt{2}.`$
If now the atom is detected by a state-sensitive detector and the $`|g`$ state is detected, the two cavities are projected on $`\alpha _A|0_A\overline{|\psi }_B+\beta _A|1_A|\psi _B`$ and one has implemented just the desired C-NOT-INV gate. On the contrary, if the atom is found in the excited state $`|e`$, the state of $`A`$ and $`B`$ becomes $`\alpha _A|0_A\overline{|\psi }_B\beta _A|1_A|\psi _B`$ and the C-NOT-INV gate is obtained once that cavity $`A`$ is subject to a $`\pi `$-phase shift, which can be realized by means of two not-phase gates $`N(\theta )`$, the first with $`\theta =\pi /2`$ and the second with $`\theta =0`$. In this way the atom is disentangled by the measurement. However, the practical application of this scheme is seriously limited by the quantum efficiency of atomic detectors, which is usually far from $`100\%`$.
## VI Many-qubit gates
### A The Toffoli gate
We have shown how to implement a set of universal quantum gates with the proposed cavity QED scheme. Therefore, in principle, the most general quantum operation involving $`n`$ qubits can be realized in terms of the one and two-qubit operations described above. This decomposition however implies a degree of network complexity, and a number of resources and steps which is rapidly increasing with the number of qubits $`n`$. One of the main advantages of the present proposal is that it is particularly suited for the efficient implementation of many-qubit quantum gates, which, in many cases, can be realized with the same number of steps of the two-qubit C-NOT gate of Section III.
A particularly clear example of the possibilities of the proposed scheme is provided by the Toffoli gate
$$|x_A|y_B|z_C|x_A|y_B|\left[z+(xy)\right]_{mod2}_C,$$
(24)
in which the target qubit $`C`$ is controlled by the first two, $`A`$ and $`B`$. The effect of the Toffoli gate on the generic three qubit state
$$|\mathrm{\Psi }_0=\alpha _1|000+\alpha _2|001+\alpha _3|010+\alpha _4|011+\alpha _5|100+\alpha _6|101+\alpha _7|110+\alpha _8|111$$
(25)
($`|n,m,l=|n_A|m_B|l_C`$ are the tensor product of the cavity mode Fock states) is to exchange the last two components $`|110`$ and $`|111`$. The implementation of this gate needs the same arrangement of aligned cavities crossed by Rydberg atoms used for the C-NOT gate of Fig. $`3`$, except that now one has three cavity qubits (with the corresponding classical sources $`S_A`$, $`S_B`$ and $`S_C`$) instead of two. The auxiliary cavity $`M`$ is again needed for the atomic mirror scheme used to disentangle the atom. The atom is initially prepared in state $`|g`$ and when it is in the first cavity $`A`$, is subject to a $`\pi `$ pulse between the dressed state $`|𝒱_+^{(0)}(0)`$ and $`|g,0`$. This pulse creates atom-cavity entanglement and the state of the total system becomes
$$\left[\alpha _1|000+\alpha _2|001+\alpha _3|010+\alpha _4|011\right]|e+\left[\alpha _5|100+\alpha _6|101+\alpha _7|110+\alpha _8|111\right]|g.$$
(26)
Then, when the atom reaches the second cavity $`B`$, it undergoes another $`\pi `$ pulse, at the new frequency $`\omega _2`$ corresponding to the transition between $`|g,0`$ and $`|𝒱_+^{(2)}(0)`$, so that the transformation
$$|0_B|g|2_B|e,$$
(27)
is realized. This means temporarily leaving the logical subspace, even though this allows us to realize a significant simplification of the scheme. The state after this second step is therefore
$$\left[\alpha _1|000+\alpha _2|001+\alpha _3|010+\alpha _4|011+\alpha _5|120+\alpha _6|121\right]|e+\left[\alpha _7|110+\alpha _8|111\right]|g.$$
(28)
When the atom enters in $`C`$, the classical field $`S_C`$ is applied so to realize the C-NOT(atom $``$ cavity $`C`$) of Sec. III and the state of Eq. (28) becomes
$$\left[\alpha _1|000+\alpha _2|001+\alpha _3|010+\alpha _4|011+\alpha _5|120+\alpha _6|121\right]|e+\left[\alpha _7|111+\alpha _8|110\right]|g.$$
(29)
At this point one has to disentangle the atom from the three cavities and also to adjust the state components in which the cavity $`C`$ contains two photons. Both problems can be solved using again the auxiliary cavity $`M`$ and a second atom crossing the apparatus in the opposite direction as in the “atomic mirror” configuration of Sec. III. The cavity $`M`$ transfers the entanglement with the cavities from the first to the second atom, which is not subject to any classical pulse in $`C`$. Then the second atom enters $`B`$, where it undergoes a $`\pi `$ pulse at the frequency $`\omega _2`$, which simply inverts the transformation of Eq. (27) (thanks to the fact that no $`|0_B|g`$ term is present), correcting in this way the terms of Eq. (29) in which the second cavity contains two photons. As a consequence, the state of the atom-cavities system becomes
$$\left[\alpha _1|000+\alpha _2|001+\alpha _3|010+\alpha _4|011\right]|e+\left[\alpha _5|100+\alpha _6|101+\alpha _7|111+\alpha _8|110\right]|g.$$
(30)
Finally the atom enters $`A`$, where it is subjected to a $`\pi `$ pulse resonant with the transition $`|𝒱_+^{(0)}(0)|g,0`$, exchanging $`|0_A|g`$ with $`|0_A|e`$, so that the second atom is disentangled from the cavities and one gets the desired generic output of a Toffoli gate, i.e.,
$$\left[\alpha _1|000+\alpha _2|001+\alpha _3|010+\alpha _4|011+\alpha _5|100+\alpha _6|101+\alpha _7|111+\alpha _8|110\right]|g.$$
(31)
Notice that in this way we have implemented the Toffoli gate with two atoms only, as in the C-NOT gate of Sec. III. Moreover this scheme can be easily extended to the case of $`n4`$ cavity qubits, for the implementation of the n-qubit generalization of the Toffoli gate. We need only two atoms crossing the aligned cavities in opposite directions also in this more general case. The pulse sequence is similar to that discussed above: both atoms undergo a $`\pi `$ pulse resonant with the transition $`|𝒱_+^{(0)}(0)|g,0`$ in the first cavity, while in the following $`n2`$ cavities they are submitted to a $`\pi `$ pulse at the frequency $`\omega _2`$. In the last cavity, the target qubit, the first atom experiences a C-NOT(atom $``$ cavity) while the second atom crosses it undisturbed. In the scheme proposed here, the target qubit is necessarily the last cavity. However, it is always possible to realize the n-qubit generalized Toffoli gate with the target qubit in a generic position of the string of cavities, by simply applying a two-qubit operation to the above scheme. In our scheme this means using four atoms at most, and in any case this is much more convenient than realizing this generic n-qubit gate from one and two qubits operation.
### B Encoding and decoding in quantum error correction codes
Other examples for which the present cavity QED scheme offers the possibility of an efficient implementation of operations involving many qubits are the encoding and decoding processes used in quantum error correction schemes . Errors in quantum information processing are due to the interaction with uncontrolled degrees of freedom (the environment), yielding an entanglement of the quantum state of the register with some environmental states. The main idea of quantum error correction is to combat “bad” entanglement with “good” entanglement, that is, protecting quantum information by storing it not in a single qubit but in an entangled state of $`n`$ qubits. This is the encoding process; if the error rate is not too large, it is possible to recover the original quantum information by using a suitable decoding procedure, because the eventual data corruption can be revealed by a measurement on the auxiliary qubits and information can be finally restored with single qubit operations. The more general one qubit error (flip error, phase error, or a combination of the two) can always be corrected using a five-qubit encoding and decoding procedure . However, if one considers one specific form of error only, three-qubit encoding is sufficient for the implementation of quantum error correction codes. For simplicity, let us consider this latter case.
Let us assume that we want to protect a generic state $`\alpha |0_A+\beta |1_A`$ of the cavity $`A`$. For the encoding process one needs two other ancilla qubits, cavity $`B`$ and $`C`$, and one has to realize the following transformation into a maximally entangled, GHZ state of three cavities
$$\left[\alpha |0_A+\beta |1_A\right]|0_B|0_C\alpha |000+\beta |111.$$
(32)
This encoding process can be realized using a scheme analogous to those discussed above for the C-NOT and Toffoli gates. Again, only two atoms are needed, with the second one, crossing the aligned cavities in opposite direction, which serves the purpose of disentangling the first atom, with the help of the auxiliary cavity $`M`$ in the “atomic mirror” scheme described in Sec. III.
The initial state of the system is
$$\left[\alpha |0_A+\beta |1_A\right]|0_B|0_C|e\left[\alpha |000+\beta |100\right]|e.$$
(33)
When the atom enters $`A`$ it undergoes the C-NOT(cavity $`A`$ $``$ atom) described by Eq. (5); when the atom arrives in $`B`$, the classical field $`S_B`$ is switched on in order to realize the C-NOT(atom $``$ cavity $`B`$) transformation (see Sec. III) and the same C-NOT(atom $``$ cavity $`C`$) operation is applied when the atom is in $`C`$. Using Eqs. (5) and (6), one can show that the state of the total system becomes
$$\left[\alpha |000|e+\beta |111|g\right].$$
(34)
Except for the entanglement with the atom, the state of $`A`$,$`B`$ and $`C`$ is of the desired form, and therefore the situation is analogous to that of the C-NOT gate of Sec. III (see Eq. (8)). Atom disentanglement can be obtained using again the atomic mirror scheme of Sec. III, or eventually, the atomic detection scheme discussed in Sec. IV, which is however seriously limited by detector inefficiencies. The logical transformations we have implemented in this section is schematically described in Fig. 6 (the dotted line box represents the atomic mirror)
Decoding is obtained by repeating exactly the same procedure adopted for encoding the state and assuming as initial condition for the qubits set the encoded state (32). It is straightforward to check that this amounts to realize the inverse transformation
$$\alpha |000+\beta |111\alpha |000+\beta |100=\left[\alpha |0_A+\beta |1_A\right]|0_B|0_C.$$
(35)
It is also easy to see that the encoding scheme described here can be extended for the controlled preparation of maximally entangled states of $`n`$ cavities. One has to consider $`n`$ cavities (plus the auxiliary cavity $`M`$) and, as in the $`n=3`$ case, one needs only two atoms crossing in opposite directions the $`n+1`$ aligned cavities. In analogy with the description above, in the first cavity the two atoms undergo the C-NOT(cavity $`A`$ $``$ atom) transformation of Eq. (5), while in all the others $`n1`$ cavities the first atom undergoes to a C-NOT(atom $``$ cavity) gate and the second one crosses them with no classical field applied. In this way, thanks to the disentanglement action of the second atom, the following maximally entangled state of $`n`$ cavities is prepared
$$\left[\alpha |0_1+\beta |1_1\right]|0_2\mathrm{}|0_n\alpha |00\mathrm{}0+\beta |11\mathrm{}1.$$
(36)
## VII Conclusions
In this paper we have presented a scheme for implementing quantum logic operations within a cavity QED configuration. The quantum register is composed by a series of high-Q cavities and information is encoded in the two lowest cavity Fock states. Both the preparation and the detection of the quantum state of individual qubits, which is an essential ingredient in quantum algorithms, can be easily performed. In particular the detection of the two Fock states could be even performed in a quantum non-demolition way, as recently demonstrated . Both one-qubit and two-qubit operations can be performed sending appropriately prepared atoms through the cavities. An important advantage of the scheme is that it is particularly suitable for the direct implementation of some useful many-qubit quantum gates, as for example the Toffoli gate and its n-qubit generalization, or the encoding-decoding network of quantum error correction codes. The scheme could be implemented in a generic cavity QED scheme, even if here we have specialized to the case of microwave cavities and circular Rydberg atoms, for which entanglement manipulation has been already demonstrated . For example, the same scheme could be adapted to the optical frequency domain, by using high-Q optical cavities, as for example the whispering gallery modes of silica microspheres , in which one can have a miniaturization of the scheme and, therefore, a faster gate operation. The scheme proposed here is in principle scalable, even if in practice its scalability will be limited by various facts. In the microwave case one is limited by the spontaneous emission from the circular Rydberg levels ($`30`$ ms) and by the fact that all the apparatus has to be cooled at cryogenic temperatures to avoid the thermal radiation. In the optical case cooling is no more needed and also the limitations due to the spontaneous emission could be avoided in principle by using atomic $`\mathrm{\Lambda }`$ transitions and adiabatic passage through a dark state . However in the microwave case, the proposed quantum gates could be implemented using available technology, and therefore proof of principle demonstrations of quantum computation with, say, ten qubits could be achievable. Instead, even if whispering gallery modes in microspheres with $`Q10^9`$ has been already realized , entanglement manipulation in this case have not been experimentally demonstrated yet.
Recently, a similar proposal of using the two lowest Fock states of a high-Q cavity as logical qubit has been presented, involving an engineered network of defects in a photonic band-gap material . This proposal is promising with respect to scalability since, using atoms travelling along engineered waveguides in the photonic band-gap material, spontaneous emission could be completely eliminated. However, even if this proposal is promising in terms of the technological realization, in this case, as in the silica microsphere case, entanglement manipulation has not yet been experimentally demonstrated.
From a general point of view the scheme proposed here is analogous to the linear ion trap scheme, except that now the high-Q cavities play the role of the ions, and the atoms, and not the collective center-of-mass motion, play the role of the quantum bus. At first sight, it may seem unpratical to reverse the role of atoms and photons as we have done here, since the common wisdom is that atoms and ions are suitable for storing informations while photons are best suited to transfer quantum information between different sites. However, the practical implementation of quantum algorithms on linear ion traps is presently limited by the heating of the center-of mass motion . On the contrary, in the case of photonic qubits discussed here, once the limitations due to the spontaneous emission are eliminated (as in the photonic band-gap case and in the microsphere case with dark state transitions), the scheme is then only limited by the decoherence due to the finite $`Q`$ of the cavities, which could reach however values of the order of $`10^{10}`$, allowing therefore a sufficient number of gate operations. Moreover, in view of the fact that photons are in any case the best tools for quantum information transport, it may be nonetheless useful to have schemes able to process and temporarily store quantum information using photons.
## A A realization of the Deutsch algorithm
As an example of a simple algorithm which can be implemented with the quantum gates presented before we consider the Deutsch algorithm in the improved version of Ref. .
Consider a generic Boolean function $`f(x)`$ mapping $`\{0,1\}`$ into $`\{0,1\}`$. There are four different possibilities, the two constant functions $`f_1(x)=0`$, $`f_2(x)=1`$, and the two balanced functions $`f_3(x)=x`$, $`f_4(x)=1x`$. Using classical algorithms, the distinction between these two different classes of functions necessarily requires that both values of $`f(x)`$ have to be evaluated. On the contrary, the Deutsch quantum algorithm solves the problem in just one step. We have to consider the two cavities quantum circuit of Fig. 7. Initially both cavities are in the vacuum state; then they are submitted to the Hadamard-phase transform of equation (21), with phases equal to $`0`$ and $`\pi `$ respectively (step $`a)`$ of Fig. 7). As shown in Sec. IV, this transformations can be implemented with a single atom prepared in the ground state $`|g`$ crossing both cavities. At this point the system undergoes the transformation of step $`b)`$, namely the following “$`f`$” gate:
$$|x_A|y_B|x_A|\left[y+f_i(x)\right]_{mod2}_B,$$
(A1)
where $`f_i(x)`$, $`i=1,2,3,4`$ are the four functions defined above, and $`\left[y+f_i(x)\right]_{mod2}`$ means addition modulo $`2`$. For $`i=1`$, this means that nothing happens to the system; on the contrary for $`i=2`$ $`A`$ remains unchanged but the second qubit in cavity $`B`$ undergoes a NOT transformation. Finally it is easy to verify that for $`i=3`$ and for $`i=4`$ the $`f`$ gate of Eq. (A1) is equivalent to the C-NOT gate and to the C-NOT-INV gate respectively. In the preceding sections we have seen how to implement all these transformations. At stage $`c)`$ we have then to implement another Hadamard transformation with zero phase on $`A`$ and measure the state of this qubit. According to the logical network of Fig. 7, it is possible to show that if the $`f_i(x)`$ of the $`f`$ gate is a constant function, then the cavity $`A`$ must be found in $`|0_A`$, otherwise, if $`f_i(x)`$ is a balanced function, the cavity $`A`$ will be in $`|1_A`$: in this way, one can establish if the function $`f_i`$ is constant or balanced using a single function evaluation. The physical implementation of the Deutsch problem in terms of the cavities is sketched in Fig. 8. In the case of the constant function it is possible to only use two atoms, because for $`i=2`$ one atom is needed for the NOT transformation on the cavity $`B`$ of stage $`b)`$ and another atom is needed for the Hadamard transform on the cavity $`A`$ at step $`c)`$. For the balanced functions, the number of atoms is instead equal to four, because, besides the two atoms for the implementation of the Hadamard transformations, one has to use two atoms for the C-NOT (if $`i=3`$) or the C-NOT-INV (if $`i=4`$).
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# The classical massive Thirring system revisited
## 1. Introduction
Ever since its publication in 1958, the Thirring model kept its fascination as is witnessed by the incredible amount of attention paid to it since then and by the interest it continues to generate (see, e.g., for a recent review). In the present paper we are not concerned with its importance as a solvable quantum field theory model but rather restrict our attention to its complete integrability aspects from a classical point of view. Thirring’s classical $`(1+1)`$-dimensional model equations in appropriate light cone coordinates, and after appropriate rescaling of the mass and coupling constant parameters, etc., can be cast in the form
$`iu_x+2v+2|v|^2u`$ $`=0,`$
$`iv_t+2u+2|u|^2v`$ $`=0.`$ (1.1)
Formal integrability of (1.1) was originally established by Mikhailov in 1976 by establishing a corresponding commutator representation (cf. (2.8) and (2.9)). In fact, one can replace (1.1) by a more general system, without identifying $`u^{}`$ and $`v^{}`$ with the complex conjugates $`\overline{u}`$ and $`\overline{v}`$ of $`u`$ and $`v`$, respectively,
$`iu_x+2v+2vv^{}u`$ $`=0,`$
$`iu_x^{}+2v^{}+2vv^{}u^{}`$ $`=0,`$ (1.2)
$`iv_t+2u+2uu^{}v`$ $`=0,`$
$`iv_t^{}+2u^{}+2uu^{}v^{}`$ $`=0,`$
without losing formal integrability, and we will actually investigate (1.2) rather than (1.1). Both (1.1) and (1.2) have been studied by numerous authors, who derived the inverse scattering approach , , , , considered soliton solutions , , , , , , , , investigated Bäcklund transformations and close connections with other integrable equations (especially, the sine-Gordon equation) , , , , , , , , , and considered monodromy deformations .
In the present paper we focus on algebro-geometric solutions of the classical massive Thirring system (1.2). The first attempt to derive algebro-geometric solutions of (1.1) is due to Date in 1978 and almost simultaneously to Prikarpatskii and Holod (see also ). Both papers are remarkably similar in strategy, in fact, they are nearly identical. In particular, both discuss theta function representations for symmetric functions of appropriate symmetric functions associated with auxiliary divisors, but neither derives explicit theta function representations of $`u`$ and $`v`$. The first theta function representations of $`u`$, $`v`$, $`u^{}`$, $`v^{}`$ for the general massive Thirring system (1.2) were derived by Bikbaev , however, with insufficient care paid to details. (In fact, his terms $`e^w`$ and $`e^{\widehat{w}}`$ on p. 581 are not defined, and in his formula (29), $`(x,t)`$-dependent terms are missing.) More recently, algebro-geometric solutions of (1.1) were also briefly considered by Wisse , again without explicitly deriving theta function representations for $`u`$ and $`v`$.
In Section 2 we follow Date’s explicit realization of Mikhailov’s commutator representation in terms of polynomials in the spectral parameter. In Section 3 we develop the basic algebro-geometric formalism for (1.2), and from that point on we deviate from previous investigations and focus on a different approach based on the solution $`\varphi `$ of a Riccati-type equation associated with the Thirring system (1.2). We consider Dubrovin-type equations for auxiliary divisors and define the Baker–Akhiezer vector associated with the system (1.2) in terms of the fundamental function $`\varphi `$ on $`𝒦_n`$, the underlying hyperelliptic curve of genus $`n_0`$. We also study the algebro-geometric initial value problem in detail. Our principal results, the theta function representations of $`u`$, $`v`$, $`u^{}`$, $`v^{}`$, and $`\varphi `$ are derived in detail in Section 4. Finally, Appendix A collects some basic results on compact Riemann surfaces and introduces the terminology freely used in Sections 3 and 4.
## 2. The basic polynomial setup
In this section we start from Mikhailov’s commutator representation of the classical massive Thirring system in a form used by Date (see also , , which contain similar material) in his analysis of quasi periodic solutions of this model.
Assuming $`u,v,u^{},v^{}:^2`$ to satisfy
$`u(,t),u^{}(,t)C^1(),v(,t),v^{}(,t)C^{\mathrm{}}(),t,`$
$`u(x,),u^{}(x,)C^1(),_x^kv(x,),_x^kv^{}(x,)C^1(),k_0,x,`$ (2.1)
we introduce the $`2\times 2`$ matrices
$`U(\zeta ,x,t)`$ $`=i\left(\begin{array}{cc}zv(x,t)v^{}(x,t)& 2\zeta v(x,t)\\ 2\zeta v^{}(x,t)& z+v(x,t)v^{}(x,t)\end{array}\right),`$ (2.2)
$`V_{n+1}(\zeta ,x,t)`$ $`=i\left(\begin{array}{cc}G_{n+1}(z,x,t)& \zeta F_n(z,x,t)\\ \zeta H_n(z,x,t)& G_{n+1}(z,x,t)\end{array}\right),n_0,`$ (2.3)
$`\stackrel{~}{V}(\zeta ,x,t)`$ $`=i\left(\begin{array}{cc}z^1u(x,t)u^{}(x,t)& 2\zeta ^1u(x,t)\\ 2\zeta ^1u^{}(x,t)& z^1+u(x,t)u^{}(x,t)\end{array}\right),`$ (2.4)
$`\zeta \{0\},z=\zeta ^2,(x,t)^2,`$
where $`F_n(z,x,t)`$, $`H_n(z,x,t)`$, and $`G_{n+1}(z,x,t)`$ are polynomials with respect to $`z`$ of degree $`n`$ and $`n+1`$, respectively, that is, they are of the type
$`F_n(z,x,t)`$ $`={\displaystyle \underset{j=0}{\overset{n}{}}}f_{nj}(x,t)z^j=f_0(x,t){\displaystyle \underset{j=1}{\overset{n}{}}}(z\mu _j(x,t)),`$ (2.5)
$`G_{n+1}(z,x,t)`$ $`={\displaystyle \underset{j=0}{\overset{n+1}{}}}g_{n+1j}(x,t)z^j,g_0(x,t)=1,`$ (2.6)
$`H_n(z,x,t)`$ $`={\displaystyle \underset{j=0}{\overset{n}{}}}h_{nj}(x,t)z^j=h_0(x,t){\displaystyle \underset{j=1}{\overset{n}{}}}(z\nu _j(x,t)).`$ (2.7)
The classical massive Thirring system is then defined by demanding the zero-curvature representation
$`V_{n+1,x}(\zeta ,x,t)+[U(\zeta ,x,t),V_{n+1}(\zeta ,x,t)]`$ $`=0,(\zeta ,x,t)\{0\}\times ^2,`$ (2.8)
$`V_{n+1,t}(\zeta ,x,t)+[\stackrel{~}{V}(\zeta ,x,t),V_{n+1}(\zeta ,x,t)]`$ $`=0,(\zeta ,x,t)\{0\}\times ^2.`$ (2.9)
Explicitly, equations (2.8) and (2.9) yield
$`F_{n,x}(z,x,t)`$ $`=2i(v(x,t)v^{}(x,t)z)F_n(z,x,t)+4iv(x,t)G_{n+1}(z,x,t),`$ (2.10)
$`G_{n+1,x}(z,x,t)`$ $`=2izv^{}(x,t)F_n(z,x,t)+2izv(x,t)H_n(z,x,t),`$ (2.11)
$`H_{n,x}(z,x,t)`$ $`=2i(v(x,t)v^{}(x,t)z)H_n(z,x,t)+4iv^{}(x,t)G_{n+1}(z,x,t),`$ (2.12)
$`\begin{array}{cc}\hfill F_{n,t}(z,x,t)& =2i(u(x,t)u^{}(x,t)z^1)F_n(z,x,t)\hfill \\ & +4iz^1u(x,t)G_{n+1}(z,x,t),\hfill \end{array}`$ (2.13)
$`G_{n+1,t}(z,x,t)`$ $`=2iu^{}(x,t)F_n(z,x,t)+2iu(x,t)H_n(z,x,t),`$ (2.14)
$`\begin{array}{cc}\hfill H_{n,t}(z,x,t)& =2i(u(x,t)u^{}(x,t)z^1)H_n(z,x,t)\hfill \\ & +4iz^1u^{}(x,t)G_{n+1}(z,x,t).\hfill \end{array}`$ (2.15)
By (2.10)–(2.14) one infers that
$$\left(G_{n+1}^2zF_nH_n\right)_x=\left(G_{n+1}^2zF_nH_n\right)_t=0$$
(2.16)
and hence
$$G_{n+1}(z,x,t)^2zF_n(z,x,t)H_n(z,x,t)=R_{2n+2}(z),$$
(2.17)
where the integration constant $`R_{2n+2}(z)`$ is a monic polynomial in $`z`$ of degree $`2n+2`$, that is,
$$R_{2n+2}(z)=\underset{m=0}{\overset{2n+1}{}}(zE_m),\{E_m\}_{m=0,\mathrm{},2n+1},$$
(2.18)
since we chose $`g_0=1`$. Moreover, (2.17) implies
$$g_{n+1}(x,t)^2=\underset{m=0}{\overset{2n+1}{}}E_m$$
(2.19)
and we will choose
$$g_{n+1}0,\text{ that is, }E_m0\text{}m=0,\mathrm{},2n+1.$$
(2.20)
The actual sign of $`g_{n+1}`$ will be determined later (cf. (3.9), (3.10)). A comparison of coefficients of $`z^k`$ in (2.10)–(2.15) then yields
$`f_0`$ $`=2v,`$
$`f_1`$ $`=iv_x+2v^2v^{}+c_1(2v),`$
$`f_n`$ $`=2g_{n+1}u,`$
$`g_0`$ $`=1,`$
$`g_1`$ $`=2vv^{}+c_1,`$ (2.21)
$`g_{n+1}`$ $`=\left({\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m\right)^{1/2},`$
$`h_0`$ $`=2v^{},`$
$`h_1`$ $`=iv_x^{}2v(v^{})^2+c_12v^{},`$
$`h_n`$ $`=2g_{n+1}u^{},\text{ etc.,}`$
where $`\{c_{\mathrm{}}\}_{\mathrm{}}`$ denote integration constants, and
$`iu_x(x,t)+2v(x,t)+2v(x,t)v^{}(x,t)u(x,t)`$ $`=0,`$ (2.22)
$`iu_x^{}(x,t)+2v^{}(x,t)+2v(x,t)v^{}(x,t)u^{}(x,t)`$ $`=0,`$ (2.23)
$`iv_t(x,t)+2u(x,t)+2u(x,t)u^{}(x,t)v(x,t)`$ $`=0,`$ (2.24)
$`iv_t^{}(x,t)+2u^{}(x,t)+2u(x,t)u^{}(x,t)v^{}(x,t)`$ $`=0.`$ (2.25)
Equations (2.22)–(2.25) represent the classical massive Thirring system in light cone coordinates. It should be emphasized that the original Thirring equations are given by (2.22), (2.24) imposing the constraints
$$u^{}(x,t)=\overline{u(x,t)},v^{}(x,t)=\overline{v(x,t)},$$
(2.26)
where the bar denotes the operation of complex conjugate. In this paper, however, we will not impose the constraints (2.26) but rather study the system (2.22)–(2.25).
Given (2.22)–(2.25), a straightforward computation verifies the commutator relation
$$U_t(\zeta ,x,t)\stackrel{~}{V}_x(\zeta ,x,t)+[U(\zeta ,x,t),\stackrel{~}{V}(\zeta ,x,t)]=0,(\zeta ,x,t)\{0\}\times ^2,$$
(2.27)
complementing (2.8) and (2.9).
This concludes our brief review of the polynomial setup by Date , and for the remainder of this paper we will deviate from his strategy and focus on an approach based on the solution $`\varphi `$ of a Riccati-type equation associated with the Thirring system. This will enable us to employ a formalism previously applied to the KdV, AKNS, Toda, Boussinesq, and the combined sine-Gordon–mKdV hierarchies , , , , , , .
We conclude this section by mentioning the elementary fact that the Thirring system (2.22)–(2.25) is invariant under the scaling transformation,
$$(u,v,u^{},v^{})(Au,Av,A^1u^{},A^1v^{}),A\{0\}.$$
(2.28)
In the special case where $`u^{}=\overline{u},v^{}=\overline{v}`$, $`A`$ in (2.28) is further constrained by $`|A|=1`$.
## 3. The basic algebro-geometric formalism
Introducing the (possibly singular) hyperelliptic curve $`𝒦_n`$ of (arithmetic) genus $`n_0`$,
$`𝒦_n:_n(z,y)`$ $`=y^2R_{2n+2}(z)=0,`$ (3.1)
$`R_{2n+2}(z)`$ $`={\displaystyle \underset{m=0}{\overset{2n+1}{}}}(zE_m),\{E_m\}_{m=0,\mathrm{},2n+1}\{0\},`$ (3.2)
we denote points $`P`$ on $`𝒦_n`$ by $`P=(z,y)`$ and compactify $`𝒦_n`$ by joining two points at infinity $`P_\mathrm{}_+`$, $`P_{\mathrm{}_{}}`$, $`P_\mathrm{}_+P_{\mathrm{}_{}}`$, still denoting the compactified curve by $`𝒦_n`$. Moreover, we recall the hyperelliptic involution (sheet exchange map) $``$ on $`𝒦_n`$,
$$:𝒦_n𝒦_n,P=(z,y)P^{}=(z,y),P_{\mathrm{}_+}^{}{}_{}{}^{}=P_{\mathrm{}_{}}.$$
(3.3)
For additional facts on $`𝒦_n`$ and further notation freely employed throughout this paper, the reader may want to consult Appendix A.
Next, we define the fundamental meromorphic function $`\varphi (,x,t)`$ on $`𝒦_n`$ by
$`\varphi (P,x,t)`$ $`={\displaystyle \frac{y(P)+G_{n+1}(z,x,t)}{F_n(z,x,t)}}`$ (3.4)
$`={\displaystyle \frac{zH_n(z,x,t)}{y(P)G_{n+1}(z,x,t)}},P=(z,y)𝒦_n,(x,t)^2,`$ (3.5)
where we used (2.17) to obtain (3.5). Introducing
$`\widehat{\mu }_j(x,t)`$ $`=(\mu _j(x,t),G_{n+1}(\mu _j(x,t),x,t))𝒦_n,j=1,\mathrm{},n,(x,t)^2,`$ (3.6)
$`\widehat{\nu }_j(x,t)`$ $`=(\nu _j(x,t),G_{n+1}(\nu _j(x,t),x,t))𝒦_n,j=1,\mathrm{},n,(x,t)^2,`$ (3.7)
and
$$P_{0,\pm }=(0,\pm G_{n+1}(0))=(0,\pm g_{n+1})𝒦_n,$$
(3.8)
we fix the branch of $`y(P)`$ near $`P_\mathrm{}_\pm `$ according to
$$\underset{|z|\mathrm{}}{lim}\frac{y(P)}{G_{n+1}(z,x,t)}=\underset{|z|\mathrm{}}{lim}\frac{y(P)}{z^{n+1}}=1\text{ as }PP_\mathrm{}_\pm $$
(3.9)
and consequently determine the sign of $`g_{n+1}`$,
$$g_{n+1}=\left(\underset{m=0}{\overset{2n+1}{}}E_m\right)^{1/2},$$
(3.10)
by compatibility of all local charts on $`𝒦_n`$. We note that $`P_{0,\pm }`$ and $`P_\mathrm{}_\pm `$ are not necessarily on the same sheet of $`𝒦_n`$. The actual sheet on which $`P_{0,\pm }`$ lie depends on the sign of $`g_{n+1}`$ and hence on the location of all $`E_m`$.
Given these conventions, the divisor $`(\varphi (,x,t))`$ of $`\varphi (,x,t)`$ then reads
$$(\varphi (,x,t))=𝒟_{P_{0,}\underset{¯}{\overset{^}{\nu }}(x,t)}𝒟_{P_{\mathrm{}_{}}\underset{¯}{\overset{^}{\mu }}(x,t)},(x,t)^2.$$
(3.11)
Next we collect a few characteristic properties of $`\varphi `$.
###### Lemma 3.1.
Assume (2.1), (2.8), (2.9), and (3.2) and let $`P=(z,y)𝒦_n`$, $`(x,t)^2`$. Then $`\varphi `$ satisfies the Riccati-type equations
$`\begin{array}{cc}\hfill \varphi _x(P,x,t)& +2iv(x,t)\varphi (P,x,t)^2\hfill \\ & +2i(zv(x,t)v^{}(x,t))\varphi (P,x,t)=2izv^{}(x,t),\hfill \end{array}`$ (3.12)
$`\begin{array}{cc}\hfill \varphi _t(P,x,t)& +2iz^1u(x,t)\varphi (P,x,t)^2\hfill \\ & +2i(z^1u(x,t)u^{}(x,t))\varphi (P,x,t)=2iu^{}(x,t).\hfill \end{array}`$ (3.13)
Moreover,
$`\varphi (P,x,t)\varphi (P^{},x,t)`$ $`=zH_n(z,x,t)/F_n(z,x,t),`$ (3.14)
$`\varphi (P,x,t)+\varphi (P^{},x,t)`$ $`=2G_{n+1}(z,x,t)/F_n(z,x,t),`$ (3.15)
$`\varphi (P,x,t)\varphi (P^{},x,t)`$ $`=2y(P)/F_n(z,x,t).`$ (3.16)
###### Proof.
Equation (3.12) follows from (2.10)–(2.11), (2.17), and (3.4). Similarly, (3.13) follows from (2.13)–(2.14), (2.17), and (3.4). Relations (3.14)–(3.16) are obvious from (2.17) and (3.4). ∎
Given $`\varphi (P,x,t)`$, we can define the Baker–Akhiezer vector $`\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0)`$ by
$`\begin{array}{cc}\hfill \mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0)& =\left(\begin{array}{c}\psi _1(P,x,x_0,t,t_0)\\ \psi _2(P,\zeta ,x,x_0,t,t_0)\end{array}\right),\hfill \\ & P=(z,y)𝒦_n\{P_\mathrm{}_\pm \},z=\zeta ^2,(x,t),(x_0,t_0)^2,\hfill \end{array}`$ (3.17)
$`\begin{array}{cc}\hfill \psi _1(& P,x,x_0,t,t_0)\hfill \\ & =\mathrm{exp}(i{\displaystyle _{t_0}^t}ds(z^1u(x_0,s)u^{}(x_0,s)+2z^1u(x_0,s)\varphi (P,x_0,s))\hfill \\ & +i{\displaystyle _{x_0}^x}dx^{}(zv(x^{},t)v^{}(x^{},t)+2v(x^{},t)\varphi (P,x^{},t))),\hfill \end{array}`$ (3.18)
$`\psi _2(`$ $`P,\zeta ,x,x_0,t,t_0)=\zeta ^1\psi _1(P,x,x_0,t,t_0)\varphi (P,x,t).`$ (3.19)
Properties of $`\mathrm{\Psi }`$ are summarized in the following result.
###### Lemma 3.2.
Assume (2.1), (2.8), (2.9), and (3.2) and let $`P=(z,y)𝒦_n\{P_\mathrm{}_\pm \}`$, $`(x,t),(x_0,t_0)^2`$. Then $`\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0)`$ satisfies
$`\mathrm{\Psi }_x(P,\zeta ,x,x_0,t,t_0)`$ $`=U(\zeta ,x,t)\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0),`$ (3.20)
$`\mathrm{\Psi }_t(P,\zeta ,x,x_0,t,t_0)`$ $`=\stackrel{~}{V}(\zeta ,x,t)\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0),`$ (3.21)
$`iy(P)\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0)`$ $`=V_{n+1}(\zeta ,x,t)\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0).`$ (3.22)
Moreover, if the zeros of $`F_n(,x,t)`$ are all simple for $`(x,t)\mathrm{\Omega }`$, $`\mathrm{\Omega }^2`$ open and connected, then $`\psi _1(,x,x_0,t,t_0)`$, $`(x,t),(x_0,t_0)\mathrm{\Omega }`$, is meromorphic on $`𝒦_n\{P_\mathrm{}_\pm \}`$. In addition,
$$\begin{array}{c}\psi _1(P,x,x_0,t,t_0)=\left(\frac{F_n(z,x,t)}{F_n(z,x_0,t_0)}\right)^{1/2}\times \hfill \\ \hfill \times \mathrm{exp}\left(2iy(P)z^1_{t_0}^t𝑑s\frac{u(x_0,s)}{F_n(z,x_0,s)}+2iy(P)_{x_0}^x𝑑x^{}\frac{v(x^{},t)}{F_n(z,x^{},t)}\right),\end{array}$$
(3.23)
$`\psi _1(P,x,x_0,t,t_0)\psi _1(P^{},x,x_0,t,t_0)`$ $`=F_n(z,x,t)/F_n(z,x_0,t_0),`$ (3.24)
$`\psi _2(P,\zeta ,x,x_0,t,t_0)\psi _2(P^{},\zeta ,x,x_0,t,t_0)`$ $`=H_n(z,x,t)/F_n(z,x_0,t_0),`$ (3.25)
$`\psi _1(P,x,x_0,t,t_0)\psi _2(P^{},\zeta ,x,x_0,t,t_0)+`$ $`\psi _1(P^{},x,x_0,t,t_0)\psi _2(P,\zeta ,x,x_0,t,t_0)`$
$`=2\zeta ^1G_{n+1}(z,x,t)/F_n(z,x_0,t_0).`$ (3.26)
###### Proof.
Equations (3.20), (3.21) are verified using (2.10)–(2.14), (3.12), (3.13), (3.18), and (3.19). (3.22) follows by combining (2.3), (3.4), (3.5), (3.18), and (3.19). Clearly $`\psi _1`$ is meromorphic on $`𝒦_n\{P_\mathrm{}_\pm ,\widehat{\mu }_1(x,t),\mathrm{},\widehat{\mu }_n(x,t)\}`$ by (3.18). Since
$`2iv(x^{},t)\varphi (P,x^{},t)`$ $`\underset{P\widehat{\mu }_j(x^{},t)}{=}{\displaystyle \frac{}{x^{}}}\mathrm{ln}\left(F_n(z,x^{},t)\right)+O(1)\text{ as }z\mu _j(x^{},t),`$ (3.27)
$`2iz^1u(x_0,s)\varphi (P,x_0,s)`$ $`\underset{P\widehat{\mu }_j(x_0,s)}{=}{\displaystyle \frac{}{s}}\mathrm{ln}\left(F_n(z,x_0,s)\right)+O(1)\text{ as }z\mu _j(x_0,s),`$ (3.28)
one infers that $`\psi _1`$ is meromorphic on $`𝒦_n\{P_\mathrm{}_\pm \}`$ if the zeros of $`F_n(,x,t)`$ are all simple. This follows from (3.18) by restricting $`P`$ to a sufficiently small neighborhood $`𝒰_j(x_0)`$ of $`\{\widehat{\mu }_j(x_0,s)𝒦_n|(x_0,s)\mathrm{\Omega },s[t_0,t]\}`$ such that $`\widehat{\mu }_k(x_0,s)𝒰_j(x_0)`$ for all $`s[t_0,t]`$ and all $`k\{1,\mathrm{},n\}\{j\}`$, and similarly, by restricting $`P`$ to a sufficiently small neighborhood $`𝒰_j(t)`$ of $`\{\widehat{\mu }_j(x^{},t)𝒦_n|(x^{},t)\mathrm{\Omega },x^{}[x_0,x]\}`$ such that $`\widehat{\mu }_k(x^{},t)𝒰_j(t)`$ for all $`x^{}[x_0,x]`$ and all $`k\{1,\mathrm{},n\}\{j\}`$. Equation (3.23) follows from (3.18) after replacing $`\varphi `$ by the right-hand side of (3.4) and utilizing (2.10) in the $`x^{}`$-integral and (2.13) in the $`s`$-integral. Equations (3.24)–(3.26) immediately follow from (3.14)–(3.16), and (3.19). ∎
Next we discuss the asymptotic behavior of $`\varphi (P,x,t)`$ as $`PP_{0,\pm },P_\mathrm{}_\pm `$ in some detail since this will turn out to be a crucial ingredient for the theta function representation to be derived in Section 4.
###### Lemma 3.3.
Assume (2.1), (2.8), (2.9), and (3.2). Then
$`\varphi (P,x,t)`$ $`\underset{z\mathrm{}}{=}{\displaystyle \frac{1}{v(x,t)}}z+{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{1}{v(x,t)}}\right)_x+O\left({\displaystyle \frac{1}{z}}\right)\text{ as }P=(z,y)P_{\mathrm{}_{}},`$ (3.29)
$`\varphi (P,x,t)`$ $`\underset{z\mathrm{}}{=}v^{}(x,t)+{\displaystyle \frac{i}{2}}v_x^{}(x,t){\displaystyle \frac{1}{z}}+O\left({\displaystyle \frac{1}{z^2}}\right)\text{ as }P=(z,y)P_\mathrm{}_+,`$ (3.30)
$`\varphi (P,x,t)`$ $`\underset{z0}{=}u^{}(x,t)z+{\displaystyle \frac{i}{2}}u_t^{}(x,t)z^2+O(z^3)\text{ as }P=(z,y)P_{0,},`$ (3.31)
$`\varphi (P,x,t)`$ $`\underset{z0}{=}{\displaystyle \frac{1}{u(x,t)}}+{\displaystyle \frac{i}{2}}\left({\displaystyle \frac{1}{u(x,t)}}\right)_tz+O(z^2)\text{ as }P=(z,y)P_{0,+}.`$ (3.32)
###### Proof.
The existence of these asymptotic expansions (in terms of local coordinates $`\zeta =1/z`$ near $`P_\mathrm{}_\pm `$ and local coordinate $`\zeta =z`$ near $`P_{0,\pm }`$) is clear from the explicit form of $`\varphi `$ in (3.4). Insertion of the polynomials $`F_n`$, $`H_n`$, and $`G_{n+1}`$, then in principle, yields the explicit expansion coefficients in (3.29)–(3.32). However, this is a cumbersome procedure, especially with regard to the next to leading coefficients in (3.29)–(3.32). Much more efficient is the actual computation of these coefficients utilizing the Riccati-type equations (3.12) and (3.13). Indeed, inserting the ansatz
$$\varphi \underset{z\mathrm{}}{=}z\varphi _1+\varphi _0+O\left(\frac{1}{z}\right)$$
(3.33)
into (3.12) and comparing the first two leading powers of $`z`$ immediately yields (3.29). Similarly, the ansatz
$$\varphi \underset{z\mathrm{}}{=}\varphi _0+\varphi _1\frac{1}{z}+O\left(\frac{1}{z^2}\right)$$
(3.34)
inserted into (3.12) immediately produces (3.30). In exactly the same manner, inserting the ansatz
$$\varphi \underset{z0}{=}\varphi _1z+\varphi _2z^2+O(z^3)$$
(3.35)
and the ansatz
$$\varphi \underset{z0}{=}\varphi _0+\varphi _1z+O(z^2)$$
(3.36)
into (3.13) immediately yields (3.31) and (3.32), respectively. ∎
We follow up with a similar asymptotic analysis of $`\psi _1(P,x,x_0,t,t_0)`$.
###### Lemma 3.4.
Assume (2.1), (2.8), (2.9), and (3.2). Then
$`\psi _1(P,x,x_0,t,t_0)`$ $`\underset{z\mathrm{}}{=}\mathrm{exp}\left(iz(xx_0)+O(1)\right)\text{ as }P=(z,y)P_{\mathrm{}_{}},`$ (3.37)
$`\psi _1(P,x,x_0,t,t_0)`$ $`\underset{z0}{=}\mathrm{exp}\left(\pm iz^1(tt_0)+O(1)\right)\text{ as }P=(z,y)P_{0,}.`$ (3.38)
###### Proof.
Equations (3.37) and (3.38) follow from (3.18) noting
$`\begin{array}{cc}\hfill i(zv(x,t)v^{}(x,t))+2iv(x,t)\varphi (& P,x,t)\underset{z\mathrm{}}{=}iz+O(1)\hfill \\ & \text{ as }P=(z,y)P_{\mathrm{}_{}},\hfill \end{array}`$ (3.39)
$`\begin{array}{cc}\hfill i(z^1u(x_0,s)u^{}(x_0,s))+2iz^1u(x_0,& s)\varphi (P,x_0,s)\underset{z\mathrm{}}{=}O(1)\hfill \\ & \text{ as }P=(z,y)P_{\mathrm{}_{}},\hfill \end{array}`$ (3.40)
$`\begin{array}{cc}\hfill i(zv(x,t)v^{}(x,t))+2iv(& x,t)\varphi (P,x,t)\underset{z0}{=}O(1)\hfill \\ & \text{ as }P=(z,y)P_{0,},\hfill \end{array}`$ (3.41)
$`\begin{array}{cc}\hfill i(z^1u(x_0,s)u^{}(x_0,s))+2iz^1u(x_0,s)\varphi (P,x_0& ,s)\underset{z0}{=}\pm iz^1+O(1)\hfill \\ & \text{ as }P=(z,y)P_{0,}.\hfill \end{array}`$ (3.42)
In some of the following considerations it is appropriate to assume that $`𝒦_n`$ is nonsingular and hence we then assume
$$E_mE_m^{}\text{ for }mm^{},m,m^{}=0,\mathrm{},2n+1$$
(3.43)
in addition to (3.2).
Next, we turn to Dubrovin-type equations for $`\mu _j(x,t)`$, $`\nu _j(x,t)`$, $`j=1,\mathrm{},n`$, that is, we derive the nonlinear first-order system of partial differential equations governing their $`(x,t)`$-variation.
###### Lemma 3.5.
Let $`n`$. Assume (2.1), (2.8), (2.9), and (3.2) and suppose that the zeros $`\{\mu _j(x,t)\}_{j=1,\mathrm{},n}`$ of $`F_n(,x,t)`$ remain distinct for $`(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu `$, where $`\stackrel{~}{\mathrm{\Omega }}_\mu ^2`$ is open and connected. Then $`\{\mu _j(x,t)\}_{j=1,\mathrm{},n}`$ satisfies the following system of differential equations
$`\mu _{j,x}(x,t)`$ $`=2iy(\widehat{\mu }_j(x,t)){\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}}(\mu _j(x,t)\mu _{\mathrm{}}(x,t))^1,`$ (3.44)
$`\mu _{j,t}(x,t)`$ $`=(1)^ng_{n+1}^1\left({\displaystyle \underset{\begin{array}{c}k=1\\ kj\end{array}}{\overset{n}{}}}\mu _k(x,t)\right)2iy(\widehat{\mu }_j(x,t)){\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}}(\mu _j(x,t)\mu _{\mathrm{}}(x,t))^1,`$
$`j=1,\mathrm{},n,(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu .`$ (3.45)
Next, assume $`𝒦_n`$ to be nonsingular and introduce the initial condition
$$\{\widehat{\mu }_j(x_0,t_0)\}_{j=1,\mathrm{},n}𝒦_n,$$
(3.46)
where $`\{\mu _j(x_0,t_0)\}_{j=1,\mathrm{},n}`$ remain distinct and distinct from zero. Then there exists an open and connected set $`\mathrm{\Omega }_\mu ^2`$, with $`(x_0,t_0)\mathrm{\Omega }_\mu `$, such that the initial value problem (3.44)–(3.46) has a unique solution $`\{\widehat{\mu }_j(x,t)\}_{j=1,\mathrm{},n}`$ satisfying
$$\widehat{\mu }_jC^{\mathrm{}}(\mathrm{\Omega }_\mu ,𝒦_n),j=1,\mathrm{},n.$$
(3.47)
For the zeros $`\{\nu _j(x,t)\}_{j=1,\mathrm{},n}`$ of $`H_n(,x,t)`$ identical statements hold with $`\mu `$ replaced by $`\nu `$, $`\stackrel{~}{\mathrm{\Omega }}_\mu `$ by $`\stackrel{~}{\mathrm{\Omega }}_\nu `$, etc. In particular, $`\{\widehat{\nu }_j(x,t)\}_{j=1,\mathrm{},n}`$ satisfies
$`\nu _{j,x}(x,t)`$ $`=2iy(\widehat{\nu }_j(x,t)){\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}}(\nu _j(x,t)\nu _{\mathrm{}}(x,t))^1,`$ (3.48)
$`\nu _{j,t}(x,t)`$ $`=(1)^ng_{n+1}^1\left({\displaystyle \underset{\begin{array}{c}k=1\\ kj\end{array}}{\overset{n}{}}}\nu _k(x,t)\right)2iy(\widehat{\nu }_j(x,t)){\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}}(\nu _j(x,t)\nu _{\mathrm{}}(x,t))^1,`$
$`j=1,\mathrm{},n,(x,t)\stackrel{~}{\mathrm{\Omega }}_\nu .`$ (3.49)
###### Proof.
Equations (2.5), (2.10), and (3.6) imply
$`F_{n,x}(\mu _j)=f_0(\mu _{j,x}){\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}}(\mu _j\mu _{\mathrm{}})=4ivG_{n+1}(\mu _j)=4ivy(\widehat{\mu }_j).`$ (3.50)
Using $`f_0=2v`$ by (2.21), one concludes (3.44). Similarly, one derives from (2.5), (2.13), and (3.6),
$$F_{n,t}(\mu _j)=f_0(\mu _{j,t})\underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}(\mu _j\mu _{\mathrm{}})=(4iu/\mu _j)G_{n+1}(\mu _j)=(4iu/\mu _j)y(\widehat{\mu }_j).$$
(3.51)
Since
$$4iu/f_0=2if_n/(f_0g_{n+1})=2i(1)^n\left(\underset{k=1}{\overset{n}{}}\mu _k\right)/g_{n+1}$$
(3.52)
by (2.21) and (2.5), one arrives at (3.45). Equations (3.48) and (3.49) are derived analogously. In order to conclude (3.47), one first needs to investigate the case where $`\widehat{\mu }_j(x,t)`$ hits one of the branch points $`(E_m,0)(𝒦_n)`$ and hence the right-hand sides of (3.44) and (3.45) vanish. Thus we suppose that
$$\mu _{j_0}(x,t)E_{m_0}\text{ as }(x,t)(\stackrel{~}{x}_0,\stackrel{~}{t}_0)$$
(3.53)
for some $`j_0\{1,\mathrm{},n\}`$, $`m_0\{0,\mathrm{},2n+1\}`$ and some $`(\stackrel{~}{x}_0,\stackrel{~}{t}_0)\mathrm{\Omega }_\mu `$. Introducing
$$\zeta _{j_0}(x,t)=(\mu _{j_0}(x,t)E_{m_0})^{1/2},\mu _{j_0}(x,t)=E_{m_0}+\zeta _{j_0}(x,t)^2$$
(3.54)
for $`(x,t)`$ in an open neighborhood of $`(\stackrel{~}{x}_0,\stackrel{~}{t}_0)\mathrm{\Omega }_\mu `$, equations (3.44) and (3.45) become
$`\zeta _{j_0,x}(x,t)\underset{(x,t)(\stackrel{~}{x}_0,\stackrel{~}{t}_0)}{=}`$ $`2i\left({\displaystyle \underset{\begin{array}{c}m=0\\ mm_0\end{array}}{\overset{2n+1}{}}}(E_{m_0}E_m)\right)^{1/2}\left({\displaystyle \underset{\begin{array}{c}k=1\\ kj_0\end{array}}{\overset{n}{}}}(E_{m_0}\mu _k(x,t))^1\right)\times `$
$`\times (1+O(\zeta _{j_0}(x,t)^2)),`$ (3.55)
$`\zeta _{j_0,t}(x,t)\underset{(x,t)(\stackrel{~}{x}_0,\stackrel{~}{t}_0)}{=}`$ $`2i\left({\displaystyle \underset{\begin{array}{c}m=0\\ mm_0\end{array}}{\overset{2n+1}{}}}(E_{m_0}E_m)\right)^{1/2}\left({\displaystyle \underset{\begin{array}{c}k=1\\ kj_0\end{array}}{\overset{n}{}}}(E_{m_0}\mu _k(x,t))^1\right)\times `$
$`\times \left({\displaystyle \underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j_0\end{array}}{\overset{n}{}}}\mu _{\mathrm{}}(x,t)\right)(1+O(\zeta _{j_0}(x,t)^2)).`$ (3.56)
Since by hypothesis the right-hand sides of (3.55) and (3.56) are nonvanishing, one arrives at (3.47). ∎
Next we derive a few trace formulas involving $`u,v,u^{},v^{}`$ and some of their $`x`$-derivatives in terms of $`\mu _j(x,t)`$ and $`\nu _j(x,t)`$.
###### Lemma 3.6.
Let $`n`$ and assume (2.1), (2.8), (2.9), and (3.2). Then
$`i{\displaystyle \frac{v_x(x,t)}{v(x,t)}}+2v(x,t)v^{}(x,t)2c_1=2{\displaystyle \underset{j=1}{\overset{n}{}}}\mu _j(x,t),`$ (3.57)
$`i{\displaystyle \frac{v_x(x,t)}{v(x,t)}}2v(x,t)v^{}(x,t)=i{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\mu _{j,x}(x,t)}{\mu _j(x,t)}}+{\displaystyle \frac{2(1)^ng_{n+1}}{_{j=1}^n\mu _j(x,t)}},`$ (3.58)
$`{\displaystyle \frac{v(x,t)}{u(x,t)}}={\displaystyle \frac{(1)^ng_{n+1}}{_{j=1}^n\mu _j(x,t)}},`$ (3.59)
$`i{\displaystyle \frac{v_x^{}(x,t)}{v^{}(x,t)}}2v(x,t)v^{}(x,t)+2c_1=2{\displaystyle \underset{j=1}{\overset{n}{}}}\nu _j(x,t),`$ (3.60)
$`i{\displaystyle \frac{v_x^{}(x,t)}{v^{}(x,t)}}+2v(x,t)v^{}(x,t)=i{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\nu _{j,x}(x,t)}{\nu _j(x,t)}}{\displaystyle \frac{2(1)^ng_{n+1}}{_{j=1}^n\nu _j(x,t)}},`$ (3.61)
$`{\displaystyle \frac{v^{}(x,t)}{u^{}(x,t)}}={\displaystyle \frac{(1)^ng_{n+1}}{_{j=1}^n\nu _j(x,t)}}.`$ (3.62)
Here
$$c_1=\frac{1}{2}\underset{m=0}{\overset{2n+1}{}}E_m$$
(3.63)
and $`g_{n+1}=\left(_{m=0}^{2n+1}E_m\right)^{1/2}`$ has been introduced in (3.9) and (3.10).
###### Proof.
Equations (3.57) and (3.60) follow from (2.5), (2.7) by comparing powers of $`z^n`$ and $`z^{n1}`$, using (2.21). (3.58) and (3.61) follow from taking $`z=0`$ in (2.10) and (2.12), using again (2.21). Finally, (3.59) and (3.62) follow from $`f_n=f_0_{j=1}^n(\mu _j)`$, $`h_n=h_0_{j=1}^n(\nu _j)`$ and (2.21). ∎
While we are not explicitly introducing the hierarchy of massive Thirring equations in this paper, we note that Dubrovin-type equations such as (3.44), (3.45) combined with trace formulas for $`u,v,u^{},v^{}`$ in terms of $`\mu _j(x,t)`$, enable one to discuss such a hierarchy following the approach outlined in .
Up to this point we assumed the zero curvature equations (2.8) and (2.9), or equivalently, (2.10)–(2.15) and as a consequence, derived the corresponding algebro-geometric formalism. In the remainder of this section we will study the algebro-geometric initial value problem, that is, starting from the Dubrovin equations (3.44)–(3.46) and the trace formulas (3.57)–(3.59), derive (2.10)–(2.15), and hence the zero curvature equations (2.8) and (2.9).
We start with an elementary result extending the scaling transformation mentioned in (3.29).
###### Lemma 3.7.
Assume (2.1) and suppose $`u,v,u^{},v^{}`$ satisfy the Thirring system (2.22)–(2.25). Assume $`A(t)=\mathrm{exp}\left(^t𝑑sa(s)\right)`$, $`t`$, with $`aC()`$ and consider the time-dependent scaling transformation
$$(u,v,u^{},v^{})(\stackrel{˘}{u},\stackrel{˘}{v},\stackrel{˘}{u}^{},\stackrel{˘}{v}^{})=(Au,Av,A^1u^{},A^1v^{}).$$
(3.64)
Then $`\stackrel{˘}{u},\stackrel{˘}{v},\stackrel{˘}{u}^{},\stackrel{˘}{v}^{}`$ satisfy the corresponding extended massive Thirring system
$`i\stackrel{˘}{u}_x(x,t)+2\stackrel{˘}{v}(x,t)+2\stackrel{˘}{v}(x,t)\stackrel{˘}{v}^{}(x,t)\stackrel{˘}{u}(x,t)`$ $`=0,`$ (3.65)
$`i\stackrel{˘}{u}_x^{}(x,t)+2\stackrel{˘}{v}^{}(x,t)+2\stackrel{˘}{v}(x,t)\stackrel{˘}{v}^{}(x,t)\stackrel{˘}{u}^{}(x,t)`$ $`=0,`$ (3.66)
$`i\stackrel{˘}{v}_t(x,t)+2\stackrel{˘}{u}(x,t)+2\stackrel{˘}{u}(x,t)\stackrel{˘}{u}^{}(x,t)\stackrel{˘}{v}(x,t)+ia(t)\stackrel{˘}{v}(x,t)`$ $`=0,`$ (3.67)
$`i\stackrel{˘}{v}_t^{}(x,t)+2\stackrel{˘}{u}^{}(x,t)+2\stackrel{˘}{u}(x,t)\stackrel{˘}{u}^{}(x,t)\stackrel{˘}{v}^{}(x,t)ia(t)\stackrel{˘}{v}^{}(x,t)`$ $`=0.`$ (3.68)
###### Proof.
It suffices to insert (3.64) into the system (2.22)–(2.25). ∎
In the special case where $`u^{}(x,t)=\overline{u(x,t)},v^{}(x,t)=\overline{v(x,t)}`$, $`A(t)`$ in Lemma 3.7 is further constrained by $`|A(t)|=1`$, $`t`$.
Next we provide the basic setup for the algebro-geometric initial value problem. We start from the following assumptions.
###### Hypothesis 3.8.
Given the hyperelliptic curve $`𝒦_n`$ in (3.2), and the proper choice of the branch of $`g_{n+1}`$ defined by $`g_{n+1}=\left(_{m=0}^{2n+1}E_m\right)^{1/2}`$, according to (3.9) (i.e., according to $`lim_{|z|\mathrm{}}y(P)z^{n1}=\mathrm{}`$ as $`PP_\mathrm{}_\pm `$), consider the Dubrovin-type system of differential equations (3.44), (3.45) on $`\mathrm{\Omega }_\mu `$, for some intitial conditions (3.46). Here $`\mathrm{\Omega }_\mu ^2`$ is assumed to be open and connected, and such that the projections $`\mu _j(x,t)`$ of $`\widehat{\mu }_j(x,t)`$ onto $``$ remain distinct and distinct from zero for $`(x,t)\mathrm{\Omega }_\mu `$, that is,
$`\mu _j(x,t)\mu _j^{}(x,t)\text{ for }jj^{},j,j^{}=1,\mathrm{},n,(x,t)\mathrm{\Omega }_\mu ,`$ (3.69)
$`\{\mu _j(x,t)\}_{j=1,\mathrm{},n}\{0\}=\mathrm{},(x,t)\mathrm{\Omega }_\mu .`$ (3.70)
Assuming Hypothesis 3.8 in the following, we will next define $`u,v,u^{},v^{}`$ and the polynomials $`F_n,G_{n+1},H_n`$ in the following steps (S1)–(S4).
(S1). Use the trace formulas (3.57)–(3.59) on $`\mathrm{\Omega }_\mu `$, that is,
$`i{\displaystyle \frac{v_x(x,t)}{v(x,t)}}+2v(x,t)v^{}(x,t)2c_1=2{\displaystyle \underset{j=1}{\overset{n}{}}}\mu _j(x,t),`$ (3.71)
$`i{\displaystyle \frac{v_x(x,t)}{v(x,t)}}2v(x,t)v^{}(x,t)=i{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\mu _{j,x}(x,t)}{\mu _j(x,t)}}+{\displaystyle \frac{2(1)^ng_{n+1}}{_{j=1}^n\mu _j(x,t)}},`$ (3.72)
$`u(x,t)=(1)^ng_{n+1}^1v(x,t){\displaystyle \underset{j=1}{\overset{n}{}}}\mu _j(x,t),(x,t)\mathrm{\Omega }_\mu ,`$ (3.73)
to define $`u(x,t),v(x,t),v^{}(x,t)`$ on $`\mathrm{\Omega }_\mu `$ up to a possibly $`t`$-dependent multiple factor according to the scale transformation described in Lemma 3.7.
(S2). Define the polynomial $`F_n(z,x,t)`$ on $`\times \mathrm{\Omega }_\mu `$ of degree $`n`$ with respect to $`z`$ by
$$F_n(z,x,t)=2v(x,t)\underset{j=1}{\overset{n}{}}(z\mu _j(x,t)),(z,x,t)\times \mathrm{\Omega }_\mu $$
(3.74)
and define the polynomial $`G_{n+1}(z,x,t)`$ on $`\times \mathrm{\Omega }_\mu `$ of degree $`n+1`$ with respect to $`z`$ by
$`F_{n,x}(z,x,t)=2i(v(x,t)v^{}(x,t)z)F_n(z,x,t)+4iv(x,t)G_{n+1}(z,x,t),`$ (3.75)
$`(z,x,t)\times \mathrm{\Omega }_\mu .`$
One then verifies from
$$2iy(\widehat{\mu }_j)=\mu _{j,x}\underset{\begin{array}{c}k=1\\ kj\end{array}}{\overset{n}{}}(\mu _j\mu _k)=\frac{F_{n,x}(\mu _j)}{2v},j=1,\mathrm{},n$$
(3.76)
and (3.75) that
$$y(\widehat{\mu }_j(x,t))=\frac{F_{n,x}(\mu _j(x,t)),x,t)}{4iv(x,t)}=G_{n+1}(\mu _j(x,t),x,t),j=1,\mathrm{},n,(x,t)\mathrm{\Omega }_\mu $$
(3.77)
and hence
$$\left(G_{n+1}(z,x,t)^2R_{2n+2}(z)\right)|_{z=\mu _j(x,t)}=0,j=1,\mathrm{},n,(x,t)\mathrm{\Omega }_\mu .$$
(3.78)
(S3). Taking $`z=0`$ in (3.75), using (3.74), results in
$$\frac{2(1)^nG_{n+1}(0)}{_{j=1}^n\mu _j}=i\frac{v_x}{v}2vv^{}+i\underset{j=1}{\overset{n}{}}\frac{\mu _{j,x}}{\mu _j}$$
(3.79)
and hence a comparison with (3.72) yields
$$G_{n+1}(0,x,t)=g_{n+1}=\left(\underset{m=0}{\overset{2n+1}{}}E_m\right)^{1/2}$$
(3.80)
and thus,
$$\left(G_{n+1}(z,x,t)^2R_{2n+2}(z)\right)|_{z=0}=0,(x,t)\mathrm{\Omega }_\mu .$$
(3.81)
Because of (3.78) and (3.81) we can define a polynomial $`H_n(z,x,t)`$ on $`\times \mathrm{\Omega }_\mu `$ of degree $`n`$ with respect to $`z`$ by
$$G_{n+1}(z,x,t)^2R_{2n+2}(z)=zF_n(z,x,t)H_n(z,x,t),(z,x,t)\times \mathrm{\Omega }_\mu .$$
(3.82)
(S4). Given $`H_n(z,x,t)`$ we finally define $`u^{}(x,t)`$ on $`\mathrm{\Omega }_\mu `$ by
$$u^{}(x,t)=\frac{H_n(0,x,t)}{2g_{n+1}},(x,t)\mathrm{\Omega }_\mu ,$$
(3.83)
Again $`u^{}(x,t)`$ is defined up to a possibly $`t`$-dependent factor in accordance with Lemma 3.7.
The algebro-geometric initial value problem now can be solved as follows.
###### Theorem 3.9.
Assume Hypothesis 3.8, define $`u,v,u^{},v^{}`$ and $`F_n,G_{n+1},H_n`$ as in $`(S1)(S4)`$ and let $`(x,t)\mathrm{\Omega }_\mu `$. Then there exists a function $`aC^{\mathrm{}}(\mathrm{\Omega }_\mu )`$, independent of $`x`$ ($`a_x|_{\mathrm{\Omega }_\mu }=0`$), such that
$`F_{n,x}(z,x,t)`$ $`=2i(v(x,t)v^{}(x,t)z)F_n(z,x,t)+4iv(x,t)G_{n+1}(z,x,t),`$ (3.84)
$`G_{n+1,x}(z,x,t)`$ $`=2izv^{}(x,t)F_n(z,x,t)+2izv(x,t)H_n(z,x,t),`$ (3.85)
$`H_{n,x}(z,x,t)`$ $`=2i(v(x,t)v^{}(x,t)z)H_n(z,x,t)+4iv^{}(x,t)G_{n+1}(z,x,t),`$ (3.86)
$`\begin{array}{cc}\hfill F_{n,t}(z,x,t)& =2i(u(x,t)u^{}(x,t)z^1)F_n(z,x,t)+a(t)F_n(z,x,t)\hfill \\ & +4iz^1u(x,t)G_{n+1}(z,x,t),\hfill \end{array}`$ (3.87)
$`G_{n+1,t}(z,x,t)`$ $`=2iu^{}(x,t)F_n(z,x,t)+2iu(x,t)H_n(z,x,t),`$ (3.88)
$`\begin{array}{cc}\hfill H_{n,t}(z,x,t)& =2i(u(x,t)u^{}(x,t)z^1)H_n(z,x,t)a(t)H_n(z,x,t)\hfill \\ & +4iz^1u^{}(x,t)G_{n+1}(z,x,t).\hfill \end{array}`$ (3.89)
In particular, $`u,v,u^{},v^{}`$ satisfy the extended massive Thirring system (3.65)–(3.68) on $`\mathrm{\Omega }_\mu `$,
$`iu_x(x,t)+2v(x,t)+2v(x,t)v^{}(x,t)u(x,t)`$ $`=0,`$ (3.90)
$`iu_x^{}(x,t)+2v^{}(x,t)+2v(x,t)v^{}(x,t)u^{}(x,t)`$ $`=0,`$ (3.91)
$`iv_t(x,t)+2u(x,t)+2u(x,t)u^{}(x,t)v(x,t)+ia(t)v(x,t)`$ $`=0,`$ (3.92)
$`iv_t^{}(x,t)+2u^{}(x,t)+2u(x,t)u^{}(x,t)v^{}(x,t)ia(t)v^{}(x,t)`$ $`=0.`$ (3.93)
###### Proof.
Define the polynomial
$`P_n(z,x,t)=2izv^{}(x,t)F_n(z,x,t)+2izv(x,t)H_n(z,x,t)G_{n+1,x}(z,x,t),`$
$`(z,x,t)\times \mathrm{\Omega }_\mu .`$ (3.94)
Using (3.77) and $`2G_{n+1}G_{n+1,x}=z(F_{n,x}H_n+F_nH_{n,x})`$ (by differentiating (3.82) with respect to $`x`$) one then computes
$`G_{n+1}(\mu _j)P_n(\mu _j)`$ $`=2i\mu _jvH_n(\mu _j)G_{n+1}(\mu _j)G_{n+1}(\mu _j)G_{n+1,x}(\mu _j)`$
$`={\displaystyle \frac{1}{2}}\mu _jH_n(\mu _j)F_{n,x}(\mu _j){\displaystyle \frac{1}{2}}\mu _jF_{n,x}(\mu _j)H_n(\mu _j)=0,`$ (3.95)
$`j=1,\mathrm{},n.`$
In order to investigate the leading-order term with respect to $`z`$ of $`P_n(z)`$ we first study the leading-order $`z`$-behavior of $`F_n(z),G_{n+1}(z)`$, and $`H_n(z)`$. Writing (cf. (2.5)–(2.7))
$`F_n(z)={\displaystyle \underset{j=0}{\overset{n}{}}}f_{nj}z^j,H_n(z)={\displaystyle \underset{j=0}{\overset{n}{}}}h_{nj}z^j,G_{n+1}(z)={\displaystyle \underset{j=0}{\overset{n+1}{}}}g_{n+1j}z^j,g_0=1,`$ (3.96)
a comparison of leading powers with respect to $`z`$ in (3.74), (3.75), and (3.82) yields
$`f_0`$ $`=2v,`$ (3.97)
$`g_0`$ $`=1,`$ (3.98)
$`v_x+2iv^2v^{}+if_1+2ig_1`$ $`=0,`$ (3.99)
$`2g_1+2vh_0+{\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m`$ $`=0.`$ (3.100)
Since (3.71) can be rewritten in the form
$$f_1=iv_x+2v^2v^{}+v\underset{m=0}{\overset{2n+1}{}}E_m,$$
(3.101)
a comparison of (3.99) and (3.101) then yields
$$g_1=2vv^{}\frac{1}{2}\underset{m=0}{\overset{2n+1}{}}E_m$$
(3.102)
and hence
$$h_0=2v^{}.$$
(3.103)
Insertion of (3.97), (3.98), and (3.103) into (3.94) then yields
$$P_n(z,x,t)=O(z^n)\text{ as }|z|\mathrm{}.$$
(3.104)
Thus, (3.95) and (3.104) prove
$$P_n(z,x,t)=b(x,t)F_n(z,x,t),(z,x,t)\times \mathrm{\Omega }_\mu $$
(3.105)
for some $`bC^{\mathrm{}}(\mathrm{\Omega }_\mu )`$ (independent of $`z`$), implying
$`G_{n+1,x}(z,x,t)=2izv^{}(x,t)F_n(z,x,t)+2izv(x,t)H_n(z,x,t)b(x,t)F_n(z,x,t),`$
$`(z,x,t)\times \mathrm{\Omega }_\mu .`$ (3.106)
Taking $`z=0`$ in (3.106), observing that $`G_{n+1}(0,x,t)`$ is independent of $`(x,t)\mathrm{\Omega }_\mu `$ by (3.80), then shows that
$$0=b(x,t)F_n(0,x,t),(x,t)\mathrm{\Omega }_\mu ,$$
(3.107)
and hence $`b=0`$ on $`\mathrm{\Omega }_\mu `$ because of (3.70). Thus,
$$G_{n+1,x}(z,x,t)=2izv^{}(x,t)F_n(z,x,t)+2izv(x,t)H_n(z,x,t),(z,x,t)\times \mathrm{\Omega }_\mu .$$
(3.108)
Differentiating (3.82) with respect to $`x`$, inserting (3.75) and (3.108), then yields
$`H_{n,x}(z,x,t)=2i(v(x,t)v^{}(x,t)z)H_n(z,x,t)+4iv^{}(x,t)G_{n+1}(z,x,t),`$ (3.109)
$`(z,x,t)\times \mathrm{\Omega }_\mu `$
and we proved (3.84)–(3.86).
Next, combining (3.45), (3.73), and (3.77) one computes
$`F_{n,t}(\mu _j)`$ $`=2v{\displaystyle \frac{(1)^n}{g_{n+1}}}\left({\displaystyle \underset{\begin{array}{c}k=1\\ kj\end{array}}{\overset{n}{}}}\mu _k\right)2iy(\widehat{\mu }_j)={\displaystyle \frac{(1)^n}{g_{n+1}}}\left({\displaystyle \underset{\begin{array}{c}k=1\\ kj\end{array}}{\overset{n}{}}}\mu _k\right){\displaystyle \frac{4iv}{\mu _j}}G_{n+1}(\mu _j)`$
$`={\displaystyle \frac{4iu}{\mu _j}}G_{n+1}(\mu _j),j=1,\mathrm{},n.`$ (3.110)
Since clearly
$$F_{n,t}(z)\left(2i(uu^{}z^1)F_n(z)+4iz^1uG_{n+1}(z)\right)=O(z^n)\text{ as }|z|\mathrm{},$$
(3.111)
a comparison of (3.110) and (3.111) yields
$`F_{n,t}(z,x,t)(2i(u(x,t)u^{}(x,t)z^1)F_n(z,x,t))+4iz^1uG_{n+1}(z,x,t))`$
$`=a(x,t)F_n(z,x,t),(z,x,t)\times \mathrm{\Omega }_\mu `$ (3.112)
for some $`aC^{\mathrm{}}(\mathrm{\Omega }_\mu )`$ (independent of $`z`$), and hence (3.87) (except for $`a_x=0`$). A comparison of powers of $`z^n`$ in (3.112) then yields (3.92).
Next, we restrict $`\mathrm{\Omega }_\mu `$ a bit further and introduce $`\stackrel{~}{\mathrm{\Omega }}_\mu \mathrm{\Omega }_\mu `$ by the requirement that $`\mu _j(x,t)`$ remain distinct and also distinct from $`\{E_m\}_{m=0,\mathrm{},2n+1}\{0\}`$ for $`(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu `$, that is, we suppose
$`\mu _j(x,t)\mu _j^{}(x,t)\text{ for }jj^{},j,j^{}=1,\mathrm{},n,(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu ,`$ (3.113)
$`\{\mu _j(x,t)\}_{j=1,\mathrm{},n}\{\{E_m\}_{m=0,\mathrm{},2n+1}\{0\}\}=\mathrm{},(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu .`$ (3.114)
Differentiating (3.82) with respect to $`t`$ inserting (3.112) then yields
$`2G_{n+1}(z)G_{n+1,t}(z)`$ $`=zF_n(z)\left(2i(uu^{}z^1)H_n(z)+aH_n(z)+H_{n,t}(z)\right)`$
$`+4iuG_{n+1}H_n(z).`$ (3.115)
Since the zeros of $`F_n`$ and $`G_{n+1}`$ are disjoint by hypothesis (3.114) (cf. also (3.77)), $`zH_{n,t}(z)`$ necessarily must be of the form
$`zH_{n,t}(z,x,t)`$ $`=2i(zu(x,t)u^{}(x,t)1)H_n(z,x,t)a(x,t)zH_n(z,x,t)`$
$`+4id(x,t)G_{n+1}(z,x,t),(z,x,t)\times \stackrel{~}{\mathrm{\Omega }}_\mu `$ (3.116)
for some $`dC^{\mathrm{}}(\stackrel{~}{\mathrm{\Omega }}_\mu )`$ (independent of $`z`$) and (3.116) inserted into (3.115) then yields
$$G_{n+1}(z,x,t)=2iu(x,t)H_n(z,x,t)+2id(x,t)F_n(z,x,t),(z,x,t)\times \stackrel{~}{\mathrm{\Omega }}_\mu .$$
(3.117)
Since
$$u(x,t)=\frac{F_n(0,x,t)}{2g_{n+1}},(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu ,$$
(3.118)
combining (3.73) and (3.74), taking $`z=0`$ in (3.117), observing (3.80) and (3.83), results in
$$0=2iu2g_{n+1}u^{}+2id(2g_{n+1}u)$$
(3.119)
and hence in
$$d(x,t)=u^{}(x,t),(x,t)\stackrel{~}{\mathrm{\Omega }}_\mu .$$
(3.120)
Using property (3.47), (3.116)–(3.120) then extend by continuity from $`\stackrel{~}{\mathrm{\Omega }}_\mu `$ to $`\mathrm{\Omega }_\mu `$. This proves (3.88) and (3.89) (except for $`a_x=0`$). A comparison of powers of $`z^n`$ in (3.89) then yields (3.93). Taking $`z=0`$ in (3.84) and (3.86), observing (3.83) and (3.118), then proves (3.90) and (3.91). Finally, computing the partial $`t`$-derivative of $`F_{n,x}`$ and separately the partial $`x`$-derivative of $`F_{n,t}`$, utilizing (3.84), (3.85), (3.87), (3.88), and (3.90)–(3.93) then shows
$$F_{n,xt}(z,x,t)F_{n,tx}(z,x,t)=a_x(x,t)F_n(z,x,t),(z,x,t)\times \mathrm{\Omega }_\mu $$
(3.121)
and hence
$$a_x(x,t)=0,(x,t)\mathrm{\Omega }_\mu .$$
(3.122)
###### Remark 3.10.
(i) The fact that the system of Dubrovin equations (3.44)–(3.46) cannot uniquely determine the solutions $`u,v,u^{},v^{}`$ of the massive Thirring system (2.22)–(2.25), as is evident from the occurrence of $`a(t)`$ in (3.92), (3.93), is of course due to the scale covariance displayed explicitly in Lemma 3.7. In particular, once a certain $`a(t)`$ has been identified, a scaling transformation of the type (3.64) (with $`A(t)`$ replaced by $`1/A(t)`$) will restore the extended massive Thirring system (3.90)–(3.93) to it’s original form in (2.22)–(2.25).
(ii) For simplicity we formulated Theorem 3.9 in terms of $`\{\widehat{\mu }_j\}_{j=1,\mathrm{},n}`$ and (3.44)–(3.46) only. Of course there exists a completely analogous approach starting with $`\{\widehat{\nu }_j\}_{j=1,\mathrm{},n}`$ and the system (3.48), (3.49) instead.
(iii) Invoking the explicit theta function representations for $`u,v,u^{},v^{}`$ to be proven in Section 4 next (this approach is independent of that used to prove Theorem 3.10), one can extend the principal assertions (3.84)–(3.93) of Theorem 3.10 by continuity to $`(x,t)`$ lying in a larger set $`\mathrm{\Omega }^2`$ as long as the divisors $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ remain nonspecial for $`(x,t)\mathrm{\Omega }`$ (cf. Theorem 4.4 and Theorem A.7).
## 4. Theta function representations
In our final section we now derive theta function representations for the principal objects of Section 3, including $`\varphi `$, $`\psi _1`$, $`u`$, $`v`$, $`u^{}`$, $`v^{}`$. These representations complement the papers by Date and Prikarpatskii and Golod , where theta function representations were derived for appropriate symmetric functions associated with auxiliary divisors, but not explicitly for $`u`$, $`v`$, $`u^{}`$, $`v^{}`$. Moreover, we correct some inaccuracies of such formulas in a paper by Bikbaev (which follows a different strategy than ours).
According to our shift in emphasis from the Baker–Akhiezer vector $`\mathrm{\Psi }`$ to our fundamental meromorphic function $`\varphi `$ on $`𝒦_n`$, we next aim at the theta function representation of $`\varphi `$.
Assuming $`𝒦_n`$ to be nonsingular for the remainder of this section (i.e., $`E_mE_m^{}`$ for $`mm^{}`$, $`m,m^{}=0,\mathrm{},2n+1`$) and $`n`$ for simplicity (to avoid repeated case distinctions), we next recall the formula for a normal differential of the third kind, which has simple poles at $`P_{0,}`$ and $`P_{\mathrm{}_{}}`$, corresponding residues $`+1`$ and $`1`$, vanishing $`a`$-periods, and is holomorphic otherwise on $`𝒦_n`$. One computes
$$\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}=\frac{y+y_{0,}}{2z}\frac{dz}{y}+\frac{_{j=1}^n(z\lambda _j)dz}{2y},P_{0,}=(0,y_{0,})=(0,g_{n+1}),$$
(4.1)
where $`\{\lambda _j\}_{j=1,\mathrm{},n}`$ are uniquely determined by the normalization
$$_{a_j}\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}=0,j=1,\mathrm{},n.$$
(4.2)
The explicit formula (4.1) then implies (using the local coordinate $`\zeta =z`$ near $`P_{0,}`$)
$$\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}(P)\underset{\zeta 0}{=}\left\{\begin{array}{c}\zeta ^1\\ 0\end{array}\right\}d\zeta \pm \left(\underset{q=0}{\overset{\mathrm{}}{}}(q+1)\omega _{q+1}^0\zeta ^q\right)d\zeta \text{ as }PP_{0,},$$
(4.3)
and similarly (using the local coordinate $`\zeta =1/z`$ near $`P_{\mathrm{}_{}}`$),
$$\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}(P)\underset{\zeta 0}{=}\left\{\begin{array}{c}\zeta ^1\\ 0\end{array}\right\}d\zeta \pm \left(\underset{q=0}{\overset{\mathrm{}}{}}(q+1)\omega _{q+1}^{\mathrm{}}\zeta ^q\right)d\zeta \text{ as }PP_{\mathrm{}_{}}.$$
(4.4)
In particular,
$`{\displaystyle _{Q_0}^P}\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}`$ $`\underset{\zeta 0}{=}\left\{\begin{array}{c}\mathrm{ln}(\zeta )\\ 0\end{array}\right\}+\omega _0^{0,}\pm \omega _1^0\zeta \pm \omega _2^0\zeta ^2+O(\zeta ^3)\text{ as }PP_{0,},`$ (4.5)
$`{\displaystyle _{Q_0}^P}\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}`$ $`\underset{\zeta 0}{=}\left\{\begin{array}{c}\mathrm{ln}(\zeta )\\ 0\end{array}\right\}+\omega _0^{\mathrm{}_{}}\pm \omega _1^{\mathrm{}}\zeta \pm \omega _2^{\mathrm{}}\zeta ^2+O(\zeta ^3)\text{ as }PP_{\mathrm{}_{}}.`$ (4.6)
Here $`Q_0(𝒦_n)`$ is an appropriate base point and we agree to choose the same path of integration from $`Q_0`$ to $`P`$ in all Abelian integrals in this section.
A comparison of (4.3), (4.4) with (4.1), (A.12), and (A.14) then yields
$`\omega _1^0`$ $`={\displaystyle \frac{1}{4}}{\displaystyle \underset{m=0}{\overset{2n+1}{}}}{\displaystyle \frac{1}{E_m}}{\displaystyle \frac{(1)^n}{2g_{n+1}}}{\displaystyle \underset{j=1}{\overset{n}{}}}\lambda _j,`$ (4.7)
$`\omega _1^{\mathrm{}}`$ $`={\displaystyle \frac{1}{4}}{\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m+{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{n}{}}}\lambda _j.`$ (4.8)
Next, we intend to go a step further and derive alternative expressions for the expansion coefficients $`\omega _0^{0,\pm }`$, $`\omega _1^0`$, $`\omega _0^\mathrm{}_\pm `$, and $`\omega _1^{\mathrm{}}`$ in (4.5) and (4.6). To begin these calculations we first recall the notion of a nonsingular odd half-period $`\mathrm{{\rm Y}}`$ defined by
$$2\underset{¯}{\mathrm{{\rm Y}}}=0(modL_n),\theta (\underset{¯}{\mathrm{{\rm Y}}})=0,\frac{\theta (\underset{¯}{z})}{z_j}|_{\underset{¯}{z}=\underset{¯}{\mathrm{{\rm Y}}}}0\text{ for some }j\{1,\mathrm{},n\}.$$
(4.9)
Discussions of even and odd half-periods (singular and nonsingular ones) can be found, for instance, in \[16, p. 12–15\], . In addition, it is convenient to introduce the notation
$`\underset{¯}{\mathrm{\Delta }}_0`$ $`=\underset{¯}{A}_{Q_0}(P_{0,+}),\underset{¯}{\mathrm{\Delta }}_{\mathrm{}}=\underset{¯}{A}_{Q_0}(P_\mathrm{}_+),`$ (4.10)
$`\underset{¯}{W}_1^0`$ $`=(W_{1,1}^0,\mathrm{},W_{1,n}^0),W_{1,j}^0={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{b_j}}\omega _{P_{0,+},0}^{(2)}={\displaystyle \frac{c_j(1)}{g_{n+1}}},j=1,\mathrm{},n,`$ (4.11)
$`\underset{¯}{W}_2^0`$ $`=(W_{2,1}^0,\mathrm{},W_{2,n}^0),W_{2,j}^0={\displaystyle \frac{c_j(1)}{4g_{n+1}}}{\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m^1+{\displaystyle \frac{c_j(2)}{2g_{n+1}}},j=1,\mathrm{},n,`$ (4.12)
$`\underset{¯}{W}_1^{\mathrm{}}`$ $`=(W_{1,1}^{\mathrm{}},\mathrm{},W_{1,n}^{\mathrm{}}),W_{1,j}^{\mathrm{}}={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{b_j}}\omega _{P_\mathrm{}_+,0}^{(2)}=c_j(n),j=1,\mathrm{},n,`$ (4.13)
$`\underset{¯}{W}_2^{\mathrm{}}`$ $`=(W_{2,1}^{\mathrm{}},\mathrm{},W_{2,n}^{\mathrm{}}),W_{2,j}^{\mathrm{}}={\displaystyle \frac{c_j(n)}{4}}{\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m+{\displaystyle \frac{c_j(n1)}{2}},j=1,\mathrm{},n.`$ (4.14)
Moreover, we abbreviate directional derivatives of $`f`$ in the direction of $`\underset{¯}{W}=(W_1,\mathrm{},W_n)^n`$ by
$`(_{\underset{¯}{W}}f)(\underset{¯}{z})={\displaystyle \underset{j=1}{\overset{n}{}}}W_j{\displaystyle \frac{f}{z_j}}(\underset{¯}{z}),(_{\underset{¯}{W}}^2f)(\underset{¯}{z})={\displaystyle \underset{j,k=1}{\overset{n}{}}}W_jW_k{\displaystyle \frac{^2f}{z_jz_k}}(\underset{¯}{z}),`$ (4.15)
$`\underset{¯}{z}=(z_1,\mathrm{},z_n)^n.`$
Then one obtains the following result.
###### Lemma 4.1.
Given (4.1)–(4.13) one obtains
$`\omega _0^{0,+}`$ $`=\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}2\underset{¯}{\mathrm{\Delta }}_0)\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}}\right),`$ (4.16)
$`\omega _0^{0,}`$ $`=\mathrm{ln}\left({\displaystyle \frac{(_{\underset{¯}{W}_1^0}\theta )(\underset{¯}{\mathrm{{\rm Y}}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}{\theta (\underset{¯}{\mathrm{{\rm Y}}}+\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}}\right),`$ (4.17)
$`\omega _1^0`$ $`=_{\underset{¯}{W}_1^0}\mathrm{ln}(\theta (\underset{¯}{\mathrm{{\rm Y}}}+\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}}))+{\displaystyle \frac{(_{\underset{¯}{W}_2^0}\theta )(\underset{¯}{\mathrm{{\rm Y}}})+2^1(_{\underset{¯}{W}_1^0}^2\theta )(\underset{¯}{\mathrm{{\rm Y}}})}{(_{\underset{¯}{W}_1^0}\theta )(\underset{¯}{\mathrm{{\rm Y}}})}}`$ (4.18)
$`=_{\underset{¯}{W}_1^0}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}2\underset{¯}{\mathrm{\Delta }}_0)}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}}\right),`$ (4.19)
$`\omega _0^\mathrm{}_+`$ $`=\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}}}{\theta (\underset{¯}{\mathrm{{\rm Y}}}2\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}}\right),`$ (4.20)
$`\omega _0^{\mathrm{}_{}}`$ $`=\mathrm{ln}\left({\displaystyle \frac{(_{\underset{¯}{W}_1^{\mathrm{}}}\theta )(\underset{¯}{\mathrm{{\rm Y}}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0+\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}}\right),`$ (4.21)
$`\omega _1^{\mathrm{}}`$ $`=_{\underset{¯}{W}_1^{\mathrm{}}}\mathrm{ln}(\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0+\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})){\displaystyle \frac{(_{\underset{¯}{W}_2^{\mathrm{}}}\theta )(\underset{¯}{\mathrm{{\rm Y}}})+2^1(_{\underset{¯}{W}_1^{\mathrm{}}}^2\theta )(\underset{¯}{\mathrm{{\rm Y}}})}{(_{\underset{¯}{W}_1^{\mathrm{}}}\theta )(\underset{¯}{\mathrm{{\rm Y}}})}}`$ (4.22)
$`=_{\underset{¯}{W}_1^{\mathrm{}}}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}{\theta (\underset{¯}{\mathrm{{\rm Y}}}2\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}}\right).`$ (4.23)
###### Proof.
Abbreviating
$$\underset{¯}{w}(P,Q_0)=\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{A}_{Q_0}(Q)(modL_n),$$
(4.24)
one infers from
$`\underset{¯}{A}_{Q_0}(P)`$ $`\underset{\zeta 0}{=}\underset{¯}{A}_{Q_0}(P_{0,\pm })\pm \underset{¯}{W}_1^0\zeta \pm \underset{¯}{W}_2^0\zeta ^2+O(\zeta ^3)\text{ as }PP_{0,\pm },`$ (4.25)
$`\underset{¯}{A}_{Q_0}(P)`$ $`\underset{\zeta 0}{=}\underset{¯}{A}_{Q_0}(P_\mathrm{}_\pm )\pm \underset{¯}{W}_1^{\mathrm{}}\zeta \pm \underset{¯}{W}_2^{\mathrm{}}\zeta ^2+O(\zeta ^3)\text{ as }PP_\mathrm{}_\pm `$ (4.26)
(cf. (A.11) and (A.13)), and (4.11), (4.13), that
$`\theta (\underset{¯}{w}(P,Q))\underset{\zeta 0}{=}`$ $`\theta (\underset{¯}{w}(P_{0,\pm },Q))(_{\underset{¯}{W}_1^0}\theta )(\underset{¯}{w}(P_{0,\pm },Q))\zeta (_{\underset{¯}{W}_2^0}\theta )(\underset{¯}{w}(P_{0,\pm },Q))\zeta ^2`$
$`+2^1(_{\underset{¯}{W}_1^0}^2\theta )(\underset{¯}{w}(P_{0,\pm },Q))\zeta ^2+O(\zeta ^3)\text{ as }PP_{0,\pm },`$ (4.27)
$`\theta (\underset{¯}{w}(P,Q))\underset{\zeta 0}{=}`$ $`\theta (\underset{¯}{w}(P_\mathrm{}_\pm ,Q))(_{\underset{¯}{W}_1^{\mathrm{}}}\theta )(\underset{¯}{w}(P_\mathrm{}_\pm ,Q))\zeta (_{\underset{¯}{W}_2^{\mathrm{}}}\theta )(\underset{¯}{w}(P_{0,\pm },Q))\zeta ^2`$
$`+2^1(_{\underset{¯}{W}_1^{\mathrm{}}}^2\theta )(\underset{¯}{w}(P_{0,\pm },Q))\zeta ^2+O(\zeta ^3)\text{ as }PP_\mathrm{}_\pm .`$ (4.28)
Next, observing the fact that
$$\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}=d\text{log}\left(\frac{\theta (\underset{¯}{w}(,P_{0,}))}{\theta (\underset{¯}{w}(,P_{\mathrm{}_{}}))}\right),$$
(4.29)
it becomes a straightforward matter deriving (4.16)–(4.23). For simplicity we just focus on the expansion of $`_{Q_0}^P\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}`$ as $`PP_{0,\pm }`$, the rest is completely analogous. Using
$`\underset{¯}{w}(Q_0,P_{0,\pm })`$ $`=\underset{¯}{\mathrm{{\rm Y}}}\pm \underset{¯}{\mathrm{\Delta }}_0,\underset{¯}{w}(Q_0,P_\mathrm{}_\pm )=\underset{¯}{\mathrm{{\rm Y}}}\pm \underset{¯}{\mathrm{\Delta }}_{\mathrm{}},\underset{¯}{w}(Q,Q)=\underset{¯}{\mathrm{{\rm Y}}},Q𝒦_n,`$
$`\underset{¯}{w}(P_\mathrm{}_\sigma ,P_{0,\sigma ^{}})`$ $`=\underset{¯}{\mathrm{{\rm Y}}}+\sigma ^{}\underset{¯}{\mathrm{\Delta }}_0\sigma \underset{¯}{\mathrm{\Delta }}_{\mathrm{}},\underset{¯}{w}(P_{0,\sigma ^{}},P_\mathrm{}_\sigma )=\underset{¯}{\mathrm{{\rm Y}}}\sigma ^{}\underset{¯}{\mathrm{\Delta }}_0+\sigma \underset{¯}{\mathrm{\Delta }}_{\mathrm{}},`$
$`\underset{¯}{w}(P_{0,\sigma },P_{0,\sigma ^{}})`$ $`=\underset{¯}{\mathrm{{\rm Y}}}+(\sigma ^{}\sigma )\underset{¯}{\mathrm{\Delta }}_0,\underset{¯}{w}(P_{0,\sigma ^{}},P_{0,\sigma })=\underset{¯}{\mathrm{{\rm Y}}}+(\sigma \sigma ^{})\underset{¯}{\mathrm{\Delta }}_0,`$
$`\underset{¯}{w}(P_\mathrm{}_\sigma ,P_\mathrm{}_\sigma ^{})`$ $`=\underset{¯}{\mathrm{{\rm Y}}}+(\sigma ^{}\sigma )\underset{¯}{\mathrm{\Delta }}_{\mathrm{}},\underset{¯}{w}(P_\mathrm{}_\sigma ^{},P_\mathrm{}_\sigma )=\underset{¯}{\mathrm{{\rm Y}}}+(\sigma \sigma ^{})\underset{¯}{\mathrm{\Delta }}_{\mathrm{}},`$ (4.30)
$`\sigma ,\sigma ^{}\{1,1\},`$
and (4.27)–(4.29), one computes by comparison with (4.5),
$`{\displaystyle _{Q_0}^P}\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}={\displaystyle _{Q_0}^P}𝑑\text{log}\left({\displaystyle \frac{\theta (\underset{¯}{w}(P^{},P_{0,}))}{\theta (\underset{¯}{w}(P^{},P_{\mathrm{}_{}}))}}\right)`$
$`=\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{w}(P,P_{0,}))}{\theta (\underset{¯}{w}(P,P_{\mathrm{}_{}}))}}\right)\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{w}(Q_0,P_{0,}))}{\theta (\underset{¯}{w}(Q_0,P_{\mathrm{}_{}}))}}\right)`$
$`=\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{w}(P,P_{0,}))}{\theta (\underset{¯}{w}(P,P_{\mathrm{}_{}}))}}\right)\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}}\right)`$
$`=\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}2\underset{¯}{\mathrm{\Delta }}_0)\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}}\right)_{\underset{¯}{W}_1^0}\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}2\underset{¯}{\mathrm{\Delta }}_0)}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}}\right)\zeta +O(\zeta ^2)`$
$`=\omega _0^{0,+}\omega _1^0\zeta +O(\zeta ^2)\text{ as }PP_{0,+}.`$ (4.31)
This proves (4.16) and (4.19). Similarly, one calculates,
$`{\displaystyle _{Q_0}^P}\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}=\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{w}(P,P_{0,}))}{\theta (\underset{¯}{w}(P,P_{\mathrm{}_{}}))}}\right)\mathrm{ln}\left({\displaystyle \frac{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}{\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}}\right)`$
$`=\mathrm{ln}(\zeta )+\mathrm{ln}\left({\displaystyle \frac{(_{\underset{¯}{W}_1^0}\theta )(\underset{¯}{\mathrm{{\rm Y}}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})}{\theta (\underset{¯}{\mathrm{{\rm Y}}}+\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}})\theta (\underset{¯}{\mathrm{{\rm Y}}}\underset{¯}{\mathrm{\Delta }}_0)}}\right)_{\underset{¯}{W}_1^0}\mathrm{ln}(\theta (\underset{¯}{\mathrm{{\rm Y}}}+\underset{¯}{\mathrm{\Delta }}_0\underset{¯}{\mathrm{\Delta }}_{\mathrm{}}))\zeta `$
$`+{\displaystyle \frac{(_{\underset{¯}{W}_2^0}\theta )(\underset{¯}{\mathrm{{\rm Y}}})+2^1(_{\underset{¯}{W}_1^0}^2\theta )(\underset{¯}{\mathrm{{\rm Y}}})}{(_{\underset{¯}{W}_1^0}\theta )(\underset{¯}{\mathrm{{\rm Y}}})}}\zeta +O(\zeta ^2)`$
$`=\mathrm{ln}(\zeta )+\omega _0^{0,}+\omega _1^0\zeta +O(\zeta ^2)\text{ as }PP_{0,},`$ (4.32)
proving (4.17) and (4.18). ∎
The results of Lemma 4.1 can conveniently be reformulated in terms of theta functions with characteristics associated with the vector $`\underset{¯}{\mathrm{{\rm Y}}}`$, but we omit further details at this point.
Combining (3.11) and Theorem A.5, the theta function representation of $`\varphi `$ must be of the form
$$\begin{array}{c}\varphi (P,x,t)=C(x,t)\frac{\theta \left(\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)})\right)}{\theta \left(\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)})\right)}\mathrm{exp}\left(_{Q_0}^P\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}\right),\hfill \\ \hfill P𝒦_n,(x,t)\mathrm{\Omega },\end{array}$$
(4.33)
assuming $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$to be nonspecial for $`(x,t)\mathrm{\Omega }`$, where $`\mathrm{\Omega }^2`$ is open and connected. We refer to Appendix A for our notational conventions concerning Abel maps $`\underset{¯}{A}_{Q_0}`$, $`\underset{¯}{\alpha }_{Q_0}`$ and $`\theta `$-functions. Here $`Q_0𝒦_n\{P_{0,\pm },P_\mathrm{}_\pm \}`$ is a fixed base point which we will always choose among the branch points of $`𝒦_n`$ (e.g., $`Q_0=(E_0,0)`$). Indeed, by (3.11), (4.5), (4.6), and Theorem A.5, $`\varphi (P,x,t)`$ and
$$\frac{\theta \left(\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)})\right)}{\theta \left(\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)})\right)}\mathrm{exp}\left(_{Q_0}^P\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}\right)$$
(4.34)
have the same singularity structure with respect to $`P𝒦_n`$. Moreover, by (A.20), (A.29), and (A.30), the expression (4.34) is single-valued and hence meromorphic on $`𝒦_n`$. Nonspecialty of $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ then yields (4.33).
It remains to analyze the function $`C(x,t)`$ in (4.33) (which is $`P`$-independent) and in the course of that we will also obtain the theta function representations of $`u`$, $`u^{}`$, $`v`$, $`v^{}`$. (The strategy to follow parallels the one used in in connection with algebro-geometric solutions of the AKNS hierarchy.)
In the following it will occasionally be convenient to use a short-hand notation for the arguments of the theta functions in (4.33) and hence we introduce the abbreviation
$`\underset{¯}{z}(P,\underset{¯}{Q})=\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{Q}}),\underset{¯}{Q}=(Q_1,\mathrm{},Q_n)\sigma ^n𝒦_n.`$ (4.35)
Next we show that the Abel maps linearizes the auxiliary divisors $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$.
###### Lemma 4.2.
Assume (2.1), (2.8), (2.9), and (3.2), and $`(x,t),(x_0,t_0)\mathrm{\Omega }`$, where $`\mathrm{\Omega }^2`$ is open and connected. Moreover, suppose $`𝒦_n`$ is nonsingular and $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ are nonspecial for $`(x,t)\mathrm{\Omega }`$. Then
$`\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)})`$ $`=\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(x_0,t_0)})+2i\underset{¯}{c}(n)(xx_0)2i\underset{¯}{c}(1)g_{n+1}^1(tt_0),`$ (4.36)
$`\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)})`$ $`=\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}(x_0,t_0)})+2i\underset{¯}{c}(n)(xx_0)2i\underset{¯}{c}(1)g_{n+1}^1(tt_0).`$ (4.37)
###### Proof.
Given the expansions (A.11) and (A.13) of $`\omega `$ near $`P_\mathrm{}_\pm `$ and $`P_{0,\pm }`$, (4.36) and (4.37) are standard facts following from Lagrange interpolation results of the type (see, e.g., )
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\mu _j^{k1}}{_{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}^n(\mu _j\mu _{\mathrm{}})}}`$ $`=\delta _{k,n},`$
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\mu _j^{k1}\left(_{\begin{array}{c}m=1\\ mj\end{array}}^n\mu _m\right)}{_{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}^n(\mu _j\mu _{\mathrm{}})}}`$ $`=(1)^{n+1}\delta _{k,1},k=1,\mathrm{},n.`$ (4.38)
In Lemma 3.3 we determined the asymptotic behavior of $`\varphi (P,x,t)`$ as $`PP_\mathrm{}_\pm ,P_{0,\pm }`$ comparing (3.4) with (3.12) and (3.13). Now we will recompute the asymtotics of $`\varphi `$ starting from (4.33).
###### Lemma 4.3.
Assume (2.1), (2.8), (2.9), and (3.2). Moreover, suppose $`𝒦_n`$ is nonsingular and $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ are nonspecial for $`(x,t)\mathrm{\Omega }`$, where $`\mathrm{\Omega }^2`$ is open and connected. Then
$`\begin{array}{cc}\hfill \varphi (P,x,t)& \underset{\zeta 0}{=}C(x,t)e^{\omega _0^{\mathrm{}_{}}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\zeta ^1\hfill \\ & +C(x,t)\omega _1^{\mathrm{}}e^{\omega _0^{\mathrm{}_{}}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\hfill \\ & C(x,t)e^{\omega _0^{\mathrm{}_{}}}{\displaystyle \frac{i}{2}}{\displaystyle \frac{}{x}}\left({\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\right)+O(\zeta )\text{ as }PP_{\mathrm{}_{}},\hfill \end{array}`$ (4.39)
$`\begin{array}{cc}\hfill \varphi (P,x,t)& \underset{\zeta 0}{=}C(x,t)e^{\omega _0^\mathrm{}_+}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\hfill \\ & C(x,t)\omega _1^{\mathrm{}}e^{\omega _0^\mathrm{}_+}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\zeta \hfill \\ & +C(x,t)e^{\omega _0^\mathrm{}_+}{\displaystyle \frac{i}{2}}{\displaystyle \frac{}{x}}\left({\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\right)\zeta +O(\zeta ^2)\text{ as }PP_\mathrm{}_+,\hfill \end{array}`$ (4.40)
$`\begin{array}{cc}\hfill \varphi (P,x,t)& \underset{\zeta 0}{=}C(x,t)e^{\omega _0^{0,}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\zeta \hfill \\ & +C(x,t)\omega _1^0e^{\omega _0^{0,}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\zeta ^2\hfill \\ & +C(x,t)e^{\omega _0^{0,}}{\displaystyle \frac{i}{2}}{\displaystyle \frac{}{t}}\left({\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\right)\zeta ^2+O(\zeta ^3)\text{ as }PP_{0,},\hfill \end{array}`$ (4.41)
$`\begin{array}{cc}\hfill \varphi (P,x,t)& \underset{\zeta 0}{=}C(x,t)e^{\omega _0^{0,+}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\hfill \\ & C(x,t)\omega _1^0e^{\omega _0^{0,+}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\zeta \hfill \\ & C(x,t)e^{\omega _0^{0,+}}{\displaystyle \frac{i}{2}}{\displaystyle \frac{}{t}}\left({\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\right)\zeta +O(\zeta ^2)\text{ as }PP_{0,+}.\hfill \end{array}`$ (4.42)
###### Proof.
Using (4.25) and (4.26) (cf. (A.11) and (A.13)) one obtains
$`\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))`$ $`\underset{\zeta 0}{=}\underset{¯}{z}(P_\mathrm{}_\pm ,\underset{¯}{\overset{^}{\mu }}(x,t))\underset{¯}{c}(n)\zeta +O(\zeta ^2)\text{ as }PP_\mathrm{}_\pm ,`$ (4.43)
$`\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))`$ $`\underset{\zeta 0}{=}\underset{¯}{z}(P_{0,\pm },\underset{¯}{\overset{^}{\mu }}(x,t))\underset{¯}{c}(1)g_{n+1}^1\zeta +O(\zeta ^2)\text{ as }PP_{0,\pm }`$ (4.44)
and hence
$`\begin{array}{cc}\hfill \theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))\right)& \underset{\zeta 0}{=}\theta \left(\underset{¯}{z}(P_\mathrm{}_\pm ,\underset{¯}{\overset{^}{\mu }}(x,t))\right)\hfill \\ & \pm {\displaystyle \frac{i}{2}}{\displaystyle \frac{}{x}}\theta \left(\underset{¯}{z}(P_\mathrm{}_\pm ,\underset{¯}{\overset{^}{\mu }}(x,t))\right)\zeta +O(\zeta ^2)\text{ as }PP_\mathrm{}_\pm ,\hfill \end{array}`$ (4.45)
$`\begin{array}{cc}\hfill \theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))\right)& \underset{\zeta 0}{=}\theta \left(\underset{¯}{z}(P_{0,\pm },\underset{¯}{\overset{^}{\mu }}(x,t))\right)\hfill \\ & {\displaystyle \frac{i}{2}}{\displaystyle \frac{}{t}}\theta \left(\underset{¯}{z}(P_{0,\pm },\underset{¯}{\overset{^}{\mu }}(x,t))\right)\zeta +O(\zeta ^2)\text{ as }PP_{0,\pm }.\hfill \end{array}`$ (4.46)
Here we used (4.36) to convert the directional derivatives $`_{j=1}^nc_j(n)/w_j`$ and $`_{j=1}^nc_j(1)/w_j`$, $`\underset{¯}{w}=(w_1,\mathrm{},w_n)^n`$ into $`/x`$ and $`/t`$ derivatives. Since by (4.37) exactly the same formulas (4.45) and (4.46) apply to $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$, insertion of (4.5), (4.6), (4.45), and (4.46) (and their $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ analogs) into (4.33) proves (4.39)–(4.42). ∎
Lemma 4.3 may seem to be just another asymptotic result, however, a comparison with Lemma 3.3 reveals that in passing we have actually derived the theta function representations for $`u`$, $`v`$, $`u^{}`$, and $`v^{}`$.
###### Theorem 4.4.
Assume (2.1), (2.8), (2.9), and (3.2), and suppose $`𝒦_n`$ is nonsingular. In addition, let $`P𝒦_n`$ and $`(x,t)\mathrm{\Omega }`$, where $`\mathrm{\Omega }^2`$ is open and connected. Moreover, suppose that $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ are nonspecial for $`(x,t)\mathrm{\Omega }`$. Then $`\varphi (P,x,t)`$ admits the representation
$$\varphi (P,x,t)=C_0e^{2i(\omega _1^{\mathrm{}}x\omega _1^0t)}\frac{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}\mathrm{exp}\left(_{Q_0}^P\omega _{P_{0,},P_{\mathrm{}_{}}}^{(3)}\right)$$
(4.47)
for some constant $`C_0\{0\}`$ and the theta function representations for the algebro-geometric solutions $`u`$, $`u^{}`$, $`v`$, and $`v^{}`$ of the classical massive Thirring system (2.22)–(2.25) read
$`u(x,t)`$ $`=C_0^1e^{\omega _0^{0,+}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}}e^{2i(\omega _1^{\mathrm{}}x\omega _1^0t)},`$ (4.48)
$`v(x,t)`$ $`=C_0^1e^{\omega _0^{\mathrm{}_{}}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}}e^{2i(\omega _1^{\mathrm{}}x\omega _1^0t)},`$ (4.49)
$`u^{}(x,t)`$ $`=C_0e^{\omega _0^{0,}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}e^{2i(\omega _1^{\mathrm{}}x\omega _1^0t)},`$ (4.50)
$`v^{}(x,t)`$ $`=C_0e^{\omega _0^\mathrm{}_+}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}e^{2i(\omega _1^{\mathrm{}}x\omega _1^0t)},`$ (4.51)
with $`\omega _0^{0,\pm },\omega _1^0,\omega _0^\mathrm{}_\pm `$, and $`\omega _1^{\mathrm{}}`$ given by (4.16)–(4.23) (cf. also (4.5)–(4.8)).
###### Proof.
A comparison of (3.29)–(3.32) and (4.39)–(4.42) yields
$$C_x(x,t)=2i\omega _1^{\mathrm{}}C(x,t),C_t(x,t)=2i\omega _1^0C(x,t),$$
(4.52)
and hence
$$C(x,t)=C_0e^{2i(\omega _1^{\mathrm{}}x\omega _1^0t)},$$
(4.53)
proves (4.47). Insertion of (4.53) into the leading asymptotic term of (4.39)–(4.42) then yields (4.48)–(4.51). ∎
###### Remark 4.5.
(i) The constant $`C_0`$ in (4.47)–(4.51) remains open due to the scaling invariance (2.28) of the Thirring system. One can rewrite (4.48)–(4.51) in the form
$`\begin{array}{cc}\hfill u(x,t)& =u(x_0,t_0){\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}}\times \hfill \\ & \times \mathrm{exp}\left(2i(\omega _1^{\mathrm{}}(xx_0)\omega _1^0(tt_0))\right),\hfill \end{array}`$ (4.54)
$`\begin{array}{cc}\hfill v(x,t)& =v(x_0,t_0){\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}}\times \hfill \\ & \times \mathrm{exp}\left(2i(\omega _1^{\mathrm{}}(xx_0)\omega _1^0(tt_0))\right),\hfill \end{array}`$ (4.55)
$`\begin{array}{cc}\hfill u^{}(x,t)& =u^{}(x_0,t_0){\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\nu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\times \hfill \\ & \times \mathrm{exp}\left(2i(\omega _1^{\mathrm{}}(xx_0)\omega _1^0(tt_0))\right),\hfill \end{array}`$ (4.56)
$`\begin{array}{cc}\hfill v^{}(x,t)& =v^{}(x_0,t_0){\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}}\times \hfill \\ & \times \mathrm{exp}\left(2i(\omega _1^{\mathrm{}}(xx_0)\omega _1^0(tt_0))\right),\hfill \end{array}`$ (4.57)
where
$`\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\mu }}(x,t))`$ $`=\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))+2i\underset{¯}{c}(n)(xx_0)+2i\underset{¯}{c}(1)g_{n+1}^1(tt_0),`$ (4.58)
$`\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\nu }}(x,t))`$ $`=\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\nu }}(x_0,t_0))+2i\underset{¯}{c}(n)(xx_0)+2i\underset{¯}{c}(1)g_{n+1}^1(tt_0),`$ (4.59)
by (4.35), (4.36), and (4.37).
(ii) Since the divisors $`𝒟_{P_{0,}\underset{¯}{\overset{^}{\nu }}(x,t)}`$ and $`𝒟_{P_{\mathrm{}_{}}\underset{¯}{\overset{^}{\mu }}(x,t)}`$ are linearly independent by (3.11), one infers
$`\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)})=\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)})+\underset{¯}{\mathrm{\Delta }},(x,t)\mathrm{\Omega },\underset{¯}{\mathrm{\Delta }}=\underset{¯}{A}_{P_{0,}}(P_{\mathrm{}_{}}).`$ (4.60)
Hence on can replace $`\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\nu }}(x,t))`$ in (4.47)–(4.51), (4.54)–(4.57), (4.59) in terms of $`\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\mu }}(x,t))`$ according to
$$\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\nu }}(x,t))=\underset{¯}{z}(Q,\underset{¯}{\overset{^}{\mu }}(x,t))+\underset{¯}{\mathrm{\Delta }}.$$
(4.61)
In principle, Theorem 4.4 completes the primary aim of this paper, the derivation of the theta function representation of algebro-geometric solutions of the classical massive Thirring system (2.22)–(2.25). The reader will have noticed that our approach thus far is nontraditional in the sense that we did not use the Baker–Akhiezer vector $`\mathrm{\Psi }`$ at all, but instead put all emphasis on the meromorphic $`\varphi `$ on $`𝒦_n`$. Just for completeness we finally derive the theta function representation for $`\psi _1`$ in (3.18).
The singularity structure of $`\psi _1(P,x,x_0,t,t_0)`$ near $`P_\mathrm{}_\pm `$ displayed in Lemma 3.4 suggests introducing Abelian differentials $`\omega _{Q,0}^{(2)}`$ of the second kind, normalized by the vanishing of their $`a`$-periods,
$$_{a_j}\omega _{Q,0}^{(2)}=0,j=1,\mathrm{},n,$$
(4.62)
with a second-order pole at $`Q`$ of the type
$$\omega _{Q,0}^{(2)}\underset{\zeta 0}{=}(\zeta ^2+O(1))d\zeta \text{ as }PQ,$$
(4.63)
and holomorphic on $`𝒦_n\{Q\}`$. More precisely, we introduce
$`\mathrm{\Omega }_{\mathrm{},0}^{(2)}`$ $`=\omega _{P_\mathrm{}_+,0}^{(2)}\omega _{P_{\mathrm{}_{}},0}^{(2)},`$ (4.64)
$`\mathrm{\Omega }_{0,0}^{(2)}`$ $`=\omega _{P_{0,+},0}^{(2)}\omega _{P_{0,},0}^{(2)},`$ (4.65)
and note that
$`{\displaystyle _{Q_0}^P}\mathrm{\Omega }_{\mathrm{},0}^{(2)}`$ $`\underset{\zeta 0}{=}\pm (\zeta ^1+e_{\mathrm{},0}+e_{\mathrm{},1}\zeta +O(\zeta ^2))\text{ as }PP_{\mathrm{}_{}},`$ (4.66)
$`{\displaystyle _{Q_0}^P}\mathrm{\Omega }_{0,0}^{(2)}`$ $`\underset{\zeta 0}{=}\pm (\zeta ^1+e_{0,0}+e_{0,1}\zeta +O(\zeta ^2))\text{ as }PP_{0,}.`$ (4.67)
###### Theorem 4.6.
Assume (2.1), (2.8), (2.9), (3.2), and suppose $`𝒦_n`$ is nonsingular. In addition, let $`P𝒦_n\{P_{0,\pm },P_\mathrm{}_\pm \}`$ and $`(x,t),(x_0,t_0)\mathrm{\Omega }`$, where $`\mathrm{\Omega }^2`$ is open and connected. Moreover, suppose that $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ are nonspecial for $`(x,t)\mathrm{\Omega }`$. Then $`\psi _1(P,x,x_0,t,t_0)`$ admits the representation
$$\begin{array}{c}\psi _1(P,x,x_0,t,t_0)=\left(\frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}\right)^{1/2}\times \hfill \\ \hfill \times \frac{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}\times \\ \hfill \times \mathrm{exp}\left(i(xx_0)\left(\omega _1^{\mathrm{}}+_{Q_0}^P\mathrm{\Omega }_{\mathrm{},0}^{(2)}\right)+i(tt_0)\left(\omega _1^0+_{Q_0}^P\mathrm{\Omega }_{0,0}^{(2)}\right)\right),\end{array}$$
(4.68)
or, equivalently,
$$\begin{array}{c}\psi _1(P,x,x_0,t,t_0)=\left(\frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}\right)^{1/2}\times \hfill \\ \hfill \times \frac{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}\times \\ \hfill \times \mathrm{exp}\left(i(xx_0)\left(\omega _1^{\mathrm{}}+_{Q_0}^P\mathrm{\Omega }_{\mathrm{},0}^{(2)}\right)+i(tt_0)\left(\omega _1^0+_{Q_0}^P\mathrm{\Omega }_{0,0}^{(2)}\right)\right).\end{array}$$
(4.69)
###### Proof.
Introducing
$$\begin{array}{c}\widehat{\psi }_1(P,x,x_0,t,t_0)=\frac{C(x,t)}{C(x_0,t_0)}\frac{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}\times \hfill \\ \hfill \times \mathrm{exp}\left(i(xx_0)_{Q_0}^P\mathrm{\Omega }_{\mathrm{},0}^{(2)}+i(tt_0)_{Q_0}^P\mathrm{\Omega }_{0,0}^{(2)}\right),\\ \hfill P𝒦_n\{P_{0,\pm },P_\mathrm{}_\pm \},(x,t),(x_0,t_0)\mathrm{\Omega },\end{array}$$
(4.70)
with an appropriate normalization $`C(x,t)`$ (which is $`P`$-independent) to be determined later, we next intend to prove that
$$\psi _1(P,x,x_0,t,t_0)=\widehat{\psi }_1(P,x,x_0,t,t_0),P𝒦_n\{P_{0,\pm },P_\mathrm{}_\pm \},(x,t),(x_0,t_0)\mathrm{\Omega }.$$
(4.71)
A comparison of (3.23), (3.37), (3.38), (4.66), (4.67), and (4.70) shows that $`\psi _1`$ and $`\widehat{\psi }_1`$ share the identical essential singularity near $`P_\mathrm{}_\pm `$. Next we turn to the local bahavior of $`\psi _1(P,x,x_0,t,t_0)`$ with respect to its zeros and poles. We temporarily restrict $`\mathrm{\Omega }`$ to $`\stackrel{~}{\mathrm{\Omega }}\mathrm{\Omega }`$ such that for all $`(x^{},s)\stackrel{~}{\mathrm{\Omega }}`$, $`\mu _j(x^{},s)\mu _k(x^{},s)`$ for all $`jk`$, $`j,k=1,\mathrm{},n`$. Then arguing as in the paragraph following (3.28) one infers from (3.18) that
$`\psi _1(P,x,x_0,t,t_0)=\{\begin{array}{cc}(\mu _j(x,t)z)O(1)\hfill & \text{ as }P\widehat{\mu }_j(x,t)\widehat{\mu }(x_0,t_0),\hfill \\ O(1)\hfill & \text{ as }P\widehat{\mu }_j(x,t)=\widehat{\mu }(x_0,t_0),\hfill \\ (\mu _j(x_0,t_0)z)^1O(1)\hfill & \text{ as }P\widehat{\mu }_j(x_0,t_0)\widehat{\mu }(x,t),\hfill \end{array}`$
$`P=(z,y)𝒦_n,(x,t),(x_0,t_0)\stackrel{~}{\mathrm{\Omega }},`$ (4.72)
where $`O(1)0`$. Applying Lemma A.6 then proves (4.71) for $`(x,t),(x_0,t_0)\stackrel{~}{\mathrm{\Omega }}`$. By continuity this extends to $`(x,t),(x_0,t_0)\mathrm{\Omega }`$ as long as $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}\sigma ^n𝒦_n`$ remains nonspecial. Finally we determine $`C(x,t)/C(x_0,t_0)`$. A comparison of (2.5), (2.21), (3.24), (4.54), (4.55), and (4.70) yields
$`\psi _1(P,x,x_0,t,t_0)\psi _1(P^{},x,x_0,t,t_0)`$
$`\underset{z\mathrm{}}{=}{\displaystyle \frac{C(x,t)^2}{C(x_0,t_0)^2}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}}\times `$
$`\times \mathrm{exp}\left(2i((xx_0)e_{\mathrm{},0}(tt_0)e_{0,0})\right)`$
$`={\displaystyle \frac{v(x,t)}{v(x_0,t_0)}}`$ (4.73)
$`={\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}}\mathrm{exp}\left(2i((xx_0)\omega _1^{\mathrm{}}(tt_0)\omega _1^0)\right)`$
and
$`\psi _1(P,x,x_0,t,t_0)\psi _1(P^{},x,x_0,t,t_0)`$
$`\underset{z0}{=}{\displaystyle \frac{C(x,t)^2}{C(x_0,t_0)^2}}{\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)}}\times `$
$`\times \mathrm{exp}\left(2i((xx_0)e_{\mathrm{},0}(tt_0)e_{0,0})\right)`$
$`={\displaystyle \frac{u(x,t)}{u(x_0,t_0)}}`$ (4.74)
$`={\displaystyle \frac{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}}\times `$
$`\times \mathrm{exp}\left(2i((xx_0)\omega _1^{\mathrm{}}(tt_0)\omega _1^0)\right).`$
Thus, (4.73) implies
$$\begin{array}{c}\frac{C(x,t)^2}{C(x_0,t_0)^2}=\frac{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}\times \hfill \\ \hfill \times \mathrm{exp}\left(2i((xx_0)\omega _1^{\mathrm{}}(tt_0)\omega _1^0)\right)\end{array}$$
(4.75)
and (4.74) yields
$$\begin{array}{c}\frac{C(x,t)^2}{C(x_0,t_0)^2}=\frac{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x_0,t_0))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x_0,t_0))\right)}{\theta \left(\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))\right)\theta \left(\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t))\right)}\times \hfill \\ \hfill \times \mathrm{exp}\left(2i((xx_0)\omega _1^{\mathrm{}}(tt_0)\omega _1^0)\right).\end{array}$$
(4.76)
In order to reconcile the two expressions (4.75) and (4.76) for $`C(x,t)^2/C(x_0,t_0)^2`$ it suffices to recall the linear dependence of the divisors $`𝒟_{P_{\mathrm{}_{}}\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{P_{0,}\underset{¯}{\overset{^}{\nu }}(x,t)}`$, that is,
$`\underset{¯}{A}_{Q_0}(P_{\mathrm{}_{}})`$ $`+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)})=\underset{¯}{A}_{Q_0}(P_{0,})+\underset{¯}{\alpha }_{Q_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}),`$ (4.77)
and
$`\underset{¯}{A}_{Q_0}(P_{0,})`$ $`=\underset{¯}{A}_{Q_0}(P_{0,+}),\underset{¯}{A}_{Q_0}(P_{\mathrm{}_{}})=\underset{¯}{A}_{Q_0}(P_\mathrm{}_+),`$ (4.78)
to conclude that
$$\underset{¯}{z}(P_\mathrm{}_+,\underset{¯}{\overset{^}{\mu }}(x,t))=\underset{¯}{z}(P_{0,+},\underset{¯}{\overset{^}{\nu }}(x,t)),\underset{¯}{z}(P_{0,},\underset{¯}{\overset{^}{\mu }}(x,t))=\underset{¯}{z}(P_{\mathrm{}_{}},\underset{¯}{\overset{^}{\nu }}(x,t))$$
(4.79)
and hence equality of the right-hand sides of (4.75) and (4.76). This proves (4.68) and (4.69). ∎
The explicit representation (4.68) for $`\psi _1`$ complements Lemma 3.2 and shows that $`\psi _1`$ stays meromorphic on $`𝒦_n\{P_\mathrm{}_\pm \}`$ as long as $`𝒟_{\underset{¯}{\overset{^}{\mu }}(x,t)}`$ and $`𝒟_{\underset{¯}{\overset{^}{\nu }}(x,t)}`$ are nonspecial (assuming $`𝒦_n`$ to be nonsingular). An analogous theta function derivation can be performed for $`\zeta \psi _2(P,\zeta ,x,x_0,t,t_0)`$, but we omit further details at this point.
We emphasize that $`\varphi (P,x,t)`$ and $`\psi _1(P,x,x_0,t,t_0)`$ are naturally defined on the two-sheeted Riemann surface $`𝒦_n`$, whereas $`\psi _2(P,\zeta ,x,x_0,t,t_0)`$ requires a four-sheeted Riemann surface due to the additional factor $`1/z^{1/2}`$ in (3.19). In particular, the Baker–Akhiezer vector $`\mathrm{\Psi }(P,\zeta ,x,x_0,t,t_0)`$ in (3.17) requires a four-sheeted Riemann surface, clearly a disadvantage when compared to our use of $`\varphi (P,x,t)`$. Finally, we note that reality constraints of the type (2.26) and their effects on algebro-geometric quantities, such as pairs of real and complex conjugate branch points of $`𝒦_n`$, etc., are discussed in (see also ).
We conclude with the elementary genus zero example (i.e., $`n=0`$), a case thus far excluded in this section.
###### Example 4.7.
Assume $`n=0`$. Then
$`𝒦_0:_0(z,y)`$ $`=y^2R_2(z)=y^2(zE_0)(zE_1)=0,`$ (4.80)
$`c_1`$ $`=(E_0+E_1)/2,g_1=(E_0E_1)^{1/2},`$ (4.81)
$`\omega _1^{\mathrm{}}`$ $`=(g_1+c_1)/2,\omega _1^0=\omega _1^{\mathrm{}}/(E_0E_1),`$ (4.82)
$`v(x,t)`$ $`=v(x_0,t_0)\mathrm{exp}(2i(\omega _1^{\mathrm{}}(xx_0)\omega _1^0(tt_0)))`$
$`=g_1u(x,t),`$ (4.83)
$`v^{}(x,t)`$ $`=v^{}(x_0,t_0)\mathrm{exp}(2i(\omega _1^{\mathrm{}}(xx_0)\omega _1^0(tt_0)))`$
$`=g_1u^{}(x,t),`$ (4.84)
$`v(x,t)v^{}(x,t)`$ $`=(c_1g_1)/2=g_1^2u(x,t)u^{}(x,t),`$ (4.85)
$`\varphi (P,x,t)`$ $`={\displaystyle \frac{y(P)+z+g_1}{2v(x,t)}}={\displaystyle \frac{2v^{}(x,t)z}{y(P)zg_1}},`$ (4.86)
$`\psi _1(P,x,x_0,t,t_0)`$ $`=\mathrm{exp}(i(xx_0)(y(P)+\omega _1^{\mathrm{}})i(tt_0)(g_1^1z^1y(P)\omega _1^0)).`$ (4.87)
## Appendix A Hyperelliptic curves and their theta functions
We give a brief summary of some of the fundamental properties and notations needed from the theory of hyperelliptic curves. More details can be found in some of the standard textbooks and , as well as monographs dedicated to integrable systems such as , Ch. 2, , App. A–C.
Fix $`n`$. The hyperelliptic curve $`𝒦_n`$ of genus $`n`$ used in Sections 3 and 4 is defined by
$$\begin{array}{c}𝒦_n:_n(z,y)=y^2R_{2n+2}(z)=0,R_{2n+2}(z)=\underset{m=0}{\overset{2n+1}{}}(zE_m),\hfill \\ \hfill \{E_m\}_{m=0,\mathrm{},2n+1},E_mE_n\text{ for }mn,m,n=0,\mathrm{},2n+1.\end{array}$$
(A.1)
The curve (A.1) is compactified by adding the points $`P_\mathrm{}_+`$ and $`P_{\mathrm{}_{}}`$, $`P_\mathrm{}_+P_{\mathrm{}_{}}`$, at infinity. One then introduces an appropriate set of $`n+1`$ nonintersecting cuts $`𝒞_j`$ joining $`E_{m(j)}`$ and $`E_{m^{}(j)}`$. We denote
$$𝒞=\underset{j\{1,\mathrm{},n+1\}}{}𝒞_j,𝒞_j𝒞_k=\mathrm{},jk.$$
(A.2)
Define the cut plane
$$\mathrm{\Pi }=𝒞,$$
(A.3)
and introduce the holomorphic function
$$R_{2n+2}()^{1/2}:\mathrm{\Pi },z\left(\underset{m=0}{\overset{2n+1}{}}(zE_m)\right)^{1/2}$$
(A.4)
on $`\mathrm{\Pi }`$ with an appropriate choice of the square root branch in (A.4). Define
$$_n=\{(z,\sigma R_{2n+2}(z)^{1/2})z,\sigma \{\pm 1\}\}\{P_\mathrm{}_+,P_{\mathrm{}_{}}\}$$
(A.5)
by extending $`R_{2n+2}()^{1/2}`$ to $`𝒞`$. The hyperelliptic curve $`𝒦_n`$ is then the set $`_n`$ with its natural complex structure obtained upon gluing the two sheets of $`_n`$ crosswise along the cuts. The set of branch points $`(𝒦_n)`$ of $`𝒦_n`$ is given by
$$(𝒦_n)=\{(E_m,0)\}_{m=0,\mathrm{},2n+1}$$
(A.6)
and finite points $`P`$ on $`𝒦_n`$ are denoted by $`P=(z,y)`$, where $`y(P)`$ denotes the meromorphic function on $`𝒦_n`$ satisfying $`_n(z,y)=y^2R_{2n+2}(z)=0`$. Local coordinates near $`P_0=(z_0,y_0)𝒦_n\{(𝒦_n)\{P_\mathrm{}_+,P_{\mathrm{}_{}}\}\}`$ are given by $`\zeta _{P_0}=zz_0`$, near $`P_\mathrm{}_\pm `$ by $`\zeta _{P_\mathrm{}_\pm }=1/z`$, and near branch points $`(E_{m_0},0)(𝒦_n)`$ by $`\zeta _{(E_{m_0},0)}=(zE_{m_0})^{1/2}`$. The Riemann surface $`𝒦_n`$ defined in this manner has topological genus $`n`$.
One verifies that $`dz/y`$ is a holomorphic differential on $`𝒦_n`$ with zeros of order $`n1`$ at $`P_\mathrm{}_\pm `$ and hence
$$\eta _j=\frac{z^{j1}dz}{y},j=1,\mathrm{},n$$
(A.7)
form a basis for the space of holomorphic differentials on $`𝒦_n`$. Introducing the invertible matrix $`C`$ in $`^n`$,
$`\begin{array}{cc}\hfill C& =(C_{j,k})_{j,k=1,\mathrm{},n},C_{j,k}={\displaystyle _{a_k}}\eta _j,\hfill \\ \hfill \underset{¯}{c}(k)& =(c_1(k),\mathrm{},c_n(k)),c_j(k)=C_{j,k}^1,\hfill \end{array}`$ (A.8)
the corresponding basis of normalized holomorphic differentials $`\omega _j`$, $`j=1,\mathrm{},n`$ on $`𝒦_n`$ is given by
$$\omega _j=\underset{\mathrm{}=1}{\overset{n}{}}c_j(\mathrm{})\eta _{\mathrm{}},_{a_k}\omega _j=\delta _{j,k},j,k=1,\mathrm{},n.$$
(A.9)
Here $`\{a_j,b_j\}_{j=1,\mathrm{},n}`$ is a homology basis for $`𝒦_n`$ with intersection matrix of the cycles satisfying
$$a_jb_k=\delta _{j,k},j,k=1,\mathrm{},n.$$
(A.10)
Near $`P_\mathrm{}_\pm `$ one infers
$`\underset{¯}{\omega }`$ $`=(\omega _1,\mathrm{},\omega _n)=\pm \left({\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{\underset{¯}{c}(j)\zeta ^{nj}}{\left(_{m=0}^{2n+1}(1E_m\zeta )\right)^{1/2}}}\right)d\zeta `$
$`\underset{\zeta 0}{=}\pm \left(\underset{¯}{c}(n)+\left({\displaystyle \frac{1}{2}}\underset{¯}{c}(n){\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m+\underset{¯}{c}(n1)\right)\zeta +O(\zeta ^2)\right)d\zeta ,\zeta =1/z,`$ (A.11)
and
$$y(P)\underset{\zeta 0}{=}\left(1\frac{1}{2}\left(\underset{m=0}{\overset{2n+1}{}}E_m\right)\zeta +O(\zeta ^2)\right)\zeta ^{n1}\text{ as }PP_\mathrm{}_\pm .$$
(A.12)
Similarly, near $`P_{0,\pm }`$ one computes
$`\underset{¯}{\omega }\underset{\zeta 0}{=}\pm {\displaystyle \frac{1}{g_{n+1}}}\left(\underset{¯}{c}(1)+\left({\displaystyle \frac{1}{2}}\underset{¯}{c}(1){\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m^1+\underset{¯}{c}(2)\right)\zeta +O(\zeta ^2)\right)d\zeta ,`$ (A.13)
$`g_{n+1}=\left({\displaystyle \underset{m=0}{\overset{2n+1}{}}}E_m\right)^{1/2},\zeta =z,`$
using
$$y(P)\underset{\zeta 0}{=}\pm g_{n+1}+O(\zeta )\text{ as }PP_{0,\pm },$$
(A.14)
with the sign of $`g_{n+1}`$ determined by the compatibility of charts.
Associated with the homology basis $`\{a_j,b_j\}_{j=1,\mathrm{},n}`$ we also recall the canonical dissection of $`𝒦_n`$ along its cycles yielding the simply connected interior $`\widehat{𝒦}_n`$ of the fundamental polygon $`\widehat{𝒦}_n`$ given by
$$\widehat{𝒦}_n=a_1b_1a_1^1b_1^1a_2b_2a_2^1b_2^1\mathrm{}a_n^1b_n^1.$$
(A.15)
Let $`(𝒦_n)`$ and $`^1(𝒦_n)`$ denote the set of meromorphic functions (0-forms) and meromorphic differentials (1-forms) on $`𝒦_n`$. The residue of a meromorphic differential $`\nu ^1(𝒦_n)`$ at a point $`Q𝒦_n`$ is defined by
$$\text{res}_Q(\nu )=\frac{1}{2\pi i}_{\gamma _Q}\nu ,$$
(A.16)
where $`\gamma _Q`$ is a counterclockwise oriented smooth simple closed contour encircling $`Q`$ but no other pole of $`\nu `$. Holomorphic differentials are also called Abelian differentials of the first kind (dfk). Abelian differentials of the second kind (dsk) $`\omega ^{(2)}^1(𝒦_n)`$ are characterized by the property that all their residues vanish. They are normalized, for instance, by demanding that all their $`a`$-periods vanish, that is,
$$_{a_j}\omega ^{(2)}=0,j=1,\mathrm{},n.$$
(A.17)
If $`\omega _{P_1,n}^{(2)}`$ is a dsk on $`𝒦_n`$ whose only pole is $`P_1\widehat{𝒦}_n`$ with principal part $`\zeta ^{n2}d\zeta `$, $`n_0`$ near $`P_1`$ and $`\omega _j=(_{m=0}^{\mathrm{}}d_{j,m}(P_1)\zeta ^m)d\zeta `$ near $`P_1`$, then
$$\frac{1}{2\pi i}_{b_j}\omega _{P_1,m}^{(2)}=\frac{d_{j,m}(P_1)}{m+1},m=0,1,\mathrm{}$$
(A.18)
Any meromorphic differential $`\omega ^{(3)}`$ on $`𝒦_n`$ not of the first or second kind is said to be of the third kind (dtk). A dtk $`\omega ^{(3)}^1(𝒦_n)`$ is usually normalized by the vanishing of its $`a`$-periods, that is,
$$_{a_j}\omega ^{(3)}=0,j=1,\mathrm{},n.$$
(A.19)
A normal dtk $`\omega _{P_1,P_2}^{(3)}`$ associated with two points $`P_1`$, $`P_2\widehat{𝒦}_n`$, $`P_1P_2`$ by definition has simple poles at $`P_j`$ with residues $`(1)^{j+1}`$, $`j=1,2`$ and vanishing $`a`$-periods. If $`\omega _{P,Q}^{(3)}`$ is a normal dtk associated with $`P`$, $`Q\widehat{𝒦}_n`$, holomorphic on $`𝒦_n\{P,Q\}`$, then
$$\frac{1}{2\pi i}_{b_j}\omega _{P,Q}^{(3)}=_Q^P\omega _j,j=1,\mathrm{},n,$$
(A.20)
where the path from $`Q`$ to $`P`$ lies in $`\widehat{𝒦}_n`$ (i.e., does not touch any of the cycles $`a_j`$, $`b_j`$). Explicitly, one obtains
$`\omega _{P_\mathrm{}_+,P_{\mathrm{}_{}}}^{(3)}`$ $`={\displaystyle \frac{\stackrel{~}{\pi }^nd\stackrel{~}{\pi }}{y}}+{\displaystyle \underset{j=1}{\overset{n}{}}}d_j\omega _j={\displaystyle \frac{_{j=1}^n(\stackrel{~}{\pi }\lambda _j)d\stackrel{~}{\pi }}{y}},`$ (A.21)
$`\omega _{P_1,P_\mathrm{}_+}^{(3)}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{y+y_1}{\stackrel{~}{\pi }z_1}}{\displaystyle \frac{d\stackrel{~}{\pi }}{y}}{\displaystyle \frac{_{j=1}^n(\stackrel{~}{\pi }\widehat{\lambda }_j)d\stackrel{~}{\pi }}{2y}},`$ (A.22)
$`\omega _{P_1,P_{\mathrm{}_{}}}^{(3)}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{y+y_1}{\stackrel{~}{\pi }z_1}}{\displaystyle \frac{d\stackrel{~}{\pi }}{y}}+{\displaystyle \frac{_{j=1}^n(\stackrel{~}{\pi }\stackrel{~}{\lambda }_j)d\stackrel{~}{\pi }}{2y}},`$ (A.23)
$`\omega _{P_1,P_2}^{(3)}`$ $`={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{y+y_1}{\stackrel{~}{\pi }z_1}}{\displaystyle \frac{y+y_2}{\stackrel{~}{\pi }z_2}}\right){\displaystyle \frac{d\stackrel{~}{\pi }}{y}},P_1,P_2𝒦_n\{P_\mathrm{}_+,P_{\mathrm{}_{}}\},`$ (A.24)
where $`d_j,\lambda _j,\widehat{\lambda }_j,\stackrel{~}{\lambda }_j`$, $`j=1,\mathrm{},n`$ are uniquely determined by the requirement of vanishing $`a`$-periods and we abbreviated $`P_j=(z_j,y_j)`$, $`j=1,2`$. (If $`n=0`$, we use the standard convention that the product over an empty index set is replaced by $`1`$.)
We shall always assume (without loss of generality) that all poles of dsk’s and dtk’s on $`𝒦_n`$ lie on $`\widehat{𝒦}_n`$ (i.e., not on $`\widehat{𝒦}_n`$).
Define the matrix $`\tau =(\tau _{j,\mathrm{}})_{j,\mathrm{}=1,\mathrm{},n}`$ by
$$\tau _{j,\mathrm{}}=_{b_j}\omega _{\mathrm{}},j,\mathrm{}=1,\mathrm{},n.$$
(A.25)
Then
$$\text{Im}(\tau )>0,\text{and}\tau _{j,\mathrm{}}=\tau _{\mathrm{},j},j,\mathrm{}=1,\mathrm{},n.$$
(A.26)
Associated with $`\tau `$ one introduces the period lattice
$$L_n=\{\underset{¯}{z}^n\underset{¯}{z}=\underset{¯}{m}+\tau \underset{¯}{n},\underset{¯}{m},\underset{¯}{n}^n\}$$
(A.27)
and the Riemann theta function associated with $`𝒦_n`$ and the given homology basis $`\{a_j,b_j\}_{j=1,\mathrm{},n}`$,
$$\theta (\underset{¯}{z})=\underset{\underset{¯}{n}^n}{}\mathrm{exp}\left(2\pi i(\underset{¯}{n},\underset{¯}{z})+\pi i(\underset{¯}{n},\tau \underset{¯}{n})\right),\underset{¯}{z}^n,$$
(A.28)
where $`(\underset{¯}{u},\underset{¯}{v})=_{j=1}^n\overline{u}_jv_j`$ denotes the scalar product in $`^n`$. It has the fundamental properties
$`\theta (z_1,\mathrm{},z_{j1},z_j,z_{j+1},\mathrm{},z_n)=\theta (\underset{¯}{z}),`$ (A.29)
$`\theta (\underset{¯}{z}+\underset{¯}{m}+\tau \underset{¯}{n})=\mathrm{exp}\left(2\pi i(\underset{¯}{n},\underset{¯}{z})\pi i(\underset{¯}{n},\tau \underset{¯}{n})\right)\theta (\underset{¯}{z}),\underset{¯}{m},\underset{¯}{n}^n.`$ (A.30)
Next, fix a base point $`Q_0𝒦_nP_{0,\pm },P_\mathrm{}_\pm \}`$, denote by $`J(𝒦_n)=^n/L_n`$ the Jacobi variety of $`𝒦_n`$, and define the Abel map $`\underset{¯}{A}_{Q_0}`$ by
$$\underset{¯}{A}_{Q_0}:𝒦_nJ(𝒦_n),\underset{¯}{A}_{Q_0}(P)=(_{Q_0}^P\omega _1,\mathrm{},_{Q_0}^P\omega _n)(modL_n),P𝒦_n.$$
(A.31)
Similarly, we introduce
$$\underset{¯}{\alpha }_{Q_0}:\mathrm{Div}(𝒦_n)J(𝒦_n),𝒟\underset{¯}{\alpha }_{Q_0}(𝒟)=\underset{P𝒦_n}{}𝒟(P)\underset{¯}{A}_{Q_0}(P),$$
(A.32)
where $`\mathrm{Div}(𝒦_n)`$ denotes the set of divisors on $`𝒦_n`$. Here $`𝒟:𝒦_n`$ is called a divisor on $`𝒦_n`$ if $`𝒟(P)0`$ for only finitely many $`P𝒦_n`$. (In the main body of this paper we will choose $`Q_0`$ to be one of the branch points, i.e., $`Q_0(𝒦_n)`$, and for simplicity we will always choose the same path of integration from $`Q_0`$ to $`P`$ in all Abelian integrals.)
In connection with divisors on $`𝒦_n`$ we shall employ the following (additive) notation,
$$\begin{array}{c}𝒟_{Q_0\underset{¯}{Q}}=𝒟_{Q_0}+𝒟_{\underset{¯}{Q}},𝒟_{\underset{¯}{Q}}=𝒟_{Q_1}+\mathrm{}+𝒟_{Q_n},\hfill \\ \hfill \underset{¯}{Q}=(Q_1,\mathrm{},Q_n)\sigma ^n𝒦_n,Q_0𝒦_n,\end{array}$$
(A.33)
where for any $`Q𝒦_n`$,
$$𝒟_Q:𝒦_n_0,P𝒟_Q(P)=\{\begin{array}{cc}1\hfill & \text{for }P=Q,\hfill \\ 0\hfill & \text{for }P𝒦_n\{Q\},\hfill \end{array}$$
(A.34)
and $`\sigma ^n𝒦_n`$ denotes the $`n`$th symmetric product of $`𝒦_n`$. In particular, $`\sigma ^m𝒦_n`$ can be identified with the set of nonnegative divisors $`0𝒟\mathrm{Div}(𝒦_n)`$ of degree $`m`$.
For $`f(𝒦_n)\{0\}`$, $`\omega ^1(𝒦_n)\{0\}`$ the divisors of $`f`$ and $`\omega `$ are denoted by $`(f)`$ and $`(\omega )`$, respectively. Two divisors $`𝒟`$, $`\mathrm{Div}(𝒦_n)`$ are called equivalent, denoted by $`𝒟`$, if and only if $`𝒟=(f)`$ for some $`f(𝒦_n)\{0\}`$. The divisor class $`[𝒟]`$ of $`𝒟`$ is then given by $`[𝒟]=\{\mathrm{Div}(𝒦_n)𝒟\}`$. We recall that
$$\mathrm{deg}((f))=0,\mathrm{deg}((\omega ))=2(n1),f(𝒦_n)\{0\},\omega ^1(𝒦_n)\{0\},$$
(A.35)
where the degree $`\mathrm{deg}(𝒟)`$ of $`𝒟`$ is given by $`\mathrm{deg}(𝒟)=_{P𝒦_n}𝒟(P)`$. It is customary to call $`(f)`$ (respectively, $`(\omega )`$) a principal (respectively, canonical) divisor.
Introducing the complex linear spaces
$`(𝒟)`$ $`=\{f(𝒦_n)f=0\text{ or }(f)𝒟\},r(𝒟)=dim_{}(𝒟),`$ (A.36)
$`^1(𝒟)`$ $`=\{\omega ^1(𝒦_n)\omega =0\text{ or }(\omega )𝒟\},i(𝒟)=dim_{}^1(𝒟),`$ (A.37)
($`i(𝒟)`$ the index of speciality of $`𝒟`$) one infers that $`\mathrm{deg}(𝒟)`$, $`r(𝒟)`$, and $`i(𝒟)`$ only depend on the divisor class $`[𝒟]`$ of $`𝒟`$. Moreover, we recall the following fundamental facts.
###### Theorem A.1.
Let $`𝒟\mathrm{Div}(𝒦_n)`$, $`\omega ^1(𝒦_n)\{0\}`$. Then
$$i(𝒟)=r(𝒟(\omega )),n_0.$$
(A.38)
The Riemann-Roch theorem reads
$$r(𝒟)=\mathrm{deg}(𝒟)+i(𝒟)n+1,n_0.$$
(A.39)
By Abel’s theorem, $`𝒟\mathrm{Div}(𝒦_n)`$, $`n`$ is principal if and only if
$$\mathrm{deg}(𝒟)=0\text{ and }\underset{¯}{\alpha }_{Q_0}(𝒟)=\underset{¯}{0}.$$
(A.40)
Finally, assume $`n`$. Then $`\underset{¯}{\alpha }_{Q_0}:\mathrm{Div}(𝒦_n)J(𝒦_n)`$ is surjective (Jacobi’s inversion theorem).
Next we introduce
$$\underset{¯}{W}_0=\{0\}J(𝒦_n),\underset{¯}{W}_m=\underset{¯}{\alpha }_{Q_0}(\sigma ^m𝒦_n),m$$
(A.41)
and note that while $`\sigma ^m𝒦_n\sigma ^n𝒦_n`$ for $`m<n`$, one has $`\underset{¯}{W}_m\underset{¯}{W}_n`$ for $`m<n`$. Thus $`\underset{¯}{W}_m=J(𝒦_n)`$ for $`mn`$ by Jacobi’s inversion theorem.
Denote by $`\underset{¯}{\mathrm{\Xi }}_{Q_0}=(\mathrm{\Xi }_{Q_{0,1}},\mathrm{},\mathrm{\Xi }_{Q_{0,n}})`$ the vector of Riemann constants,
$$\mathrm{\Xi }_{Q_{0,j}}=\frac{1}{2}(1+\tau _{j,j})\underset{\begin{array}{c}\mathrm{}=1\\ \mathrm{}j\end{array}}{\overset{n}{}}_a_{\mathrm{}}\omega _{\mathrm{}}(P)_{Q_0}^P\omega _j,j=1,\mathrm{},n.$$
(A.42)
###### Theorem A.2.
The set $`\underset{¯}{W}_{n1}+\underset{¯}{\mathrm{\Xi }}_{Q_0}J(𝒦_n)`$ is the complete set of zeros of $`\theta `$ on $`J(𝒦_n)`$, that is,
$$\theta (X)=0\text{ if and only if }X\underset{¯}{W}_{n1}+\underset{¯}{\mathrm{\Xi }}_{Q_0}$$
(A.43)
(i.e., if and only if $`X=(\underset{¯}{\alpha }_{Q_0}(𝒟)+\underset{¯}{\mathrm{\Xi }}_{Q_0}(modL_n)`$ for some $`𝒟\sigma ^{n1}𝒦_n`$). The set $`\underset{¯}{W}_{n1}+\underset{¯}{\mathrm{\Xi }}_{Q_0}`$ has complex dimension $`n1`$.
###### Theorem A.3.
Let $`𝒟_{\underset{¯}{Q}}\sigma ^n𝒦_n`$, $`\underset{¯}{Q}=(Q_1,\mathrm{},Q_n)`$. Then
$$1i(𝒟_{\underset{¯}{Q}})=s(n/2)$$
(A.44)
if and only if there are $`s`$ pairs of the type $`(P,P^{})\{Q_1,\mathrm{},Q_n\}`$ (this includes, of course, branch points for which $`P=P^{}`$).
###### Remark A.4.
While $`\theta (\underset{¯}{z})`$ is well-defined (in fact, entire) for $`\underset{¯}{z}^n`$, it is not well-defined on $`J(𝒦_n)=^n/L_n`$ because of (A.30). Nevertheless, $`\theta `$ is a “multiplicative function” on $`J(𝒦_n)`$ since the multipliers in (A.30) cannot vanish. In particular, if $`\underset{¯}{z}_1=\underset{¯}{z}_2(modL_n)`$, then $`\theta (\underset{¯}{z}_1)=0`$ if and only if $`\theta (\underset{¯}{z}_2)=0`$. Hence it is meaningful to state that $`\theta `$ vanishes at points of $`J(𝒦_n)`$. Since the Abel map $`\underset{¯}{A}_{Q_0}`$ maps $`𝒦_n`$ into $`J(𝒦_n)`$, the function $`\theta (\underset{¯}{A}_{Q_0}(P)\underset{¯}{\xi })`$ for $`\underset{¯}{\xi }^n`$, becomes a multiplicative function on $`𝒦_n`$. Again it makes sense to say that $`\theta (\underset{¯}{A}_{Q_0}()\underset{¯}{\xi })`$ vanishes at points of $`𝒦_n`$.
###### Theorem A.5.
Let $`\underset{¯}{Q}=(Q_1,\mathrm{},Q_n)\sigma ^n𝒦_n`$ and assume $`𝒟_{\underset{¯}{Q}}`$ to be nonspecial, that is, $`i(𝒟_{\underset{¯}{Q}})=0`$. Then
$$\theta (\underset{¯}{\mathrm{\Xi }}_{Q_0}\underset{¯}{A}_{Q_0}(P)+\alpha _{Q_0}(𝒟_{\underset{¯}{Q}}))=0\text{ if and only if }P\{Q_1,\mathrm{},Q_n\}.$$
(A.45)
###### Lemma A.6.
\[7, Lemmas 5.4 and 6.1\] Let $`(x,t),(x_0,t_0)\mathrm{\Omega }`$ for some $`\mathrm{\Omega }^2`$. Assume $`\psi (,x,t_r)`$ to be meromorphic on $`𝒦_n\{P_\mathrm{}_+,P_{\mathrm{}_{}},P_{0,+},P_{0,}\}`$ with essential singularities at $`P_\mathrm{}_\pm `$, $`P_{0,\pm }`$ such that $`\stackrel{~}{\psi }(,x,t)`$ defined by
$$\stackrel{~}{\psi }(P,x,t)=\psi (P,x,t)\mathrm{exp}\left(i(xx_0)_{Q_0}^P\mathrm{\Omega }_{\mathrm{},0}^{(2)}i(tt_0)_{Q_0}^P\mathrm{\Omega }_{0,0}^{(2)}\right)$$
(A.46)
is meromorphic on $`𝒦_n`$ and its divisor satisfies
$$(\stackrel{~}{\psi }(,x,t))𝒟_{\underset{¯}{\overset{^}{\mu }}(x_0,t_0)}.$$
(A.47)
Here $`\mathrm{\Omega }_{\mathrm{},0}^{(2)}`$ and $`\mathrm{\Omega }_{0,0}^{(2)}`$ are defined in (4.64) and (4.65) and the path of integration is chosen identical to that in the Abel maps (A.31) and (A.32)<sup>1</sup><sup>1</sup>1This is to avoid multi-valued expressions and hence the use of the multiplicative Riemann–Roch theorem in the proof of Lemma A.6.. Define a divisor $`𝒟_0(x,t)`$ by
$$(\stackrel{~}{\psi }(,x,t))=𝒟_0(x,t)𝒟_{\underset{¯}{\overset{^}{\mu }}(x_0,t_0)}.$$
(A.48)
Then
$$𝒟_0(x,t)\sigma ^n𝒦_n,𝒟_0(x,t)>0,\mathrm{deg}(𝒟_0(x,t))=n.$$
(A.49)
Moreover, if $`𝒟_0(x,t)`$ is nonspecial for all $`(x,t)\mathrm{\Omega }`$, that is, if
$$i(𝒟_0(x,t))=0,(x,t)\mathrm{\Omega },$$
(A.50)
then $`\psi (,x,t)`$ is unique up to a constant multiple (which may depend on $`x`$ and $`t`$).
###### Theorem A.7.
Suppose $`𝒟_{\underset{¯}{\overset{^}{\mu }}}\sigma ^n𝒦_n`$ is nonspecial, $`\underset{¯}{\overset{^}{\mu }}=(\widehat{\mu }_1,\mathrm{},\widehat{\mu }_n)`$, and $`\widehat{\mu }_{n+1}𝒦_n`$ with $`\widehat{\mu }_{n+1}^{}\{\widehat{\mu }_1,\mathrm{},\widehat{\mu }_n\}`$. Let $`\{\widehat{\lambda }_1,\mathrm{},\widehat{\lambda }_{n+1}\}𝒦_n`$ with $`𝒟_{\underset{¯}{\overset{^}{\lambda }}\widehat{\lambda }_{n+1}}𝒟_{\underset{¯}{\overset{^}{\mu }}\widehat{\mu }_{n+1}}`$ (i.e., $`𝒟_{\underset{¯}{\overset{^}{\lambda }}\widehat{\lambda }_{n+1}}[𝒟_{\underset{¯}{\overset{^}{\mu }}\widehat{\mu }_{n+1}}]`$). Then any $`n`$ points $`\widehat{\nu }_j\{\widehat{\lambda }_1,\mathrm{},\widehat{\lambda }_{n+1}\}`$, $`j=1,\mathrm{},n`$ define a nonspecial divisor $`𝒟_{\underset{¯}{\overset{^}{\nu }}}\sigma ^n𝒦_n`$, $`\underset{¯}{\overset{^}{\nu }}=(\widehat{\nu }_1,\mathrm{},\widehat{\nu }_n)`$.
###### Proof.
Since $`i(𝒟_P)=0`$ for all $`P𝒦_1`$, there is nothing to prove in the special case $`n=1`$. Hence we assume $`n2`$. Let $`P_0(𝒦_n)`$ be a fixed branch point of $`𝒦_n`$ and suppose that $`𝒟_{\underset{¯}{\overset{^}{\nu }}}`$ is special. Then by Theorem A.3 there is a pair $`\{\widehat{\nu },\widehat{\nu }^{}\}\{\widehat{\nu }_1,\mathrm{},\widehat{\nu }_n\}`$ such that
$$\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}})=\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\underset{¯}{\overset{^}{\nu }}}}),$$
(A.51)
where $`\underset{¯}{\underset{¯}{\overset{^}{\nu }}}=(\widehat{\nu }_1,\mathrm{},\widehat{\nu }_n)\{\widehat{\nu },\widehat{\nu }^{}\}\sigma ^{n2}𝒦_n`$. Let $`\widehat{\nu }_{n+1}\{\widehat{\lambda }_1,\mathrm{},\widehat{\lambda }_{n+1}\}\{\widehat{\nu }_1,\mathrm{},\widehat{\nu }_n\}`$ so that $`\{\widehat{\nu }_1,\mathrm{},\widehat{\nu }_{n+1}\}=\{\widehat{\lambda }_1,\mathrm{},\widehat{\lambda }_{n+1}\}\sigma ^{n+1}𝒦_n`$. Then
$`\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\underset{¯}{\overset{^}{\nu }}}\widehat{\nu }_{n+1}})`$ $`=\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\overset{^}{\nu }}\widehat{\nu }_{n+1}})=\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\overset{^}{\lambda }}\widehat{\lambda }_{n+1}})`$
$`=\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}\widehat{\mu }_{n+1}})=\underset{¯}{A}_{P_0}(\widehat{\mu }_{n+1}^{})+\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}}),`$ (A.52)
and hence by Theorem A.2 and (A.52),
$$0=\theta (\underset{¯}{\mathrm{\Xi }}_{P_0}+\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\underset{¯}{\overset{^}{\nu }}}\widehat{\nu }_{n+1}}))=\theta (\underset{¯}{\mathrm{\Xi }}_{P_0}\underset{¯}{A}_{P_0}(\widehat{\mu }_{n+1}^{})+\underset{¯}{\alpha }_{P_0}(𝒟_{\underset{¯}{\overset{^}{\mu }}})).$$
(A.53)
Since by hypothesis $`𝒟_{\underset{¯}{\overset{^}{\mu }}}`$ is nonspecial and $`\widehat{\mu }_{n+1}^{}\{\widehat{\mu }_1,\mathrm{},\widehat{\mu }_n\}`$, (A.53) contradicts Theorem A.5. Thus, $`𝒟_{\underset{¯}{\overset{^}{\nu }}}`$ is nonspecial. ∎
Acknowledgments. It is a great pleasure to dedicate this paper to Sergio Albeverio on the occasion of his 60th birthday. F. G. and H. H. gratefully acknowledge many years of close interaction with an extraordinary mentor and friend, who profoundly affected our early scientific endeavors.
In addition, we also dedicate this paper with admiration to Walter E. Thirring, to honor his influence on the physical sciences in general, and on one of us (F. G.), in particular.
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# From GM law to A powerful mean field scheme
## 1 Introduction
For many years basic mean-field theory has been applied to a huge variety of problems. It is a very simple way to tackle collective phenomena . Most of the time it yieds a correct qualitative description. However, quantitatively the results are very poor. In particular, all aspects of the critical behavior are .grossly misrepresented . Moreover some results are even wrong like for instance the existence of long range order for the Ising system at both one dimension for ferromagnets and at two dimensions for the triangular antiferromagnets.
The Bethe scheme was then introduced to extend the Weiss one spin approach to a cluster of fluctuating spins. For decades it has been looked upon as a solid imrovement. It yields no long range order at one dimension and even becomes exact for the hypercube at infinite dimensions . However quantitative results are yet rather poor. Moreover it was demonstrated recently that the Bethe scheme violates systematically translational invariance . As such it is forbidden by symmetry.
In this paper we present a new and powerful mean field scheme. It embodies the Bethe idea of including a fluctuating spin cluster yet preserving the overall lattice translational invariance. In addition, the connectivity between fluctuating clusters is rescaled according to the GM law introduced few years ago in percolation theory .
Associated critical temperatures are calculated for a large variety of lattices and at several dimensions. Discrepancies with available exact estimates are only within few percent. The lower critical dimension for the onset of phase transitions is found to be $`d_l=1+\frac{2}{q}`$ for Ising systems. It turns to he Golden number $`d_l=\frac{1+\sqrt{5}}{2}`$ for hypercubes ($`q=2d`$). In the case of the triangular Ising antiferromagnet the exact result of no long range order is reproduced
## 2 Revisiting the Bethe aprroximation
Mean-field theory is a one-site approach which first breaks the lattice symmetry by discriminating bewteen fluctuating degrees of freedom and averaged ones . Two interpenetrated lattices are thus defined. Equating the thermal average of fluctuating degree of freedom to the already averaged ones restores the inital lattice symmetry. Simultaneously a self-consistent equation of state is obtained.
To implement a Bethe scheme on a lattice 3 distinct interpenetrated sublattices (A, B, C) must be introduced. First the fluctuating center (A), then the fluctuating nearest neigbhors (B) and last the mean field nearest neigbhors (nn) of the nn (B) not including the center (A). From the A-spin plus these B and C shells, a cell is constituted to pave the whole space and reproduce the full lattice topology.
Cluster center (A) has thus all its nn spins (B) which are fluctuating while surface cluster spins (B) have mean-field nn spins (C) and one nn fluctuating spin (A). Simultaneously mean-field spins (C) have all their nn which are fluctuating spins, making their environnement identical to the cluster center.
On this basis the Bethe requiremnt $`<S_A>=<S_B>`$ is not compatible with the equality $`m_C=<S_B>`$ which should also hold to ensure translational invariance. The Bethe topology is therefore forbidden by symmetry. It is not the case for one-site Weiss theory. For a detailled demonstration see
Last but not least, it is worth noticing it is indeed this very symmetry problem which makes the Bethe approach exact on the Caley tree lattice. This lattice does not exhibit translational invariance by construction. This symmetry breaking was overlooked for several decades.
## 3 A new powerful mean field schme
From the discovery of a systematic Bethe induced symmetry breaking arises the question of the possibility to indeed extend a mean field treatment to more than one site.
Above analysis of the Bethe scheme emphazises the role of the cluster center in the irreversible breaking of the symmetry. It hints to avoid such a fluctuating center. One way to achieve this constraint is to use compact closed linear loops within the lattice topology. For instance compact 4-spin squares and 3-spin triangles for respectively square and triangular lattices.
Each one of these plaquettes is then set respectively as A-species (fluctuating) and B-species (mean field) with a staggered-like coverage pattern. A-plaquettes (B-plaquettes) have thus all their nn plaquettes as B-plaquettes (A-plaquettes).
For a given plaquette, each spin has two nn spins of the same species within the plaquette itself and $`(q2)`$ nn spins of the other species belonging to nn plaquettes. At this stage we have a series of fluctuating one-dimensional closed chains in an external field $`h`$. The number of spins $`N`$ in each chain is determined from lattice topology. It is $`N=4`$, $`N=3`$, $`N=3`$ and $`N=6`$ for respectively square, triangular, Kagom$`\stackrel{´}{e}`$ and Honeycomb lattices.
it is the interactions with nn mean-field spin plaquettes which produce the field $`h`$. We have $`h=\delta Jm`$ where $`\delta `$ accounts for connectivety to B-sublattices, $`J`$ is the nn coupling constant and $`m`$ the averaged magnetization on the B-sublattice. The problem can now be solved exactly. In particular, the chain site magnetization is ,
$$<S_i>=\beta \mathrm{exp}2K\{\frac{(1\mathrm{tanh}(K)^N)}{(1+\mathrm{tanh}(K)^N)}\}h,$$
(1)
at order one in $`h`$. Here $`iA`$-plaquettes. $`K\frac{J}{k_BT}`$ where $`k_B`$ is the Boltzman constant and $`T`$ the temperature.
Putting $`<S_i>=m`$ restores the initial lattice symmetry. It is indeed possible since only two sublattices were involved which was not the case for the 3 sublattice Bethe scheme. The self-consistent equation of state is,
$$m=\delta K\mathrm{exp}2K\{\frac{(1\mathrm{tanh}(K)^N)}{(1+\mathrm{tanh}(K)^N)}\}m+\mathrm{},$$
(2)
at order one in $`m`$ and using $`h=\delta Jm`$. To solve Eq. (2) needs to determine the value of $`\delta `$.
It is then worth to evoke a recent work on percolation thresholds, the GM law . It shows that relevant connectivity variables for site and bond dilution are repectively $`(d1)(q1)`$ and $`\frac{(d1)(q1)}{d}`$. In other words, for site percolation, the number of possible directions $`(q1)`$ from a given site, has to be multiplied by $`(d1)`$. For bond percolation this effective number of site directions has to be divided by dimension $`d`$. Using these variables, the GM law was found to yield all percolation thresholds for all Bravais lattices at all dimensions .
This percolation finding suggests to consider here a rescaled connectivity between closed loops instead of $`\delta =q2`$. Using above counting, we first start with $`q`$ instead of $`(q1)`$ since now dealing with pair exchange interactions and not percolation. Second we renormalize $`q`$ by $`(d1)`$ giving $`q(d1)`$. However the $`2`$ neighboring sites which are treated exactly within the closed loop have to be substracted from the effective number of sites which gives $`q(d1)2`$. Moreover, interactions being related to bonds, we divide this number by $`d`$ as for bond percolation. These considerations lead to a connectivity,
$$\delta =\frac{q(d1)2}{d}.$$
(3)
## 4 Results
We can now check the validity of our simple symmetry preserving model with respect to critical temperatures. From Eq. (2) we get,
$$\delta K_c^G\mathrm{exp}2K_c^G\{\frac{(1\mathrm{tanh}(K_c^G)^N)}{(1+\mathrm{tanh}(K_c^G)^N)}\}=1.$$
(4)
The trivial connectivity counting $`\delta =q2`$ already improves Weiss model. For instance $`K_c^G=0.29`$ in the square case and $`T_c=0`$ at $`d=1`$. We now proceed using Eq. (3) for connectivty.
For the square case ($`q=4,N=4`$), $`\delta =1`$ which gives $`K_c^G=0.4399`$. Exact result is $`K_c^e=0.4407`$. In the case of triangular lattice ($`q=6,N=3`$), $`K_c^G=0.2919`$ with $`\delta =2`$ while the exact estimate is $`K_c^e=0.2746`$. For Kagom$`\stackrel{´}{e}`$ ($`q=4,N=3`$) $`\delta =1`$ yielding $`K_c^G=0.4649`$ for an exact estimate of $`K_c^e=0.4666`$. And $`K_c^G=0.6160`$ for the honeycomb lattice ($`q=3,N=6`$) where $`\delta =\frac{1}{2}`$ for an exact estimate of $`K_c^e=0.6585`$ (see Table I).
Going to $`d=3`$ imposes to restrict the plaquette size to $`N=4`$ since a one-dimensional loop cannot embody a three-dimensional topology. However there exits a $`d`$dependence through Eq. (3). We get $`\delta =\frac{10}{3}`$, $`\delta =\frac{14}{3}`$ and $`\delta =\frac{22}{3}`$ for respectively cubic, fcc and bcc lattices. Corresponding critical temperatures are given by $`K_c^G=0.2012,0.1568,0.1096`$ respectively for exact estimates of $`0.2217,0.1575,`$ and $`0.1021`$ (see Table I).
Critical temperature estimates are available for the hypercube at $`d=5,6,7`$. These are $`K_c^e=0.1139,0.0923,0.0777`$ respectively.
To get the $`d\mathrm{}`$ asymptotic limit of our model we take both $`q\mathrm{}`$ and $`J0`$ under the constraint $`qJ=cst`$. From Eq. (3) connectivity limit is $`\delta q(1\frac{1}{d})`$ which gives always,
$$\delta q$$
(5)
at leading order. Indeed $`q`$ diverges always quicker than $`\frac{q}{d}`$.even for $`fcc`$-lattices where $`q=2d(d1)`$. In turn Eq. (4) becomes,
$$K_c^G=\frac{1}{q},$$
(6)
which is the mean-field result as expected in the $`d\mathrm{}`$ limit.
To evaluate the sensibility on the loop size, it is fruitful to expand Eq. (4) in powers of $`K`$. It gives
$$K_c^G(1+2K_c^G+\mathrm{}+\frac{(2K_c^G)^N}{N!})(1(K_c^G)^N+\mathrm{})(1(K_c^G)^N+\mathrm{})=\frac{1}{\delta }.$$
(7)
Since $`N3`$, a simple analytic expression is obtained only at order one,
$$K_c^G=\frac{1}{\delta }.$$
(8)
At two dimensions Eq. (8) gives $`K_c^G=1,\frac{1}{2},1,2`$ for respectively the square, triangular, Kagom$`\stackrel{´}{e}`$ and for honeycomb lattices. These results are rather poor and shows the importance of the finite value of $`N`$ which embodies part of the lattice topology.
## 5 Conclusion
We have presented a very simple self-consistent model which yields rather good values for critical temperatures within a few percent of exact results. Besides a rescaled lattice connectivity, the finite length of the loops is also taken into account. This new scheme represents a substantial improvement over existing mean-field cluster approximations.
We can also determine from our model a lower critical dimension for phase transitions. It comes from the condition $`h=0`$ for which we have a one-dimensional finite system. Such a system has no long range order at $`T0`$. Phase transitions are thus obtained only in the range $`h0`$ which gives $`q(d1)>2`$ leading to,
$$d_l=1+\frac{2}{q}.$$
(9)
For the Ising hypercube ($`q=2d`$) it becomes the Golden number $`d_l=\frac{1+\sqrt{5}}{2}`$, which excludes the $`d=1`$ case and contains $`d=2`$ as it should be. Last but not least, applying our scheme to the Ising triangular antiferromagnet we do revover the exact result of no long range order at non-zero temperatures; contrary to usual mean field approaches.
References
1. R. K. Pathria, Statistical Mechanics, Pergamon Press (1972)
2. Sh-k Ma, Modern Theory of Critical Phenomena, The Benjamin Inc.: Reading MA (1976)
3. H. A. Bethe, Proc. Roy. Soc. London A150, 552 (1935)
4. P. Weiss, J. Phys. Radium, Paris 6, 667 (1907)
5. S. Galam, Phys. Rev. B54, 1599 (1996)
6. S. Galam and A. Mauger, Phys. Rev. B53, 2171 (1996)
7. S. Galam and P. V. Koseleff, to be published (2000)
8. M. E. Fisher, Repts. Prog. Phys. VXXX (II), 671 (1967)
9. J. Adler, in “Recent developments in computer simulation studies in Condensed matter physics”, VIII, edited by D. P. Landau, Springer (1995)
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# Density dependent strong coupling constant of QCD derived from compact star data
Abstract
The present work is an endeavour to connect the properties of tiny nearly massless objects with those of some of the most massive ones, the compact stars.
Since 1996 there is major influx of X-ray and $`\gamma `$ ray data from binary stars, one or both of which are compact objects that are difficult to explain as neutron stars since they contain a mass M in too small a radius R . The suggestion has been put forward that these are strange quark stars (SS) explainable in a simple model with chiral symmetry restoration (CSR) for the quarks and the M, R and other properties like QPOs (quasi periodic oscillations) in their X-ray power spectrum.
It would be nice if this astrophysical data could shed some light on fundamental properties of quarks obeying QCD. One can relate the strong coupling constant of QCD, $`\alpha _s`$ to the quark mass through the Dyson-Schwinger gap equation using the real time formalism of Dolan and Jackiw. This enables us to obtain the density dependence of $`\alpha _s`$ from the simple CSR referred to above. This way fundamental physics, difficult to extract from other models like for example lattice QCD, can be constrained from present - day compact star data and may be put back to modelling the dense quark phase of early universe. (1) IUCAA, Ganeshkhind, Pune 411 007, India
(2) Azad Physics Centre, Dept. of Physics, Maulana Azad College, Calcutta 700 013, India
(3) Dept. of Physics, Presidency College, Calcutta 700 073, India
$``$ permanent address; 1/10 Prince Golam Md. Road, Calcutta 700 026, India, IUCAA Senior Associate, e-mail : deyjm@giascl01.vsnl.net.in
\** Suported by DST Grant no. SP/S2/K18/96, Govt. of India.
A calculation was done for cold stars in and important conclusions drawn from there about chiral symmetry restoration in QCD when the EOS was used to get SS fitting definite mass-radius (M-R) relations . The empirical M-R relations were derived from astrophysical observations like luminosity variation and some properties of QPO-s in the X-ray power spectrum of these compact stars. The calculations are compared to these stars which emits the X-rays, generated from accretion from their binary partner. The different QPO-s show a correlation for not only these stars but black holes also so that one cannot explain these as being due to magnetic fields or other properties of the stars but some scale-dependent phenomenon in the orbits of the particles which are accreted . One of the compact objects, the SAX J1808.4 - 3658 with period 2.49 millisecond, has been called the holy grail of X-ray astronomy . Its discovery was anticipated for nearly 20 years because magnetospheric disk accretion theory as well as evolutionary ideas concerning the genesis of millisecond radio pulsars strongly suggested that such rapid spin frequencies must occur in accreting magnetic field neutron stars.
The interesting point made by Dey et al. is that starting from an empirical form for the density dependent masses of the up (u), down (d) and strange (s) quarks (q in short) given below, one can constrain the parameter of the form of this mass from recent astronomical data<sup>1</sup><sup>1</sup>1 $`\rho _B=(\rho _u+\rho _d+\rho _s)/3`$ is the baryon number density, $`\rho _0=0.17fm^3`$ is the normal nuclear matter density, and $`\nu `$ is a numerical parameter. The current quark masses $`m_i`$ used in the following are 4, 7 and 150 MeV for u, d and s respectively.:
$$M_i=m_i+M_Qsech(\nu \frac{\rho _B}{\rho _0}),i=u,d,s.$$
(1)
For three values of $`\nu =0.286,0.333`$ and $`0.400`$ the q-star can have a sequence of masses obtained from the standard TOV equations for different choices of the central density. This is shown in Fig. (1).
In other words the masses of stars in units of solar mass, ($`M/M_{}`$), found as a function of the star radius R, calculated using the above eqn.(1), - produces constraints which enable us to restrict the parameter $`\nu `$. At high $`\rho _B`$ the q- mass $`M_i`$ falls from its constituent value $`M_Q`$ to its current one $`m_i`$. The other parameter $`M_Q`$ was taken to be 0.31 $`GeV`$ to match up with constituent quark masses assuming the known fact that the hadrons have very little potential energy. The results are not very sensitive in so far as changing $`M_Q`$ to 0.32 $`GeV`$ changes the maximum mass of the star from 1.43735 M to 1.43738 M and the corresponding radius changes from 7.0553 kms to 7.0558 kms.
It is interesting to plot the up (u), down (d) and strange (s) quark masses at various radii in a star. This is done for a particular value of the parameter $`\nu =1/3`$ already discussed in which gives a sequence of stars falling right into the allowed region of the M-R curve. Fig (2) shows that the quarks do not have the constant current masses assumed in the bag model nor do they have the constituent masses of zero density hadrons. Upto a radius about 2 kms the quarks have their chiral mass but in the major portion of the star their masses are substantially higher. At the surface the strange q- mass is about 0.278 $`GeV`$ and the u,d q- masses $`0.13`$ $`GeV`$.
Given the form for the mass eq. (1) we have to look for a formalism to calculate this mass in QCD. This can be done very conveniently using the real time formalism of Dolan and Jackiw (DJ in short) since here one does not encounter the problem of imaginary chemical potential. As is well known usually people work in the Euclidean space for lattice and other gauge invariant theories. But this excludes finite density calculations since that involves imaginary chemical potentials that makes the action unbounded from below. Using DJ, the price one has to pay is that this is a gauge dependent formalism. We work in the well established Landau gauge used by . We believe in the physical nature of our results and wish to impress upon the reader the robustness of this physicality. In other words we argue that the modelling that we have to invoke does not take away the basic nature of the results that we obtain, namely that the strong coupling constant $`\alpha _s`$ decreases with increasing density and in principle this can be constrained by compact star data if one believes them to be SS.
The Fermion propagator in the DJ formalism is as follows :
$$S_F(p)=\frac{i(p/+m)}{p^2m^2+iϵ}\frac{2\pi \delta (p^2m^2)}{exp[p_0ϵ_F/T]+1}$$
(2)
where $`T`$ is the temperature and $`ϵ_F`$, the chemical potential. There is a similar propagator for the gluon involving the Bose function instead of Fermi but with no chemical potential. One has to calculate the self energy of the quark with this propagator including a gluon loop of four momentum $`k`$ in Fig(3).
Using $`\gamma ^\mu D_{\mu \nu }\gamma ^\nu =3i/k^2`$ and a colour factor 4/3 for the colour group SU(3) one gets
$$m(\rho )=4g_s^2\frac{m(\rho )}{(2\pi )^4}d^4k[\frac{i}{k^2[(pk)^2m(\rho )^2]}\frac{2\pi \delta [(pk)^2m(\rho )^2}{k^2}f(p_0k_0)]$$
(3)
where we have neglected a Boson finite $`T`$-term which deos not contribute and $`f`$ is the Fermi function
$$f(E)=\frac{1}{1+exp[(EE_F)/T]}.$$
(4)
In the rest frame, $`\stackrel{}{p}=0`$. The Fermi function becomes a step function for our case, since the temperature T = 0. On closing the contours in the lower half $`k_0`$ plane, first term reduces to
$$\frac{d^3k}{(2\pi )^3}[\frac{1}{2E[(mE)^2\stackrel{}{k}^2]}\frac{1}{4m\stackrel{}{k}^2}]=\frac{4g_s^2m}{8\pi ^2}\frac{dk}{E}$$
(5)
while the second term becomes, on using the variable $`pkq`$,
$$4g_s^2m\frac{d^4q}{(2\pi )^2}\frac{\delta (q^2m(\rho )^2)}{(pq)^2.(exp[(q_0\mu )\beta ]+1}=\frac{4g_s^2m}{8\pi ^2}\frac{dq}{E}\frac{2}{exp[(q_0E_F)\beta ]+1}$$
(6)
The first integral is logarithmically divergent and therefore must be set to renormalize the $`\rho =0,T1/beta=0`$ quark mass. We do not use a high value of $`g_s^2/4\pi \alpha _s=0.75`$ like since it gives a cut-off which is too small, of magnitude $`1.29GeV`$ only. We prefer a lower $`\alpha _s=0.55`$ and this also gives us a high cut-off $`k_D=2.86GeV`$. We use $`m(\rho =0)=0.33GeV.`$
Now we introduce the second term for finite $`\rho `$ but $`T=0`$ so that the Fermi function reduces to a step function.
$$m(\rho )=m(0)\frac{\alpha _s}{\pi }_0^{E_F}\frac{dE}{\sqrt{E^2m(\rho )^2}}.$$
(7)
Given our form for $`m(\rho )`$, which is $`M_Q`$ of eq.(1), dropping the current quark mass $`m_i`$, $`\alpha _s`$ is now evaluated for all densities. For small density its value increases but this is unnecessary for our model. We get very reasonable values of $`\alpha _s`$ for $`\rho =5\rho _0`$ upwards as can be seen in Table 1 and also in Fig(4). The calculation serves double purposes, it shows that for the concerned densities our choice of the mass function, namely eq.(1), is reasonable and also as already stressed it gives a novel shape for the variation of $`\alpha _s`$ with density not easy to obtain otherwise.
The variation with $`\nu `$ is very small at the surface but deep inside the star when the density increases the values of $`\alpha _s`$ goes down substantially and the change assumes significance. The physicality of the results will enable us to extend our calculation now to $`T0`$ and apply to early universe where temperature of order $`100MeV`$ is expected.
In summary we have shown both in principle and in practice that the strong coupling constant can be derived at various densities once one knows the behaviour of the chiral symmetry restoration for the quark mass constrained from stellar data. The practical model that we have chosen uses an empirical form for chiral restoration at high density proposed and tested against various star properties through the allowed mass radius regions of the compact objects. We can therefore urge the astrophysics data-analysts to pinpoint the mass and radius of compact objects like SAX J1808.8 to a greater precision to help efforts like ours.
It is a pleasure to thank Dr. Arun Thampan for helpful discussions. The stimulation for the work came from a discussion with Prof. Donald Lynden-Bell.
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# Growth enhanced surface diffusion and elastic instability on amorphous solids
## Abstract
A continuum model for growth of solids is developed, considering adatom deposition, surface diffusion, and configuration dependent incorporation rate. For amorphous solids it is related to surface energy densities. The high adatom density leads to growth enhanced dynamics of (a) Mullins’ classical equation \[J. Appl. Phys. 28, 333 (1957)\] without, and (b) of the Asaro-Tiller-Grinfeld-Srolovitz instability with lateral stress in the growing film. The latter mechanism is attributed to morphologies found in recent experiments.
The theoretical approach to kinetic roughening , due to its roots in Statistical Mechanics, relates various growth processes on solids to so called universality classes. Their distinction allows to identify essential properties for the large scale morphology beyond the microscopic mechanisms, say certain symmetries or conservation laws, which manifest themselves in terms of universal scaling exponents.
One such important property is mass conservation suppressing adatom loss to the surrounding space, which allows for surface relaxation only through transport along the surface. The total solid mass then is always equal to its initial value plus the integrated incoming flux.
Molecular Beam Epitaxy (MBE) garantees such conditions, and growth parameters like for instance flux intensity, temperature, chemical composition of the adsorbate can be easily and precisely controlled . It has been widely used and refined for growth of semiconductor and metal crystals. However, under typical conditions the diffusion length and average crystal step distance are quite large, imposing their own features to the morphology , such that any theoretically predicted asymptotic scaling regime lies beyond the reach of experimental observation.
On the other hand it may become observable in growing amorphous substances where it seems plausible that intrinsic lengthscales remain small. Some experiments have been performed in recent years looking for kinetic roughening on amorphous solids . It turns out that their theoretical interpretation in terms of standard continuum and discrete models leaves two main open questions which are addressed in this letter. First, far from equilibrium the kinetics on the surface are different from relaxation by thermally activated adatoms leading to Mullins’ continuum equation . Second, in some experiments there appears a very pronounced intermediate lateral lengthscale together with a mounding instability whose origin was unclear. Obviously it cannot be due to step edge barriers as is the case on crystal surfaces . Another suggested mechanism, deflexion of incoming particle trajectories or “steering” , is in some cases relevant at grazing incidence but negligible for normal beams .
Here a new continuum approach to growth with surface diffusion is developed resulting in a flux dependent coefficient for surface evolution (see Eqs. (8) and (16)), the counterpart to equilibrium relaxation Eq. (1) . The mounding instability is shown to be due to lateral elastic stress in the growing solid , which is responsible for the well known Asaro-Tiller-Grinfeld-Srolovitz (ATGS) instability . Surface modulations allow for relaxation of a laterally stressed solid. For large enough wavelengths this energy gain overcomes the cost in additional capillary energy, disfavouring and destabilizing a flat surface. Due to the nonequilibrium kinetics it appears in a new light: Its driving force is an equilibrium, the dynamical evolution a nonequilibrium phenomenon. In the following calculations the surface energetics account well for the observed wavelength (the size of emerging mounds), and the nonequilibrium kinetics yield the linear growth rate of the instabilty.
To begin, the continuum model of equilibrium relaxation is briefly recalled . Surface configurations are described by a space and time dependent height field $`H(𝐱,t)`$, neglecting overhangs and voids. The driving force for equilibrium relaxation is the surface free energy, $`=d^2x\gamma \sqrt{1+(H)^2}`$. If an atom is added somewhere at the surface this total area may change, resulting in a extra energetic contribution $`\mu =\mathrm{\Omega }\gamma ^2H`$ . $`\mu `$ is often called chemical potential and $`\mathrm{\Omega }=a^3`$ denotes the atomic volume. Spatial variations in $`\mu `$ bias the adatom diffusion and create a macroscopic mass exchange. Local mass (and volume) conservation impose a continuity equation for the surface dynamics, which for $`|\mu |k_BT`$ reads
$$_tH=aD\frac{\mathrm{\Omega }\gamma }{k_BT}^2H\kappa (^2)^2H.$$
(1)
A reasonable assumption for the adatom diffusion coefficient is $`D=a^2/\tau _0e^{\epsilon _{\mathrm{eff}}/k_BT}`$, an Arrhenius form with some activation barrier $`\epsilon _{\mathrm{eff}}`$ (which on a random surface comes from an average over waiting times at each site ) and an “attempt frequency” $`1/\tau _0`$ . Eq. (1) predicts an exponential decay with rate $`\kappa k^4`$ for fluctuations with wavenumber $`k`$, which has been veryfied experimentally . On a phenomenological basis Eq. (1) has also been applied to growth with relaxation by curvature driven surface diffusion, where under deposition (“shot”) noise it results in a power spectrum of height fluctuations $`S(𝐤,t)|H(𝐤,t)|^2|𝐤|^4`$ . Of course under growth the coefficient $`\kappa `$ takes a different form as above in Eq. (1) which is adressed next.
On a growing solid there is a constant supply of mobile adatoms out of the beam hitting the surface. The number of thermally activated adatoms, responsible for equilibrium relaxation, now becomes negligible. Here an ansatz is presented where the incorporation rate $`I\{H\}`$ depends on the surface configuration $`H(𝐱,t)`$ and connects it thus to the adatom density $`\rho (𝐱,t)`$,
$`_tH`$ $`=`$ $`I\{H\}\mathrm{\Omega }\rho `$ (2)
$`_t\rho `$ $`=`$ $`D\{H,\rho \}\rho I\{H\}\rho +F/\mathrm{\Omega }.`$ (3)
Note that $`H+\mathrm{\Omega }\rho `$ is a conserved quantity, increasing by the average growth velocity $`F`$. The problem is to find the dependence $`I\{H\}`$, but before turning to that point its rôle in a phenomenological equation as (1) shall be worked out.
Eqs. (2) are expanded around the “flat” growing surface $`H_0(t)H_0+Ft`$, on which the incorporation rate takes some value $`I_0`$, and around the average adatom density $`\rho _0F/(\mathrm{\Omega }I_0)`$. Small deviations $`h(𝐱,t)H(𝐱,t)H_0(t)`$ and $`c(𝐱,t)\mathrm{\Omega }(\rho (𝐱,t)\rho _0)`$ obey to linear order
$`_th`$ $`=`$ $`{\displaystyle \frac{F}{I_0}}I_1h+I_0c`$ (4)
$`_tc`$ $`=`$ $`{\displaystyle \frac{F}{I_0}}I_1h+(D^2I_0)c,`$ (5)
where $`DD\{H_0(t),\rho _0\}`$. $`I\{H\}I_0I_1h`$ is approximated linearly (not necessarily local in space) and will be evaluated in its Fourier transform $`I_1(𝐤)`$.
In the long wavelength limit, $`Dk^2I_0`$ (beyond the diffusion length $`\mathrm{}_d`$, the typical distance an adatom moves before incorporation), and also $`|I_1(𝐤)|F/I_0I_0`$ , the eigenvalues of (4) are
$`\lambda _1`$ $`=`$ $`{\displaystyle \frac{DF}{I_0^2}}k^2I_1(𝐤)`$ (6)
$`\lambda _2`$ $`=`$ $`I_0`$ (7)
with relative errors $`O(Dk^2/I_0)+O(F|I_k(𝐤)|/I_0^2)`$. $`\lambda _2`$, remaining finite in the long wavelength limit, rules the enhanced (diminished) incorporation under higher (lower) adatom density. The eigenvector of $`\lambda _1`$ mainly lies in direction of $`h`$, so in the considered large scale limit the linearization of Eqs. (2) turns into
$$_th(𝐤,t)=\frac{DF}{I_0^2}k^2I_1(𝐤)h(𝐤,t),$$
(8)
a non-equilibrium version of Eq. (1).
On crystal surfaces $`I_1(𝐤)`$ depends essentially on the step configurations, but on amorphous solids it is proportional to variations in energy density at the surface, and can be derived in a mean field type approach. For a diffusing adatom the amorphous surface consists of potential wells of various depths. They have an average distance $`a`$ from each other and a probability distribution of depths $`n(\epsilon /\epsilon _0)/\epsilon _0`$, where the energy scale $`\epsilon _0`$ is explicitely included in the notation. In a simple view of the nonequilibrium growth process an atom gets deposited, diffuses, and is finally incorporated into the solid in a sufficiently deep potential well. Dimer formation or other types of nucleation can be neglected if sticky sites are denser than adatoms — a condition to be verified a posteriori. An adatom “sticks” to a site, if it cannot escape until it is buried by the further growing solid. Thus the depths $`\epsilon `$ of “sticky” potential wells fulfill $`\tau _0e^{\epsilon /k_BT}>a/F`$. Assuming an exponential distribution of energy depths, $`n(\epsilon /\epsilon _0)=\mathrm{exp}(\epsilon /\epsilon _0)`$ , this yields a relative proportion of sticky sites
$$\frac{a^2}{\mathrm{}_d^2}=\left(\frac{F\tau _0}{a}\right)^{k_BT/\epsilon _0},$$
(9)
which also defines the diffusion length $`\mathrm{}_d`$, because the fractal dimension of a random walk is two and the adatom can “fully explore” its surroundings on the surface. The average incorporation rate is the inverse of the time needed to hit a sticky site, so $`I_0=D/\mathrm{}_d^2`$.
As seen in the context of Eq. (1) deviations from a flat surface change the energy density and the chemical potential near the surface. A natural way to account for this is shifting the energy scale $`\epsilon _0\epsilon _0\mu =\epsilon _0\mathrm{\Omega }`$. A positive additional energy density $``$ lowers the depth of potential wells encountered by the adatoms, and changes the incorporation rate by
$$II_0=I_0\mathrm{log}\frac{a}{F\tau _0}\frac{k_BT}{\epsilon _0}\frac{\mathrm{\Omega }}{\epsilon _0}=2I_0\mathrm{log}\frac{\mathrm{}_d}{a}\frac{\mathrm{\Omega }}{\epsilon _0}.$$
(10)
The change is proportional to the relative energy change for adatoms. Flux and temperature enter via $`\mathrm{}_d`$ and $`I_0`$.
Now the surface energetics have to be evaluated in order to obtain $`I_1(𝐤)`$ in Eq. (8) via (10). As above in Eq. (1) also here the change in surface free energy enters, so there must be a capillary contribution
$$_{\mathrm{cap}}^{\mathrm{lin}}(𝐱,t)=\gamma ^2h(𝐱,t).$$
(11)
The other important part of $``$ comes from elastic deformations of the growing solid. In the experiments considered here amorphous metallic glasses (i.e. the alloy Zr<sub>65</sub>Al<sub>7.5</sub>Cu<sub>27.5</sub>) grow under lateral expansive stress the origin of which is not fully clarified . It builds up during growth and reaches a constant level at a film thickness of about 50 nm. The order of magnitude is then 1 GPa, corresponding to a lateral deformation of $`\alpha 0.1`$ to $`1`$ %, given Young’s modulus to be roughly $`E10^2`$ GPa . Applying the linear relations of stress and strain for small deformations one can calculate the strain and stress tensors in the film with a flat surface, $`u_{ij}^0`$ and $`\sigma _{ij}^0`$, as well as the elastic energy density $`_{\mathrm{el}}^0=\sigma _{ij}^0u_{ij}^0/2`$.
Changes in shape are slow compared to mechanical balancing inside the body, so the strain and stress fields follow the surface configuration quasistatically. Surface variations can be seen as perturbing $`\sigma _{ij}^0`$ by an additional stress field $`\tau _{ij}`$ with boundary conditions $`\tau _{iz}=_ihE\alpha /(1\sigma )`$ for $`i=x,y`$ and $`\tau _{zz}=0`$ to linear order in $`h`$ ($`\sigma `$ without indices denotes the Poisson number). Green’s function corresponding to the geometry of the body yields the perturbative $`\tau _{ij}`$ and corresponding strain $`w_{ij}`$ throughout the solid. In view of Eq. (10) only the corrections to the elastic energy density $`_{\mathrm{el}}^0`$ at the surface to linear order in $`\tau _{ij}`$ and $`w_{ij}`$ (and therefore $`h`$) are needed,
$$_{\mathrm{el}}^{\mathrm{lin}}(𝐱,t)=\frac{\alpha ^2}{\pi }\frac{E}{12\sigma }d^2x^{}\frac{(𝐱𝐱^{})h(𝐱^{},t)}{|𝐱𝐱^{}|^3}.$$
(12)
Curved parts of the surface also compress or elongate the solid, contributing to the energy density by
$$_{\mathrm{el}^{}}^{\mathrm{lin}}(𝐱,t)=\sigma \frac{12\sigma }{1\sigma }\alpha \gamma ^2h(𝐱,t).$$
(13)
The derivation of these elastic energy densities will be presented elsewhere in more detail . Here the next step is to gather Eqs. (11), (12), and (13) to the Fourier transform of the energy density variations at the surface,
$`^{\mathrm{lin}}(𝐤,t)`$ $`=`$ $`B(𝐤)h(𝐤,t)`$ (14)
$`=[(1`$ $`+`$ $`\sigma \alpha {\displaystyle \frac{12\sigma }{1\sigma }})\gamma k^2{\displaystyle \frac{E}{\pi }}\alpha ^2{\displaystyle \frac{1+\sigma }{1\sigma }}|𝐤|]h(𝐤,t),`$ (15)
which is local in $`𝐤`$. With $`B(𝐤)`$ defined above, the linear surface evolution in Eq. (8) including capillary and elastic effects becomes
$$_th(𝐤,t)=2\mathrm{}_d^2\mathrm{log}\frac{\mathrm{}_d}{a}F\frac{\mathrm{\Omega }B(𝐤)}{\epsilon _0}k^2h(𝐤,t).$$
(16)
In particular, without any elastic effects the kinetic coefficient in Eq. (1) is $`\kappa =2\mathrm{}_d^2\mathrm{log}(\mathrm{}_d/a)F(\mathrm{\Omega }\gamma /\epsilon _0)`$. The full $`B(𝐤)`$ reflects the ATGS instability: It is positive for large $`|𝐤|`$, but negative below a critical wavenumber $`k_c`$, attaining a minimum at the wavevectors $`|𝐤|k_{}=3/4\alpha ^2/\pi (1+\sigma )/(1\sigma )E/\gamma `$, where $`B(𝐤)`$ takes the value $`\gamma k_{}^2/3`$. In Eq. (8) because of (even small) initial roughness and noise in the deposition and diffusion processes from this linear instability random patterns of buckles with typical scale $`\lambda _{}=2\pi /k_{}`$ emerge growing exponentially in amplitude with a rate $`1/\tau _{}2/3\mathrm{}_d^2\mathrm{log}(\mathrm{}_d/a)F(\mathrm{\Omega }\gamma /\epsilon _0)k_{}^4=\kappa k_{}^4/3`$.
So far a nonequilibirum dynamic equation for growing surfaces, Eq. (16), has been derived perturbatively close to a horizontal interface. Now it is compared to experimental results . The instability has been carefully observed e.g. for Zr<sub>65</sub>Al<sub>7.5</sub>Cu<sub>27.5</sub> , but the lateral stress is documented only for related materials . The best test would be to measure both simultaneously, but already with the given information an order of magnitude estimate for $`\lambda _{}`$ and $`1/\tau _{}`$ is possible.
Here are the experimental parameters: The elastic constants were given above, the surface tension is $`\gamma 2`$ J/m<sup>2</sup> , temperature $`k_BT4\times 10^{21}`$ J, growth velocity $`F=8`$ Å/s, atomic size $`a=\mathrm{\Omega }^{1/3}3`$ Å. Order of magnitude estimates are used for the time between adatom hop attempts, $`\tau _010^{13}`$ s , and surface energetics, $`\epsilon _010^{19}`$ J (about 1 eV). This yields a diffusion length $`\mathrm{}_d1`$ nm, which suits well the observation that atoms move a few diameters before incorporation . Indeed this is smaller than the average distance between mobile adatoms $`\sqrt{\mathrm{\Omega }/\rho _0}\mathrm{}_d\sqrt{D/F}a\sqrt{\epsilon _0/k_BT}`$ (D is calculated from the distribution of waiting times in “unsticky” sites ), so dimer formation is suppressed. Given these values the theoretical predictions are $`\lambda _{}25`$ nm and $`1/\tau _{}10^2`$ 1/s.
This fits fairly well to experimental observations, as illustrated in Fig. 1. The growing film develops buckles, which after a thickness of about 30 nm take a constant lateral size of $`R_c17`$ nm. From 30 to 240 nm film thickness their vertical amplitude increases exponentially with rate $`6.5\times 10^3`$ 1/s. Surprisingly this quantitative expression of a linear growth instability is not adressed explicitly in the original work , where the authors focus on an observed early time algebraic increase for both quantities. It is caused by kinetic roughening of large $`k`$ modes at times before $`\tau _{}`$ . Besides, in experiments on related materials the lateral homogeneous strain $`\alpha `$ reaches a constant level only after about 50 nm film thickness . So only then the linear instability as described by Eq. (16) with $`B(𝐤)`$ constant in time becomes visible.
Some remarks comparing to different interpretations of the observed instability are in order. First, the elastic energy density in Eq. (12) is calculated for film and substrate of the same material . This does apply to the experimental system, where the Zr alloy film was deposited on a previous 100 nm thick layer of the same material. In particular, effects of perfectly rigid substrates won’t be observed. Second, the partial relaxation of the film close to the modulated surface will not produce a measurable relief of total stress in the layer. Even close to the surface the stress is lowered only by a factor $`1O(|h|)`$ and the method of substrate deformation measures only an average stress across the whole film . Third, compression of convex expansion of concave parts by surface tension should not be measurable in the total stress changes, since positive and negative curvature compensate each other. Besides it is only a minor effect (compare Eq. (13) to (11) and (12)). Forth, recently a local continuum equation with linear terms $`_th=(\nu k^2\kappa k^4)h`$ has been fitted to experimental results . It would be interesting to see whether a destabilizing term $`|𝐤|^3`$ as in Eq. (8) can give a better description.
In conclusion, in this letter a theoretical framework for a continuum thery of surface growth with diffusion has been constructed. As in the fundamental lattice models the basic processes are particle deposition and diffusion until an energetically favourable site is reached. Surface free energy stabilizes the interface by a configuration dependent attachment rate. For amorphous solids a mean field type of approach yields Eq. (10), a configuration dependence through the energy density near the surface, resulting in a nonequilibrium evolution equation (16). Second, an experimentally observed growth instability on amorphous films has been shown to be stress induced. Its spatial properties can be explained by standard energy arguments of the Asaro-Tiller-Grinfeld-Srolovitz instability, its temporal evolution needs the above nonequilibrium framework. Far from equilibrium the instability is enhanced. A challenge is the extension to crystal growth, where steps act as sticky sites and cannot be treated in a simple mean field way, and where island nucleation becomes important. Random nucleation leads to different step configurations on maxima compared to minima, so it contributes to $`I_1(𝐤)`$ which may be a way to understand the relation $`\kappa \mathrm{}_D^4F`$ obtained from dimensional analysis , where $`\mathrm{}_D`$ denotes the average distance between island nuclei on crystal terraces, conceptually different from $`\mathrm{}_d`$ here. Also the elastic interactions should be worked out for crystals and for heteroepitaxy with different rigidities of substrate and growing film, which would enable very important applications .
It is a pleasure to thank Joachim Krug for encouraging discussions and helpful comments. This work has been supported by the Academy of Finland.
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# Contents
## 1 Introduction
Finite type invariants have received a lot of attention over the past decade. One reason for this is that they provide a common framework for many of the most powerful knot invariants, such as the Conway, Jones, HOMFLYPT and Kauffman invariants. The framework also allows us to study these invariants using elementary combinatorics, by looking at associated functionals (called weight systems) on spaces of chord diagrams. This provides a new ways of describing the invariants.
The modest goal of this paper is to define a few weight systems in terms of the adjacency matrix of the intersection graph of the chord diagrams, and to show that among these weight systems are those associated with the Conway, HOMFLYPT and Kauffman polynomials, in both their framed and unframed incarnations. This gives us new formulas for the weight systems associated to these important knot invariants. We build on ideas of Bar-Natan and Garoufalides , who first found the formula we give for the Conway polynomial.
In section 2 we will review the necessary background for the paper, including finite type invariants, the 2-term relations introduced by Bar-Natan and Garoufalides, intersection graphs of chord diagrams and Lando’s graph bialgebra. In section 3 we will study the adjacency matrix of the intersection graph, and show that the weight systems associated with the Conway and HOMFLYPT polynomials can be defined in terms of the determinant and rank of this matrix. In section 4 we look at marked chord diagrams and define an extended set of 2-term relations on these diagrams. We give an explicit set of generators for the space of marked chord diagrams modulo these relations. Finally, we show that the weight system associated with the Kauffman polynomial can be defined in terms of the rank of the adjacency matrix of marked chord diagrams.
Remark: The result for the Conway polynomial (Theorem 4) has been previously proved by Bar-Natan and Garoufalidis , but is included here for completeness and to place it in the context of Lando’s bialgebra. After distributing the first version of this paper , the author discovered that the adjacency matrix of an intersection graph has also been studied by Soboleva , who has also proven Theorem 5, and a weaker version of Theorem 11. The intersection graphs we study are also related to the trip matrix of a knot, studied by Zulli .
## 2 Preliminaries
### 2.1 Finite Type Invariants
In 1990, V.A. Vassiliev introduced the idea of Vassiliev or finite type knot invariants, by looking at certain groups associated with the cohomology of the space of knots. Shortly thereafter, Birman and Lin gave a combinatorial description of finite type invariants. We will give a brief overview of this combinatorial theory. For more details, see Bar-Natan .
A knot is an embedding of the circle $`S^1`$ into the 3-sphere $`S^3`$. A knot invariant is a map from these embeddings to some set which is invariant under isotopy of the embedding. We will also consider invariants of regular isotopy, where the isotopy preserves the framing of the knot (i.e. a chosen section of the normal bundle of the knot in $`S^3`$). We first note that we can extend any knot invariant to an invariant of singular knots, where a singular knot is an immersion of $`S^1`$ in 3-space which is an embedding except for a finite number of isolated double points. Given a knot invariant $`v`$, we extend it via the relation:
An invariant $`v`$ of singular knots is then said to be of finite type, specifically of type n, if $`v`$ is zero on any knot with more than $`n`$ double points (where $`n`$ is a finite nonnegative integer). The smallest such $`n`$ is called the order of $`v`$. We denote by $`V_n`$ the vector space over $``$ generated by (framing-independent) finite type invariants of type $`n`$ (i.e., whose order is $``$ $`n`$). We can completely understand the space of finite type invariants by understanding all of the vector spaces $`V_n/V_{n1}`$. An element of this vector space is completely determined by its behavior on knots with exactly $`n`$ singular points. In addition, since such an element is zero on knots with more than $`n`$ singular points, any other (non-singular) crossing of the knot can be changed without affecting the value of the invariant. This means that elements of $`V_n/V_{n1}`$ can be viewed as functionals on the space of chord diagrams:
###### Definition 1
A chord diagram of degree n is an oriented circle, together with $`n`$ chords of the circles, such that all of the $`2n`$ endpoints of the chords are distinct. The circle represents a knot, the endpoints of a chord represent 2 points identified by the immersion of this knot into 3-space. The diagram is determined by the order of the $`2n`$ endpoints.
Functionals on the space of chord diagrams which are derived from finite type knot invariants will satisfy certain relations. This leads us to the definition of a weight system:
###### Definition 2
A weight system of degree n is a linear functional $`W`$ on the space of chord diagrams of degree $`n`$ (with values in an associative commutative ring $`𝐊`$ with unity) which satisfies the 1-term and 4-term relations, shown in Figure 1.
It can be shown (see ) that the space of weight systems of degree $`n`$ is isomorphic to $`V_n/V_{n1}`$. For convenience, we take the dual approach, and simply study the space of chord diagrams of degree $`n`$ modulo the 1-term and 4-term relations. The 1-term relation is occasionally referred to as the ”framing-independence” relation, because it arises from the framing-independence of the invariants in $`V_n`$ (essentially, from the first Reidemeister move). Since most of the interesting structure of the vector spaces arises from the 4-term relation, it is common to look at the more general setting of invariants of regular isotopy, and consider the vector space $`A_n`$ of chord diagrams of degree $`n`$ modulo the 4-term relation alone. We will call the space $`W_n`$ of linear functionals on $`A_n`$ the space of regular weight systems of degree $`n`$. We will let $`\widehat{A_n}`$ denote the vector space of chord diagrams modulo both the 1-term and 4-term relations, and $`\widehat{W_n}`$ denote the space of functionals on $`\widehat{A_n}`$, the unframed weight systems.
It is useful to combine all of these spaces into a graded module $`A=_{n1}A_n`$ via direct sum. We can give this module a bialgebra (or Hopf algebra) structure by defining an appropriate product and co-product:
* We define the product $`D_1D_2`$ of two chord diagrams $`D_1`$ and $`D_2`$ as their connect sum. This is well-defined modulo the 4-term relation (see ).
* We define the co-product $`\mathrm{\Delta }(D)`$ of a chord diagram $`D`$ as follows:
$$\mathrm{\Delta }(D)=\underset{J}{}D_J^{}D_J^{\prime \prime }$$
where $`J`$ is a subset of the set of chords of $`D`$, $`D_J^{}`$ is $`D`$ with all the chords in $`J`$ removed, and $`D_J^{\prime \prime }`$ is $`D`$ with all the chords not in $`J`$ removed.
It is easy to check the compatibility condition $`\mathrm{\Delta }(D_1D_2)=\mathrm{\Delta }(D_1)\mathrm{\Delta }(D_2)`$.
There is a natural deframing map $`\varphi :AAA`$, defined by:
$$\varphi (D_1D_2)=(\mathrm{\Theta })^{deg(D_1)}D_2$$
Here $`\mathrm{\Theta }`$ represents the chord diagram consisting of a single chord. This map gives a canonical projection $`\widehat{}:W_n\widehat{W_n}`$, defined by $`\widehat{W}(D)=W(\varphi (\mathrm{\Delta }(D)))`$ (see , Exercise 3.16).
### 2.2 2-Term relations
Of course, any particular weight system will satisfy relations in addition to the 1-term and 4-term relations, and it can be useful to look at weight systems which lie in the subspaces determined by these additional relations. In particular, Bar-Natan and Garoufalides noted that the weight system associated with the Conway polynomial satisfies the 2-term relations in Figure 2.
Clearly, these relations imply that the weight system satisfies the 4-term relation as well. As a result, the product and coproduct of section 2.1 are still well-defined. So we can give the vector space of chord diagrams modulo the 2-term relations the structure of a bialgebra. We will denote this bialgebra (and the underlying vector space) by $`B`$. There is a natural projection from $`A`$ to $`B`$.
Bar-Natan and Garoufalides also showed that $`B`$ is generated (as a vector space) by ($`m_1,m_2`$)-caravans of $`m_1`$ ”one-humped camels” (isolated chords which intersect no other chords) and $`m_2`$ ”two-humped camels” (pairs of chords which intersect each other, but no other chords). An example of such a caravan is shown in Figure 3.
### 2.3 Intersection Graphs
###### Definition 3
Given a chord diagram $`D`$, we define its intersection graph $`\mathrm{\Gamma }(D)`$ as the graph such that:
* $`\mathrm{\Gamma }(D)`$ has a vertex for each chord of $`D`$.
* Two vertices of $`\mathrm{\Gamma }(D)`$ are connected by an edge if and only if the corresponding chords in $`D`$ intersect, i.e. their endpoints on the bounding circle alternate.
For example:
Note that these graphs are simple (i.e. every edge has two distinct endpoints, and there is at most one edge connecting any two vertices). These graphs are also known as circle graphs, and have been studied extensively by graph theorists. A combinatorial classification of circle graphs has been given by Bouchet .
A circle graph can be the intersection graph for more than one chord diagram. For example, there are three different chord diagrams with the following intersection graph:
However, these chord diagrams are all equivalent modulo the 4-term relation. Chmutov, Duzhin and Lando conjectured that intersection graphs actually determine the chord diagram, up to the 4-term relation. In other words, they proposed:
###### Conjecture 1
If $`D_1`$ and $`D_2`$ are two chord diagrams with the same intersection graph, i.e. $`\mathrm{\Gamma }(D_1)`$ = $`\mathrm{\Gamma }(D_2)`$, then for any weight system $`W`$, $`W(D_1)=W(D_2)`$.
This Intersection Graph Conjecture is now known to be false in general. Morton and Cromwell found a finite type invariant of type 11 which can distinguish some mutant knots, and Le and Chmutov and Duzhin have shown that mutant knots cannot be distinguished by intersection graphs. However, the conjecture is true in many special cases, and the exact extent to which it fails is still unknown, and potentially very interesting.
The conjecture is known to hold in the following cases :
* For chord diagrams with 8 or fewer chords (checked via computer calculations);
* For the weight systems coming from the defining representations of Lie algebras gl(N) or so(N) as constructed by Bar-Natan in ;
* When $`\mathrm{\Gamma }(D_1)=\mathrm{\Gamma }(D_2)`$ is a tree (or, more generally, a linear combination of forests);
* When $`\mathrm{\Gamma }(D_1)=\mathrm{\Gamma }(D_2)`$ has a single loop (see ).
The second item above includes the weight systems arising from the Conway, HOMFLYPT and Kauffman polynomials. A main goal of this paper is to find explicit formulas for these weight systems in terms of intersection graphs.
### 2.4 Lando’s graph bialgebra
Lando has given more structure to the questions surrounding intersection graphs by extending the map $`\mathrm{\Gamma }`$ to a homomorphism between the bialgebra $`A`$ of chord diagrams and a particular bialgebra of graphs. Lando’s bialgebra of graphs is constructed by defining an analogue of the 4-term relation for graphs, as follows:
###### Definition 4
Consider the graded vector space (over $``$) of formal linear combinations of graphs, graded by the number of vertices in the graphs. For any graph G and vertices A and B in V(G) we impose on the vector space the relation:
$$GG_{AB}^{}\stackrel{~}{G}_{AB}+\stackrel{~}{G}_{AB}^{}=0$$
where $`G_{AB}^{}`$ is the result of complementing the edge AB in G (i.e. adding or removing it), $`\stackrel{~}{G}_{AB}`$ is the result of complementing the edge AC for every vertex C in V(G) which is adjacent to B and $`\stackrel{~}{G}_{AB}^{}`$ is the result of complementing the edge AB in $`\stackrel{~}{G}_{AB}`$. Here is an example of such a relation:
The bialgebra $`F`$ is defined as this graded vector space, together with a product and a coproduct. The product is simply disjoint union of graphs. The coproduct is a map $`\mu :FFF`$, defined as follows. For any graph G, and subset $`JV(G)`$ of its vertices, let $`G_J`$ denote the subgraph induced by $`J`$. Then:
$$\mu (G)=\underset{JV(G)}{}G_JG_{V(G)\backslash J}$$
An example is shown below:
It is now easy to show that $`\mathrm{\Gamma }`$ extends to a bialgebra homomorphism from $`A`$ to $`F`$ (see ).
We can easily extend Lando’s results to include the 1-term relation and framing-independent invariants. We define the algebra $`\widehat{F}`$ to be simply $`F`$ modulo graphs with isolated vertices (these correspond to the isolated chords of the 1-term relation for chord diagrams). It is then trivial to show that $`\mathrm{\Gamma }`$ extends to a bialgebra homomorphism from $`\widehat{A}`$ to $`\widehat{F}`$.
A regular graph weight system is a linear functional $`\gamma :F`$ (Lando called these functionals 4-invariants) . Then, given any regular graph weight system $`\gamma `$, $`\gamma \mathrm{\Gamma }:A`$ is a regular weight system. Similarly, if we define a graph weight system to be a linear functional of $`\widehat{F}`$, then for any graph weight system $`\alpha `$, $`\alpha \mathrm{\Gamma }`$ will be a weight system.
Just as for chord diagrams, there is a natural deframing map $`\varphi :FFF`$, defined by:
$$\varphi (G_1G_2)=()^{deg(G_1)}G_2$$
Here $``$ represents the trivial graph consisting of a single vertex and no edges. This map gives a canonical projection $`\widehat{}:F^{}\widehat{F}^{}`$, defined by $`\widehat{\gamma }(G)=\gamma (\varphi (\mu (G)))`$.
## 3 The Adjacency Matrix of an Intersection Graph
In this section we will show that the determinant and rank of the adjacency matrix of a graph (over $`𝐙_2`$) are regular graph weight systems, and that the determinant is, in addition, a graph weight system. We will do this by showing that the isomorphism class of the adjacency matrix (as a symmetric bilinear form over $`𝐙_2`$) satisfies 2-term relations analagous to those in section 2.2. We will then show that these weight systems are essentially the same as those associated with the Conway and HOMFLYPT polynomials.
### 3.1 Graph Weight Systems from the Adjacency Matrix
We begin by recalling the definition of the adjacency matrix of a graph.
###### Definition 5
Given a graph G with n vertices, labeled $`\{v_1,\mathrm{},v_n\}`$, the adjacency matrix of G, or adj(G), is the symmetric $`n\times n`$ matrix defined by:
$$adj(G)_{ij}=\{\begin{array}{c}1ifv_iandv_jareconnectedbyanedgeinG\\ 0otherwise\end{array}$$
In the case of a simple graph, the diagonal entries of the matrix will all be 0.
This matrix can be viewed as a symmetric bilinear form over $`𝐙_2`$. If we permute the labels on the vertices of $`G`$, we change the matrix $`adj(G)`$ by the corresponding permutations of the rows and columns. But this does not change the isomorphism class of the form (see ). So, as an isomorphism class of symmetric bilinear forms, the adjacency matrix of an unlabeled graph is well-defined. From Milnor and Husemoller , we know that the determinant and rank of the matrix are invariants of the isomorphism class of the form, and hence are well-defined invariants of the graph. This leads us to define the following functions on graphs:
###### Definition 6
Given a graph G, we define the determinant of G and the rank of G as follows:
$$det(G)=det(adj(G))𝐙_2$$
$$rank(G)=rank(adj(G))$$
We extend these functions linearly to get Z-valued functionals on the space of graphs. We will also call these extensions the determinant and rank. We will see that the determinant gives a Z-valued graph weight system, and the rank gives a Z-valued regular graph weight system (the rank does not satisfy the 1-term relation). We first show that both functionals are regular graph weight systems. To do this, we will show that they satisfy 2-term relations, analogous to those in section 2.2, defined as follows. Consider graphs $`G,G_{AB}^{},\stackrel{~}{G}_{AB},\stackrel{~}{G}_{AB}^{}`$ as in section 2.4. Then the 2-term relations are:
$$G\stackrel{~}{G}_{AB}=0$$
$$G_{AB}^{}\stackrel{~}{G}_{AB}^{}=0$$
It is clear that any functional which satisfies these 2-term relations will also satisfy the 4-term relation. So the vector space $`E`$ of graphs modulo the 2-term relations can be given the structure of a bialgebra, using the same product and coproduct as for $`F`$. There is a natural projection from $`F`$ to $`E`$. Moreover, the pullback by $`\mathrm{\Gamma }`$ of any functional on $`E`$ will be a functional on $`B`$ (defined in section 2.2).
###### Theorem 1
The isomorphism class of the adjacency matrix of a graph satisfy the 2-term relations above.
Proof: Consider two vertices $`A`$ and $`B`$, giving rise to the four graphs $`G,G_{AB}^{},\stackrel{~}{G}_{AB},\stackrel{~}{G}_{AB}^{}`$. We want to show that $`adj(G)adj(\stackrel{~}{G}_{AB})`$ and $`adj(G_{AB}^{})adj(\stackrel{~}{G}_{AB}^{})`$. The easiest way to do this is simply to write down the matrices explicitly. The vertices of $`G`$ other than $`A`$ and $`B`$ can be partitioned into four sets $`S_{AB},S_A,S_B,\mathrm{and}S_0`$, where $`S_{AB}`$ contains those vertices adjacent to both $`A`$ and $`B`$ in $`G`$, $`S_A`$ contains those vertices adjacent to $`A`$ but not $`B`$ in $`G`$, $`S_B`$ contains those vertices adjacent to $`B`$ but not $`A`$ in $`G`$, and $`S_0`$ contains those vertices adjacent to neither $`A`$ nor $`B`$ in $`G`$.
The adjacency matrices for the four graphs, with respect to the basis $`\{A,B,S_{AB},S_A,S_B,S_0\}`$, are shown below. We assume that $`A`$ and $`B`$ are connected by an edge in $`G`$ (if not, simply interchange $`G`$ and $`G^{}`$). Here $`I`$ and $`O`$ represent a row or column of 1’s and 0’s respectively:
$$adj(G)=\left[\begin{array}{cccccc}0& 1& I& I& O& O\\ 1& 0& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}0& 1& O& I& I& O\\ 1& 0& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj(\stackrel{~}{G}_{AB})$$
$$adj(G_{AB}^{})=\left[\begin{array}{cccccc}0& 0& I& I& O& O\\ 0& 0& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}0& 0& O& I& I& O\\ 0& 0& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj(\stackrel{~}{G}_{AB}^{})$$
The isomorphisms are just the result of adding the second row (and column) of the matrix on the left to its first row (and column), modulo 2. So the isomorphism classes of the adjacency matrices of the graphs satisfy the 2-term relations. $`\mathrm{}`$
###### Corollary 1
The (linear extensions of) the rank and determinant of a graph are regular graph weight systems.
###### Theorem 2
The (linear extension of) the determinant of a graph is a graph weight system.
Proof: We need to show that the determinant of a graph satisfies the 1-term relation - i.e. that it is trivial on graphs with isolated vertices. Let $`G`$ be a graph with an isolated vertex, $`v`$, and let $`G=G\{v\}`$. Then the adjacency matrix for $`G`$ can be represented:
$$adj(G)=\left[\begin{array}{cc}0& 0\\ 0& adj(G)\end{array}\right]$$
Since there is a row (and column) of 0’s, $`det(G)=0`$, so the determinant satisfies the 1-term relation. $`\mathrm{}`$
As we mentioned earlier, the rank of a graph does not satisfy the 1-term relation, so it is not a graph weight system. However, we can use the canonical projection from section 2.4 to construct a graph weight system from the rank. In fact, we will construct a polynomial graph weight system by beginning with the invariant $`R(G)(x)=x^{rank(G)}`$, whose linear extension is clearly also a regular graph weight system. To apply our projection, it suffices to note that $`rank(G_1G_2)=rank(G_1)+rank(G_2)`$, and so $`rank(^{deg(G_1)}G_2)=deg(G_1)rank()+rank(G_2)=rank(G_2)`$.
###### Theorem 3
Given a graph G, we define a polynomial $`\widehat{R}(G)(x)`$ as follows. Here J is a subset of the vertices of G, $`|J|`$ is the size of J, n is the total number of vertices in G, and $`G_J`$ is the subgraph of G induced by J:
$$\widehat{R}(G)(x)=\underset{J}{}(1)^{n|J|}x^{rank(G_J)}$$
This polynomial is the canonical projection of R(G), and so its linear extension to a Z\[x\]-valued functional on the space of graphs is a graph weight system.
### 3.2 The Conway and HOMFLYPT weight systems
The Conway polynomial $`\mathrm{\Delta }`$ of a link is a power series $`\mathrm{\Delta }(L)=_{n0}a_n(L)z^n`$. It can be computed via the skein relation (where $`L_+,L_{},L_0`$ are as in Figure 4):
$$\mathrm{\Delta }(L_+)\mathrm{\Delta }(L_{})=z\mathrm{\Delta }(L_0)$$
$$\mathrm{\Delta }(unlinkofkcomponents)=\{\begin{array}{c}1ifk=1\\ 0ifk>1\end{array}$$
The coefficient $`a_n`$ is a finite type invariant of type $`n`$ (see , ), and therefore defines a weight system $`b_n`$ of degree $`n`$. The collection of all these weight systems is called the Conway weight system, denoted $`C`$. Consider a chord diagram $`D`$, together with a chord $`v`$. Let $`D_v`$ be the result of surgery on v, i.e. replacing $`v`$ by an untwisted band, and then removing the interior of the band and the intervals where it is attached to $`D`$, as shown in Figure 5
(so $`D_v`$ may have multiple boundary circles). The skein relations for the Conway polynomial give rise to the following relations for $`C`$:
$$C(D)=C(D_v)$$
$$C(unlinkofkcomponents)=\{\begin{array}{c}1ifk=1\\ 0ifk>1\end{array}$$
It is easy to show (see ) that this weight system satisfies the 2-term relations of section 2.2. Simply surger the two chords; the 2-term relation then says just that one band can be ”slid” over the other, which doesn’t change the topology of the diagram. We will show that this weight system is the same as the the determinant of the intersection graph of the diagram. Our proof is essentially the same as that in ; we include it for completeness.
###### Theorem 4
For any chord diagram D, $`C(D)=det(\mathrm{\Gamma }(D))`$.
Proof: Since both of these weight systems satisfy the 2-term relations, it suffices to show that they agree on caravans. Consider a caravan $`D`$ with $`m_1`$ one-humped camels and $`m_2`$ two-humped camels, as shown in Figure 3. Then $`adj(\mathrm{\Gamma }(D))[0]^{m_1}\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]^{m_2}`$. So $`det(\mathrm{\Gamma }(D))=0^{m_1}1^{m_2}=\{\begin{array}{c}1ifm_1=0\\ 0otherwise\end{array}`$.
On the other hand, if we surger all the chords of the diagram, we obtain an unlink with $`m_1+1`$ components, which means that $`C(D)=\{\begin{array}{c}1ifm_1=0\\ 0otherwise\end{array}`$. So the two weight systems agree. $`\mathrm{}`$
We now turn to the HOMFLYPT polynomial. We will begin by considering a framed version of the HOMFLYPT polynomial - i.e. an invariant of regular isotopy, rather than isotopy. This invariant is the Laurent polynomial $`P(l,m)𝐙[l^{\pm 1},m^{\pm 1}]`$ defined by the following skein relations (see ) ($`L^+`$ is the result of adding a positive kink to the link $`L`$):
$$P(L_+)P(L_{})=mP(L_0)$$
$$P(L^+)=lP(L)$$
$$P(LO)=\frac{ll^1}{m}P(L)$$
$$P(O)=1$$
If we make the substitutions $`m=e^{ax/2}e^{ax/2}`$ and $`l=e^{bx/2}`$, and expand the resulting power series, we transform the HOMFLYPT polynomial into a power series in $`x`$, whose coefficients are finite type invariants (of regular isotopy). These invariants give rise to regular weight systems which we can collect together as the HOMFLYPT regular weight system $`H`$. The skein relations above give rise to the following relations for $`H`$, by looking at the first terms of the power series (as before, $`D_v`$ is the result of surgering the chord $`v`$ in $`D`$):
$$H(D)=aH(D_v)$$
$$H(DO)=bH(D)$$
$$H(O)=1$$
So if $`D`$ is an unlink of $`k`$ components, $`H(D)=b^{k1}`$. Since the first of these relations is almost the same as for the Conway weight system $`C`$, the same argument shows that $`H`$ satisfies the 2-term relations. We will use this to show that the HOMFLYPT regular weight system is equivalent to the rank of the intersection graph of the diagram. (This result was found independently by Soboleva, for the case $`a=1`$ .)
###### Theorem 5
For any chord diagram D of degree k, $`H(D)=a^kb^{krank(\mathrm{\Gamma }(D))}=(ab)^kR(\mathrm{\Gamma }(D))(b^1)`$
Proof: As with Theorem 4, it suffices to show that the weight systems agree on caravans. Let $`D`$ be the caravan with $`m_1`$ one-humped camels and $`m_2`$ two-humped camels, as in Figure 3 (so the degree of $`D`$ is $`m_1+2m_2`$). As before, $`adj(\mathrm{\Gamma }(D))[0]^{m_1}\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]^{m_2}`$, so the rank is $`2m_2`$. On the other hand, if we surger all the chords (each time multiplying $`H`$ by $`a`$), the resulting link has $`m_1+1`$ components, so $`H(D)=a^kb^{m_1}=a^kb^{krank(\mathrm{\Gamma }(D))}`$. $`\mathrm{}`$
###### Corollary 2
If D is a chord diagram of degree k, and $`L_D`$ is the link with c components obtained by surgering all of the chords of D, then $`rank(\mathrm{\Gamma }(D))=kc+1`$.
We can also consider the unframed HOMFLYPT polynomial $`\widehat{P}(l,m)`$, defined by $`\widehat{P}(L)=l^{writhe(L)}P(L)`$ (see ). This invariant is determined by the skein relations:
$$l\widehat{P}(L_+)l^1\widehat{P}(L_{})=m\widehat{P}(L_0)$$
$$\widehat{P}(LO)=\frac{ll^1}{m}\widehat{P}(L)$$
$$\widehat{P}(O)=1$$
After making the same substitutions as before, we again obtain a power series whose coefficients are finite type invariants (this time of isotopy). The collection of the associated weight systems $`\widehat{H}`$ was described by Meng (here $`D_v`$ is the result of surgery on the chord $`v`$, and $`D\backslash v`$ is the result of removing the chord $`v`$):
$$\widehat{H}(D)=a\widehat{H}(D_v)b\widehat{H}(D\backslash v)$$
$$\widehat{H}(DO)=b\widehat{H}(D)$$
$$\widehat{H}(O)=1$$
It is easy to see that this weight system is simply the canonical projection of $`H`$, and so we can conclude that:
###### Theorem 6
For any chord diagram D of degree k, $`\widehat{H}(D)=(ab)^k\widehat{R}(\mathrm{\Gamma }(D))(b^1)`$
Proof: Both weight systems are the canonical projections of $`H`$. $`\mathrm{}`$
Remark: Rather than considering the rank of the adjacency matrix, we could as easily have studied its nullity. If we define $`N(G)(x)=x^{nullity(adj(G))}`$, and let $`\widehat{N}(G)`$ be its canonical projection, then Theorems 5 and 6 imply that $`H(D)=a^kN(\mathrm{\Gamma }(D))(b)`$ and $`\widehat{H}(D)=a^k\widehat{N}(\mathrm{\Gamma }(D))(b)`$.
## 4 Marked Chord Diagrams and the Kauffman weight system
In this section we will look at marked chord diagrams, motivated by the Kauffman polynomial. The idea is that, where we replaced a chord with a band in the previous section, a marked chord will be replaced by a twisted band, as in Figure 6.
The two different surgeries correspond to the two resolutions of a crossing, $`L_0`$ and $`L_{\mathrm{}}`$, in Figure 4.
We will begin by defining marked chord diagrams and graphs, together with a natural map from the space of chord diagrams (graphs) to the space of marked chord diagrams (graphs). We will define an expanded set of 2-term relations on these spaces, and show that a modification of the adjacency matrix is invariant under these relations. We will use this to construct regular graph weight systems, and show that one of these systems is equivalent to the regular weight system associated with the (framed) Kauffman polynomial.
### 4.1 Marked Chord Diagrams and Graphs
A marking of a chord diagram $`D`$ (respectively, graph $`G`$) is simply a partition of the set of chords $`J(D)`$ (resp. vertices $`V(G)`$) into two disjoint subsets $`J_m`$ and $`J_u`$ ($`V_m`$ and $`V_u`$), where $`J_m`$ ($`V_m`$) is the set of marked chords (vertices), and $`J_u`$ ($`V_u`$) is the set of unmarked chords (vertices). We will typically denote a marked chord by labeling it with a pound sign (#).
There is a natural map from the vector space of chord diagrams to the vector space of marked chord diagrams, simply taking a diagram to the sum (with signs) of all possible ways of marking it.
###### Definition 7
Consider a chord diagram D, and a subset J of the set of chords of D. Let $`D^J`$ denote the marked chord diagram obtained by marking all the chords in J. Then we define a map M from the vector space of chord diagrams to the vector space of marked chord diagrams by: $`M(D)=_{JJ(D)}(1)^{|J|}D^J`$.
We can define a similar map (which we will also denote $`M`$) from the vector space of graphs to the space of marked graphs. There is an obvious lifting of the map $`\mathrm{\Gamma }`$ from a chord diagram to its intersection graph to a map from a marked chord diagram to its marked intersection graph, which we will also denote $`\mathrm{\Gamma }`$ (simply mark the vertices corresponding to the marked chords). Clearly, for any chord diagram $`D`$, $`M(\mathrm{\Gamma }(D))=\mathrm{\Gamma }(M(D))`$.
### 4.2 Extended 2-term relations
We can extend the 2-term relations from section 2.2 to a set of 2-term relations on the space of marked chord diagrams. Just as the original 2-term relations were motivated by the idea of replacing chords by bands, the extensions are motivated by the idea of replacing marked chords by twisted bands.
The extended 2-term relations are shown in Figure 7.
The space of marked chord diagrams modulo these relations can be given the structure of a bialgebra by using the same product and coproduct as in section 2.1. We will denote this bialgebra (and the underlying graded vector space) $`B^m`$. The only point which needs to be checked is that the product is well-defined modulo the 2-term relations. This verification is very similar to the corresponding proof for chord diagrams in , and is left as an exercise for the reader.
###### Proposition 1
If $`\beta `$ is a functional on $`B^m`$, then $`\beta M`$ is a regular weight system (a functional on $`A`$).
Proof: It is easy to check that the image of a 4-term relation under $`M`$ is a linear combination of 2-term relations, so $`M`$ is a bialgebra homomorphism from $`A`$ to $`B^m`$. $`\mathrm{}`$
We want to find a set of generators for the vector space of marked chord diagrams modulo the 2-term relations. One such spanning set is a generalization of the caravans of the original 2-term relations.
###### Definition 8
A marked ($`n_1,n_2,n_3`$)-caravan is a marked chord diagram $`(\mathrm{\Theta }_m)^{n_1}\mathrm{\Theta }^{n_2}X^{n_3}`$, where $`\mathrm{\Theta }_m`$ is the chord diagram consisting of a single marked chord (a marked one-humped-camel), $`\mathrm{\Theta }`$ is the diagram consisting of a single unmarked chord (a one-humped-camel), and X is the diagram consisting of two intersecting unmarked chords (a two-humped-camel).
An example of a marked caravan is shown in Figure 8.
We will now show that these caravans span the space of marked chord diagrams, modulo the extended 2-term relations, using an argument similar to that in . In fact, we will show a slightly stronger fact:
###### Theorem 7
Any marked chord diagram is equivalent to a marked caravan, modulo the 2-term relations.
Proof: Let $`D`$ be a marked chord diagram. To begin with, assume that $`D`$ has two intersecting chords $`c_1`$ and $`c_2`$ (possibly marked). There are four possibilities: both chords are unmarked, only $`c_1`$ is marked, only $`c_2`$ is marked, or both chords are marked. In each case, the pair of chords can be slid to the right using the 2-term relations, as in Figure 9,
until a (possibly marked) two-humped-camel is factored out. Continuing inductively, we can factor out (possibly marked) two-humped camels until there are no remaining pairs of intersecting chords. Then, among the remaining chords, there will be a ”smallest” chord, whose endpoints are not separated by the endpoints of any other chord. This chord, whether marked or unmarked, can be slid to the right as in Figure 10,
until a (possibly marked) one-humped-camel is factored out. Continuing inductively, we can reduce the remaining chords to a series of marked and unmarked one-humped-camels. Finally, we can reduce the marked two-humped-camels to pairs of one-humped-camels as in Figure 11.
We are left with a product of marked and unmarked one-humped-camels and unmarked two-humped-camels, which is a marked caravan. This completes the proof. $`\mathrm{}`$
We can define similar 2-term relations for marked graphs. Consider a marked graph $`G`$ with vertices $`A,B`$. Let $`(G)_{AB}`$ be this graph with both $`A`$ and $`B`$ unmarked, $`(G)_{AB}`$ be the graph with $`A`$ marked and $`B`$ unmarked, $`(G)_{AB}`$ be the graph with $`A`$ unmarked and $`B`$ marked, and $`(G)_{AB}`$ be the graph with both $`A`$ and $`B`$ marked. Then the 2-term relations are ($`G_{AB}^{},\stackrel{~}{G}_{AB},\stackrel{~}{G}_{AB}`$ are as defined in section 2.4):
$$(G)_{AB}(\stackrel{~}{G}_{AB})_{AB}=0$$
$$(G_{AB}^{})_{AB}(\stackrel{~}{G}_{AB}^{})_{AB}=0$$
$$(G)_{AB}(\stackrel{~}{G}_{AB})_{AB}=0$$
$$(G_{AB}^{})_{AB}(\stackrel{~}{G}_{AB}^{})_{AB}=0$$
$$(G)_{AB}(\stackrel{~}{G}_{AB}^{})_{AB}=0$$
$$(G_{AB}^{})_{AB}(\stackrel{~}{G}_{AB})_{AB}=0$$
$$(\stackrel{~}{G}_{AB})_{AB}(G_{AB}^{})_{AB}=0$$
$$(\stackrel{~}{G}_{AB}^{})_{AB}(G)_{AB}=0$$
We let $`E^m`$ denote the vector space of marked graphs modulo these relations. Then $`E^m`$ can be given the structure of a bialgebra by using the same product and coproduct as in section 2.4.
###### Proposition 2
If $`\gamma `$ is a functional on $`E^m`$, then $`\gamma M`$ is a regular graph weight system, and $`\gamma \mathrm{\Gamma }`$ is a functional on $`B^m`$.
Proof: To show the first part of the proposition, we simply need to check that the image of a 4-term relation under $`M`$ is a linear combination of 2-term relations, so $`M`$ is a bialgebra homomorphism from $`F`$ to $`E^m`$. The second part of the proposition is immediate. $`\mathrm{}`$
The commutative diagram below summarizes the maps between the various bialgebras we have discussed. All of the maps are bialgebra homomorphisms. It is worth noting that the map $`M`$ is not a homomorphism from $`B`$ to $`B^m`$, because the image of a 2-term relation in $`B`$ may not be a sum of 2-term relations in $`B^m`$. The maps $`p`$ are the natural projections from $`A`$ and $`F`$ to $`B`$ and $`E`$, respectively. The maps $`\stackrel{~}{p}`$ are projections from $`B^m`$ and $`E^m`$ to $`B`$ and $`E`$ (respectively), defined by sending all diagrams (graphs) with marked chords (vertices) to 0.
### 4.3 Marked Adjacency Matrices
Now that we have defined the algebra $`E^m`$, and shown that functionals on this algebra give rise to regular graph weight systems and hence (via the deframing map) regular weight systems, we want to construct explicit examples. Once again, we will use the adjacency matrix of a graph. The adjacency matrix of a marked graph is defined as in section 3.1, except that $`adj(G)_{ii}=1`$ if $`v_i`$ is a marked vertex. (We can visualize a marked vertex as having a small loop attached to it, so it is adjacent to itself.)
As in Section 3.1, this matrix can be viewed as a symmetric bilinear form over $`𝐙_2`$, and is well-defined up to isomorphism of forms. As before, we define the $`rank`$ (resp. $`det`$) of a marked graph as the rank (resp. determinant) of the adjacency matrix of the graph.
###### Theorem 8
The isomorphism class of the adjacency matrix of a marked graph satisfies the extended 2-term relations.
Proof: Consider a graph $`G`$ with vertices $`A`$ and $`B`$. We simply need to verify the eight 2-term relations for the adjacency matrix. We can do this by writing down the matrices explicitly, as we did in Theorem 1. As before, we write our matrices with respect to the basis $`\{A,B,S_{AB},S_A,S_B,S_0\}`$, and we assume that $`A`$ and $`B`$ are connected by an edge in $`G`$. Also as before, $`I`$ and $`O`$ represent a row or column of 1’s and 0’s respectively:
$$adj((G)_{AB})=\left[\begin{array}{cccccc}0& 1& I& I& O& O\\ 1& 0& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}0& 1& O& I& I& O\\ 1& 0& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj((\stackrel{~}{G}_{AB})_{AB})$$
$$adj((G_{AB}^{})_{AB})=\left[\begin{array}{cccccc}0& 0& I& I& O& O\\ 0& 0& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}0& 0& O& I& I& O\\ 0& 0& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj((\stackrel{~}{G}_{AB}^{})_{AB})$$
$$adj((G)_{AB})=\left[\begin{array}{cccccc}1& 1& I& I& O& O\\ 1& 0& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}1& 1& O& I& I& O\\ 1& 0& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj((\stackrel{~}{G}_{AB})_{AB})$$
$$adj((G_{AB}^{})_{AB})=\left[\begin{array}{cccccc}1& 0& I& I& O& O\\ 0& 0& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}1& 0& O& I& I& O\\ 0& 0& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj((\stackrel{~}{G}_{AB}^{})_{AB})$$
$$adj((G)_{AB})=\left[\begin{array}{cccccc}0& 1& I& I& O& O\\ 1& 1& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}1& 0& O& I& I& O\\ 0& 1& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj((\stackrel{~}{G}_{AB}^{})_{AB})$$
$$adj((G_{AB}^{})_{AB})=\left[\begin{array}{cccccc}0& 0& I& I& O& O\\ 0& 1& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}1& 1& O& I& I& O\\ 1& 1& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]=adj((\stackrel{~}{G}_{AB})_{AB})$$
$$adj((\stackrel{~}{G}_{AB}^{})_{AB})=\left[\begin{array}{cccccc}0& 0& O& I& I& O\\ 0& 1& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}1& 0& I& I& O& O\\ 0& 1& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]=adj((G)_{AB})$$
$$adj((\stackrel{~}{G}_{AB})_{AB})=\left[\begin{array}{cccccc}0& 1& O& I& I& O\\ 1& 1& I& O& I& O\\ O& I& & & & \\ I& O& & & & \\ I& I& & & & \\ O& O& & & & \end{array}\right]\left[\begin{array}{cccccc}1& 0& I& I& O& O\\ 0& 1& I& O& I& O\\ I& I& & & & \\ I& O& & & & \\ O& I& & & & \\ O& O& & & & \end{array}\right]=adj((G_{AB}^{})_{AB})$$
The isomorphisms are just the result of adding the second row (and column) of the matrix on the left to its first row (and column), modulo 2. So the isomorphism classes of the adjacency matrices of the graphs satisfy the extended 2-term relations. $`\mathrm{}`$
###### Corollary 3
The rank and determinant of a marked graph are functionals on $`E^m`$.
We can combine these functionals with $`M`$ to obtain regular graph weight systems. In order to construct polynomial-valued weight systems, we will begin with $`s(G)=x^{rank(G)}`$ and $`t(G)=x^{det(G)}`$, whose linear extensions are also functionals on $`E^m`$:
###### Theorem 9
Given an unmarked graph G, and a subset J $``$ V(G), define $`G^J`$ as the result of marking the vertices in J. Define the maps S(G) and T(G) as follows:
$$S(G)(x)=\underset{JV(G)}{}(1)^{|J|}x^{rank(G^J)}$$
$$T(G)(x)=\underset{JV(G)}{}(1)^{|J|}x^{det(G^J)}$$
Then these maps are regular graph weight systems. Moreover, the map S(G) is multiplicative: $`S(G_1G_2)=S(G_1)S(G_2)`$
Proof: $`S(G)=s(M(G))`$ and $`T(G)=t(M(G))`$ (where $`s`$ and $`t`$ are extended linearly), so these are regular graph weight systems by Proposition 2. Since $`rank(G_1G_2)=rank(G_1)+rank(G_2)`$, we can see that $`s(G_1G_2)=s(G_1)s(G_2)`$. It is easy to check that $`M`$ is also multiplicative. Therefore $`S(G)`$ is multiplicative. $`\mathrm{}`$
Moreover, we can obtain graph weight systems by applying the canonical projection from section 2.4.
###### Theorem 10
Given an unmarked graph G with n vertices, a subset J $``$ V(G), and a subset $`J_mJ`$, define $`G_J^{J_m}`$ as the subgraph induced by J, with the vertices in $`J_m`$ marked. Then we define $`\widehat{S}(G)`$ and $`\widehat{T}(G)`$ as follows:
$$\widehat{S}(G)(x)=\underset{JV(G)}{}(x1)^{n|J|}S(G_J)=\underset{JV(G)}{}\underset{J_mJ}{}(1)^{|J_m|}(x1)^{n|J|}x^{rank(G_J^{J_m})}$$
$$\widehat{T}(G)(x)=\underset{JV(G)}{}T(G_J)=\underset{JV(G)}{}\underset{J_mJ}{}(1)^{|J_m|}x^{det(G_J^{J_m})}$$
These maps are the canonical projections of S(G) and T(G), and so are graph weight systems.
Proof: Recall that the deframing map $`\varphi (G_1G_2)=()^{deg(G_1)}G_2`$. It is easy to check that the map $`M:FE^m`$ is multiplicative, i.e. $`M(G_1G_2)=M(G_1)M(G_2)`$. So $`M(\varphi (G_1G_2))=M()^{deg(G_1)}M(G_2)`$. Since $`s(G)`$ is also multiplicative, we have:
$`S(\varphi (G_1G_2))`$ $`=`$ $`s(M(\varphi (G_1G_2)))`$
$`=`$ $`s(M())^{deg(G_1)}s(M(G_2))`$
$`=`$ $`(x1)^{deg(G_1)}s(M(G_2))`$
$`=`$ $`(x1)^{deg(G_1)}S(G_2)`$
From this, it is straightforward to see that the projection of $`S(G)`$ is $`_{JV(G)}(x1)^{n|J|}S(G_J)`$, as desired.
On the other hand, the determinant of a graph with any isolated unmarked chords is 0. So for $`T`$ we have (denoting the graph consisting of a single marked vertex by $`\mathrm{\#}`$):
$`T(\varphi (G_1G_2))`$ $`=`$ $`t(M(\varphi (G_1G_2)))`$
$`=`$ $`t(M()^{deg(G_1)}M(G_2))`$
$`=`$ $`t\left((\mathrm{\#})^{deg(G_1)}{\displaystyle \underset{JV(G_2)}{}}(1)^{|J|}G_2^J\right)`$
$`=`$ $`t\left({\displaystyle \underset{J}{}}{\displaystyle \underset{k=0}{\overset{deg(G_1)}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{deg(G_1)}{k}}\right)(\mathrm{\#})^k()^{deg(G_1)k}(1)^{|J|}G_2^J\right)`$
$`=`$ $`{\displaystyle \underset{J}{}}{\displaystyle \underset{k=0}{\overset{deg(G_1)}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{deg(G_1)}{k}}\right)(1)^{|J|}(1)^{deg(G_1)k}t((\mathrm{\#})^k()^{deg(G_1)k}G_2^J)`$
Since $`det((\mathrm{\#})^k()^{deg(G_1)k}G_x^J)=\{\begin{array}{c}det(G_2^J)ifk=deg(G_1)\\ 0otherwise\end{array}`$, we know that $`t((\mathrm{\#})^k()^{deg(G_1)k}G_x^J)=\{\begin{array}{c}t(G_2^J)ifk=deg(G_1)\\ 1otherwise\end{array}`$. So our equation reduces to:
$`T(\varphi (G_1G_2))`$ $`=`$ $`{\displaystyle \underset{J}{}}(1)^{|J|}\left(t(G_2^J)+{\displaystyle \underset{k=0}{\overset{deg(G_1)1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{deg(G_1)}{k}}\right)(1)^{deg(G_1)k}\right)`$
$`=`$ $`{\displaystyle \underset{J}{}}(1)^{|J|}\left(t(G_2^J)+(11)^{deg(G_1)}1\right)`$
$`=`$ $`{\displaystyle \underset{J}{}}(1)^{|J|}t(G_2^J){\displaystyle \underset{J}{}}(1)^{|J|}`$
$`=`$ $`t(M(G_2))0=T(G_2)`$
From this, we can conclude that the projection of $`T(G)`$ is $`\widehat{T}(G)(x)=_{JV(G)}T(G_J)`$, as desired. $`\mathrm{}`$
### 4.4 The Kauffman weight system
We want to show that $`S(\mathrm{\Gamma }(D))`$ and $`\widehat{S}(\mathrm{\Gamma }(D))`$ are the weight systems associated with the Kauffman polynomial. We will begin by considering a framed version of the Kauffman polynomial $`F(y,z)`$, defined by the following skein relations ($`L_+`$, $`L_{}`$, $`L_0`$ and $`L_{\mathrm{}}`$ are as shown in Figure 4, and $`L^+`$ is the result of adding a postive kink to $`L`$):
$$F(L_+)F(L_{})=z(F(L_0)F(L_{\mathrm{}})$$
$$F(L^+)=yF(L)$$
$$F(LO)=\left(\frac{yy^1}{z}+1\right)F(L)$$
$$F(O)=1$$
To derive finite type invariants, we make the substitutions $`z=e^{ax/2}e^{ax/2}`$ and $`y=e^{(b1)x/2}`$. If we then expand the polynomial as a power series in $`x`$, the coefficients will be finite type invariants. The regular weight system associated with this collection of invariants is defined by the skein relations below. Here $`D`$ is an unmarked chord diagram, $`v`$ is a chord in $`D`$, $`D_v`$ is the result of replacing $`v`$ by an untwisted band, and $`D^v`$ is the result of replacing $`v`$ by a band with a half-twist:
$$K(D)=a(K(D_v)K(D^v))$$
$$K(DO)=bK(D)$$
$$K(O)=1$$
Note that, if $`D`$ is an unlink of $`k`$ components, then $`K(D)=b^{k1}`$.
Our first task is to show that this regular weight system factors through the algebra $`B^m`$. We define a map $`K^m:B^m𝐙[a,b]`$ recursively by the following relations, where $`D`$ is a marked chord diagram, and $`v`$ is a chord in $`D`$:
$$K^m(D)=\{\begin{array}{c}aK^m(D_v)ifvisunmarked\\ aK^m(D^v)ifvismarked\end{array}$$
$$K^m(DO)=bK^m(D)$$
$$K^m(O)=1$$
Note that, if $`D`$ is a diagram with no chords and $`k`$ components, then $`K^m(D)=b^{k1}`$.
###### Proposition 3
$`K^m`$ satisfies the extended 2-term relations.
Proof: In each of the 2-term relations of Figure 7, replace each unmarked chord by an untwisted band and each marked chord by a band with a half-twist. It is clear that the relations are simply the result of sliding one band over another, and don’t change the topology of the diagram. We need only keep in mind that when a band is slid over a half-twisted band (marked chord), it receives a half-twist itself. We view a band with a full twist as equivalent to an untwisted band, since it does not change the number of components of the diagram, which is all that matters in the base case of the definition of $`K^m`$. $`\mathrm{}`$
###### Proposition 4
$`K=K^mM`$, so K is the pullback of $`K^m`$ by M.
Proof: Consider a diagram $`D`$ in $`A`$. We will prove the proposition via induction on the number of chords of $`D`$. If $`D`$ has no chords, then $`M(D)=D`$. Since $`K`$ and $`K^m`$ differ only in their first skein relation (which only applies if there are chords), we conclude that $`K^m(M(D))=K^m(D)=K(D)`$.
For our inductive step, assume $`D`$ has a chord $`v`$. Note that $`D_v`$ and $`D^v`$ each have fewer chords than $`D`$, so $`K^m(M(D_v))=K(D_v)`$ and $`K^m(M(D^v))=K(D^v)`$. If $`J`$ is a subset of the chords of $`D`$, we let $`D^J`$ denote the marked chord diagram which results by marking all the chords in $`J`$. (However, for the single chord $`v`$, we will still let $`D^v`$ denote the result of replacing $`v`$ with a half-twisted band.) $`M(D)`$ is then given by:
$$M(D)=\underset{J}{}(1)^{|J|}D^J=\underset{Js.t.vJ}{}(1)^{|J|}(D^JD^{Jv})$$
Then:
$`K^m(M(D))`$ $`=`$ $`{\displaystyle \underset{Js.t.vJ}{}}(1)^{|J|}(K^m(D^J)K^m(D^{Jv}))`$
$`=`$ $`{\displaystyle \underset{Js.t.vJ}{}}(1)^{|J|}(aK^m((D_v)^J)aK^m((D^v)^J))`$
$`=`$ $`a(K^m(M(D_v))K^m(M(D^v)))`$
$`=`$ $`a(K(D_v)K(D^v))`$
$`=`$ $`K(D)`$
So by induction, we conclude that for any diagram $`D`$, $`K(D)=K^m(M(D))`$. $`\mathrm{}`$
###### Theorem 11
For any $`DA`$ of degree $`k`$, $`K(D)=(ab)^kS(\mathrm{\Gamma }(D))(b^1)`$.
Proof: Since $`S(\mathrm{\Gamma }(D))=s(M(\mathrm{\Gamma }(D))=s(\mathrm{\Gamma }(M(D))`$, and $`K(D)=K^m(M(D))`$, it suffices to show that $`(ab)^k(s\mathrm{\Gamma }(D))(b^1)=K^m(D)`$ for any $`DB^m`$. Since both of these maps satisfy the extended 2-term relations, it suffices to show that they agree on marked caravans, by Theorem 7.
Consider a marked ($`n_1,n_2,n_3`$)-caravan $`D`$, as shown in Figure 8. The degree of this caravan is $`k=n_1+n_2+2n_3`$. Then $`adj(\mathrm{\Gamma }(D))[1]^{n_1}[0]^{n_2}\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]^{n_3}`$. So $`rank(\mathrm{\Gamma }(D))=n_1+2n_3`$, and $`(ab)^ks(\mathrm{\Gamma }(D))(b^1)=(ab)^kb^{n_12n_3}=a^kb^{kn_12n_3}=a^kb^{n_2}`$.
On the other hand, $`K^m(D)`$ is computed by replacing all the unmarked chords with untwisted bands and all the marked chords with twisted bands (multiplying by $`a`$ each time), and then looking at the number of components of the resulting link. This link will have $`n_2+1`$ components, so $`K^m(D)=a^kb^{n_2}=(ab)^ks(\mathrm{\Gamma }(D))(b^1)`$, which completes the proof. $`\mathrm{}`$
We can also consider the unframed Kauffman polynomial $`\widehat{F}(y,z)`$, defined by $`\widehat{F}(L)=y^{writhe(L)}F(L)`$ (see ). This invariant is also determined by the skein relations:
$$y\widehat{F}(L_+)y^1\widehat{F}(L_{})=m(\widehat{F}(L_0)\widehat{F}(L_{\mathrm{}}))$$
$$\widehat{F}(LO)=\left(\frac{yy^1}{z}+1\right)\widehat{F}(L)$$
$$\widehat{F}(O)=1$$
After making the same substitutions as before, we again obtain a power series whose coefficients are finite type invariants (this time of isotopy). The collection of the associated weight systems $`\widehat{K}`$ was described by Meng (here $`D_v`$ is the result of replacing the chord $`v`$ by an untwisted band, $`D^v`$ is the result of replacing the chord $`v`$ by a half-twisted band, and $`D\backslash v`$ is the result of removing the chord $`v`$):
$$\widehat{K}(D)=a\widehat{K}(D_v)a\widehat{K}(D^v)b\widehat{K}(D\backslash v)$$
$$\widehat{K}(DO)=b\widehat{K}(D)$$
$$\widehat{K}(O)=1$$
It is easy to see that this weight system is simply the canonical projection of $`K`$, and so we can conclude that:
###### Theorem 12
For any chord diagram D of degree k, $`\widehat{K}(D)=(ab)^k\widehat{S}(\mathrm{\Gamma }(D))(b^1)`$
Proof: Both weight systems are the canonical projections of $`K`$. $`\mathrm{}`$
Remark: Rather than considering the rank of the marked adjacency matrix, we could as easily have studied its nullity. If we define $`u(G)(x)=x^{nullity(adj(G))}`$ and $`U(G)=u(M(G))`$, and let $`\widehat{U}(G)`$ be the canonical projection of $`U(G)`$, then Theorems 11 and 12 imply that $`K(D)=a^kU(\mathrm{\Gamma }(D))(b)`$ and $`\widehat{K}(D)=a^k\widehat{U}(\mathrm{\Gamma }(D))(b)`$.
We now have explicit formulas for computing the Conway, HOMFLYPT and Kauffman weight systems directly from intersection graphs. Hopefully, these interpretations will help shed some light on the geometric meanings of these polynomials.
## 5 Acknowledgements
I would like to thank Dror Bar-Natan and Louis Zulli for informing me of their previous work. I would also like to thank Sergei Lando for sending me E. Soboleva’s paper, along with his own work on obtaining Vassiliev invariants from intersection graphs. Finally, I would like to thank the anonymous reviewer who made several suggestions which led to substantial revisions of the paper.
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# Can a Large Neutron Excess Help Solve the Baryon Loading Problem in Gamma-Ray Burst Fireballs?
\[
## Abstract
We point out that the baryon-loading problem in Gamma-Ray Burst (GRB) models can be amelioriated if a significant fraction of the baryons which inertially confine the fireball are converted to neutrons. A high neutron fraction in some circumstances can result in a reduced transfer of energy from relativistic light particles in the fireball to baryons. The energy needed to produce the required relativistic flow in the GRB is consequently reduced, in some cases by orders of magnitude. This could be relevant to GRB models because a high neutron-to-proton ratio has been calculated in neutron star-merger fireball environments. Significant neutron excess also could occur near compact objects with high neutrino fluxes.
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In this Letter we show how the baryon loading problem can be alleviated in certain gamma-ray burst (GRB) models when significant numbers of baryons are converted to neutrons. Interestingly, many of the proposed GRB “central engines” involve compact objects which are themselves highly neutronized, or which are accompanied by intense neutrino fluxes. Weak interactions induced by these neutrino fluxes can result in significant proton-to-neutron conversion, especially if resonant neutrino flavor transformation takes place .
Inferences of the energetics and spectral observations of GRBs imply (i) total energies in gamma-rays approaching $`10^{53}`$ ergs for the most energetic events (in the absence of beaming), and (ii) large Lorentz factors of the progenitor fireball ($`\gamma 10^3`$) (for a recent review, see Ref. ). Excessive baryon pollution of the fireball precludes attainment of these features for many GRB models. This is a consequence of the conversion of radiation energy in the electron/positron/photon fireball to kinetic energy in baryons . However, the relatively small cross sections characterizing the interactions of neutrons with the electron/positron/photon plasma may afford a solution to this problem.
This can be seen by considering the fictitious limit of completely noninteracting neutrons. Imagine that protons inertially tether an electron/positron/photon fireball via photon Thomson drag on $`e^\pm `$, which in turn influences protons through Coulomb interactions. If these protons were suddenly converted to non-interacting “neutrons”, then the fireball would expand relativistically, leaving behind the baryonic component. Real neutrons can approximate this limit as they interact with the electron/positron/photon plasma only via the neutron magnetic dipole moment. These cross sections are small compared to the Thomson cross section $`\sigma _\mathrm{T}`$: neutron-electron (positron) scattering has $`\sigma _{\mathrm{n}e}10^7\sigma _\mathrm{T}`$ ; and neutron-photon scattering has $`\sigma _{\mathrm{n}\gamma }10^{12}\sigma _\mathrm{T}`$ .
However, the real limit on the efficacy of this mechanism is the strong interaction neutron-proton scattering which will dominate the energy transfer process when conversion of neutrons to protons is incomplete. Therefore, the degree to which the baryon loading burden can be lifted in our proposed mechanism will depend on the neutron excess in the fireball environment. Here we will measure the neutron content of the plasma in terms of the electron fraction $`Y_e`$, the net number of electrons ($`n_e^{}n_{e^+}`$) per baryon, or in terms of the neutron-to-proton ratio $`Y_e=1/(\mathrm{n}/\mathrm{p}+1)`$.
We note that although previous studies have invoked neutrino oscillations to attempt a baryon loading problem solution , none has exploited the $`Y_e`$-changing aspect of the weak interaction.
To go beyond the simplistic picture of non-interacting neutrons, we can consider a two-component ((i) neutrons, and (ii) protons/$`e^\pm `$/photons) plasma in the context of a homogeneous fireball with initial radius, temperature, Lorentz factor, and electron fraction, $`R_0,T_0,\gamma _0`$ & $`Y_{e0}`$, respectively. Numerical and analytic work have shown the following simple scaling laws for such a configuration :
$`\text{For }R<\eta R_0/\gamma _0`$ $`\{\begin{array}{cc}\gamma =\gamma _0R/R_0\hfill & \\ T=T_0R_0/R\hfill & \end{array},`$ (1)
$`\text{For }R>\eta R_0/\gamma _0`$ $`\gamma =\eta .`$ (3)
As a matter of convenience we will take the scaling to be such that $`\gamma _0=1`$. In Eq. (1), the ratio of energy in radiation $`E`$ to total baryon rest mass $`M`$ is $`\eta E/M`$.
One can relate $`\eta `$ to the entropy per baryon, $`s`$, using the number density of baryons $`N=\rho _\mathrm{b}/m_p=\rho _{\mathrm{rad}}/m_p\eta `$ where $`\rho _b`$ the baryon component rest mass energy density and $`m_p`$ is the proton rest mass. Using this relation, and noting that in terms of the proper entropy density $`S`$, the entropy-per-baryon is $`s=S/N`$, the relation between $`s`$, $`\eta `$, and the temperature $`T_0`$ is $`s1250\eta (1\mathrm{MeV}/T_0)`$.
For a large enough $`\eta `$, baryon loading is unimportant . In fact when $`\eta \begin{array}{c}>\hfill \\ \hfill \end{array}10^5(E/10^{52}\mathrm{erg})^{1/3}(10^7\mathrm{cm}/R_0)^{2/3},`$ the fireball becomes optically thin before transferring its energy to kinetic energy in baryons (here $`E`$ is the total energy of the fireball).
Written in terms of the time $`t`$ as measured in a frame comoving with the fireball the above relations imply
$$R=R_0e^{t/\tau _{\mathrm{dyn}}},\gamma =e^{t/\tau _{\mathrm{dyn}}},T=T_0e^{t/\tau _{\mathrm{dyn}}},$$
(4)
for $`\gamma <\eta `$. Here the dynamic timescale is defined to be the initial light crossing time, $`\tau _{\mathrm{dyn}}R_0/c`$. In fireballs resulting from neutron star mergers, for example, $`\tau _{dyn}2.7\times 10^5s`$, corresponding to an $`R_0`$ of 8 km .
A particle co-moving with the expanding plasma experiences a 4-acceleration $`a^\mu `$, with magnitude $`\sqrt{a^\mu a_\mu }=d\gamma /dR`$. As noted above, the force that drags the neutrons along with the expanding plasma arises princiapally from n-p collisions. The relative contribution to the total force on the neutrons from collisions with electrons and positrons is roughly $`F_{ne}/F_{np}m_en_e\sigma _{ne}/m_pn_p\sigma _{np}10^{10}\left(s/Y_e\right)`$ and is small for the conditions we consider. The neutron-photon cross section is small enough ($`\sigma _{n\gamma }10^{36}[E_\gamma /(1\mathrm{MeV})]^2\mathrm{cm}^2`$ where $`E_\gamma `$ is the photon energy in the neutron rest frame ) that n-$`\gamma `$ interactions are negligible.
The relations in Eqs. (1) imply that an inertial observer with time coordinate $`t^{}`$ initially (at $`t^{}=0`$) comoving with the plasma sees the plasma accelerate according to $`\gamma v=t^{}/\tau _{\mathrm{dyn}}`$. Hereafter we adopt natural units where $`c=1`$. If we denote by $`\tau _{\mathrm{coll}}^1`$ the frequency of neutron/proton collisions (per neutron), we expect that the two components of the plasma will achieve a relative velocity given by $`v_{\mathrm{rel}}2\tau _{\mathrm{coll}}/\tau _{\mathrm{dyn}}`$, where the factor of 2 arises from the approximate angle independence of the neutron-proton scattering cross section and the near equality of the neutron and proton masses. An equivalent expression is found if one considers the force on the neutrons from collisions with protons . It is clear then that when $`\tau _{\mathrm{dyn}}\tau _{\mathrm{coll}}`$ the neutrons are coupled to the rest of the plasma. However, decoupling occurs as these two timescales become comparable. Since the baryon number density in the plasma frame decreases as $`e^{3t/\tau _{\mathrm{dyn}}}`$, decoupling will occur quickly, i.e. on a timescale shorter than $`\tau _{\mathrm{dyn}}`$.
When significant decoupling occurs we can neglect the thermal contribution to the collision frequency and write
$$\tau _{\mathrm{coll}}^1(9\times 10^{12}\mathrm{s}^1)\frac{T_{\mathrm{MeV}}^3}{s_5}v_{\mathrm{rel}}Y_e\sigma _{10},$$
(5)
where $`s_5s/10^5`$ and $`\sigma _{10}`$ is the neutron-proton cross-section in units of 10 fm<sup>2</sup>, and $`T_{\mathrm{MeV}}`$ is the temperature in MeV. As the precise energy dependence of $`\sigma _{10}`$ is not important here, it suffices to note a few representative values: $`\sigma _{10}(v_{\mathrm{rel}}=0.1)17`$, $`\sigma _{10}(v_{\mathrm{rel}}=0.3)2`$ and $`\sigma _{10}(v_{\mathrm{rel}}=0.6)0.4`$ .
The requirement of a non-negligible relative velocity then gives the decoupling time $`t=t_{\mathrm{dec}}`$ as
$$\frac{t_{\mathrm{dec}}}{\tau _{\mathrm{dyn}}}=4.6+(1/3)\mathrm{ln}\left(\frac{\sigma _{10}\tau _6T_{0,\mathrm{MeV}}^3Y_e}{s_5}\right).$$
(6)
In the above, $`\tau _6\tau _{\mathrm{dyn}}/10^6\mathrm{sec}`$, and $`\tau _{\mathrm{coll}}`$ was evaluated at a terminal velocity of $`0.5`$. The calculated decoupling time is logarithmically sensitive to this choice. In reality the neutrons do not sharply decouple but continue to interact with the plasma over roughly a dynamical timescale. In this sense, the $`t_{\mathrm{dec}}`$ appearing in Eq. (6) is an “effective” decoupling time. An accurate determination of $`t_{\mathrm{dec}}`$ requires solving in detail the neutron and proton transport equations. However, because of the exponential decrease of density with time in the plasma frame, the number 4.6 appearing in Eq. (6) is only uncertain to approximately $`\pm (1/3)`$.
Once the neutrons decouple they will have an energy $`\gamma _{\mathrm{dec}}(1Y_e)M`$. The ratio of kinetic energy in neutrons to the total energy in the fireball is then
$$f_n(1Y_e)\frac{e^{t_{\mathrm{dec}}/\tau _{\mathrm{dyn}}}}{\eta }1.3(1Y_e)\left(\frac{Y_e\sigma _{10}\tau _6}{s_5^4}\right)^{1/3}.$$
(7)
(Here “total” energy includes both the thermal $`e^\pm /\gamma `$ energy and the bulk kinetic energy of baryons.) From this we see that for $`Y_e`$ less than
$$Y_{e,\mathrm{crit}}0.46\left(\frac{s_5^4}{\sigma _{10}\tau _6}\right)$$
(8)
at the time of decoupling the baryon loading problem is diminished. (Protons and neutrons each move with $`\gamma _{\mathrm{dec}}`$ at the decoupling point; thereafter, protons will possess larger Lorentz factors than do average neutrons). If the plasma remains optically thick to radiation until the energy in radiation is converted to kinetic energy of the remaining protons, energy conservation gives the final Lorentz factor of the protons
$$\gamma \eta \left(\frac{1f_n}{Y_e}\right).$$
(9)
A simple ansatz for the condition that the fireball remains optically thick after decoupling is $`\eta \begin{array}{c}<\hfill \\ \hfill \end{array}(Y_e)10^5(E/10^{52}\mathrm{erg})^{1/3}(10^7\mathrm{cm}/R_0)^{2/3}`$. This is obtained by applying the result from and making the replacements $`\eta (\eta /Y_e)(1f_n),E(1f_n)E`$ and $`ss/Y_e`$. Note that even for modest values of $`f_n`$, $`\gamma `$ can be increased significantly if $`Y_e`$ is low.
The above results are summarized in Fig. 1, where we have plotted the smallest entropy $`s_5`$ ($`1.25\eta _{100}/T_0`$) for which decoupling occurs as a function of $`Y_e`$. For example, if $`s_5=0.6,T_0=2\mathrm{MeV}`$ (corresponding to $`\eta 100`$), and $`Y_{e0}=0.02`$, then the final Lorentz factor of the plasma after neutron decoupling (Eqs. (7), (9)) would be $`\gamma 1500`$, which is 15 times larger than the standard case of $`\gamma =\eta `$. As another example, consider the Ref. values of $`\tau _6=27`$ and $`Y_e=0.1`$ and suppose that $`T_0=10\mathrm{MeV}`$ and $`s_5=2.5`$ (corresponding to $`\eta =2000`$). In this case we find $`\gamma =1.1\times 10^4`$, an increase by a factor of 5.7. Clearly, the importance of this effect depends on how low $`Y_e`$ can be.
Two conditions must be met in order to achieve a low $`Y_e`$ at the time of decoupling: (i) $`Y_e`$ must be low initially and (ii) $`Y_e`$ must not be unacceptably raised during the evolution of the fireball. We can divide up the discussion of $`Y_e`$ in this way because the initial electron fraction depends in detail on the GRB central engine, whereas the later evolution of the fireball is generically given by the relations in Eq. (4). Many proposed GRB central engines involve neutrino heating or are sited in environments subject to intense neutrino fluxes . General discussions of the relation between neutrino processes and the dynamics of outflow may be found in Refs. . However, the details of neutron decoupling are insensitive to how $`Y_e`$ is set and we are not arguing for a specific GRB site.
The processes which have a significant effect on $`Y_e`$ in the fireball environment are lepton capture/decay involving free nucleons and inelastic $`\mathrm{nn}\mathrm{np}\pi `$ scattering (charged pion-nucleon bremsstrahlung),
$`\nu _e+\mathrm{n}`$ $``$ $`\mathrm{p}+\mathrm{e}^{}`$ (11)
$`\overline{\nu }_e+\mathrm{p}`$ $``$ $`\mathrm{n}+\mathrm{e}^+`$ (12)
$`\mathrm{n}`$ $``$ $`\mathrm{p}+e^{}+\overline{\nu }_e`$ (13)
$`\mathrm{n}+\mathrm{n}`$ $``$ $`\mathrm{n}+\mathrm{p}+\pi ^{}.`$ (14)
In general, $`Y_e`$ is set by the competition between the above processes . For the range of fireball parameters of interest to us, free neutron decay (13) is unimportant as the fraction of neutrons decaying during the evolution of the fireball is $`10^9\tau _6\mathrm{ln}\eta `$. Furthermore, as lepton capture is only important during the early, hot, evolution of the fireball and inelastic nucleon-nucleon scattering only occurs after neutron decoupling, the lepton capture and pion bremsstrahlung processes may be considered separately.
In environments where neutrino heating is important the forward reactions (11) and (12) can dominate in setting the electron fraction . Integration of the rate equations corresponding to the lepton capture processes gives $`\mathrm{n}/\mathrm{p}\lambda _{\overline{\nu }_e\mathrm{p}}/\lambda _{\nu _e\mathrm{n}}(L_{\overline{\nu }_e}E_{\overline{\nu }_e})/(L_{\nu _e}E_{\nu _e})`$, where $`\lambda _{\overline{\nu }_e\mathrm{p}}`$ and $`\lambda _{\nu _e\mathrm{n}}`$ are the rates for the reactions in Eqs. (11) and (12), $`E_{\overline{\nu }_e}`$ and $`E_{\nu _e}`$ are the average energies characterizing the energy spectra of the $`\overline{\nu }_e`$ and $`\nu _e`$ neutrinos, respectively, while $`L_{\overline{\nu }_e}`$ and $`L_{\nu _e}`$ are the corresponding energy luminosities. Absent neutrino oscillations and flavor/type mixings, any thermal neutrino emission scenario from a compact object will yield a characteristic average neutrino energy heirarchy for solar mass scale objects: $`E_{\nu _\mu }E_{\overline{\nu }_\mu }E_{\nu _\tau }E_{\overline{\nu }_\tau }>E_{\nu _{\overline{e}}}>E_{\nu _e}`$. These considerations are consistent with findings in Ref. in which a hard $`\overline{\nu }_e`$ spectrum from a collapsing neutron star leads to an electron fraction in the fireball of $`Y_e0.1`$.
If the $`\nu _e`$ component of the neutrino emission were to disappear or be greatly reduced, then the competition inherent in the above equations would be unbalanced in favor of the reaction $`\overline{\nu }_e+\mathrm{p}\mathrm{n}+e^+`$. This, in turn, would result in the wholesale production of neutrons. In fact, several schemes involving matter-enhanced active-sterile neutrino transformation have been proposed as a way of enabling $`r`$-process nucleosynthesis in neutrino-heated supernova ejecta: one of these involves matter-enhanced $`\nu _e\nu _s`$ and $`\overline{\nu }_e\overline{\nu }_s`$ ; the other involves matter-enhanced conversion $`\nu _{\mu ,\tau }\nu _s`$ followed by an active-active matter-enhanced conversion $`\nu _{\mu ,\tau }\nu _e`$ . In either case, the intial $`\nu _e`$ flux can be reduced by more than an order of magnitude and, in turn, this can translate into a substantial decrease in the initial $`Y_e`$. (Just how low depends on central engine outflow hydrodynamics and on neutrino background effects .)
If we demand that an initially low $`Y_e`$ not be raised above $`Y_{e,\mathrm{crit}}`$, consideration of lepton capture on neutrons allows us to place rough constraints on the fireball and neutrino parameters. We incorporate the uncertainty in the initial fireball evolution by supposing that the relations in Eq. (4) are valid only after the fireball has a Lorentz factor $`\gamma _i`$ and temperature $`T_i`$. Consideration of positron capture after $`\gamma =\gamma _i,T=T_i`$ then leads to $`T_i<(22\mathrm{MeV})\left(Y_{e,\mathrm{crit}}/\tau _6\right)^{1/5}`$ Similarily, consideration of $`\nu _e`$ capture on neutrons leads to $`T_{\nu _e}<(40\mathrm{MeV})\gamma _i\left(Y_{e,\mathrm{crit}}/\tau _6\right)^{1/5}.`$ In deriving this limit we have taken the $`\nu _e`$ spectrum to be a Fermi-Dirac blackbody with temperature $`T_{\nu _e}`$ and zero chemical potential. This limit could be modified or weakened if $`\nu _e`$ flavor transformation occurs.
Determining the increase in $`Y_e`$ due to pion production requires a proper treatment of neutron transport in the plasma. However, an upper limit on the increase is readily obtained by considering the extreme case where (i) the protons are frozen into the accelerating plasma (ii) non-forward n-p collissions are assumed to result in maximal momentum exchange (iii) n-n collissions are ignored except as a post-processing step to determine $`\pi `$ production and (iv) the change in $`Y_e`$ due to inelastic n-p and inelastic p-p scatterings is ignored. This simple picture gives an upper limit on the increase in $`Y_e`$ because an exchange of any of the assumptions (i)-(iii) for more realistic ones has the effect of decreasing the velocity dipersion of the neutrons. By a calculation with the above assumptions we obtain the upper limit on the increase in $`Y_e`$ to be $`\mathrm{\Delta }Y_e\begin{array}{c}<\hfill \\ \hfill \end{array}10^3/Y_{e0}`$, where $`Y_{e0}`$ is the initial electron fraction. Fig. (2) displays the evolution of the neutron distribution function as calculated with the above assumptions. The distribution function drops sharply at the instantaneous plasma velocity because there is no mechanism for boosting neutrons to higher velocity.
The ineffectiveness of pion bremsstrahlung in increasing $`Y_e`$ may be attributed to the fact that the processes which increase $`Y_e`$ result from a two step process (i.e. an n-p scattering boosts a neutron which then inelastically scatters with another neutron), the decrease of $`\sigma (v)v`$ with increasing velocity (the product of cross section and relative velocity decreases by a factor of 5 as $`v`$ increases from 0 to c ), and the late onset of this process in n-n scattering. Pion production does not begin until the pion mass threshold is reached at a relative velocity of $`0.645`$ and even at a relative velocity of $`0.728`$ (center of mass energy 2.08 GeV), the inelastic contribution to the cross section is only 8% of the total . For $`\mathrm{\Delta }Y_e>Y_{e0}`$ our perturbative approach to calculating $`\mathrm{\Delta }Y_e`$ breaks down and the increase in $`Y_e`$ may cause a recoupling of the proton and neutron flows. Note that our calculation only gives an upper limit on the increase in $`Y_e`$ due to inelastic n-n scattering. A more careful transport calculation will likely show a smaller increase in $`Y_e`$. It is consistent then to discuss decoupling at low $`Y_e`$ for a wide range of fireball parameters. The final Lorentz factor of the plasma may then be substantially increased for a given energy input and baryon load.
Aside from the fireball energetics versus $`Y_e`$ issue addressed here, neutron-proton separation recently has been shown to have implications for the electromagnetic and high-energy neutrino signatures of GRBs . Future large volume detectors such as AMANDA/ICECUBE will be able to provide high energy neutrino data on GRBs . Perhaps the details of neutron decoupling and the associated electromagnetic/neutrino signature could allow a diagnostic of the weak interaction physics deep in GRB central engine environments.
It is a pleasure to acknowledge discussions with N. Dalal, M. Patel, X. Shi and J. R. Wilson. This research was supported in part by NSF Grant PHY98-00980, an IGPP grant, and a NASA GSRP for KA.
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# The crumpling transition of membranes driven by quantum fluctuations in a 𝐷=ϵ expansion
## 1 Introduction
Fluid membranes, like biomembranes, belong to the classical world— there is no need to take quantum fluctuations into consideration since typical temperatures are high and sizes are large. Nevertheless, we might reduce temperature down to zero and increase $`\mathrm{}`$ in a Gedankenexperiment and study the quantum fluctuations of a (non relativistic) flexible membrane. The most important feature of the resulting model is a second-order quantum transition at a finite $`\mathrm{}`$ from an almost flat to a crumpled phase, contrary to the thermal case, where the membrane remains always in the crumpled phase . Ingredients of the action describing the quantum membrane are a kinetic energy term, surface tension and the Canham-Helfrich curvature energy . The latter elastic term disfavours curved configuration of the membrane and is proportional to the square of the mean curvature, integrated over the surface of the membrane. The action should not depend on the internal coordinates which are used to parametrize the surface—it should be a functional of the geometry of the membrane only.
Our model serves as a toy model for quantum interfaces like the Helium liquid–vapour interface at very low temperatures . It is remarkable, that He<sup>3</sup> enriches at a He<sup>4</sup> interface and acts as a surfactant, lowering the surface tension of the He<sup>4</sup> interface . Therefore, the next-leading curvature terms become relevant. The He<sup>3</sup> film on top of the He<sup>4</sup> bulk could be in fact a suitable candidate for the quantum membrane under investigation. It is, however, unlikey that the predicted second-order phase transition from a flat to a crumpled interface (at finite $`\mathrm{}`$) can be observed experimentally. The surface tension of the He-system is still non-zero, and moreover, there is no direct way to change $`\mathrm{}`$ in an experiment. Nevertheless, it might be possible by changing the He<sup>3</sup> concentration to tune the system closer to the transition point.
So far, nonrelativistic quantum membranes with curvature stiffness and surface tension were discussed in at fixed dimension $`D=2`$ of the membrane and $`d=3`$ of the embedding space and in for an infinite dimensional embedding space $`d\mathrm{}`$. Extending these results we study the quantum fluctuations of a membrane in a systematic $`D=ϵ`$ expansion and find a fixed point at zero temperature and finite $`\mathrm{}`$ (for $`D>0`$), which is not seen by .
## 2 The model
Following we start from the (imaginary time) action
$$S_0=t^D\sigma \sqrt{g}\left(\frac{1}{2\nu }(_t\mathrm{X}\mathbf{)}^2\mathbf{+}\mathrm{r}\mathbf{+}\frac{2}{𝜶}\mathrm{H}^2\right),$$
(2.1)
where $`\mathrm{X}\mathbf{(}𝝈\mathbf{,}\mathrm{t}\mathbf{)}`$ describes the time-dependent $`D`$-dimensional surface embedded in a $`D+1`$-dimensional, Euclidean space. Hence $`\sigma `$ is a $`D`$-dimensional set of internal coordinates which parametrize the surface. $`^D\sigma \sqrt{g}`$ is the invariant element of area, $`H`$ is the mean curvature, $`1/\nu `$ is the mass density, $`r`$ is the surface tension and $`1/\alpha `$ is the bending rigidity. The action (2.1) only makes sense in Lagrangian coordinates, i.e. coordinates which follow the elements of fluid. Then, $`\mathrm{X}\mathbf{(}𝝈\mathbf{,}\mathrm{t}\mathbf{)}`$ is (for fixed $`\sigma `$) the trajectory of an element of fluid and $`_t\mathrm{X}`$ the corresponding velocity. The velocity might be decomposed into a normal part $`v_{}=\mathrm{N}\mathbf{}\mathbf{}_\mathrm{t}\mathrm{X}`$ and the tangential parts $`u_i=_i\mathrm{X}\mathbf{}\mathbf{}_\mathrm{t}\mathrm{X}`$, where $`\mathrm{N}`$ is the normal vector and $`_i\mathrm{X}\mathbf{=}\mathbf{}\mathrm{X}\mathbf{/}\mathbf{}𝝈^\mathrm{i}`$ are the tangential vectors (see e. g. ). The normal velocity $`v_{}`$ allows for a geometrical interpretation—it encodes the variations of the shape of the membrane, whereas the tangential velocity simply generates reparametrizations of the surface (coordinate transformations). Indeed, a purely tangential velocity field leaves the shape of the membrane unchanged and, therefore, belongs to degrees of freedom beyond the geometric ones<sup>1</sup><sup>1</sup>1 The divergence of the tangential velocity field, however, is fixed by the (geometrical) condition $`2Hv_{}=D^iu_i`$ in case of an incompressible membrane..
Instead of using action (2.1) and taking care of the tangential degrees of freedom, we write down a simplified action, depending only on the shape of the membrane and its variation (with conveniently chosen coupling constants)
$$S_0/\mathrm{}=t^D\sigma \sqrt{g}\left(\frac{\lambda }{2\mathrm{}}v_{}^2+\frac{r}{\mathrm{}}+\frac{2}{\mathrm{}\lambda }H^2\right).$$
(2.2)
It can be seen easily that $`v_{}`$ is really a scalar under general time-dependent coordinate transformations $`\sigma \sigma (\sigma ^{},t)`$. Consequently, the action $`S_0`$ itself is invariant under any time-dependent coordinate transformation. To calculate the partition function $`𝒵=D[\mathrm{X}\mathbf{]}\mathrm{𝐞𝐱𝐩}\mathbf{(}\mathbf{}\mathrm{S}_0\mathbf{/}\mathbf{}\mathbf{)}`$ and related expectation values, we have to sum over all physically distinct surfaces $`\mathrm{X}\mathbf{(}𝝈\mathbf{,}\mathrm{t}\mathbf{)}`$ in a reparametrization-invariant manner. Being far from trivial (see , the quantum case does not pose an extra complication), we have to restrict the discussion of the measure problem to few remarks. At first, we have to choose a certain representation of the surfaces (gauge fixing) in order to avoid over-counting of surfaces with identical shapes, but different coordinate systems. A common and practical choice is a representation of the surface in terms of a (time dependent) heigth-variable $`f(𝐱,t)`$ \- the Monge representation ($`𝐱`$ is a $`D`$dimensional Euclidean vector)
$$\mathrm{X}\mathbf{(}𝐱\mathbf{,}\mathrm{t}\mathbf{)}\mathbf{=}\mathbf{(}𝐱\mathbf{,}\mathrm{f}\mathbf{(}𝐱\mathbf{,}\mathrm{t}\mathbf{)}\mathbf{)}\mathbf{,}$$
(2.3)
which is connected to the Lagrangian coordinates by a particular time-dependent coordinate transformation. We convert the action (2.2) into the Monge representation using $`v_{}=_tf/\sqrt{g}`$ (where $`g=1+(f)^2`$) and obtain
$`S_0/\mathrm{}`$ $`=`$ $`{\displaystyle }t^Dx[{\displaystyle \frac{\lambda }{2\mathrm{}}}{\displaystyle \frac{(_tf)^2}{\sqrt{g}}}+{\displaystyle \frac{r}{\mathrm{}}}\sqrt{g}`$
$`+{\displaystyle \frac{1}{2\mathrm{}\lambda }}\sqrt{g}(({\displaystyle \frac{1}{\sqrt{g}}}f))^2]`$
$`=`$ $`{\displaystyle }t^Dx[{\displaystyle \frac{\lambda }{2\mathrm{}}}(_tf)^2+{\displaystyle \frac{r}{2\mathrm{}}}_if_if+{\displaystyle \frac{1}{2\mathrm{}\lambda }}(^2f)^2`$
$`{\displaystyle \frac{\lambda }{4\mathrm{}}}(_tf)^2_if_if{\displaystyle \frac{r}{8\mathrm{}}}_if_if_jf_jf`$
$`{\displaystyle \frac{1}{4\mathrm{}\lambda }}_if_if(^2f)^2{\displaystyle \frac{1}{\mathrm{}\lambda }}_if_jf_i_jf^2f]+\mathrm{O}(f^6),`$
where $`i,j=1\mathrm{}D`$ <sup>2</sup><sup>2</sup>2The action (2) differs from the one used in by a wrong sign in front of the first vertex. The wrong sign, however, does not affect any consequence drawn by (for zero temperature), since the authors study the field theory not at the lower critical dimension, but at the dimension $`D=2`$, where the flow of the coupling constants is mainly determined by their naive dimensions.. The corresponding invariant measure $`D[f]`$ differs from the naive measure $`_𝐱f(𝐱)`$ by the so-called Fadeev-Popov determinant and the Liouville term, which, however, contribute to two loop and higher orders only. To one loop order, we may safely use the naive measure instead . The lower critical dimension of the theory is $`D=0`$, where the coupling constants in front of the kinetic energy and in front of the curvature energy become marginal, as can be seen from the dimensions of the coupling constants ($`L`$=length), which are $`\mathrm{}L^D`$ and $`rL^2`$ ($`\lambda `$ is rendered dimensionless, $`tL^2`$). Therefore, we have to calculate the quantum fluctuations of (2) in a double $`\mathrm{}`$ and $`D=ϵ`$ expansion, which is done here to one loop order with the help of dimensional regularization and the minimal subtraction scheme in analogy with . The surface tension $`r`$ is a relevant parameter of the theory and is zero right at the critical point. In fact, a non-zero $`r`$ imposes a finite correlation length $`\xi =r^{1/2}`$ on the propagator of action (2). Not included in the action (2) are the integral over the scalar curvature $`R`$ (which does not yield a contribution to one loop order and which is a topological invariant for $`D=2`$) and a boundary term—the cross term $`2t^Dx\sqrt{g}v_{}H=t_tA`$, where $`A`$ is the surface area.
## 3 Field theory
The bare $`T=0`$ two-point vertex function $`\mathrm{\Gamma }_{0,2}(q,\omega )`$ reads (denoting from now on bare quantities with a subscript 0)
$`\mathrm{\Gamma }_{0,2}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{\lambda _0}{\mathrm{}_0}}\omega ^2+{\displaystyle \frac{r_0}{\mathrm{}_0}}q^2+{\displaystyle \frac{1}{\mathrm{}_0\lambda _0}}q^4`$ (3.1)
$`+{\displaystyle \frac{\omega ^2}{2}}\lambda _0(r_0\lambda _0)^{ϵ/2}I_ϵ{\displaystyle \frac{q^2}{2+ϵ}}r_0(r_0\lambda _0)^{ϵ/2}I_ϵ`$
$`+q^4\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{2}{ϵ}}\right){\displaystyle \frac{1}{\lambda _0}}(r_0\lambda _0)^{ϵ/2}I_ϵ,`$
where $`I_ϵ=(4\pi )^{ϵ/21/2}\mathrm{\Gamma }(1ϵ/2)\mathrm{\Gamma }(ϵ/2+1/2)/\mathrm{\Gamma }(1+ϵ/2)`$ ($`I_ϵ`$ is finite in the limit $`ϵ0`$). We introduce renormalized couplings $`\mathrm{},\lambda ,r`$ by $`\mathrm{}_0=\mu ^ϵ\mathrm{}Z_{\mathrm{}}/I_ϵ`$, $`r_0=Z_rr`$ and $`\lambda _0=Z_\lambda \lambda `$ ($`\mu `$ is an arbitrary momentum scale) and require the renormalized vertex function
$`\mathrm{\Gamma }_2(q,\omega )`$ $`=`$ $`{\displaystyle \frac{\mu ^ϵI_ϵ}{\mathrm{}}}({\displaystyle \frac{Z_\lambda }{Z_{\mathrm{}}}}\lambda \omega ^2+{\displaystyle \frac{Z_r}{Z_{\mathrm{}}}}rq^2+{\displaystyle \frac{1}{\lambda Z_\lambda Z_{\mathrm{}}}}q^4`$ (3.2)
$`+{\displaystyle \frac{\omega ^2}{2}}\lambda \mathrm{}Z_\lambda {\displaystyle \frac{q^2}{2+ϵ}}r\mathrm{}Z_r+q^4({\displaystyle \frac{1}{2}}+{\displaystyle \frac{2}{ϵ}}){\displaystyle \frac{\mathrm{}}{\lambda Z_\lambda }})`$
to be finite in the limit $`ϵ0`$. We find $`Z_{\mathrm{}}=Z_\lambda =Z_r=1+\mathrm{}/ϵ+\mathrm{O}(\mathrm{}^2)`$. The RG-equation is readily derived from the fact, that the bare quantities do not depent on the momentum scale $`\mu `$, yielding
$$\left(\mu _\mu +\beta (\mathrm{})_{\mathrm{}}+\gamma r_r+\zeta \lambda \lambda \right)\mathrm{\Gamma }_2(q,\omega )=0$$
(3.3)
with the beta-function $`\beta (\mathrm{})=\mathrm{}(ϵ\mathrm{})`$. The theory has an ultraviolet stable fixed point (not seen by ) $`\mathrm{}^{}=ϵ`$ which corresponds to a quantum-phase transition at a finite $`\mathrm{}`$ (for $`D>0`$ and $`T=0`$) from a smooth phase for small $`\mathrm{}`$ to a crumpled phase for $`\mathrm{}>\mathrm{}^{}`$. The main effect of quantum fluctuations on the statistics of membranes is a shift of the lower critical dimension from $`D=2`$ to $`D=0`$. For slightly larger dimensions a quantum-phase transition takes place at a finite $`\mathrm{}^{}`$. We expect that this picture remains true up to the physical dimension $`D=2`$, although the numerical accuracy for the one-loop critical exponents (not presented in this short communication) should be quite limited.
I thank R Bausch, H K Janssen and R Blossey for useful discussions. This work has been supported by the Deutsche Forschungsgemeinschaft under SFB 237.
## Appendix A Details of the calculation
The bare two-point vertex function reads
$`\mathrm{\Gamma }_{0,2}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{\lambda }{\mathrm{}}}\omega ^2+{\displaystyle \frac{r}{\mathrm{}}}q^2+{\displaystyle \frac{1}{\mathrm{}\lambda }}q^4`$ (1.1)
$`{\displaystyle \frac{\lambda }{2\mathrm{}}}\omega ^2I_2{\displaystyle \frac{\lambda }{2\mathrm{}}}q^2I_1{\displaystyle \frac{r}{2\mathrm{}}}q^2I_2{\displaystyle \frac{r}{\mathrm{}ϵ}}q^2I_2`$
$`{\displaystyle \frac{1}{2\mathrm{}\lambda }}q^2I_3{\displaystyle \frac{1}{2\mathrm{}\lambda }}q^4I_2{\displaystyle \frac{2}{\mathrm{}\lambda ϵ}}q^2I_3{\displaystyle \frac{2}{\mathrm{}\lambda ϵ}}q^4I_2`$
with the Feynman-integrals
$`I_1`$ $`=`$ $`\mathrm{}{\displaystyle \frac{\omega }{2\pi }\frac{^ϵq}{(2\pi )^ϵ}\frac{\omega ^2}{\lambda \omega ^2+rq^2+\lambda ^1q^4}},`$
$`I_2`$ $`=`$ $`\mathrm{}{\displaystyle \frac{\omega }{2\pi }\frac{^ϵq}{(2\pi )^ϵ}\frac{q^2}{\lambda \omega ^2+rq^2+\lambda ^1q^4}},`$
$`I_3`$ $`=`$ $`\mathrm{}{\displaystyle \frac{\omega }{2\pi }\frac{^ϵq}{(2\pi )^ϵ}\frac{q^4}{\lambda \omega ^2+rq^2+\lambda ^1q^4}}.`$ (1.2)
Within dimensional regularization we have ($`\omega ^ϵq0`$)
$`I_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{\lambda }}{\displaystyle \frac{\omega }{2\pi }\frac{^ϵq}{(2\pi )^ϵ}\frac{\lambda \omega ^2+rq^2+\lambda ^1q^4rq^2\lambda ^1q^4}{\lambda \omega ^2+rq^2+\lambda ^1q^4}}`$ (1.3)
$`=`$ $`{\displaystyle \frac{r}{\lambda }}I_2{\displaystyle \frac{1}{\lambda ^2}}I_3.`$
To evaluate $`I_2`$, we substitute $`\omega \omega q^2`$ and find
$$I_2=\mathrm{}\frac{^ϵq}{(2\pi )^ϵ}\frac{\omega }{2\pi }\frac{q^2}{\lambda \omega ^2q^2+r+\lambda ^1q^2}.$$
(1.4)
Now we are able to perform the $`q`$-integration
$$I_2=\frac{\mathrm{}}{r}\frac{\omega }{2\pi }\left(\frac{r}{\lambda \omega ^2+\lambda ^1}\right)^{1+ϵ/2}I,$$
(1.5)
where
$$I=\frac{^ϵk}{(2\pi )^ϵ}\frac{k^2}{k^2+1}=\frac{1}{(4\pi )^{ϵ/2}}\mathrm{\Gamma }(1ϵ/2).$$
(1.6)
With the help of $`(2\pi )^1s(s^2+1)^{1ϵ/2}=(4\pi )^{1/2}\mathrm{\Gamma }(ϵ/2+1/2)/\mathrm{\Gamma }(1+ϵ/2)`$ we obtain
$$I_2=\mathrm{}(r\lambda )^{ϵ/2}I_ϵ,$$
(1.7)
where
$$I_ϵ=\frac{1}{(4\pi )^{ϵ/2+1/2}}\frac{\mathrm{\Gamma }(1ϵ/2)\mathrm{\Gamma }(ϵ/2+1/2)}{\mathrm{\Gamma }(1+ϵ/2)}.$$
(1.8)
An analogous calculation yields
$$I_3=\lambda r\frac{1+ϵ}{2+ϵ}I_2=\mathrm{}\lambda r(r\lambda )^{ϵ/2}\frac{1+ϵ}{2+ϵ}I_ϵ.$$
(1.9)
Within a cut-off regularization scheme instead of dimensional regularization additional divergent terms show up in equation (1.3) and equation (1.9) which can be absorbed by an additive renormalization of the surface tension $`r`$.
## References
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# 1 Introduction
## 1 Introduction
The $`q`$-state Potts antiferromagnet (AF) exhibits nonzero ground state entropy, $`S_0>0`$ (without frustration) for sufficiently large $`q`$ on a given lattice $`\mathrm{\Lambda }`$ or, more generally, on a graph $`G`$. This is equivalent to a ground state degeneracy per site $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. Such nonzero ground state entropy is important as an exception to the third law of thermodynamics. There is a close connection with graph theory here, since the zero-temperature partition function of the above-mentioned $`q`$-state Potts antiferromagnet on a graph $`G`$ satisfies
$$Z(G,q,T=0)_{PAF}=P(G,q)$$
(1.1)
where $`P(G,q)`$ is the chromatic polynomial expressing the number of ways of coloring the vertices of the graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color (for reviews, see -). The minimum number of colors necessary for such a coloring of $`G`$ is called the chromatic number, $`\chi (G)`$. Thus<sup>1</sup><sup>1</sup>1At certain special points $`q_s`$ (typically $`q_s=0,1,..,\chi (G)`$), one has the noncommutativity of limits $`lim_{qq_s}lim_n\mathrm{}P(G,q)^{1/n}lim_n\mathrm{}lim_{qq_s}P(G,q)^{1/n}`$, and hence it is necessary to specify the order of the limits in the definition of $`W(\{G\},q_s)`$ . We use the first order of limits here; this has the advantage of removing certain isolated discontinuities in $`W`$.
$$W(\{G\},q)=\underset{n\mathrm{}}{lim}P(G,q)^{1/n}$$
(1.2)
where $`n=v(G)`$ is the number of vertices of $`G`$ and $`\{G\}=lim_n\mathrm{}G`$. Since $`P(G,q)`$ is a polynomial, one can generalize $`q`$ from $`_+`$ to $``$. The zeros of $`P(G,q)`$ in the complex $`q`$ plane are called chromatic zeros; a subset of these may form an accumulation set in the $`n\mathrm{}`$ limit, denoted $``$, which is the continuous locus of points where $`W(\{G\},q)`$ is nonanalytic. <sup>2</sup><sup>2</sup>2Although it does not happen in the cases considered here, for some families of graphs $``$ may be null, and $`W`$ may also be nonanalytic at certain discrete points. The maximal region in the complex $`q`$ plane to which one can analytically continue the function $`W(\{G\},q)`$ from physical values where there is nonzero ground state entropy is denoted $`R_1`$. The maximal value of $`q`$ where $``$ intersects the (positive) real axis is labelled $`q_c(\{G\})`$. This point is important since it separates the interval $`q>q_c(\{G\})`$ on the positive real $`q`$ axis where the Potts model (with $`q`$ extended from $`_+`$ to $``$) exhibits nonzero ground state entropy (which increases with $`q`$, asymptotically approaching $`S_0=k_B\mathrm{ln}q`$ for large $`q`$, and which for a regular lattice $`\mathrm{\Lambda }`$ can be calculated approximately via large–$`q`$ series expansions) from the interval $`0qq_c(\{G\})`$ in which $`S_0`$ has a different analytic form.
In the present work we report exact calculations of $`P(G,q)`$, $`W(\{G\},q)`$, and $``$ for families of triangular lattice strips of fixed width $`L_y=3`$ vertices and arbitrary length $`L_x=m`$ vertices. The longitudinal and transverse directions are taken to be $`\widehat{x}`$ and $`\widehat{y}`$). In Fig. 1 we display some illustrative examples. A generic form for chromatic polynomials for recursively defined families of graphs, of which strip graphs $`G_s`$ are special cases, is
$$P((G_s)_m,q)=\underset{j=1}{\overset{N_\lambda }{}}c_j(q)(\lambda _j(q))^m$$
(1.3)
where $`c_j(q)`$ and the $`N_\lambda `$ terms $`\lambda _j(q)`$ depend on the type of strip graph $`G_s`$ but are independent of $`m`$.
Let us comment on the motivations for the current study and how it extends previous work. Calculations of the chromatic polynomials for the cyclic and Möbius strip graphs of the square lattice were carried out for $`L_y=2`$ in (see also ). Subsequently, these calculations were extended to the width $`L_y=3`$ for the cyclic and Möbius square lattice strips. After studies of chromatic zeros in and, for $`L_y=2`$ strips in , the $`W`$ and $``$ defined in the infinite-length limit were determined for $`L_y=2`$ in and for $`L_y=3`$ in . Similar calculations were carried out for strips of various lattices with widths up to $`L_y=4`$ and free boundary conditions in , and for cyclic and Möbius graphs involving homeomorphic expansions of square strips and cyclic polygon chains . An important question concerns the effect of boundary conditions (BC’s), and hence graph topology, on $`P`$, $`W`$, and $``$. We use the symbols FBC<sub>y</sub> and PBC<sub>y</sub> for free and periodic transverse boundary conditions and FBC<sub>x</sub>, PBC<sub>x</sub>, and TPBC<sub>x</sub> for free, periodic, and twisted periodic longitudinal boundary conditions. The term “twisted” means that the longitudinal ends of the strip are identified with reversed orientation. These strip graphs can be embedded on surfaces with the following topologies: (i) (FBC<sub>y</sub>,FBC<sub>x</sub>): open strip; (ii) (PBC<sub>y</sub>,FBC<sub>x</sub>): cylindrical; (iii) (FBC<sub>y</sub>,PBC<sub>x</sub>): cylindrical, denoted cyclic here; (iv) (FBC<sub>y</sub>,TPBC<sub>x</sub>): Möbius; (v) (PBC<sub>y</sub>,PBC<sub>x</sub>): torus; and (vi) (PBC<sub>y</sub>,TPBC<sub>x</sub>): Klein bottle.<sup>3</sup><sup>3</sup>3These BC’s can all be implemented in a manner that is uniform in the length $`L_x`$; the case (vii) (TPBC<sub>y</sub>,TPBC<sub>x</sub>) with the topology of the projective plane requires different identifications as $`L_x`$ varies and will not be considered here. (We recall that that unlike graphs of type (i)-(v), the Klein bottle surface cannot be embedded without self-intersection in $`^3`$.)
Several specific effects of different boundary conditions have been found (for a general discussion, see ). First, from comparisons of exact calculations of $`W`$ for the infinite-length limit of lattice strips with boundary conditions of types (i)-(iv), it was observed that those with a global circuit<sup>4</sup><sup>4</sup>4A global circuit is a route following a lattice direction which has the topology of $`S^1`$ and a length $`\mathrm{}_{g.c.}`$ that goes to infinity as $`n\mathrm{}`$. lead to a locus $``$ that separates the $`q`$ plane into different regions, and it was inferred that the presence of global circuits is a sufficient condition for $``$ to have this property .<sup>5</sup><sup>5</sup>5This is not a necessary condition, as was shown in . For strip graphs of regular lattices, this is equivalent to having $`PBC_x`$ or $`TPBC_x`$, i.e., periodic or twisted periodic boundary conditions in the direction in which the strip length goes to infinity as $`L_x\mathrm{}`$. Further support for this inference was adduced from the calculation of $`P`$, $`W`$, and $``$ for the $`L_y=3`$ strip of the square lattice with torus and Klein bottle boundary conditions . Applications of a coloring matrix method for these calculations has been discussed in . For the (homogeneous) strip graphs with (FBC<sub>y</sub>,FBC<sub>x</sub>) in , $``$ consists of arcs that do not separate regions of the $`q`$ plane. As the width of the strip increases, these arcs tend to elongate, and their ends tend to move toward each other, thereby suggesting that if one considered the sequence of strip graphs of this type with width $`L_y`$ (having taken the limit $`L_x\mathrm{}`$ to obtain a locus $``$ for each member of this sequence), then in the limit $`L_y\mathrm{}`$, the arcs would close to form a $``$ that separated the $`q`$ plane into two or more regions and passed through $`q=0`$. The interesting feature of the families of graphs with global circuits is that when the graphs contain a global circuit, this separation of the $`q`$ plane into regions (with $``$ passing through $`q=0`$) already occurs for finite $`L_y`$. This means that the $`W`$ functions of the infinite-length limit of these graphs with cyclic and twisted cyclic longitudinal boundary conditions already exhibit a feature which is expected to occur for the $``$ for the $`W`$ function of the full two-dimensional lattice. This expectation is supported by the calculation of $`W`$ and $``$ for the 2D triangular lattice with cylindrical boundary conditions by Baxter . Thus, although calculations of chromatic polynomials for lattice strips with global circuits are, in general, more difficult than for the corresponding strips with free boundary conditions, such calculations are worthwhile, since in the $`L_x\mathrm{}`$ limit, the resultant locus $``$ embodies the analytic property of $``$ expected for the full 2D lattice, viz., that it crosses the real axis at least at $`q=0`$ and a maximal point $`q_c`$, and separates the $`q`$ plane into different regions.
Second, it has been found that for a given type of strip graph $`G_s`$ with FBC<sub>y</sub>, the chromatic polynomials for PBC<sub>x</sub> and TPBC<sub>x</sub> boundary conditions (i.e., cyclic and Möbius strips) have the same $`\lambda _j`$, although in general different $`c_j`$. It follows that the loci $``$ are the same for these two different longitudinal boundary condition choices . In the case of PBC<sub>y</sub>, the reversal of orientation involved in going from PBC<sub>x</sub> to TPBC<sub>x</sub> longitudinal boundary conditions (i.e. from torus to Klein bottle topology) can lead to the removal of some of the $`\lambda _j`$’s that were present; i.e., $`P`$ for the (PBC<sub>y</sub>,TPBC<sub>x</sub>) strip may involve only a subset of the $`\lambda _j`$’s that are present for the (PBC<sub>y</sub>,PBC<sub>x</sub>) strip. For example, for the $`L_y=3`$ strips of the square lattice with (PBC<sub>y</sub>, PBC<sub>x</sub>) boundary conditions, there are $`N_\lambda =8`$ $`\lambda _j`$’s, but for the strip with (PBC<sub>y</sub>, TPBC<sub>x</sub>) boundary conditions only a subset of $`N_\lambda =5`$ of these terms occurs in $`P`$ . None of the three $`\lambda _j`$’s that are absent from $`P`$ in the TPBC<sub>x</sub> case is leading, so that $``$ is the same for both of these families.
Third, for a given type of strip graph $`G_s`$ containing a global circuit, it has been found that, in the infinite-length limit where the locus $``$ is defined, it not only passes through the origin, $`q=0`$, but always crosses the positive real axis at one or more points, the maximal one being denoted $`q_c(\{G\})`$, as mentioned above. In contrast, for strip graphs that do not contain global circuits, $``$ may not cross the real axis. For example, for the (infinite-length limit of the) strip of the triangular lattice with free transverse and longitudinal boundary conditions, $`(FBC_y,FBC_x)`$, the locus $``$ crosses the real axis for width $`L_y=3`$ but does not cross it for $`L_y=4`$ (see Fig. 5 of ). For square strips with $`(FBC_y,(T)PBC_x)`$ boundary conditions, it has been found that $`q_c(\{G\})`$ is a nondecreasing function of $`L_y`$ for the cases studied so far, namely $`q_c=2`$ for $`L_y=1`$ and $`L_y=2`$, and $`q_c2.33654`$ for $`L_y=3`$ . For the strips of the triangular lattice with $`(FBC_y,(T)PBC_x)`$ boundary conditions, we only have one width for which $``$ has been determined, namely, $`L_y=2`$, for which $`q_c=3`$. It is thus of interest to explore what the value of $`q_c`$ is for cyclic and Möbius strips of the triangular lattice with larger widths. For the $`L_y=3`$ torus and Klein bottle strips of the square lattice, very interestingly, $`q_c=3`$, which is the same value as for the 2D square lattice . Clearly it is worthwhile to investigate what the value of $`q_c`$ is for torus and Klein bottle strips of the triangular lattice, and we shall do this here.
Fourth, for the strip graph $`(G_s)_m`$ with a given type of transverse boundary conditions BC<sub>y</sub>, the chromatic polynomial for PBC<sub>x</sub> has a larger value of $`N_\lambda `$ than the chromatic polynomial for FBC<sub>x</sub>, and the corresponding loci $``$ are different.
Fifth, one may ask how the $`W`$ functions in region $`R_1`$ compare for different boundary conditions. It has been found that, for a given type of strip graph $`G_s`$, in the region $`R_1`$ (with its left-hand boundary on the positive real $`q`$ axis at $`q_c`$) being defined for the PBC<sub>x</sub> boundary conditions, the $`W`$ function is the same for FBC<sub>x</sub>, PBC<sub>x</sub>, and TPBC<sub>x</sub>. This includes the region of real $`q`$ greater than this value of $`q_c`$, so that this result is somewhat reminiscent of the statement of the existence of the thermodynamic limit in statistical mechanics, i.e. the fact that in the disordered phase of a statistical mechanical system the thermodynamic functions are independent of the boundary conditions and if an ordered phase exists, then the role of the boundary conditions is only to set a preferred direction for the symmetry-breaking order parameter.
Our new results on $`P`$, $`W`$, and $``$ for the $`L_y=3`$ strip of the triangular lattice with the various boundary conditions of type (iii)-(vi) give insight concerning these five items, as do our additional results for $`L_y=4`$ cyclic strips and $`L_y=5,6`$ cylindrical strips. Calculations of these quantities for the $`L_y=2`$ cyclic and Möbius triangular strips were presented in . An early study of $`P`$ for the $`L_y=2`$ strip is in <sup>6</sup><sup>6</sup>6We note the following corrections in Sand’s result for $`P`$ for the $`L_y=2`$ cyclic strip of the triangular lattice, given on p. 88 of : the coefficients of the last two terms should be reversed in sign to read ($`zq`$ in his notation) $`(z1)\left[[(1/2)(52z+\sqrt{94z})]^m+[(1/2)(52z\sqrt{94z})]^m\right]`$, as in eq. (5.23) of . and a coloring matrix method to obtain $`P`$ has been given in . $`P`$, $`W`$, and $``$ have been calculated for triangular lattice strips with free boundary conditions (i) for $`L_y`$ up to 4 . (Some related calculations including work on cases with noncompact $``$, are , -.)
For strips with global circuits, and $`L_x`$ above the lowest values involving degenerate cases, the triangular lattice strips with boundary conditions of type (i)-(vi) have $`n=L_yL_x`$ vertices; the cyclic and Möbius triangular strips have $`e=(3L_y2)L_x`$ edges while the torus and Klein bottle triangular strips have $`e=3L_yL_x`$ edges. Define $`N_t`$ to be the number of triangles in the strips (i.e., excluding the triangles forming the cross sections of the strips in the case $`L_y=3`$); then, again for $`L_x`$ and $`L_y`$ large enough to avoid degenerate cases, the cyclic and Möbius strips have $`N_t=2(L_y1)L_x`$ such triangles, while the torus and Klein bottle strips have $`2L_yL_x`$ such triangles. For strips of type (i), $`e=3L_xL_y2L_x2L_y+1`$ and $`N_t=2(L_y1)(L_x1)`$ while for strips of type (ii), $`e=3L_xL_y2L_y`$ and $`N_t=2L_y(L_x1)`$.
Again, for $`L_x=m`$ and $`L_y`$ large enough to avoid degenerate cases, the chromatic number for the triangular lattice strips with cyclic or torus boundary conditions is
$$\chi (tri(L_y)_m,FBC_y,PBC_x)=\chi (tri(L_y)_m,PBC_y,PBC_x)=\{\begin{array}{cc}3\hfill & \text{if }m=0\text{ mod 3}\hfill \\ 4\hfill & \text{if }m=1\text{ or 2 mod 3}\hfill \end{array}$$
(1.4)
and for the triangular lattice strips with Möbius or Klein bottle boundary conditions,
$$\chi (tri(L_y)_m,FBC_y,TPBC_x)=\chi (tri(L_y)_m,PBC_y,TPBC_x)=4m3.$$
(1.5)
For strips with $`(FBC_y,FBC_x)`$ or $`(PBC_y,FBC_x)`$, $`\chi =3`$.
## 2 $`L_y=3`$ Cyclic and Möbius Triangular Strips
We calculate the chromatic polynomials by iterated use of the deletion-contraction theorem, via a generating function approach . In general, for a strip graph of type $`G_s`$ we have
$$\mathrm{\Gamma }(G_s,q,x)=\underset{m=m_0}{\overset{\mathrm{}}{}}P((G_s)_m,q)x^{mm_0}$$
(2.1)
where, as before , we take $`m_0=2`$ for strips with periodic or twisted periodic longitudinal boundary conditions. The generating functions $`\mathrm{\Gamma }(G_s,q,x)`$ are rational functions of the form
$$\mathrm{\Gamma }(G_s,q,x)=\frac{𝒩(G_s,q,x)}{𝒟(G_s,q,x)}$$
(2.2)
with
$$𝒩(G_s,q,x)=\underset{j=0}{\overset{d_𝒩}{}}A_{G_s,j}(q)x^j$$
(2.3)
and
$$𝒟(G_s,q,x)=1+\underset{j=1}{\overset{d_𝒟}{}}b_{G_s,j}(q)x^j$$
(2.4)
where the $`A_{G_s,i}`$ and $`b_{G_s,i}`$ are polynomials in $`q`$ (with no common factors) and
$$d_𝒩=deg_x(𝒩)$$
(2.5)
and
$$d_𝒟=deg_x(𝒟).$$
(2.6)
The denominator $`𝒟`$ can be written in factorized form as
$$𝒟(G_s,q,x)=\underset{j=1}{\overset{d_𝒟}{}}(1\lambda _{G_s,j}x).$$
(2.7)
Equivalently, the $`\lambda _{G_s,j}`$ are roots of the equation
$$\xi ^{d_𝒟}𝒟(G_s,q,1/\xi )=\xi ^{d_𝒟}+\underset{j=1}{\overset{d_𝒟}{}}b_{G_s,j}\xi ^{d_𝒟j}.$$
(2.8)
From the generating function, we calculate the chromatic polynomials in the form of eq. (1.3) by the use of the general formula, eq. (2.14), in ,
$$P(G_m,q)=\underset{j=1}{\overset{d_𝒟}{}}\left[\underset{s=0}{\overset{d_𝒩}{}}A_s\lambda _j^{d_𝒟s1}\right]\left[\underset{1id_𝒟;ij}{}\frac{1}{(\lambda _j\lambda _i)}\right]\lambda _j^{mm_0}$$
(2.9)
(where the conventions in eq. (2.1) were chosen such that $`m_0=0`$ in ). Note that some of the coefficients $`c_{G_s,j}`$ may vanish, so that not all of the $`\lambda _{G_s,j}`$’s in $`𝒟(G_s,q,x)`$ contribute to the sum in (1.3) . The formula (1.3) for the chromatic polynomial has the virtue of being a closed-form expression. However, for a given type of strip graph, for values of the width $`L_y`$ greater than the first one or two values, the set of $`\lambda _j`$’s include nonpolynomial algebraic roots. These occur, of course, as symmetric functions of the roots, so that, by a theorem on symmetric functions of roots of algebraic equations, their resultant contributions to $`P`$ are expressible in terms of the coefficients of these equations (which are polynomials in $`q`$) and hence are polynomials in $`q`$ . However, in calculating chromatic polynomials and the resultant chromatic zeros for large values of $`L_x=m`$, it can be convenient to use directly the expression for the generating function, since both the numerator and denominator of this function, eqs. (2.3), (2.4), are polynomials in $`q`$ so that one does not have to rely upon cancellations of nonpolynomial algebraic expressions at intermediate stages in the calculation.
The coefficients $`c_{tri(L_y),j}`$ that enter into the expressions for the chromatic polynomial (1.3) for the cyclic triangular ($`t`$) strip of width $`L_y`$ are (like those for the cyclic square strip) certain polynomials that we denote $`c^{(d)}`$, given by
$$c^{(d)}=\underset{k=1}{\overset{d}{}}(qq_{d,k})$$
(2.10)
where
$$q_{d,k}2+2\mathrm{cos}\left(\frac{2\pi k}{2d+1}\right)\mathrm{for}k=1,2,\mathrm{}d$$
(2.11)
with $`0dL_y`$. We list below the specific $`c^{(d)}`$’s that appear in our results for the $`L_y=3`$ and $`L_y=4`$ cyclic strips of the triangular lattice
$$c^{(0)}=1,c^{(1)}=q1,c^{(2)}=q^23q+1,$$
(2.12)
$$c^{(3)}=q^35q^2+6q1,$$
(2.13)
and
$$c^{(4)}=(q1)(q^36q^2+9q1).$$
(2.14)
Where the coefficient $`c^{(0)}`$ appears in chromatic polynomials, we shall simply write it as unity.
For the cyclic $`L_y=3`$ triangular strip of length $`L_x=m`$, we obtain the result
$`P(tri(3\times m),FBC_y,PBC_x,q)={\displaystyle \underset{j=1}{\overset{10}{}}}c_{t3,j}(\lambda _{t3,j})^m`$ (2.15)
$`=`$ $`c^{(3)}(1)^m+c^{(2)}\left[(q2)^m+(\lambda _{t3,3})^m+(\lambda _{t3,4})^m\right]`$ (2.17)
$`+`$ $`c^{(1)}\left[[(q2)(q3)]^m+(\lambda _{t3,6})^m+(\lambda _{t3,7})^m+(\lambda _{t3,8})^m\right]+\left[(\lambda _{t3,9})^m+(\lambda _{t3,10})^m\right].`$ (2.19)
where the short notation $`t3`$ is used in subscripts for $`tri(L_y=3)`$. The labelling of $`c_{t3,j}`$ and $`\lambda _{t3,j}`$ in eq. (LABEL:ptcyclic) is consecutive, so that $`c_{t3,1}=c^{(3)}`$, $`\lambda _{t3,1}=1`$; $`c_{t3,2}=c^{(2)}`$, $`\lambda _{t3,2}=q2`$, and so forth:
$$\lambda _{t3,(3,4)}=\frac{1}{2}\left(7+2q\pm \sqrt{258q}\right),$$
(2.22)
$`\lambda _{t3,j}`$, $`j=6,7,8`$ are the roots of the cubic equation
$$\xi ^3+b_{t3,1}\xi ^2+b_{t3,2}\xi +b_{t3,3}=0$$
(2.23)
where
$$b_{t3,1}=2q^212q+19$$
(2.24)
$$b_{t3,2}=(q2)(q^39q^2+28q29)$$
(2.25)
$$b_{t3,3}=(q2)^4(q3),$$
(2.26)
and
$$\lambda _{t3,(9,10)}=\frac{1}{2}\left[q^37q^2+18q17\pm (q^614q^5+81q^4250q^3+442q^2436q+193)^{1/2}\right].$$
(2.27)
Note that some $`\lambda _{t3,j}`$’s vanish if $`q=2`$ or $`q=3`$; specifically, if $`q=2`$, then $`\lambda _{t3,3}=0`$ and two of the roots of the cubic equation (2.23) vanish, while for $`q=3`$, $`\lambda _{t3,3}=0`$ and one of the roots of eq. (2.23) vanish. The vanishing of the chromatic polynomial at $`q=0,1,2`$ and (i) $`q=3`$ for the cyclic strip with $`m0`$ mod 3 or (ii) $`q=3`$ for the Möbius strip discussed below involves both some terms $`\lambda _j`$ vanishing and cancellations among others. This is also true of chromatic polynomials for strips of other lattices.
Parenthetically, we note that the denominator of the generating function for this strip graph actually has $`d_𝒟=16`$:
$$𝒟(tri(L_y=3),FBC_y,PBC_x,q,x)=\underset{j=1}{\overset{16}{}}(1\lambda _{t3,j}(q)x)$$
(2.28)
where the $`\lambda _{t,j}`$ were given in eq. (LABEL:ptcyclic) and
$$\lambda _{t,11}=\lambda _{t,1}=1$$
(2.29)
$$\lambda _{t,12}=\lambda _{t,13}=q3$$
(2.30)
and $`\lambda _{t,j}`$, $`j=14,15,16`$ are the roots of the cubic equation
$$\xi ^3+(2q5)(q3)\xi ^2+(q2)^2(q3)(q4)\xi (q2)^4(q3)=0.$$
(2.31)
A new feature of these results that we have not encountered before in our calculations of chromatic polynomials and generating functions for strip graphs is the occurrence of factors in $`𝒟`$ that have multiplicity higher than unity: specifically, the factors $`(1+x)^2`$, corresponding to $`\lambda _{t3,1}=\lambda _{t3,11}=1`$, and $`[1(q3)x]^2`$, corresponding to $`\lambda _{t3,12}=\lambda _{t3,13}=q3`$. None of the $`\lambda _{t3,j}`$, $`12j16`$, contribute to $`P`$ (and $`\lambda _{t3,11}`$ is identical to $`\lambda _{t3,1}`$), so that $`P`$ only involves the ten $`\lambda _{t3,j}`$’s in eq. (LABEL:ptcyclic). We had found cases earlier of some $`\lambda _j`$’s in $`𝒟`$ not contributing to $`P`$, e.g., for the $`L_y=2`$ cyclic strip of the kagomé lattice and homeomorphic expansions of the $`L_y=2`$ cyclic strip of the square lattice . However, this is the first value of $`L_y`$ where we have encountered this type of behavior for a strip of a homopolygonal lattice.<sup>7</sup><sup>7</sup>7As in , we define a homopolygonal 2D lattice as a regular tiling of the plane involving a single type of polygon, in which all vertices are equivalent, while a heteropolygonal lattice is the same but involving two or more different types of regular polygons. In the case of the kagomé lattice, these are triangles and hexagons.
The numerator of the generating function has degree $`deg_x(𝒩)=14`$. The polynomials $`A_j`$, $`j=1,..,14`$ comprising this numerator can be calculated from the expressions for $`P`$ and $`𝒟`$, eqs. (LABEL:ptcyclic) and (2.28) and hence are not listed here.
From eq. (1.4), it follows that for $`m=1`$ or 2 mod 3 and $`q=3`$, the chromatic polynomial in eq. (LABEL:ptcyclic) vanishes; for $`m=0`$ mod 3 (and $`m3`$) and $`q=3`$ we find
$$P(tri(L_y=3)_m,FBC_y,PBC_x,q=3)=6\mathrm{for}m=0mod3.$$
(2.32)
This may be contrasted with the case of a $`k`$-critical graph $`G_k`$, where $`\chi (G)=k`$ and $`P(G,q=\chi (G))=k!`$ (e.g. $`k=3`$ for the triangular lattice).
We define
$$C(G)=\underset{j=1}{\overset{N_{\lambda _G}}{}}c_{G,j}.$$
(2.33)
where the $`G`$-dependence in the coefficients is indicated explicitly. Note that for recursive graphs like the strip graphs considered here, the $`c_{G,j}`$ depend on $`L_y`$ and the boundary conditions, but not on $`L_x`$. Our results above give
$$C=q(q1)^2\mathrm{for}G=tri(L_y=3,FBC_y,PBC_x).$$
(2.34)
in accord with the generalization $`C=P(T_{L_y},q)`$ for this strip, where $`P(T_n,q)=q(q1)^{n1}`$ and $`T_n`$ is the tree graph with $`n`$ vertices. This was also true for the $`L_y=3`$ cyclic strip of the square lattice and is in accord with the coloring matrix approach .
For the $`L_y=3`$ triangular lattice Möbius strip (denoted by the subscript $`t3Mb`$) we calculate
$$P(tri(L_y=3)_m,FBC_y,TPBC_x,q)=\underset{j=1}{\overset{10}{}}c_{t3Mb,j}(\lambda _{t3Mb,j})^m$$
(2.35)
where
$$\lambda _{t3Mb,j}=\lambda _{t3,j},j=1,..,10.$$
(2.36)
As was true in the case of the L$`{}_{y}{}^{}=2`$ Möbius strip of the triangular lattice, the coefficients $`c_j`$ are not just of the $`c^{(d)}`$ polynomial form; instead, some coefficients are rational functions of $`q`$ and others are algebraic nonpolynomial functions of $`q`$:
$$c_{t3Mb,1}=c^{(2)}$$
(2.37)
$$c_{t3Mb,2}=\frac{(q^24q+5)}{2(q2)}$$
(2.38)
$$c_{t3Mb,(3,4)}=\frac{1}{4(q2)}\left[(q3)^2\frac{(q1)(q5)}{\lambda _{t,3}\lambda _{t,4}}\right]$$
(2.39)
$$c_{t3Mb,5}=\frac{(q1)(q^24q+5)}{2(q2)}$$
(2.40)
$$c_{t3Mb,9}=c_{t3Mb,10}=c^{(0)}=1.$$
(2.41)
The $`j=3,4`$ terms can be written in a manner which manifestly exhibits their property of being a symmetric function of $`\lambda _{t3,3}`$ and $`\lambda _{t3,4}`$:
$$c_{t3Mb,3}(\lambda _{t3,3})^m+c_{t3Mb,4}(\lambda _{t3,4})^m=\frac{(q3)^2}{4(q2)}\left[(\lambda _{t3,3})^m+(\lambda _{t3,4})^m\right]\frac{(q1)(q5)}{4(q2)}\left[\frac{(\lambda _{t3,3})^m(\lambda _{t3,4})^m}{\lambda _{t3,3}\lambda _{t3,4}}\right].$$
(2.42)
It is convenient to leave the expressions for $`c_{t3Mb,j}`$, $`j=6,7,8`$, in the general forms that follow from eq. (2.9):
$$c_{t3Mb,j}=\lambda _{t3,j}^2\left[\underset{s=0}{\overset{8}{}}A_{tm3,s}(\lambda _{t3,j})^{9s}\right]\left[\underset{1i10;ij}{}\frac{1}{(\lambda _{t3,j}\lambda _{t3,i})}\right],j=6,7,8$$
(2.43)
where the $`A_{t3Mb,s}`$ are given in the appendix.
For the sum of the coefficients for the $`L_y=3`$ Möbius strip of the triangular lattice we calculate
$$C=\underset{j=1}{\overset{10}{}}c_{t3Mb,j}=q(q1)\mathrm{for}G=tri(L_y=3,FBC_y,TPBC_x)$$
(2.44)
in agreement with a general formula that we have derived elsewhere for a Möbius strip graph $`G_s`$ of width $`L_y`$ of the square or triangular lattice ,
$$C(G_s(L_y),Mb)=\{\begin{array}{cc}0\hfill & \text{for even }L_y\hfill \\ P(T_{(\frac{L_y+1}{2})},q)\hfill & \text{for odd }L_y\hfill \end{array}.$$
(2.45)
Chromatic zeros for the $`L_y=3`$, $`L_x=m=20`$ cyclic graph of the triangular lattice are shown in Fig. 2; with this value of $`m`$, the chromatic zeros for the triangular lattice Möbius graph are quite similar and both sets of chromatic zeros lie close to the boundary $``$ and indicate its position. The locus $``$ and the $`W`$ functions are the same for the cyclic and Möbius graph families. We find
$$q_c(tri(L_y=3),FBC_y,(T)PBC_x)=3.$$
(2.46)
This will be discussed further below when we give our result for $`q_c`$ for the corresponding $`L_y=4`$ strip of the triangular lattice.
The locus $``$ has support for $`Re(q)0`$ and separates the $`q`$ plane into five main regions. The outermost one, region $`R_1`$, extends to infinite $`|q|`$ and includes the intervals $`q3`$ and $`q0`$ on the real $`q`$ axis. The other regions are labelled consequtively moving leftward along the real axis from $`q=q_c=3`$ to $`q=0`$, and then including complex conjugate pairs of regions that have no overlap with the real axis. Thus, region $`R_2`$ includes the real interval $`2q3`$, while region $`R_3`$ includes the real interval $`0q2`$. There are also two complex-conjugate regions, $`R_4,R_4^{}`$ centered at approximately $`q3.35\pm 0.7i`$. The triple points occur at (i) $`q3.14+0.37i`$ and (ii) $`q3.12+1.12i`$, where $`R_1,R_2`$, and $`R_4`$ are contiguous, and (iii) $`q2.25+1.66i`$, where $`R_1`$, $`R_2`$, and $`R_3`$ are contiguous, together with the complex conjugates of these points. In addition, there are two tiny regions $`R_5,R_5^{}`$ that touch the real axis at the single point $`q_c=3`$ and extend a very short distance to the upper and lower right, with triple points at about $`q3.02\pm 0.13i`$. These are plotted in Fig. 2. Thus, $`q_c`$ is a multiple point where four branches of $``$ intersect. The occurrence of extremely small regions was also found for the ($`L_x\mathrm{}`$ limit of the) $`L_y=3`$ cyclic strip of the square lattice, as shown in Fig. 1 of . The fact that such tiny regions can occur means that when one performs the usual computer scan over the complex plane to map out the region diagram and determine the dominant $`\lambda _j`$’s in each region, since this involves a finite grid, one can only detect tiny sliver regions down to a certain resolution set by the grid of the scan. On the real axis, this is not a serious complication, since the mapping is a one-variable problem, but in the complex plane, one must a commensurately fine scanning grid to detect tiny sliver regions.
The part of $``$ and the associated chromatic zeros with largest magnitude occur at roughly $`q3.53+0.69i`$ and have $`|q|3.59`$. Comparing the present locus $``$ for $`L_y=3`$ with the one obtained in for $`L_y=2`$, we obtain further evidence supporting the conjecture that for a strip graph of a regular lattice with periodic or twisted periodic longitudinal boundary conditions ($``$ is the same for both of these), the envelope curve of $``$, i.e. the outermost portion of $``$ for a given $`L_y`$ surrounds the envelope curve for the $``$ of the same strip with a smaller value of $`L_y`$. This was found to be true for the cyclic and Möbius strips of the square lattice for the cases $`L_y=2,3`$ and 4 (and, for the cyclic case, also $`L_y=1`$) . For strips of a regular lattice $`\mathrm{\Lambda }`$, as $`L_y\mathrm{}`$, one expects that the envelope curves defined for each $`L_y`$ will approach a limiting curve, which is precisely the envelope for the boundary $``$ of the full 2D lattice $`\mathrm{\Lambda }`$ defined via this limit.
In region $`R_1`$, $`\lambda _{t3,9}`$ is the dominant $`\lambda _j`$ (with appropriate choice of branch cut for evaluation away from the negative real axis), so
$$W=(\lambda _{t3,9})^{1/3},qR_1.$$
(2.47)
The fact that this is the same as $`W`$ for the (FBC<sub>y</sub>,FBC<sub>x</sub>) case is a general result . In region $`R_2`$, $`\lambda _{t3,4}`$ is dominant, so
$$|W|=|\lambda _{t3,4}|^{1/3},qR_2$$
(2.48)
(in regions other than $`R_1`$, only $`|W|`$ can be determined unambiguously ). In region $`R_3`$, $`W`$ is given by the largest (in magnitude) of the roots of the cubic (2.23), which we label $`\lambda _{t3,6}`$:
$$|W|=|\lambda _{t3,6}|^{1/3},qR_3.$$
(2.49)
In regions $`R_4,R_4^{}`$,
$$|W|=|\lambda _{t3,8}|^{1/3},qR_4,R_4^{}.$$
(2.50)
where $`\lambda _{t3,8}`$ is the other among these cubic roots that has maximal magnitude in this region. In $`R_5,R_5^{}`$, $`|W|=|\lambda _{t3,6}|^{1/3}`$.
## 3 $`L_y=4`$ Cyclic Triangular Strip
We have succeeded in performing an exact calculation of the chromatic polynomial for the cyclic triangular strip of arbitrarily great length and of the next larger width, namely, $`L_y=4`$. The calculation is considerably more involved than for the $`L_y=3`$ cyclic strip, as is indicated by the number of $`\lambda _j`$ terms in eq. (1.3), namely, $`N_\lambda =26`$, as contrasted with the value $`L_y=10`$ for the $`L_y=3`$ cyclic strip. Elsewhere we have proved that $`N_\lambda `$ is the same for the square and triangular strips of a given width and have given a general determination of $`N_\lambda `$ as a function of $`L_y`$ . To go beyond this width for this family of cyclic strips of the triangular lattice would be increasingly difficult, since the number of terms $`N_\lambda `$ in (1.3) is 70, 192, and 534 for $`L_y=5,6,`$ and 7, respectively. For $`L_y=4`$ we find the chromatic polynomial
$$P(tri(4\times m),FBC_y,PBC_x,q)=\underset{j=1}{\overset{26}{}}c_{t4,j}(\lambda _{t4,j})^m$$
(3.1)
where one term is $`\lambda _{t4,1}=1`$ and the others, $`\lambda _{t4,j}`$, are the roots of (i) the quartic equation (9.2.1), for $`2j5`$; (ii) the quartic equation (9.2.2), for $`6j9`$; (iii) the degree-8 equation (9.2.7), for $`10j17`$; and (iv) the degree-9 equation (9.2.16), for $`18j26`$. Since these equations are somewhat lengthy, they are given in the appendix. In contrast to the $`L_y=2`$ and $`L_y=3`$ cyclic triangular strips, here it is not possible to solve for all of the terms $`\lambda _{t4,j}`$ as algebraic roots since for $`j=10`$ through $`j=26`$ these are roots of equations of degree higher than quartic. The corresponding coefficients $`c_{t4,j}`$ are
$$c_{t4,1}=c^{(4)}$$
(3.2)
$$c_{t4,j}=c^{(0)}\mathrm{for}2j5$$
(3.3)
$$c_{t4,j}=c^{(3)}\mathrm{for}6j9$$
(3.4)
$$c_{t4,j}=c^{(2)}\mathrm{for}10j17$$
(3.5)
$$c_{t4,j}=c^{(1)}\mathrm{for}18j26.$$
(3.6)
As before, one can, equivalently, list the generating function; however, since all of the relevant information is present in (3.1) with the requisite definitions of the coefficients $`c_{t4,j}`$ and terms $`\lambda _{t4,j}`$, we shall forego this.
For the Möbius $`L_y=4`$ strip of the triangular lattice, a general argument shows that the $`\lambda _j`$’s are the same as for the cyclic strip, although the coefficients $`c_j`$ are different and are more complicated than those for the cyclic strip, as was already shown by the $`L_y=2`$ triangular strip, where two coefficients for the Möbius case are algebraic nonpolynomial functions of $`q`$. As discussed before, since $``$ is determined only by the $`\lambda _j`$’s, it is the same for the cyclic and Möbius strips of a given lattice.
Chromatic zeros for the $`L_y=4`$, $`L_x=m=16`$ (hence $`n=64`$) cyclic graph of the triangular lattice are shown in Fig. 3; the complex zeros give an approximate indication of $``$. The locus $``$ separates the $`q`$ plane into six main regions, including four that contain intervals of the real axis. The outermost one, region $`R_1`$, extends to infinite $`|q|`$ and includes the intervals $`qq_c`$ and $`q0`$ on the real $`q`$ axis. Region $`R_2`$ includes the real interval $`3qq_c`$; region $`R_3`$ includes the real interval $`2q3`$, and region $`R_4`$ includes the real interval $`0q2`$. There are also complex-conjugate regions $`R_5`$ and $`R_5^{}`$ centered at approximately $`q3.5\pm 0.6i`$. There could also be other such complex-conjugate pairs of regions. The density of zeros is observed to be smaller on the parts of $``$ extending through $`q=2`$ and $`q=3`$ and also on the right-most bubble-like curves than on the rest of the of $``$. As before with the $`L_y=3`$ strip, it is straightforward to determine the various triple points. The maximal point at which the locus $``$ crosses the real axis is
$$q_c(tri(L_y=4),FBC_y,(T)PBC_x)3.228126.$$
(3.7)
Thus, for the triangular strips with $`(FBC_y,(T)PBC_x)`$ boundary conditions for which the chromatic polynomials have been calculated so far, i.e. for the widths $`L_y=2,3`$, and 4, $`q_c`$ is a nondecreasing function of $`L_y`$: $`q_c(tri(2))=3`$, $`q_c(tri(3))=3`$, and $`q_c(tri(4))=3.228\mathrm{}`$. We have found that this is also true for the strips of the square lattice with the same cyclic or Möbius boundary conditions: $`q(sq(1))=q_c(sq(2))=2`$, $`q_c(sq(3))2.34`$, and $`q_c(sq(4))2.49`$. This contrasts with the non-monotonic behavior of $`q_c`$ as a function of $`L_y`$ that we have found for strips with $`(PBC_y,FBC_x)`$ boundary conditions (see eq. (5.16) below). It is anticipated that for cyclic and Möbius strips of the triangular lattice, as $`L_y`$ increases beyond 3, $`q_c`$ will approach the 2D value $`q_c(tri)=4`$.
Another important feature of the locus $``$ for the $`L_y=4`$ cyclic triangular strip, which is the same as was true of the $`L_y=2`$ and $`L_y=3`$ cyclic triangular strips, is that it crosses the real axis at the points $`q=0,2`$ and 3. This leads us to the conjecture
$$Conjecture:\{q=0,2,3\}\mathrm{for}tri(L_y,FBC_y,(T)PBC_x)L_y2$$
(3.8)
The crossing at $`q=0`$ is present in a very wide class of families of graphs that contain global circuits .
The crossing at $`q=2`$ that we have found for $`L_y=2,3,4`$ signals a zero-temperature critical point of the Ising antiferromagnet on these triangular strips, and similarly for the generalization (3.8) to arbitrary $`L_y`$. Note that this $`T=0`$ critical point involves frustration, since not all of the spin-spin interactions around a triangle can have their energies minimized. Frustration is also responsible for the $`T=0`$ critical point of the Ising antiferromagnet on the full 2D triangular lattice , but the nature of the critical singularities in the free energy, correlation length, etc. are different: they are algebraic for the triangular lattice, but are essential singularities for the $`L_y\times \mathrm{}`$ strips . As will be shown below, this crossing of $``$ at $`q=2`$ is also found for the $`L_y=2`$ strips with torus or Klein bottle boundary conditions, i.e. $`(PBC_y,PBC_x)`$ or $`(PBC_y,TPBC_x)`$. Thus, it appears to be a general feature of the strip graphs of the triangular lattice with periodic or twisted periodic longitudinal boundary conditions. Indeed, the respective boundaries $``$ for the $`L_y=2`$ cyclic and Möbius strips of the square, triangular, and kagomé lattices, and for the $`L_y=3`$ cyclic and Möbius strips of the square lattice and the torus and Klein bottle strips of the square lattice also cross the real $`q`$ axis at $`q=2`$ (as well as $`q=0`$).
In contrast, for strips with $`(FBC_y,FBC_x)`$ (open) or $`(PBC_y,FBC_x)`$ (cylindrical) boundary conditions, $``$ does not pass through $`q=0`$ and, while $``$ for the $`L_y=3`$ open square strip crosses the real axis at $`q=2`$, the respective loci $``$ for the $`L_y=4,5`$ open and cylindrical square strips and the open and cylindrical triangular strips with $`L_y=3,4,5`$ do not, as discussed further below.
In regions $`R_j`$, $`1j4`$, the dominant terms are, respectively, the root of maximal magnitude of (1) the first quartic equation (9.2.1), the second quartic equation (9.2.2), (3) the eighth-degree equation (9.2.7), (4) the ninth-degree equation (9.2.16). In region $`R_1`$, we label the dominant term as $`\lambda _{t4,R1}`$, and so forth for the other regions. In the complex-conjugate regions $`R_5,R_5^{}`$ the dominant term is the largest-magnitude root of the ninth-degree equation (9.2.16), which we label as $`\lambda _{t4,R5}`$. We have
$$W=(\lambda _{t4,R1})^{1/4},qR_1$$
(3.9)
$$|W|=|\lambda _{t4,Rj}|^{1/4},qR_j,j=2,3,4$$
(3.10)
and
$$|W|=|\lambda _{t4,R5}|^{1/4},qR_5,R_5^{}.$$
(3.11)
As usual, the boundaries of the regions are defined by the degeneracies, in magnitude, of the leading $`\lambda _{t4,j}`$’s in these regions. For example, on the real axis, one has $`|\lambda _{t4,R1}|=|\lambda _{t4,R2}|`$ at $`q_c`$, $`|\lambda _{t4,R2}|=|\lambda _{t4,R3}|`$ at $`q=3`$, $`|\lambda _{t4,R3}|=|\lambda _{t4,R4}|`$ at $`q=2`$, and $`|\lambda _{t4,R4}|=|\lambda _{t4,R1}|`$ at $`q=0`$.
In contrast to the situation for the $`L_y=2`$ and $`L_y=3`$ cyclic strips of the triangular lattice, here $`q_c`$ is a regular instead of a multiple (= intersection) point on the algebraic curve forming $``$, in the terminology of algebraic geometry (a multiple point on an algebraic curve is a point where several branches of the curve cross each other).
Another qualitative difference is that for $`L_y=2`$ and $`L_y=3`$, $``$ contained support only for $`Re(q)0`$, and the only point on $``$ with $`Re(q)=0`$ was the origin itself; however, for $`L_y=4`$, $``$ extends slightly into the half-plane with $`Re(q)<0`$. Although these lattice strips involve global circuits, it was found in that this is not a necessary condition for chromatic zeros and loci $``$ of lattice strips (or homeomorphic expansions thereof) to have support for $`Re(q)<0`$; this was also evident in results of , ruling out a conjecture in . In the case of strips of the square lattice we had found that for $`L_y=1`$ and 2, $``$ only had support for $`Re(q)0`$, but for $`L_y=3`$ and $`L_y=4`$ it had support also for $`Re(q)<0`$.
Finally, another general feature is that the motion of the outermost curves forming part of $``$ for $`L_y=2,3,4`$ is consistent with the inference that as $`L_y`$ increases, these form a limiting curve.
## 4 $`L_y=3`$ Triangular Lattice Strips with Torus and Klein Bottle Boundary Conditions
For the $`L_y=3`$ triangular lattice strips with torus boundary conditions (denoted with the subscript $`tt`$), we calculate, by the same methods as above,
$$P(tri(L_y=3)_m,PBC_y,PBC_x,q)=\underset{j=1}{\overset{11}{}}c_{tt,j}(\lambda _{tt,j})^m$$
(4.1)
where
$$\lambda _{tt,1}=2,c_{tt,1}=\frac{1}{3}(q1)(q^25q+3),$$
(4.2)
$$\lambda _{tt,2}=q2,c_{tt,2}=\frac{1}{2}(q1)(q2),$$
(4.3)
$$\lambda _{tt,3}=3q14,c_{tt,3}=\frac{1}{2}q(q3),$$
(4.4)
$$\lambda _{tt,4}=2(q4)^2,c_{tt,4}=q1,$$
(4.5)
$$\lambda _{tt,5}=q^39q^2+29q32,c_{tt,5}=1,$$
(4.6)
$$\lambda _{tt,(6,7)}=e^{\pm 2\pi i/3},c_{tt,6}=c_{tt,7}=\frac{1}{6}q(q1)(2q7),$$
(4.7)
$$\lambda _{tt,(8,9)}=(q^25q+7)e^{\pm 2\pi i/3},c_{tt,8}=c_{tt,9}=q1,$$
(4.8)
$$\lambda _{tt,(10,11))}=(2q7)e^{\pm 2\pi i/3},c_{tt,10}=c_{tt,11}=\frac{1}{2}(q1)(q2).$$
(4.9)
Note that although several of the coefficients are complex (and have coefficients that are not integers), the chromatic polynomial itself is, of course, a polynomial with integer coefficients. The sum of the complex terms can be written as
$$\underset{j=6}{\overset{11}{}}c_{tt,j}(\lambda _{tt,j})^m=(q1)\left[\frac{q(2q7)}{3}+2(q^25q+7)^m+(q2)[(2q7)]^m\right]\mathrm{cos}\left(\frac{2m\pi }{3}\right)$$
(4.10)
None of these terms with $`6j11`$ can be dominant anywhere in the $`q`$ plane since this would imply that the limit (1.2) would not exist. As compared with previously calculated strip graphs, this exhibits a number of qualitatively new features: (i) previously, only one $`\lambda _j`$ was a constant, independent of $`q`$, and this was always equal to either 1 or $`1`$; here there are three, and none of these is equal to $`\pm 1`$; (ii) previously, all of the $`\lambda _j`$’s were either polynomials with real coefficients or algebraic functions of polynomials with real coefficients. By general arguments given before , the dominant $`\lambda _j`$ in region $`R_1`$ has as its highest power $`q^{n/m}`$, i.e., in the present case, $`q^3`$, and has coefficient 1. We find that the generating function for this case has $`deg_x(𝒩)=9`$ and $`deg_x(𝒟)=11`$, so that all of the $`\lambda _j`$’s in $`𝒟`$ contribute to $`P`$.
Given that $`\chi =4`$ for $`m=1`$ or 2 mod 3 from eq. (1.5), it follows that for these $`m`$ values, and $`q=3`$, $`P(tri(L_y=3)_m,PBC_y,PBC_x,q=3)=0`$; however, $`\chi =3`$ for $`m=0`$ mod 3; in this case, we find (for $`m3`$)
$$P(tri(L_y=3)_m,PBC_y,PBC_x,q=3)=6\mathrm{for}m=0mod3.$$
(4.11)
This is analogous to eq. (2.32) above.
For the $`L_y=3`$ triangular lattice strips with Klein bottle topology (denoted with the subscript $`tk`$) we calculate
$`P(tri(L_y=3)_m,PBC_y,TPBC_x,q)={\displaystyle \underset{j=1}{\overset{5}{}}}c_{tk,j}(\lambda _{tk,j})^m`$ (4.12)
(4.13)
$`=(q1)(2)^m{\displaystyle \frac{1}{2}}(q1)(q2)(q2)^m+{\displaystyle \frac{1}{2}}q(q3)(3q14)^m+`$ (4.14)
(4.15)
$`(q1)[2(q4)^2]^m+(q^39q^2+29q32)^m.`$ (4.16)
Thus, $`c_{tk,j}=c_{tt,j}`$, $`j=1,\mathrm{..5}`$, and the six $`\lambda _{tt,j}`$, $`j=6,\mathrm{..11}`$ that involve complex factors do not contribute to $`P(tri(L_y=3)_m,PBC_y,TPBC_x,q)`$. Since none of these six $`\lambda _{tt,j}`$’s is dominant anywhere, it follows that $``$ is the same for the $`L_y=3`$ triangular lattice strips with torus and Klein bottle topology, just as was true of the analogous square lattice strips , and in accord with the general discussion of Ref. . For this case the generating function has $`deg_x(𝒩)=2`$ and $`deg_x(𝒟)=5`$, so that, as was true of the $`L_y=3`$ Möbius and torus strips of the triangular lattice, all of the $`\lambda _j`$’s in $`𝒟`$ contribute to $`P`$.
For the sum of the coefficients we find, for the triangular $`L_y=3`$ strip with torus boundary conditions
$$C=P(K_3,q)=q(q1)(q2)\mathrm{for}G=tri(L_y=3)\mathrm{with}(PBC_y,PBC_x)$$
(4.17)
and for this strip with Klein bottle boundary conditions
$$C=0\mathrm{for}G=tri(L_y=3)\mathrm{with}(PBC_y,TPBC_x).$$
(4.18)
These are the same respective values as those calculated in for the corresponding $`L_y=3`$ strips of the square lattice with torus and Klein bottle boundary conditions.
Chromatic zeros for the $`L_y=3`$, $`L_x=m=20`$ torus graph of the triangular lattice are shown in Fig. 4; as before, with this value of $`m`$, the chromatic zeros for the torus and Klein bottle boundary conditions are generally similar and both sets of chromatic zeros lie close to the boundary $``$ and indicate its position. The locus $``$ and the $`W`$ functions are the same for the torus and Klein bottle graph families. We find
$$q_c3.7240756\mathrm{}\mathrm{for}\{G\}=tri(L_y=3)\mathrm{with}(PBC_y,PBC_x)\mathrm{or}(PBC_y,TPBC_x).$$
(4.19)
This value of $`q_c`$ is given by the real root of the degeneracy equation $`q^39q^2+32q46=0`$. It is interesting that this is just 7 % below the value for the infinite 2D triangular lattice, $`q_c(tri)=4`$. Thus, for the width $`L_y=3`$ strips where a comparison can be made, the use of both periodic transverse boundary conditions and periodic or twisted periodic longitudinal boundary conditions expedites the approach to the 2D thermodynamic limit, in the sense that it yields a value of $`q_c`$ that is substantially closer to the 2D value than was obtained for the same width strip with free transverse boundary conditions in eq. (2.46). This is expected since for the torus or Klein bottle boundary conditions (a) the resultant graphs are $`\mathrm{\Delta }`$-regular<sup>8</sup><sup>8</sup>8A $`\mathrm{\Delta }`$-regular graph is a graph all of whose vertices have the same degree, $`\mathrm{\Delta }`$, where the degree of a vertex is defined as the number of edges connected to it. with the degree (coordination number) $`\mathrm{\Delta }=6`$ of the 2D triangular lattice, and (b) there are no boundaries to the surface on which the graphs are embedded. In contrast, for the boundary conditions of types (i)-(iv) (aside from degenerate cases) the graphs are not $`\mathrm{\Delta }`$-regular and they do have boundaries.
The locus $``$ has support for $`Re(q)0`$ and separates the $`q`$ plane into three regions. The outermost one, region $`R_1`$, extends to infinite $`|q|`$ and includes the intervals $`qq_c`$ and $`q0`$ on the real $`q`$ axis. Region $`R_2`$ includes the real interval $`2qq_c`$ and extends upward and downward to the complex conjugate triple points on $``$ at $`q_t,q_t^{}4.0\pm 1.7i`$. Region $`R_3`$ is the innermost one and includes the real interval $`0q2`$. The maximum value of $`|Im(q)|`$ on the boundary between $`R_1`$ and $`R_3`$ is about 3.0, occurring at $`Re(q)2.6`$. The maximum value of $`|q|`$ on $``$ is roughly 4.6, occurring at $`q3.9\pm 2.5i`$. The boundary between $`R_2`$ and $`R_3`$ curves to the right as one increases $`|Im(q)|`$, extending from $`q=2`$ upward to $`q_t`$ and downward to $`q_t^{}`$. As is evident in Fig. 4, the density of chromatic zeros is highest on the $`R_1R_2`$ boundary near $`q_c`$. The feature that $``$ has support only for $`Re(q)0`$ is the same as was found for the $`L_y=3`$ strips of the square lattice with torus and Klein bottle topology.
In region $`R_1`$, $`\lambda _{tt,5}=\lambda _{tk,5}`$ is the dominant $`\lambda _j`$, so
$$W=(q^39q^2+29q32)^{1/3},qR_1.$$
(4.20)
The fact that this is the same as $`W`$ for the (PBC<sub>y</sub>,FBC<sub>x</sub>) case, eq. (5.3), is a general result . The importance of the PBC<sub>y</sub> is evident from the fact that for the same width of three squares, the strip with (FBC<sub>y</sub>,FBC<sub>x</sub>) yields a different $`W`$ .
In region $`R_2`$ $`\lambda _{tt,3}=\lambda _{tk,3}`$ is dominant, so
$$|W|=|3q14|^{1/3},qR_2$$
(4.21)
(in regions other than $`R_1`$, only $`|W|`$ can be determined unambiguously ). In region $`R_3`$, $`\lambda _{tt,4}=\lambda _{tk,4}`$ is dominant, so
$$|W|=|2(q4)^2|^{1/3},qR_3.$$
(4.22)
At all of the three points, $`q=0,2`$, and $`q=q_c3.72`$ where $``$ crosses the real $`q`$ axis, it does so vertically. The present results are in accord with the inference that for a recursive graph with regular lattice structure, a sufficient condition for $``$ to separate the $`q`$ plane into two or more regions is that it contains a global circuit, i.e. a path along a lattice direction whose length goes to infinity as $`n\mathrm{}`$; here this is equivalent to PBC<sub>x</sub>.
Our calculations of the zero-temperature Potts antiferromagnet partition functions (chromatic polynomials) and exponential of the entropy, $`W`$, for $`L_y=3`$ strips of the triangular lattice with various boundary conditions including free and periodic or reversed-orientation periodic thus elucidate the role that these boundary conditions and the associated topologies play. One particular feature of note is that the torus and Klein bottle graphs have interestingly different chromatic polynomials, with different $`N_\lambda `$, but the $`W`$ functions and hence the boundaries $``$ are the same.
## 5 Width $`L_y=5,6`$ Strips of the Triangular Lattice with $`(PBC_y,FBC_x)`$
We first briefly review the cases $`L_y=3,4`$ . The chromatic polynomial for the $`L_y=3`$ triangular $`(t)`$ strip with $`n=3(m+2)`$ vertices (following the labelling conventions in ) and $`(PBC_y,FBC_x)`$ boundary conditions, denoted $`t3PF`$ in the subscripts, has $`N_\lambda =1`$,
$$\lambda _{t3PF}=q^39q^2+29q32,$$
(5.1)
$$P(tri(3\times m,PBC_y,FBC_x),q)=q(q1)(q2)(\lambda _{t3PF})^{m+1},$$
(5.2)
and
$$W(tri(L_y=3),PBC_y,FBC_x,q)=(q^39q^2+29q32)^{1/3}$$
(5.3)
with $`=\mathrm{}`$.
For the $`L_y=4`$ strip with $`(PBC_y,FBC_x)`$ boundary conditions and $`n=4(m+2)`$ vertices, one had $`N_\lambda =2`$ and
$$P(tri(4\times m,PBC_y,FBC_x),q)=\underset{j=1}{\overset{2}{}}c_{t4PF,j}(\lambda _{t4PF,j})^{m+1}$$
(5.4)
where
$$\lambda _{t4PF,j}=\frac{(q3)}{2}\left[T_{t4PF}\pm (q4)\sqrt{R_{t4PF}}\right]$$
(5.5)
where
$$T_{t4PF}=q^39q^2+33q48$$
(5.6)
$$R_{t4PF}=q^410q^3+43q^2106q+129.$$
(5.7)
In this case, as is evident in Fig. 4(a) of ), $``$ includes both arcs and a closed oval invariant under complex conjugation which crossed the real axis where the prefactor of the square root vanished, i.e. at $`q_c=4`$, and at the real root of $`T_{t4PF}`$, namely $`q3.481406`$. There are two regions, $`R_1`$ and $`R_2`$, the exterior and interior of the oval. We have
$$W(tri,4\times \mathrm{},PBC_y,q)=(\lambda _{t4PF,1})^{1/4},qR_1$$
(5.8)
and
$$W(tri,4\times \mathrm{},PBC_y,q)=(\lambda _{t4PF,2})^{1/4},qR_2$$
(5.9)
with appropriate choices of branch cuts.
We have calculated new results for the width $`L_y=5`$ and $`L_y=6`$ strips with $`(PBC_y,FBC_x)`$ boundary conditions and $`n=L_y(m+2)`$ vertices. For $`L_y=5`$ we find $`N_\lambda =2`$ and
$$P(tri(5\times m,PBC_y,FBC_x),q)=\underset{j=1}{\overset{2}{}}c_{t5PF,j}(\lambda _{t5PF,j})^{m+1}$$
(5.10)
where
$$\lambda _{t5PF,j}=\frac{1}{2}\left[T_{t5PF}\pm S_{t5PF}\sqrt{R_{t5PF}}\right]$$
(5.11)
where
$$T_{t5PF}=q^515q^4+98q^3355q^2+711q614$$
(5.12)
$$R_{t5PF}=q^48q^3+26q^260q+85$$
(5.13)
and
$$S_{t5PF}=q^311q^2+43q58.$$
(5.14)
As before, the coefficients $`c_{t5PF,j}`$ can be computed using eq. (2.9) in terms of the generating function. This generating function is given in the appendix.
Just as was true for the $`L_y=4`$ triangular lattice strip with $`(PBC_y,FBC_x)`$, in the present $`L_y=5`$ case $``$ includes both a pair of complex conjugate arcs and an oval that crosses the real axis at two points. Analogously to the previous case, one of these points is the real zero of $`T_{t5PF}`$ at
$$q_{\mathrm{}}3.207224$$
(5.15)
and the other is the real zero of the prefactor $`S_{t5PF}`$ in front of the square root, at
$$q=3.284775=q_c\mathrm{for}\{G\}=tri(5\times \mathrm{},PBC_y,FBC_x).$$
(5.16)
Taking the approximate center of the oval as the average of the left and right crossings, $`q_{center}=(q_{\mathrm{}}+q_c)`$, one finds that
$$q_{center}3.2460$$
(5.17)
near to the Tutte-Beraha number $`B_73.246980..`$, where
$$B_r=2+2\mathrm{cos}\left(\frac{2\pi }{r}\right).$$
(5.18)
Chromatic zeros near $`B_7`$ were also noted in for these types of strips. Note that $`q_c`$ decreases from 4 to the above value in eq. (5.16) as $`L_y`$ increases from 4 to 5. Hence, in contrast to our calculations of strips of the square and triangular lattices with $`(FBC_y,FBC_x)`$, $`(FBC_y,(T)PBC_x)`$, and $`(PBC_y,(T)PBC_x)`$, where in all cases considered, the respective $`q_c`$’s were nondecreasing functions of $`L_y`$, here we find that this is not the case for boundary conditions of the type $`(PBC_y,FBC_x)`$.
Chromatic zeros and $``$ for this strip are shown in Fig. 5; $``$ separates the $`q`$ plane into two regions, $`R_1`$, the exterior, and $`R_2`$, the interior, of the oval shown in the figure. Comparing the ovals for the $`L_y=4`$ strip (Fig. 4(a) of ) and the $`L_y=5`$ strip (Fig. 5 here), one sees that the oval shrinks in size. We have
$$W(tri,5\times \mathrm{},PBC_y,q)=(\lambda _{t5PF,1})^{1/5},qR_1$$
(5.19)
and
$$W(tri,5\times \mathrm{},PBC_y,q)=(\lambda _{t5PF,2})^{1/5},qR_2$$
(5.20)
with appropriate choices of branch cuts.
For the strip of with $`L_y=6`$ and $`(PBC_y,FBC_x)`$ we find $`N_\lambda =5`$ and
$$P(tri(6\times m,PBC_y,FBC_x),q)=\underset{j=1}{\overset{5}{}}c_{t6PF,j}(\lambda _{t6PF,j})^{m+1}$$
(5.21)
Here the $`\lambda _{t6PF}`$’s are the roots of a fifth-order equation, so we cannot solve for them in terms of algebraic expressions as was possible for previous strips of this type with $`L_y=3,4,5`$. We give the generating function for this strip in the appendix.
Chromatic zeros for this strip are shown in Fig. 6 for $`L_x=m+2=16`$, i.e., $`n=96`$; this length is sufficiently great that the chromatic zeros give an approximate indication of the location of the locus $``$. The locus thereby inferred appears to consist of a single connected set of curves and crosses the real axis. From our exact analytic expressions, we calculate the this crossing point to be
$$q_c=3.252419\mathrm{for}\{G\}=tri(6\times \mathrm{},PBC_y,FBC_x).$$
(5.22)
The morphology of chromatic zeros for this long $`6\times 16`$ cylindrical strip is similar to that found for a $`8\times 8`$ patch with cylindrical boundary conditions in . In both cases, one can figuratively think of how the inferred loci $``$ can be modified to yield the locus for the $`L_y=\mathrm{}`$ strip with $`(PBC_y,FBC_x)`$ cylindrical boundary conditions: one (i) pulls the left-hand endpoints further to the left (into the $`Re(q)<0`$ half-plane) and around so that they meet in a cusp at the origin, $`q=0`$; (ii) pulls the right-hand endpoints of the prongs over and around so that they meet in a cusp at $`q=4`$; (iii) shifts the crossing curve slightly to the right so that it crosses the real axis at $`q3.82`$; and (iv) makes minor shifts of the rest of the locus $``$ so as to obtain $``$ for the cylindrical strip with $`L_y=\mathrm{}`$. In this respect, just as was discussed in , one sees that the periodic transverse boundary conditions help to minimize finite-size effects. Clearly, for a fixed $`L_y`$, with $`L_x\mathrm{}`$, finite-size effects are reduced the most for $`(PBC_y,(T)PBC_x)`$, next-most for $`(PBC_y,FBC_x)`$ and $`(FBC_y,(T)PBC_x)`$, and least for $`(FBC_y,FBC_x)`$. We find
$$W(tri,6\times \mathrm{},PBC_y,q)=(\lambda _{t6PF,max})^{1/6}\mathrm{for}qR_1$$
(5.23)
where $`\lambda _{t6PF,max}`$ denotes the root with the maximal magnitude.
It is of interest to use these results to study the approach of $`W`$ to the limit for the full infinite 2D triangular lattice, extending the work of . In Table 1 we list the various values of $`W(tri(L_y\times \mathrm{},PBC_y,FBC_x),q)`$, denoted as $`W(L_y,q)`$ to save space, together with the corresponding values of $`W`$ for the full 2D triangular lattice, $`W(tri,q)`$, obtained via numerical evaluation of an infinite product representation in as checked by series expansions and rigorous bounds , and the ratio $`R_W(tri(L_y\times \mathrm{}),PBC_y,FBC_x,q)`$, where
$$R_W(tri(L_y\times \mathrm{}),BC_y,BC_x,q)=\frac{W(tri(L_y\times \mathrm{}),BC_y,BC_x,q)}{W(tri,q)}.$$
(5.24)
This extends the previous study in . Recall that for $`q4`$, $`W(tri(L_y\times \mathrm{}),BC_y,BC_x,q)`$ is independent of $`BC_x`$ . Evidently, the approach of $`W`$ for the infinite-length finite-width triangular strips to the 2D thermodynamic limit as $`L_y`$ increases is quite rapid; for moderate values of $`q`$, say 6 or 7, the ratio $`R_W(tri(5\times \mathrm{}),PBC_y,FBC_x,q)`$ is equal to 1 to within approximately $`10^3`$ or better, and the ratio $`R_W(tri(6\times \mathrm{}),PBC_y,FBC_x,q)`$ is at least as close to 1. As was proved in , the approach is non-monotonic for $`PBC_y`$ (and monotonic for $`FBC_y`$).
## 6 Width $`L_y=5`$ Strips of the Triangular Lattices with $`(FBC_y,FBC_x)`$
Previous studies have been published on strips of various lattices with $`(FBC_y,FBC_x)`$ boundary conditions . In the case of the the square and triangular lattice, these went up to $`L_y=4`$ and involved $`\lambda _j`$’s that were roots of a cubic and quartic equation, respectively. Although one can calculate chromatic polynomials for wider strips, the analysis is more cumbersome if the equations defining the $`\lambda _j`$’s are higher than quartic, so that one cannot solve for these $`\lambda _j`$’s as analytic closed-form algebraic expressions. A study of these on wider open strips is in . As an illustration of this sort of situation, here we present a calculation of chromatic polynomials for the strip of the triangular lattice with width $`L_y=5`$ and $`(FBC_y,FBC_x)`$. For this width, the $`\lambda _j`$’s are solutions of a degree-9 equation and hence cannot be solved for as algebraic roots. The method of calculation is again the iterated use of the deletion-contraction theorem. In , a given strip $`(G_s)_m`$ was constructed by $`m`$ successive additions of a subgraph $`H`$ to an endgraph $`I`$; here, $`I=H`$, so that, following the notation of , the total length of the strip graph $`(G_s)_m`$ is $`m+2`$ vertices, or equivalently, $`m+1`$ edges in the longitudinal direction. The results are conveniently expressed in terms of the coefficient functions in the generating function, as discussed above. For the width $`L_y=5`$ open strip of the triangular lattice we find $`deg_z(𝒟)=N_\lambda =9`$. The coefficient functions $`b_{t5FF,j}`$ and $`A_{t5FF,j}`$ (cf. eqs. (2.4) and (2.3)) are listed in the Appendix, where $`FF`$ is short for $`(FBC_y,FBC_x)`$.
In Fig. 7 we show a plot of chromatic zeros for the open strip of the triangular lattice with $`L_y=5`$ and length $`m+2=16`$ vertices, so that the strip has $`n=80`$ vertices in all. From an analysis of the degeneracy of leading $`\lambda _j`$’s, we find the exact result
$$q_c=3\mathrm{for}sq(5\times \mathrm{},FBC_y,FBC_x).$$
(6.1)
This is in agreement with the chromatic zeros for the finite strip shown in Fig. 7. Again, we can compare Fig. 7 with the corresponding plots for $`L_y=2`$ and $`L_y=3`$ (Fig. 5(a,b) of ), and, just as was true for the corresponding three open strips of the square lattice, this comparison shows that as $`L_y`$ increases, the arcs forming $``$ elongate and the arc endpoints nearest to the origin approach more closely to the origin. We recall that for the $`L_y=4`$ open triangular strip, no $`q_c`$ is defined since $``$ does not cross the real axis.
For $`q>q_c`$, we have, for the physical ground state degeneracy per site of the $`q`$-state Potts antiferromagnet,
$$W(t(5\times \mathrm{},FBC_x,FBC_y),q)=(\lambda _{tri,j,max})^{1/5}$$
(6.2)
where $`\lambda _{sq,j,max}`$ denotes the solution of eq. (2.8) with the coefficients (9.4.1)-(9.4.29) that has the maximal magnitude in region $`R_1`$.
As before, it is of interest to use this result to study the approach of $`W`$ to the limit for the full infinite 2D triangular lattice, extending the work of for this set of boundary conditions. In Table 2 we list the various values of $`W(tri(L_y\times \mathrm{},FBC_y,BC_x),q)`$, denoted as $`W(L_y,q)`$ to save space, together with the corresponding values of $`W`$ for the full 2D triangular lattice, $`W(tri,q)`$ and the ratio $`R_W(tri(L_y\times \mathrm{}),FBC_y,BC_x,q)=W(tri(L_y\times \mathrm{}),FBC_y,BC_x,q)/W(tri,q)`$. Recall that the value of $`W`$ for $`q4`$ is independent of $`BC_x`$ . Again, the approach of $`W`$ for the infinite-length finite-width triangular strips to the 2D thermodynamic limit as $`L_y`$ increases is quite rapid.
## 7 Comparative Discussion
In this section we give a general discussion of the locus $``$. We have found several interesting features:
1. For the strips of the triangular lattice of width $`L_y=3`$ and $`L_y=4`$ with $`(FBC_y,(T)PBC_x)`$ and strips of width $`L_y=3`$ and $`(PBC_y,(T)PBC_x)`$ boundary conditions studied here, we have shown that the locus $``$ encloses regions of the $`q`$ plane including certain intervals on the real axis and passes through $`q=0`$ and 2 as well as other possible points, depending on the family. This extends the previous study of the $`L_y=2`$ strip of the triangular lattice with $`(FBC_y,(T)PBC_x)`$ . While the strips with $`(FBC_y,(T)PBC_x)`$ have a locus $``$ that passes through $`q=3`$, this is not the case with at least the $`L_y=3`$ strip with $`(PBC_y,(T)PBC_x)`$. For the strips with $`(FBC_y,(T)PBC_x)`$ considered here we have shown that $`q_c`$ is a nondecreasing function of $`L_y`$. As one increases $`L_y`$, the $`L_y=4`$ cyclic strip is the first one for which one can no longer solve for all of the $`\lambda _j`$’s as algebraic roots, since some of the equations involved are higher than quartic.
2. The crossing of $``$ at the point $`q=2`$ for the (infinite-length limit of) strips with global circuits nicely signals the property that the Ising antiferromagnet has a frustrated zero-temperature critical point on these strips. This has been discussed in in the context of exact solutions for finite-temperature Potts model partition functions on the $`L_y=2`$ cyclic and Möbius triangular strips. In contrast, this connection is not, in general, present for strips with free longitudinal boundary conditions since $``$ does not pass through $`q=2`$ (as recapitulated below). The $`q=3`$ Potts antiferromagnet also has a zero-temperature critical point on these strips, and this is similarly manifested by the crossing of the singular locus $``$ through the point $`q=3`$ for the cyclic and Möbius strips with $`(FBC_y,(T)BC_x)`$, but not for the $`L_y=3`$ torus and Klein bottle strips with $`(PBC_y,(T)PBC_x)`$.
3. An interesting feature of the cyclic strips of the triangular lattice (as well as the square and kagomé lattices) is that for all of cases that have been studied here and in , there is a correlation between the coefficient $`c_{G_s,j}`$ of the respective dominant $`\lambda _{G_s,j}`$’s in regions that include intervals of the real axis. Before, it was shown that the $`c_{G_s,j}`$ of the dominant $`\lambda _{G_s,j}`$ in region $`R_1`$ including the real intervals $`q>q_c(\{G\})`$ and $`q<0`$ is $`c^{(0)}=1`$, where the $`c^{(d)}`$ were given in eqs. (2.10), (2.11). Here we extend this, observing that the $`c_{G_s,j}`$ that multiplies the dominant $`\lambda _{G_s,j}`$ in the region containing the intervals $`0<q<2`$ is $`c^{(1)}`$. For the cyclic strips of the triangular lattice that we have studied, namely, $`L_y=2,3,4`$, the $`c_{G_s,j}`$ multiplying the dominant $`\lambda _{G_s,j}`$ in the region containing the interval $`2<q<3`$ is $`c^{(2)}`$. For the $`L_y=4`$ strip of the triangular lattice, there is another region containing the real interval $`3<q<q_c`$, where $`q_c`$ for this strip was given in (3.7), and we find that $`c^{(3)}`$ multiplies the dominant $`\lambda _{t4,j}`$ containing this interval. For the cyclic strips of the square lattice, although the values of $`q_c`$ are different, a similar correlation is observed; in particular, for the respective widths $`L_y=3,4`$, $`c^{(2)}`$ multiplies the $`\lambda _{G_s,j}`$ that is dominant in the region containing the interval $`2<q<q_c`$, where $`q_c2.34`$ and $`q_c2.49`$ for these two widths .
4. For the strip of the triangular lattice of width $`L_y=5`$ with cylindrical $`(PBC_y,FBC_x)`$ boundary conditions, we find that $``$ consists of arcs together with a (closed) oval. However, $``$ does not pass through $`q=0,2`$, or 3. This is qualitatively the same morphology that was found for the corresponding strip of width $`L_y=4`$ (see Fig. 4 of ). For comparison, the triangular strip with width $`L_y=3`$ and $`(PBC_y,FBC_x)`$ had $`=\mathrm{}`$. The point $`q_c=4`$ for $`L_y=4`$ and $`q_c3.28`$ for $`L_y=5`$, which shows that for this family of strips, in cases where there is a $`q_c`$ (there is none for $`L_y=3`$), it is not, in general, a nondecreasing function of $`L_y`$. This is somewhat reminiscent of the non-monotonicity of $`W`$ that we showed in the case of $`(PBC_y,BC_x)`$ strips, where $`BC_x=FBC_x`$ or $`(T)PBC_x`$, in . For the $`L_y=6`$ strip with $`(PBC_y,FBC_x)`$, we infer from the chromatic zeros that $``$ consists of a single component, and from the analytic results we compute that $`q_c3.25`$, which again shows the nonmonotonicity of $`q_c`$ as a function of $`L_y`$. The morphology of chromatic zeros for our $`6\times 16`$ strip is similar to that found in for $`8\times 8`$ patch, both with cylindrical boundary conditions. For strips with torus or Klein bottle boundary conditions, $`W`$ and $``$ have been calculated for only one $`L_y`$ value, namely, $`L_y=3`$ for the square lattice in and for the triangular lattice here. For the $`L_y=3`$ $`(PBC_y,(T)PBC_x)`$ square strip, $`q_c`$ is equal to the value 3 for the full 2D square lattice, but for the $`L_y=3`$ $`(PBC_y,(T)PBC_x)`$ triangular strip, our result (4.19) shows that $`q_c`$ is less than the value of 4 for the 2D triangular lattice.
5. We have included an illustrative result for a wider strip with $`(FBC_y,FBC_x)`$ boundary conditions, namely the strip with $`L_y=5`$. As one increases $`L_y`$, this is the first value at which one can no longer solve analytically for the $`\lambda _j`$’s as algebraic roots. The locus $``$ is similar to the respective loci that were found earlier in studies of open strips in that it does not pass through $`q=0`$ or $`q=2`$ and does not separate the $`q`$ plane into regions containing intervals of the real axis. (It does pass through $`q=3`$, unlike the loci for the open strips with $`L_y=3,4`$.) Our results confirm the trend that was observed in earlier work , namely that as $`L_y`$ increases, the arcs elongate and move closer together, and the arc endpoints nearest to the origin move toward this point. In contrast to the strips with global circuits, these strips do not manifest the property that the $`q=2`$ (Ising) and $`q=3`$ Potts antiferromagnets have zero-temperature critical points since $``$ does not pass through these respective points. The simplest example of this is the Ising antiferromagnet on the infinite open line; in this case, although the model has a well-known zero-temperature critical point, this is not evident in the locus $``$, which is the emptyset.
6. For the strips with $`(PBC_y,FBC_x)`$ and $`(FBC_y,FBC_x)`$, while the endpoints of the arcs on $``$ that lie closest to the origin tend to move toward the origin as $`L_y`$ increases, leading one to expect that in the limit $`L_y\mathrm{}`$, the limiting locus $``$ would pass through $`q=0`$, no such motion toward $`q=2`$ is observed in the cases so far calculated. This is in agreement with the fact that the locus found in for the triangular lattice constructed as the limit $`L_x,L_y\mathrm{}`$ with $`(PBC_y,FBC_x)`$ boundary conditions passes through $`q=0`$ and 4 (and at $`q3.82`$) but not through $`q=2`$ or $`q=3`$. Thus, assuming that our conjecture (3.8) is correct, it follows that the locus $``$ for the triangular lattice depends on the boundary conditions used to define this lattice: if one constructs as the limit $`L_y\mathrm{}`$ with $`(PBC_y,FBC_x)`$ (cylindrical) boundary conditions, then $``$ does not pass through $`q=2`$ or 3 , while if one constructs it as the limit $`L_y\mathrm{}`$ with $`(FBC_y,(T)PBC_x)`$ (cyclic or Möbius) boundary conditions, then, if the conjecture is valid, $``$ would pass through $`q=2`$ and 3 in the limit just as it does for each of the values $`L_y=2,3,4`$ studied so far. However, since the value of $`q_c`$ pertains directly to a physical quantity, as the minimal value of real $`q`$ above which $`W(q)`$ is analytic, one expects that $`q_c`$ should be independent of the boundary conditions used to define the 2D lattice. All results obtained so far are consistent with this expectation.
7. There have been a number of theorems proved concerning real chromatic zeros. An elementary result is that no chromatic zeros can lie on the negative real axis $`q<0`$, since a chromatic polynomial has alternating coefficients. It has also been proved that there are no chromatic zeros in the intervals $`0<q<1`$, and $`1<q<32/27`$ . The bound of 32/27 in has been shown to be sharp; i.e., for any $`ϵ>0`$, there exists a graph with a chromatic zero at $`q=32/27+ϵ`$ . See also and references therein. Based on our studies of strips of the triangular (and square) lattices with the various boundary conditions considered, we make the following observation: for such strips, we have not found any chromatic zeros, except for the zero at $`q=1`$, in the interior of the disk $`|q1|=1`$. This motivates the conjecture that for these strips, there are no chromatic zeros with $`|q1|<1`$ except for the zero at $`q=1`$. Assuming that this conjecture is valid, the bound would be a sharp bound, since the circuit graph with $`n`$ vertices, $`C_n`$, has chromatic zeros lying precisely on the circle $`|q1|=1`$ and at $`q=1`$ . Further work is needed to prove (or disprove) this conjecture.
Some relevant features are summarized in Table 3. The entries for $`L_y=\mathrm{}`$ with $`(PBC_y,FBC_x)`$ are from .
## 8 Conclusions
In this paper we have presented exact calculations of the zero-temperature $`q`$-state Potts antiferromagnet partition functions (equivalently, chromatic polynomials $`P`$), on strips of the triangular lattice of width $`L_y=3`$ and with boundary conditions of four types: (a) $`(FBC_y,PBC_x)=`$ cyclic, (b) $`(FBC_y,TPBC_x)=`$ Möbius, (c) $`(PBC_y,PBC_x)=`$ toroidal, and (d) $`(PBC_y,TPBC_x)=`$ Klein bottle, where $`F`$, $`P`$, and $`TP`$ denote free, periodic, and twisted periodic. In the infinite-length limits of these strips, exact results were given for the ground state degeneracy (exponential of the ground state entropy), $`W`$, and its analytic structure in the complex $`q`$ plane, in particular, the nonanalytic locus $``$, was discussed. Exact calculations of $`P`$ and $`W`$ and studies of $``$ were also presented for wider strips, including (e) cyclic, $`L_y=4`$, (f) $`(PBC_y,FBC_x)=`$ cylindrical, $`L_y=5,6`$, and an illustrative $`(FBC_y,FBC_x)=`$ open case with $`L_y=5`$. A comparative analysis of these results was included. An interesting result of our calculations of $`W`$ on infinite-length strips with different widths and transverse boundary conditions is the observation that for the cases studied, $``$ passes through $`q=2`$ (as well as $`q=0`$) for strips with periodic or twisted periodic longitudinal boundary conditions but does not for strips with free longitudinal boundary conditions. Hence, in particular, if one uses periodic or twisted periodic longitudinal boundary conditions, the locus $``$ nicely signals the existence of the zero-temperature critical point of the Ising antiferromagnet on these infinite-length, finite-width strips of the triangular lattice. Numerical values of $`W`$ were given for infinite-length strips of various widths and were shown to approach values for the 2D lattice rapidly. Some conjectures for the behavior of $``$ for arbitrarily wide strips, and for a region in the $`q`$ plane free of chromatic zeros, were stated. These exact calculations of the $`T=0`$ Potts antiferromagnet partition function and ground state degeneracy on strips of the triangular lattice give valuable analytic knowledge of properties of Potts antiferromagnets.
Note added: The original version of this paper, submitted in early Oct. 1999, contained calculations on $`L_y=3`$ strips of the triangular lattice. In response to a request by a referee to perform calculations for wider strips, we have added the results on the $`L_y=4`$ strip with $`(FBC_y,PBC_x)`$, the $`L_y=5,6`$ strips with $`(PBC_y,FBC_x)`$, and the $`L_y=5`$ strip with $`(FBC_y,FBC_x)`$.
Acknowledgment: The research of R. S. was supported in part by the U. S. NSF grant PHY-97-22101.
## 9 Appendix
### 9.1 Generating Function for $`L_y=3`$ Möbius Strip of the Triangular Lattice
As noted in the text, it is convenient to leave the coefficients $`c_{t3Mb,j}`$, $`j=6,7,8`$ in the general form (2.43). For the evaluation of these coefficients, we list here the generating function for this strip. We have $`d_𝒩=8`$ and $`d_𝒟=10`$, and
$$𝒟(tri(L_y=3),FBC_y,TPBC_x,q,x)=\underset{j=1}{\overset{10}{}}(1\lambda _{t3,j}(q)x)$$
(9.1.1)
where the $`\lambda _{t3,j}`$ were given in eq. (LABEL:ptcyclic). Since several of the $`\lambda _{t3,j}`$’s are algebraic, it is useful to display the denominator in an explicitly polynomial form:
$`𝒟(tri(L_y=3),FBC_y,TPBC_x,q,x)=(1+x)[1(q2)x]\times `$ (9.1.2)
(9.1.3)
$`\left[1(2q7)x+(q2)(q3)x^2\right]\left[1+(q2)(q3)x\right]F_{t3q3}F_{t3q2}`$ (9.1.4)
where
$$F_{t3q3}=\underset{j=6,7,8}{}(1\lambda _{t3,j}x)=1+b_{t3,1}x+b_{t3,2}x^2+b_{t3,3}x^3$$
(9.1.5)
with $`b_{t3,j}`$, $`j=1,2,3`$ given in eqs. (2.24)-(2.26), and
$$F_{t3q2}=\underset{j=9,10}{}(1\lambda _{t,j}x)=1(q^37q^2+18q17)x+(q2)^3(q3)x^2.$$
(9.1.6)
For the numerator, extracting a common factor via the definition of the reduced coefficients $`\overline{A}_{t3Mb,j}`$,
$$A_{t3Mb,j}q(q1)(q2)(q3)\overline{A}_{t3Mb,j},$$
(9.1.7)
we have
$$\overline{A}_{t3Mb,0}=q^26q+10$$
(9.1.8)
$$\overline{A}_{t3Mb,1}=(q3)(q^312q^2+45q55)$$
(9.1.9)
$$\overline{A}_{t3Mb,2}=q^620q^5+170q^4779q^3+2016q^22779q+1588$$
(9.1.10)
$$\overline{A}_{t3Mb,3}=(q3)(5q^682q^5+574q^42185q^3+4745q^25536q+2691)$$
(9.1.11)
$$\overline{A}_{t3Mb,4}=(q2)(q3)(9q^6141q^5+930q^43303q^3+6651q^27176q+3224)$$
(9.1.12)
$$\overline{A}_{t3Mb,5}=(q2)^2(q3)^3(6q^455q^3+186q^2277q+152)$$
(9.1.13)
$$\overline{A}_{t3Mb,6}=(q2)^5(q3)(q^412q^3+59q^2138q+125)$$
(9.1.14)
$$\overline{A}_{t3Mb,7}=(q2)^6(q3)^3(3q^214q+17)$$
(9.1.15)
$$\overline{A}_{t3Mb,8}=(q2)^9(q3)^3.$$
(9.1.16)
It should also be noted that there is a significant difference between the strips of the triangular lattice studied here and the analogous strips of the square lattice . In general, a strip of the square lattice of width $`L_y`$ and length $`L_x`$ with any of the boundary conditions $`(FBC_y,PBC_x)`$ (cyclic), $`(FBC_y,TPBC_x)`$ (Möbius), $`(PBC_y,PBC_x)`$ (torus), or $`(PBC_y,TPBC_x)`$ (Klein bottle) is invariant under a translation by one edge or vertex in the longitudinal direction. However, in the case of the analogous strip of the triangular lattice, this is only true of the cases with cyclic and torus boundary conditions; the strips with Möbius and Klein bottle boundary conditions have a “seam” along with the orientation of the triangles reverses. This is discussed further in the appendix. Thus, if one proceeds in a longitudinal direction along the triangular-lattice Möbius strip, starting in a manner such that the triangles are formed by edges connecting the upper left and lower right vertices of squares (relative to one’s direction of motion), then when one crosses this seam, the triangles will be formed by edges connecting the upper right and lower left vertices of the squares on the strip. Related to this, there are differences in the degrees of various vertices on these strips (where the degree of a vertex is defined as the number of edges that connect to this vertex). If one avoids the lowest few values of $`L_x`$ where the strips degenerate, then, in general, (i) for the cyclic triangular strip of width $`L_y`$ and length $`L_x`$, the $`(L_y2)L_x`$ internal vertices have degree 6 while the $`2L_x`$ vertices on the upper and lower sides have degree 4; (ii) for the same strip as in (i) but with Möbius instead of cyclic longitudinal boundary conditions, the $`(L_y2)L_x`$ internal vertices have degree 6, the $`2(L_y1)`$ vertices on the upper and lower sides except for those on the seam have degree 4, and, on the seam, the external vertices have degrees 5 and 3; (iii) for the same strip as in (i) but with torus or Klein bottle boundary conditions, all of the vertices have degree 6.
### 9.2 Equations for the Terms in the Chromatic Polynomial for the Cyclic $`L_y=4`$ Strip of the Triangular Lattice
Four of the $`\lambda _{t4,j}`$, which we label as $`j=2,3,4,5`$, are the same as for the open $`L_y=4`$ strip of the triangular lattice. These are solutions to the quartic equation
$$\xi ^4+b_{t4,1,1}\xi ^3+b_{t4,1,2}\xi ^2+b_{t4,1,3}\xi +b_{t4,1,4}=0$$
(9.2.1)
where the coefficients were given as $`b_{t(4),j}b_{t4,1,j}`$, $`j=1,..,4`$ in eqs. (B.15)-(B.18) of
The terms $`\lambda _{t4,j}`$, $`j=6,7,8,9`$ are solutions to the quartic equation
$$\xi ^4+b_{t4,2,1}\xi ^3+b_{t4,2,2}\xi ^2+b_{t4,2,3}\xi +b_{t4,2,4}=0$$
(9.2.2)
where
$$b_{t4,2,1}=4q13$$
(9.2.3)
$$b_{t4,2,2}=2(3q^218q+26)$$
(9.2.4)
$$b_{t4,2,3}=(q2)(4q^225q+38)$$
(9.2.5)
and
$$b_{t4,2,4}=(q2)^2(q3)^2.$$
(9.2.6)
Another set of $`\lambda _j`$’s for $`10j17`$ are roots of an equation of degree 8,
$$\xi ^8+\underset{k=1}{\overset{8}{}}b_{t4,3,k}\xi ^{8k}=0$$
(9.2.7)
where
$$b_{t4,3,1}=2(3q^2+19q31)$$
(9.2.8)
$$b_{t4,3,2}=15q^4186q^3+867q^21794q+1385$$
(9.2.9)
$$b_{t4,3,3}=20q^6+366q^52784q^4+11248q^325425q^2+30452q15080$$
(9.2.10)
$$b_{t4,3,4}=(q2)(15q^7334q^6+3174q^516676q^4+52294q^397852q^2+101138q44528)$$
(9.2.11)
$$b_{t4,3,5}=(q2)^2(q3)^2(6q^6126q^5+1076q^44804q^3+11861q^215378q+8185)$$
(9.2.12)
$$b_{t4,3,6}=(q2)^4(q3)^3(q^525q^4+216q^3868q^2+1670q1246)$$
(9.2.13)
$$b_{t4,3,7}=(q2)^6(q3)^4(2q^322q^2+80q97)$$
(9.2.14)
$$b_{t4,3,8}=(q2)^8(q3)^6.$$
(9.2.15)
A final set of $`\lambda _{t4,j}`$, $`18j26`$, are solutions to the equation of degree 9
$$\xi ^9+\underset{k=1}{\overset{9}{}}b_{t4,4,k}\xi ^{9k}=0$$
(9.2.16)
where
$$b_{t4,4,1}=2(q3)(2q^212q+21)$$
(9.2.17)
$$b_{t4,4,2}=(q3)(6q^590q^4+558q^31772q^2+2865q1875)$$
(9.2.18)
$`b_{t4,4,3}=4q^9111q^8+1380q^710071q^6+47476q^5149742q^4+315652q^3`$ (9.2.19)
(9.2.20)
$`428385q^2+339300q119368`$ (9.2.21)
$`b_{t4,4,4}=(q2)(q3)^2(q^934q^8+491q^74032q^6+20961q^571954q^4`$ (9.2.22)
(9.2.23)
$`+163654q^3238278q^2+201722q75672)`$ (9.2.24)
$`b_{t4,4,5}=(q2)^2(q3)^3(3q^981q^8+983q^77029q^6+32609q^5101701q^4`$ (9.2.25)
(9.2.26)
$`+213036q^3288702q^2+229385q81299)`$ (9.2.27)
$`b_{t4,4,6}=(q2)^4(q3)^4(3q^870q^7+728q^64404q^5+16929q^442286q^3`$ (9.2.28)
(9.2.29)
$`+66933q^261296q+24830)`$ (9.2.30)
$$b_{t4,4,7}=(q2)^6(q3)^6(q^620q^5+166q^4734q^3+1833q^22462q+1393)$$
(9.2.31)
$$b_{t4,4,8}=(q2)^8(q3)^7(2q^421q^3+87q^2165q+119)$$
(9.2.32)
$$b_{t4,4,9}=(q2)^{12}(q3)^8.$$
(9.2.33)
### 9.3 Generating Function for the $`L_y=5,6`$ Strips of the Triangular Lattice with $`(PBC_y,FBC_x)`$
For the $`L_y=5`$ strip we calculate a generating function of the form (2.2) with $`d_𝒟=2`$, $`d_𝒩=1`$ and, in the notation of eqs. (2.3) and (2.4),
$$b_{t5PF,1}=q^5+15q^498q^3+355q^2711q+614$$
(9.3.1)
$$b_{t5PF,2}=(q3)(3q^766q^6+619q^53205q^4+9877q^318065q^2+18078q7588)$$
(9.3.2)
$$A_{t5PF,0}=q(q1)(q2)(q3)(q^614q^5+85q^4290q^3+599q^2723q+398)$$
(9.3.3)
$`A_{t5PF,1}=q(q1)(q2)(q3)(q^22q+2)(3q^766q^6+619q^53205q^4`$ (9.3.4)
(9.3.5)
$`+9877q^318065q^2+18078q7588).`$ (9.3.6)
For the $`L_y=6`$ strip we calculate a generating function of the form (2.2) with $`d_𝒟=5`$, $`d_𝒩=4`$ and
$$b_{t6PF,1}=q^6+18q^5145q^4+680q^31980q^2+3379q2586$$
(9.3.7)
$`b_{t6PF,2}=4q^{10}128q^9+1868q^816352q^7+94977q^6382031q^5+1076317q^4`$ (9.3.8)
(9.3.9)
$`2093899q^3+2686606q^22047842q+702080`$ (9.3.10)
$`b_{t6PF,3}=2(q3)^2(q^{12}39q^{11}+695q^{10}7493q^9+54509q^8282283q^7+1068575q^6`$ (9.3.11)
(9.3.12)
$`2982861q^5+6098756q^48908956q^3+8820488q^25306146q+1462992)`$ (9.3.13)
$`b_{t6PF,4}=2(q2)(q3)^5(2q^{11}62q^{10}+882q^97601q^8+44105q^7181018q^6`$ (9.3.14)
(9.3.15)
$`+536536q^51149015q^4+1742334q^31779827q^2+1099188q309188)`$ (9.3.16)
$`b_{t6PF,5}=4(q2)^2(q3)^8(q^25q+5)(2q^736q^6+277q^51179q^4+2990q^3`$ (9.3.17)
(9.3.18)
$`4505q^2+3728q1310)`$ (9.3.19)
$`A_{t6PF,0}=q(q1)(q2)(q^921q^8+199q^71121q^6+4159q^510623q^4`$ (9.3.20)
(9.3.21)
$`+18887q^322824q^2+17177q6143)`$ (9.3.22)
$`A_{t6PF,1}=q(q1)(q2)(4q^{13}140q^{12}+2266q^{11}22416q^{10}+150973q^9`$ (9.3.23)
(9.3.24)
$`730186q^8+2607252q^76958852q^6+13899608q^520584349q^4`$ (9.3.25)
(9.3.26)
$`+22103679q^316461349q^2+7723994q1748140)`$ (9.3.27)
$`A_{t6PF,2}=2q(q1)(q2)(q3)^2(q^{15}42q^{14}+816q^{13}9734q^{12}+79793q^{11}`$ (9.3.28)
(9.3.29)
$`476549q^{10}+2144264q^97409966q^8+19852299q^741297346q^6+66301130q^5`$ (9.3.30)
(9.3.31)
$`80939629q^4+73099740q^346448750q^2+18742947q3655548)`$ (9.3.32)
$`A_{t6PF,3}=2q(q1)(q2)(q3)^5(2q^{15}72q^{14}+1212q^{13}12651q^{12}+91556q^{11}`$ (9.3.33)
(9.3.34)
$`486599q^{10}+1962326q^96116898q^8+14870220q^728223745q^6+41554711q^5`$ (9.3.35)
(9.3.36)
$`46735427q^4+39026466q^322975509q^2+8586616q1545752)`$ (9.3.37)
$`A_{t6PF,4}=4q(q1)(q2)^2(q3)^8(q^25q+5)(q^45q^3+10q^210q+5)`$ (9.3.38)
(9.3.39)
$`\times (2q^736q^6+277q^51179q^4+2990q^34505q^2+3728q1310)`$ (9.3.40)
### 9.4 Generating Function for $`L_y=5`$ Open Strip of the Triangular Lattice
For this strip we calculate a generating function of the form (2.2) with $`d_𝒟=9`$, $`d_𝒩=8`$ and, in the notation of eqs. (2.3) and (2.4),
$$b_{t5FF,1}=(q3)(q^410q^3+46q^2113q+120)$$
(9.4.1)
$`b_{t5FF,2}=6q^8131q^7+1280q^67328q^5+26930q^465081q^3+100888q^291462q+36965`$ (9.4.2)
(9.4.3)
(9.4.4)
$`b_{t5FF,3}=(q2)(q3)(15q^9378q^8+4289q^728788q^6+126096q^5374139q^4`$ (9.4.5)
(9.4.6)
$`+752541q^3989867q^2+772611q272483)`$ (9.4.7)
$`b_{t5FF,4}=(q2)(q3)^2(20q^{11}607q^{10}+8429q^970702q^8+398115q^7`$ (9.4.8)
(9.4.9)
$`1580547q^6+4515585q^59285872q^4+13471537q^313131321q^2`$ (9.4.10)
(9.4.11)
$`+7738560q2087938)`$ (9.4.12)
$`b_{t5FF,5}=(q2)^2(q3)^3(15q^{12}502q^{11}+7729q^{10}72397q^9`$ (9.4.13)
(9.4.14)
$`+459566q^82083176q^7+6915864q^616947196q^5+30430188q^4`$ (9.4.15)
(9.4.16)
$`39053679q^3+34008163q^218041392q+4408580)`$ (9.4.17)
$`b_{t5FF,6}=(q2)^3(q3)^5(6q^{12}203q^{11}+3149q^{10}29630q^9+188440q^8853749q^7`$ (9.4.18)
(9.4.19)
$`+2826657q^66893466q^5+12293272q^415636918q^3+13466958q^2`$ (9.4.20)
(9.4.21)
$`7049578q+1695556)`$ (9.4.22)
$`b_{t5FF,7}=(q2)^5(q3)^7(q^{11}34q^{10}+516q^94647q^8+27718q^7115308q^6`$ (9.4.23)
(9.4.24)
$`+342008q^5724072q^4+1072964q^31060043q^2+628196q169014)`$ (9.4.25)
$`b_{t5FF,8}=(q2)^7(q3)^9(q^923q^8+241q^71505q^6+6145q^516929q^4`$ (9.4.26)
(9.4.27)
$`+31319q^337359q^2+25972q7987)`$ (9.4.28)
$`b_{t5FF,9}=(q2)^{12}(q3)^{11}(q^38q^2+21q17)`$ (9.4.29)
We use the definition $`A_{t5FF,j}=q(q1)(q2)\overline{A}_{t5FF,j}`$ and have
$$\overline{A}_{t5FF,0}=(q2)^7$$
(9.4.30)
$`\overline{A}_{t5FF,1}=(q2)^2(6q^8113q^7+926q^64304q^5+12381q^422504q^3`$ (9.4.31)
(9.4.32)
$`+25133q^215663q+4121)`$ (9.4.33)
$`\overline{A}_{t5FF,2}=(q2)^2(q3)(15q^{10}363q^9+3936q^825141q^7+104572q^6`$ (9.4.34)
(9.4.35)
$`295342q^5+572184q^4748554q^3+630336q^2306798q+65006)`$ (9.4.36)
$`\overline{A}_{t5FF,3}=(q2)^2(q3)^2(20q^{12}587q^{11}+7867q^{10}63587q^9+344801q^8`$ (9.4.37)
(9.4.38)
$`1319670q^7+3650087q^67338694q^5+10623070q^410770336q^3+7237158q^2`$ (9.4.39)
(9.4.40)
$`2882235q+511741)`$ (9.4.41)
$`\overline{A}_{t5FF,4}=(q2)^3(q3)^3(15q^{13}487q^{12}+7266q^{11}65902q^{10}+405008q^9`$ (9.4.42)
(9.4.43)
$`1779077q^8+5739974q^713753341q^6+24433426q^531724357q^4+29191727q^3`$ (9.4.44)
(9.4.45)
$`17973781q^2+6613085q1093258)`$ (9.4.46)
$`\overline{A}_{t5FF,5}=(q2)^4(q3)^4(6q^{14}215q^{13}+3554q^{12}35907q^{11}+247594q^{10}`$ (9.4.47)
(9.4.48)
$`1231849q^9+4556733q^812718548q^7+26883798q^642760478q^5+50287310q^4`$ (9.4.49)
(9.4.50)
$`42307379q^3+24007733q^28200007q+1267630)`$ (9.4.51)
$`\overline{A}_{t5FF,6}=(q2)^6(q3)^7(q^{12}33q^{11}+485q^{10}4225q^9+24378q^898266q^7`$ (9.4.52)
(9.4.53)
$`+283651q^6589907q^5+875195q^4900611q^3+607987q^2240848q+42203)`$ (9.4.54)
(9.4.55)
(9.4.56)
$`\overline{A}_{t5FF,7}=(q2)^{10}(q3)^9(q^818q^7+145q^6674q^5+1941q^43474q^3`$ (9.4.57)
(9.4.58)
$`+3701q^22120q+499)`$ (9.4.59)
$$\overline{A}_{t5FF,8}=(q1)^3(q2)^{11}(q3)^{11}(q^38q^2+21q17).$$
(9.4.60)
### 9.5 $`L_y=2`$ Cyclic and Möbius Strips of the Triangular Lattice with Odd $`N_t`$
In this part of the appendix we shall report some new results for $`L_y=2`$ cyclic and Möbius strips of the triangular lattice with an odd number $`N_t`$ of triangles and compare these with the case of even $`N_t=L_xL_y`$. Some illustrative strip graphs are shown in Fig. 8.
For the $`L_y=2`$ cyclic strip of the triangular lattice with even $`N_t=2L_x`$ we calculated (see also )
$$P(tri(L_y=2,cyc.),N_t=2m,q)=(q^23q+1)+[(q2)^2]^m+(q1)[(\lambda _{t2,3})^m+(\lambda _{t2,4})^m]$$
(9.5.1)
where
$$\lambda _{t2,(3,4)}=\frac{1}{2}\left[52q\pm \sqrt{94q}\right].$$
(9.5.2)
For the corresponding $`L_y=2`$ Möbius strip of the triangular lattice with even $`N_t=2L_x`$, we calculated
$$P(tri(L_y=2,Mb.),N_t=2m,q)=1+[(q2)^2]^m(q1)(q3)\frac{\left[(\lambda _{t2,3})^m(\lambda _{t2,4})^m\right]}{\lambda _{t2,3}\lambda _{t2,4}}.$$
(9.5.3)
Note that the Möbius strip has a seam. The function $`W`$ and the boundary $``$ were given in ; $``$ separates the $`q`$ plane into three regions and crosses the real axis at $`q=0,2`$, and at $`q_c=3`$.
We proceed to our new results. For the $`L_y=2`$ cyclic strip of the triangular lattice containing an odd number $`N_t=2m+1`$ of triangles, we calculate
$`P(tri(L_y=2,cyc.),N_t=2m+1,q)=(q^23q+1)+(q2)[(q2)^2]^m+`$ (9.5.4)
(9.5.5)
$`{\displaystyle \frac{1}{2}}(q1)(q3)\left[\left((\lambda _{t2,3})^m+(\lambda _{t2,4})^m\right)+{\displaystyle \frac{\left((\lambda _{t2,3})^m(\lambda _{t2,4})^m\right)}{\lambda _{t2,3}\lambda _{t2,4}}}\right].`$ (9.5.6)
This strip has a seam.
For the $`L_y=2`$ Möbius strip of the triangular lattice with an odd number $`N_t=2m+1`$ of triangles (which does not have a seam) we obtain
$`P(tri(L_y=2,Mb.),N_t=2m+1,q)=1+(q2)[(q2)^2]^m+`$ (9.5.7)
(9.5.8)
$`{\displaystyle \frac{1}{2}}(1q)\left[\left((\lambda _{t2,3})^m+(\lambda _{t2,4})^m\right)+(94q){\displaystyle \frac{\left((\lambda _{t2,3})^m(\lambda _{t2,4})^m\right)}{\lambda _{t2,3}\lambda _{t2,4}}}\right].`$ (9.5.9)
Note that the sum of the coefficients $`C=0`$ for the chromatic polynomials for both of the $`L_y=2`$ strips with odd $`N_t`$, eqs. (9.5.6) and (9.5.9). For comparison, for the $`L_y=2`$ even-$`N_t`$ strips, $`C=q(q1)`$ for the cyclic case, eq. (9.5.1) and $`C=0`$ for the Möbius case, eq. (9.5.3).
Equivalently, one may write these results in terms of generating functions, with the definition analogous to (2.1):
$$\mathrm{\Gamma }(tri(L_y=2,BC_x),N_t=even,q,x)=\underset{m=2}{\overset{\mathrm{}}{}}P(tri(L_y=2,BC_x),N_t=2m,q)x^{m2}.$$
(9.5.10)
and
$$\mathrm{\Gamma }(tri(L_y=2,BC_x),N_t=odd,q,x)=\underset{m=2}{\overset{\mathrm{}}{}}P(tri(L_y=2,BC_x),N_t=2m+1,q)x^{m2}$$
(9.5.11)
where $`BC_x=PBC_x`$ or $`TPBC_x`$. The generating functions for all of these strips have the same denominator,
$`𝒟(tri(L_y=2,cyc.))={\displaystyle \underset{j=1}{\overset{4}{}}}(1\lambda _{t2,j}x)`$ (9.5.12)
(9.5.13)
$`=(1x)\left[1(q2)^2x\right]\left[1(52q)x+(q2)^2x^2\right].`$ (9.5.14)
The numerators are easily worked out from the results that we have given for the chromatic polynomials and the denominator $`𝒟`$; for example,
$$\mathrm{\Gamma }(tri(L_y=2,cyc.),N_t=even,q,x)=\frac{q(q1)(q2)\left[q3+(q2)x(q2)^3x^2\right]}{𝒟(tri(L_y=2,cyc.))}$$
(9.5.15)
$$\mathrm{\Gamma }(tri(L_y=2,cyc.),N_t=odd,q,x)=\frac{q(q1)(q2)(q3)\left[q3+(q2)^2x\right]}{𝒟(tri(L_y=2,cyc.))}.$$
(9.5.16)
and so forth for the Möbius strips.
We remark on some general features of these results. The chromatic polynomials for both even and odd $`N_t`$ and both cyclic and Möbius strips have the same four $`\lambda _j`$’s, and hence the same $`W`$ functions and boundary $``$ (given in ). For all cases, $`q_c=3`$. These properties are in agreement with the general discussion in on the effects of boundary conditions on $`W`$ and $``$.
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# Discovery of a very luminous megamaser during a radio flare in the Seyfert 2 galaxy Mrk 348
## 1 Introduction
Emission from H<sub>2</sub>O masers has been found in a few galaxies, exhibiting apparent isotropic luminosities a million times higher than in typical stellar masers Dos Santos & Lepine (1979); Gardner & Whiteoak (1982); Claussen et al. (1984); Henkel et al. (1984); Haschick & Baan (1985); Braatz et al. (1994). The detection rate of these megamasers is very low, i.e. about 5% among Seyfert galaxies Braatz et al. (1997) and almost zero among radio galaxies (e.g., Henkel et al. (1998)). The maser is associated with dense and warm material, possibly a molecular torus or disk, around an active galactic nucleus (AGN). The AGN apparently produces the seed radio photons and the X-ray photons needed to pump the masing material Neufeld et al. (1994).
With the help of Very Long Baseline Interferometry (VLBI) megamasers can be used to investigate the small-scale structure of an AGN in great detail. In the case of NGC 4258 this has helped to establish the presence of a thin, warped disk around the nucleus, to determine the black hole mass, and even to measure the precise distance to this galaxy Miyoshi et al. (1995); Herrnstein et al. (1999).
Finding new megamaser galaxies is therefore of prime interest. The only clear trend that has emerged in recent years is that megamasers are exclusively found in type 2 AGN, i.e. those Seyferts and LINERs which are expected to be obscured by a molecular torus according to the unified scheme Antonucci (1993). Many have high absorbing column depths inferred from X-ray spectroscopy. There is also an indication of an excess of megamasers in highly inclined galaxies Braatz et al. (1997); Falcke et al. (2000). Here we report the discovery of a hitherto undetected and very luminous megamaser in the Seyfert galaxy Mrk 348 during a radio flare of the AGN.
Mrk 348 (NGC 262, $`z=0.01503`$ Huchra et al. (1999), luminosity distance $`D=62.5`$ Mpc for $`z`$ converted into the Galactic Standard of Rest and $`H_0=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup>), is a well-studied Seyfert 2 galaxy with broad emission-lines in polarized light Miller & Goodrich (1990). The galaxy is classified as an S0 with a rather low inclination ($`i=16^{}`$, see Braatz et al. (1997)). Ground-based Simpson et al. (1996) and Hubble Space Telescope imaging Falcke et al. (1998) have revealed a dust lane crossing the nucleus and an excitation cone in emission-lines. Ginga observations found hard X-ray emission and a high absorbing column depth of $`N_\mathrm{H}=10^{23.1}\mathrm{cm}^2`$ towards the nucleus Warwick et al. (1989). All this suggests the presence of an obscuring torus in Mrk 348. Attempts to detect the obscuring material through radio spectroscopy have failed so far (e.g. H I, Gallimore et al. (1999)).
What makes this galaxy stand out among Seyfert galaxies is its bright and variable radio nucleus. ?) found a compact radio core on VLBI scales with a flux density of several hundred milli-Jansky and a flat to inverted spectrum. ?) presented more recent VLBI observations for two epochs, showing a two-component structure expanding with sub-relativistic speeds. They also noted a flare of the radio continuum emission at 15 GHz with the flux rising from 120 mJy to 570 mJy between 1997.10 and 1998.75. In the following we will present and discuss results of K-band radio spectroscopy of this galaxy.
## 2 Observations and Results
Data were taken in March and April 2000, using the Effelsberg 100 m telescope of the MPIfR equipped with a dual channel K-band HEMT receiver. The system temperature was of order 70 K on a $`T_\mathrm{A}^{}`$ temperature scale; the beam size was 40<sup>′′</sup>. The data were recorded using an autocorrelator with 8 $`\times `$ 256 channels and bandwidths of 80 MHz each. The eight backends were configured in two groups of four, sampling the two orthogonal linear polarizations. Frequency shifts between the four backends representing a given linear polarization were adjusted in such a way that a total velocity range of 3000 km s<sup>-1</sup> could be covered.
The measurements were carried out in a dual beam switching mode (switching frequency 1 Hz) with a beam throw of 121<sup>′′</sup> in azimuth. Only linear baselines were subtracted. Calibration was obtained by measurements of W3OH (3.2 Jy according to Mauersberger et al. (1988)). Pointing could be checked on Mrk 348 itself and was found to be stable to within 5–7<sup>′′</sup>.
The galaxy was initially observed to look for ammonia (NH<sub>3</sub>) and cyclopropenylidene (C<sub>3</sub>H<sub>2</sub>) absorption against the bright radio nucleus. No absorption features were found. The transitions, rest frequencies, and upper limits of these observations are listed in Table 1.
Since the H<sub>2</sub>O ($`6_{16}5_{23}`$) transition is very close we also observed the redshifted 22.23508 GHz line and detected an emission feature on March 17. We repeated the observations on five subsequent days and detected the line each time. The six spectra are shown in Fig. 1. We also produced a combined spectrum of all days (Fig. 2, top) and fitted the H<sub>2</sub>O line with a Gaussian profile. A one component fit yields a central velocity of $`v=4641.6\pm 1.2`$ km s<sup>-1</sup> for the line, which is redshifted by 133 km s<sup>-1</sup> from the systemic velocity. The peak flux is $`S_\nu =34`$ mJy with a full width at half maximum of $`\mathrm{\Delta }v=130\pm 3`$ km s<sup>-1</sup>. The integrated flux is 4.71$`\pm 0.1`$ Jy km s<sup>-1</sup>, yielding an apparent isotropic luminosity of 420 $`L_{}`$. A two component fit yields a narrow line at $`v=4677.7\pm 0.8`$ km s<sup>-1</sup> with $`\mathrm{\Delta }v=58.6\pm 3.2`$ km s<sup>-1</sup> and $`S_\nu =27`$ mJy plus a broad line with $`v=4609.5\pm 1.9`$ km s<sup>-1</sup>, $`\mathrm{\Delta }v=101.8\pm 1.9`$ km s<sup>-1</sup>, and $`S_\nu =26`$ mJy. The latter fit indicates that the line is asymmetric and has a pronounced blue wing. If we compare our two best spectra from 17 March and 5 April we tentatively find some variability in the shape of this blue wing. We searched for additional high velocity features and did not detect anything down to a limit of 6 mJy (1$`\sigma `$, 4 km s<sup>-1</sup>) in the range 3250 to 6200 km s<sup>-1</sup> and down to a limit of 10-15 mJy in the range $`1470`$ to 10480 km s<sup>-1</sup>. The continuum flux density we find at 22 GHz is $`0.8\pm 0.1`$ Jy, corresponding to $`410^{23}`$ Watt Hz<sup>-1</sup>.
We also reduced some archival data of earlier observations of Mrk 348 taken between October 1997 and February 1998 and combined them into one data set (Fig. 2, bottom). In this spectrum the broad line is marginally detected with a flux density approximately three times lower than in the current observations. The smaller and variable baseline in these earlier observations made a reliable identification of this line impossible without a priori information.
## 3 Discussion and Summary
We have clearly discovered a new megamaser in Mrk 348. Its luminosity is comparable to the emission from NGC 3079 which contains the second most luminous H<sub>2</sub>O maser after TXS 2226-184 Koekemoer et al. (1995). The line width is among the broadest found for a megamaser, similar to the masers in TXS 2226-184 and NGC 1052. Assuming the emission is associated with material close to the center this is the first spectroscopic evidence for molecular gas possibly obscuring the nucleus.
Despite the bright radio continuum and the molecular maser line we have found no absorption lines from either NH<sub>3</sub> or C<sub>3</sub>H<sub>2</sub>. This is in line with the earlier non-detection of H I absorption Gallimore et al. (1999). C<sub>3</sub>H<sub>2</sub> is an organic ring molecule which is widespread in the Galaxy and is associated with diffuse gas in the ISM Matthews & Irvine (1985). It was also found in absorption against the nucleus of the radio galaxy Centaurus A Seaquist & Bell (1986). From our non-detection we find a 1$`\sigma `$ upper limit for the optical depth $`\tau `$ times the covering factor $`f`$ of $`\tau f<0.0075`$. This is seven times smaller than the value found for Centaurus A and might be related to the fact that we see Mrk 348 almost face on.
On the other hand, with its type 2 AGN, polarized broad emission-lines, and a nuclear dust lane seen, Mrk 348 falls right into the roster of typical megamaser galaxies, where it is suggested that the masing material is part of the obscuring ‘torus’ in the unified scheme of AGN. The face-on orientation of Mrk 348 then would suggest that the axis of this torus and the galaxy disk axis are severely misaligned.
In October 1997 – February 1998 the flux density of the line was three times lower than it was in March 2000. In the earlier survey by ?) the maser line was not detected and given its line width and low flux density it would have resulted only in a broad feature at the $`3\sigma `$ level. Linear interpolation with time of the continuum flux density given in ?) suggests a continuum flux density for Mrk 348 around 310 mJy in October 1997 at 15 GHz. The level we measure is roughly a factor of three higher – assuming a flat spectrum – and this increase in flux density is similar to the increase in flux density of the line. This could indicate a correlation between continuum and maser flux density, implying an unsaturated maser. With its current radio luminosity the galaxy is now the most radio luminous megamaser galaxy ever discovered. The response of the line to the continuum flare within about 2 years sets an upper limit to the distance of the masers from the nucleus of $``$0.6 pc, which is of similar order as the size scale of the molecular disk found in NGC 4258.
Indeed, ?) noted a certain excess of detected megamaser galaxies with large radio powers. Mrk 348 certainly adds to this trend. The distribution of radio powers at 6 cm of the parent samples of AGN selected by ?) has a peak around $`10^{21.75}`$ Watt/Hz. If we add Mrk 348 and more recently discovered megamasers Greenhill et al. (1997); Hagiwara et al. (1997) and also complement the radio data in ?) with more recent data from the NASA Extragalactic Database, we find that 12 out of 19 detected megamaser galaxies are at or above the peak in the distribution of radio power for galaxies without megamaser detections.
Because of the possibly biased selection of the detected megamasers this is not highly significant. However, it highlights an apparently necessary prerequisite for megamaser emission, namely an AGN with a compact radio core to provide seed photons. So far all detected megamasers have compact ($`<1\mathrm{}`$), mostly flat-spectrum, radio emission at a level of a few milli-Jansky. In some cases, like Mrk 348 Neff & de Bruyn (1983), NGC 1052 Shaffer & Marscher (1979), Mrk 1210 Heisler et al. (1998), NGC 2639 Wilson et al. (1998), NGC 3079 Trotter et al. (1998), NGC 5793 Whiteoak & Gardner (1987); Gardner et al. (1992), NGC 5506 Sadler et al. (1995), and possibly NGC 4945 Elmouttie et al. (1997) the radio cores can even reach several tens to hundreds of milli-Jansky.
While compact radio cores in Seyfert and LINER galaxies are not uncommon, only very few are so prominent as those in some of the radio-bright megamaser sources. We find that all megamasers mentioned above, i.e. more than a third of known megamasers, have compact radio cores above a fiducial limit of 25 mJy at 5 GHz. On the other hand, in a survey of spiral galaxies ?) find only 3 out of 54 galaxies (22 of which are Seyfert galaxies) with compact cores above 25 mJy at 5 GHz. Similarly, in a survey of nearby AGN ?) find roughly 40% of Seyfert and LINER galaxies to contain compact flat-spectrum radio cores. However, only three out of 48 galaxies have flux densities $`>`$25 mJy. Known megamasers therefore seem to prefer galaxies with relatively bright compact radio emission.
Mrk 348 currently has the brightest and most prominent radio core among megamaser galaxies. Morphology, spectrum, and variability of the core are very similar to the radio core in III Zw 2 which was the first Seyfert galaxy discovered to contain a superluminal jet. This galaxy has a millimeter-peaked spectrum and a jet which shows a stop-and-go behavior indicative of a strong interaction with dense material on the sub-parsec scale Falcke et al. (1999); Brunthaler et al. (2000). ?) therefore speculate whether the bright inverted radio core in Mrk 348 could be interpreted similarly to those in Compact Symmetric Objects (CSOs) with a Gigahertz-Peaked-Spectrum (GPS, see O’Dea (1998)). In these galaxies bright hotspots are formed in a jet that terminates already on the parsec scale. In III Zw 2 and Mrk 348 this seems to happen on even smaller scales, leading to higher peak frequencies and could be due to frustration of the jet by a molecular cloud or even a warped or misaligned torus.
Since the masers in NGC 1052, which have similar broad line widths as in Mrk 348, are found along the radio jet Claussen et al. (1998) it should be checked whether in Mrk 348 one has an analogous situation. One can speculate that in such a case the evolution of the radio flare and the evolution of the maser flare and its blue wing could be related, possibly providing a unique diagnostic tool to study jet-ISM interactions.
In any case, with its bright radio core Mrk 348 provides an ideal opportunity to observe the maser lines in this galaxy at high resolution with VLBI during this flare even though the lines still have a rather low flux. Since radio and maser emission seem to be highly variable both should be monitored frequently. Given that Mrk 348 was not discovered in an earlier survey this finding also suggests that existing samples should be revisited to search for more flaring megamasers.
###### Acknowledgements.
We thank Alan Roy for helpful discussions. We are grateful to Jim Ulvestad for a prompt referee report and useful suggestions. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by JPL, Caltech, under contract with NASA.
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# Hadron Structure Functions within a Chiral Quark ModelTalk prepared for the QNP 2000 conference in Adelaide Feb. 2000 and to appear in the proceedings. Presentation prevented by United Airlines. This work is supported in parts by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreements #DF–FC02–94ER40818 and #DE-FG0398ER41066 and by the Deutsche Forschungsgemeinschaft (DFG) under contract We 1254/3-1.
## 1 THE CHIRAL MODEL
The bosonized action of chiral quark models can be cast in the form
$`𝒜[S,P]=iN_C\mathrm{Tr}_\mathrm{\Lambda }\mathrm{log}[i/(S+i\gamma _5P)]{\displaystyle \frac{1}{4G}}{\displaystyle }d^4x\mathrm{tr}𝒱(S,P).`$ (1)
Here $`𝒱`$ is a local potential respectively for scalar and pseudoscalar fields $`S`$ and $`P`$ which are matrices in flavor space. For example, in the Nambu–Jona–Lasinio (NJL) model one has $`𝒱=S^2+P^2+2\widehat{m}_0(S+iP)`$. From the gap–equation we obtain the VEV, $`S=m`$ which parameterizes the dynamical chiral symmetry breaking. The regularization of the quadratically divergent quark loop is indicated by the cut–off $`\mathrm{\Lambda }`$. We adjust its value as well as the coupling constant $`G`$ and the current quark mass $`\widehat{m}_0`$ to fit the phenomenological meson parameters $`m_\pi `$ and $`f_\pi `$, leaving only a single free parameter, the constituent quark mass, $`m`$. An essential feature of these models is that the derivative term in (1) is formally identical to that of a non–interacting quark model. Hence the current operator is given as $`J^\mu =\overline{q}𝒬\gamma ^\mu q`$, with $`𝒬`$ being a flavor matrix. We compute expectation values of currents by introducing pertinent sources in the bosonized action and taking appropriate derivatives.
The major concern in regularizing the functional (1) is to maintain the chiral anomaly. We achieve this goal by splitting this functional into $`\gamma _5`$–even and odd pieces and only regularize the former. The $`\gamma _5`$–odd part turns out to be conditionally finite.
Details of this presentation are published in . For related work see refs .
## 2 REGULARIZATION OF THE COMPTON TENSOR
DIS off hadrons is parameterized by the hadronic tensor $`W^{\mu \nu }(q)`$ with $`q`$ being the momentum transmitted to the hadron by the photon. $`W^{\mu \nu }(q)`$ is obtained from the hadron matrix element of the commutator $`[J^\mu (\xi ),J^\nu (0)]`$. In bosonized quark models we find it convenient to start from the absorptive part of the forward virtual Compton amplitude
$`T^{\mu \nu }(q)={\displaystyle }d^4\xi \mathrm{e}^{iq\xi }p,s|T\left(J^\mu (\xi )J^\nu (0)\right)|p,s\mathrm{and}W^{\mu \nu }(q)={\displaystyle \frac{1}{2\pi }}\mathrm{}(T^{\mu \nu }(q).`$ (2)
(We denote the momentum of the hadron by $`p`$ and eventually its spin by $`s`$.) The advantage is that the time–ordered product is unambiguously obtained from the regularized action
$`T\left(J^\mu (\xi )J^\nu (0)\right)={\displaystyle \frac{\delta ^2}{\delta v_\mu (\xi )\delta v_\nu (0)}}\mathrm{Tr}_\mathrm{\Lambda }\mathrm{log}[i/(S+i\gamma _5P)+𝒬v/]|_{v_\mu =0}.`$ (3)
In order to extract the leading twist pieces of the structure functions, we study $`W^{\mu \nu }(q)`$ in the Bjorken limit: $`q^2\mathrm{}`$ with $`x=q^2/pq`$ fixed.
We now have to specify the regularization of the functional trace in (3). We define
$`i𝐃=i/(S+i\gamma _5P)+v/𝒬\mathrm{and}i𝐃_5=i/(Si\gamma _5P)v/𝒬`$ (4)
and separate the functional trace into (un–)regularized $`\gamma _5`$–even (odd) pieces,
$`\mathrm{Tr}_\mathrm{\Lambda }\mathrm{log}[i/(S+i\gamma _5P)+𝒬v/]`$ $`=`$ $`i{\displaystyle \frac{N_C}{2}}{\displaystyle \underset{i=0}{\overset{2}{}}}c_i\mathrm{Tr}\mathrm{log}\left[\mathrm{𝐃𝐃}_5+\mathrm{\Lambda }_i^2iϵ\right]`$ (5)
$`i{\displaystyle \frac{N_C}{2}}\mathrm{Tr}\mathrm{log}\left[𝐃\left(𝐃_5\right)^1iϵ\right].`$
With the conditions $`c_0=1,\mathrm{\Lambda }_0=0,_{i=0}^2c_i=0\mathrm{and}_{i=0}^2c_i\mathrm{\Lambda }_i^2=0`$ the double Pauli–Villars regularization renders the functional in (3) finite.
## 3 PION STRUCTURE FUNCTION
DIS off pions is characterized by a single structure function, $`F(x)`$. For its computation we have to specify the pion matrix element in the Compton amplitude (2). Whence we introduce the pion field $`\stackrel{}{\pi }`$ via<sup>1</sup><sup>1</sup>1The coupling $`g`$ and the constituent quark mass $`m`$ are related by the pion decay constant. In the chiral limit the relation is linear $`m=gf_\pi `$.
$`S+iP\gamma _5=m\left(U\right)^{\gamma _5}=m\mathrm{exp}\left(i{\displaystyle \frac{g}{m}}\gamma _5\stackrel{}{\pi }\stackrel{}{\tau }\right).`$ (6)
Expanding (5,6) to linear and quadratic order in $`\stackrel{}{\pi }`$ and $`v_\mu `$, respectively yields the proper result for the anomalous decay $`\pi ^0\gamma \gamma `$. In turn we obtain the Compton amplitude for virtual pion–photon scattering by expanding (5,6) to second order in both, $`\stackrel{}{\pi }`$ and $`v_\mu `$. Due to the separation into $`𝐃`$ and $`𝐃_5`$ this calculation differs from the evaluation of the ‘handbag’ diagram because isospin violating dimension–five operators emerge. Fortunately all isospin violating pieces cancel yielding
$`F(x)={\displaystyle \frac{5}{9}}(4N_Cg^2){\displaystyle \frac{d}{dm_\pi ^2}}\left\{m_\pi ^2{\displaystyle \underset{i=0}{\overset{2}{}}}c_i{\displaystyle \frac{d^4k}{(2\pi )^4i}}\left[k^2x(1x)m_\pi ^2+m^2+\mathrm{\Lambda }_i^2iϵ\right]^2\right\}.`$ (7)
The cancellation of the isospin violating pieces is a feature of the Bjorken limit: insertions of the pion field on the propagator carrying the infinitely large photon momentum can be safely ignored. Furthermore this propagator can be taken to be the one for non–interacting massless fermions. This implies that the Pauli–Villars cut–offs can be omitted for this propagator leading to the desired scaling behavior of the structure function.
## 4 NUCLEON STRUCTURE FUNCTIONS
In the bosonized chiral quark model baryons emerge as solitons of the meson fields . We parameterize the soliton by
$`U(\stackrel{}{x},t)=A(t)\mathrm{exp}\left(i\stackrel{}{\tau }\widehat{r}\mathrm{\Theta }(r)\right)A^{}(t),`$ (8)
with the chiral angle $`\mathrm{\Theta }(r)`$ being determined from the stationary condition for constant $`A`$. Subsequently we quantize the collective coordinates $`A`$ to generate nucleon states.
As argued above we take the quark propagator with the infinite photon momentum to be free and massless. Thus, it is sufficient to differentiate
$`{\displaystyle \frac{N_C}{4i}}{\displaystyle \underset{i=0}{\overset{2}{}}}c_i\mathrm{Tr}\left\{(𝐃^{(\pi )}𝐃_5^{(\pi )}+\mathrm{\Lambda }_i^2)^1[𝒬^2v/(/)^1v/𝐃_5^{(\pi )}𝐃^{(\pi )}(v/(/)^1v/)_5𝒬^2]\right\}`$
$`+{\displaystyle \frac{N_C}{4i}}\mathrm{Tr}\left\{(𝐃^{(\pi )}𝐃_5^{(\pi )})^1[𝒬^2v/(/)^1v/𝐃_5^{(\pi )}+𝐃^{(\pi )}(v/(/)^1v/)_5𝒬^2]\right\},`$ (9)
with respect to the photon field $`v_\mu `$ as in eq (3). We have introduced the $`(\mathrm{})_5`$ description
$`\gamma _\mu \gamma _\rho \gamma _\nu =S_{\mu \rho \nu \sigma }\gamma ^\sigma iϵ_{\mu \rho \nu \sigma }\gamma ^\sigma \gamma ^5\mathrm{and}(\gamma _\mu \gamma _\rho \gamma _\nu )_5=S_{\mu \rho \nu \sigma }\gamma ^\sigma +iϵ_{\mu \rho \nu \sigma }\gamma ^\sigma \gamma ^5`$ (10)
to account for the unconventional appearance of axial sources in $`𝐃_5`$ . Upon substituting (8) into (9) and computing the functional trace, using a basis of quark states obtained from the Dirac Hamiltonian in the background of $`U(\stackrel{}{x},t)`$, we find analytical results for the structure functions. By construction their regularization is consistent with the chiral anomaly. We refer to for detailed formulae and the explicit verification of sum rules. Here we simply report the important result that the structure function entering the Gottfried sum rule is related to the $`\gamma _5`$–odd piece of the action and hence does not undergo regularization. This is surprising because in the parton model this structure function differs from the one associated with the Adler sum rule only by the sign of the anti–quark distribution. The latter structure function, however, gets regularized, in agreement with the quantization rules for the collecive coordinates. As we have consistently implemented the regularization at the level of the bosonized action this demonstrates that in effective models these structure functions are quite different from constituent quark distributions and in their description one has to go beyond identifying degrees of freedom.
## 5 NUMERICAL RESULTS FOR THE POLARIZED NUCLEON STRUCTURE FUNCTIONS
Unfortunately numerical results for the full structure functions, i.e. including the properly regularized vacuum piece are not yet available. However, we have verified that in the Pauli–Villars regularization the axial charges are saturated to 95% or more by their valence quark contributions once the self–consistent soliton is substituted. This provides sufficient justification to adopt the valence quark contribution to the polarized structure functions as a reliable approximation . In Fig. 1 we compare the model predictions for the linearly independent polarized structure functions to experimental data .
Fig. 1: Model predictions for the polarized proton structure functions $`xg_1`$ (left panel) and $`xg_2`$ (right panel). The curves labeled ‘RF’ denote the results as obtained from the valence quark contribution to (9). These undergo a projection to the infinite momentum frame ‘IMF’ and a leading order ‘LO’ DGLAP evolution . Data are from SLAC–E143 .
The evolution of the structure function $`g_2`$ to the momentum scale of the experiments requires the separation into twist–2 and twist–3 components . We observe that the model results for the polarized structure functions, which we argued to have reliably approximated, agree reasonably well with the experimental data. This encourages future work in this direction.
## ACKNOWLEDGMENTS
We thank our collaborators E. Ruiz Arriola, O. Schröder and H. Reinhardt for valuable contributions. HW is particularly grateful to the organizers of the conference for including this contribution in the proceedings despite its presentation being prevented by United Airlines.
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# Consistent Linearized Gravity in Brane Backgrounds
## 1 Introduction and Summary
There has been much interest recently in the studies of non-compact spaces containing 3-branes as domain walls as an alternative to compactification for the treatment of the hierarchy problem . These scenarios, in fact, revive older ideas described in . In a simple, non-compact scenario, which was proposed by Randall and Sundrum (RS) , the 3-brane is flat due to fine tuning of either the brane tension or the cosmological constant. Gravity on the brane resembles the usual 4-dimensional Einstein gravity for long distances, owing to a graviton bound state, whereas the Kaluza Klein modes nearly decouple from the gravity on the brane .
One of the most interesting questions in this and similar scenarios is the coupling of gravity to matter on the brane, because it represents the very mechanism of generating gravitational forces. In fact, several papers have studied this question, and the existence of long-distance Einstein gravity on the brane has been confirmed . Many of these calculations have used the RS gauge , as a result of which the brane appears bent owing to matter located on it . Moreover, calculations in global Riemannian normal coordinates are simply not consistent .
In this paper, we would like to follow up on the issue of localization of gravity on the brane and address the following points: First, we present a globally consistent formulation of linearized gravity with matter located on the brane, which is kept *straight*, both in the RS and an alternative background. Our motivation for studying the alternative background is that the brane attracts ordinary matter in that case, whereas it is repulsive in the Randall-Sundrum case . Consistency of our calculation is made possible by adopting a novel gauge fixing (cf. Secs. 2 and 3). In our gauge, the metric components transversal to the brane are non-vanishing. Fluctuations of this kind have been studied already in . Second, we determine the physical degrees of freedom of gravity by analyzing the gauge freedom for a straight brane (Sec. 4). The problem of physical degrees of freedom was studied in . Last, we find the particular solution for a matter perturbation on the brane and analyze the effective 4-dimensional gravity (Sec. 5). We would like to mention that our solution is only gauge-equivalent to the bent-brane formulation , if we are far enough from the brane, although also the results for gravity on the brane seem to be identical. Further comments can be found in Sec. 4.
Although our results are not entirely new, in that they confirm the well-established long-distance Einstein gravity, they represent, to our mind, a concise and elegant global formulation and provide important insights into the dynamics on the brane. In particular, a non-gravitational mechanism of confinement for the RS background turns out to be essential not only because of the instability of the geodesics on the brane , but also for the Newtonian dynamics on the brane, which would be spoiled by extrinsic contributions. On the other hand, Einstein’s equation does not hold for long distances for the alternative background.
Let us now summarize our results. Our metric has the form
$$ds^2=\mathrm{e}^{2k|y|}(\eta _{\mu \nu }+\gamma _{\mu \nu })dx^\mu dx^\nu +2n_\mu dx^\mu dy+(1+\varphi )dy^2,$$
(1)
where $`\eta _{\mu \nu }=\text{diag}(1,1,1,1)`$. In eqn. (1) and henceforth, positive $`k`$, $`k>0`$, corresponds to the RS background and negative $`k`$, $`k<0`$, to the alternative background.
Using the more general ansatz (1) instead of the RS gauge for the metric we are able to keep the brane straight at $`y=0`$ and obtain a globally consistent solution for the linearized Einstein equations. For a matter perturbation $`t_{\mu \nu }`$ located on the brane, our solutions for the first order quantities in eqn. (1) are
$$n_\mu =\frac{\mathrm{sgn}y}{8k}\gamma _{,\mu },\varphi =\frac{\mathrm{sgn}y}{4k}\gamma _{,y},$$
(2)
and the traceless transversal part of $`\gamma _{\mu \nu }`$ satisfies
$$\begin{array}{c}_y\left(\mathrm{e}^{2k|y|}_y\stackrel{~}{\gamma }_{\mu \nu }\right)2k\mathrm{sgn}y\mathrm{e}^{2k|y|}_y\stackrel{~}{\gamma }_{\mu \nu }+\mathrm{}\stackrel{~}{\gamma }_{\mu \nu }\hfill \\ \hfill =16\pi \delta (y)\left[t_{\mu \nu }\frac{1}{3}\left(\eta _{\mu \nu }\frac{_\nu _\mu }{\mathrm{}}\right)t\right].\end{array}$$
(3)
Moreover, the trace $`\gamma `$ satisfies the boundary conditions
$$\gamma |_{y=0}=\frac{32\pi k}{3\mathrm{}}t,\gamma |_{y=\mathrm{}}=0,$$
but can be altered by gauge transformations for $`y0`$. The step functions in eqn. (2) correspond to an apparent singularity separating two coordinate patches. We shall demonstrate explicitely in Sec. 4 how to obtain a continuous solution.
## 2 General Method
The starting point for our calculation is the action
$$S=d^4x𝑑y\sqrt{g}(R2\mathrm{\Lambda })+_{\text{brane}}d^4x\sqrt{\widehat{g}}(\sigma +_{\text{matter}}),$$
(4)
where
$$\mathrm{\Lambda }=6k^2\text{and}\sigma =12k,$$
(5)
so that the metric
$$ds^2=\mathrm{e}^{2k|y|}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2$$
(6)
is a background solutions for $`_{\text{matter}}=0`$.
In our approach, we would like to use a global coordinate system, in which the brane is located at $`y=0`$ even when a perturbation is present on the brane. Thus, it is natural to use the time-slicing formalism , where we slice with respect to $`y=const.`$ hypersurfaces. At least in the RS case ($`k>0`$), we cannot impose an a priori gauge on the metric perturbations, since we would obtain exponentially growing solutions for large $`y`$. Instead, consistency forces us to use a very particular, yet elegant, gauge. This shall become clear in the course of our calculation. Let us now give a short review of the time-slicing formalism and briefly outline the specific character of our approach.
In the time slicing formalism, one splits up the metric tensor as
$$\left(g_{ab}\right)=\left(\begin{array}{cc}\widehat{g}_{\mu \nu }& n_\mu \\ n_\nu & n_\lambda n^\lambda +n^2\end{array}\right),\left(g^{ab}\right)=\frac{1}{n^2}\left(\begin{array}{cc}n^2\widehat{g}^{\mu \nu }+n^\mu n^\nu & n^\mu \\ n^\nu & 1\end{array}\right),$$
(7)
where $`\widehat{g}_{\mu \nu }`$ is the induced metric on the hypersurfaces, $`n^\mu =\widehat{g}^{\mu \nu }n_\nu `$, and $`a,b=0,1,2,3,5`$, $`\mu ,\nu =0,1,2,3`$. Henceforth, we shall denote quantities derived from the induced metric $`\widehat{g}_{\mu \nu }`$ with a hat. $`n`$ and $`n_\mu `$ are the lapse function and shift vector, respectively.
We consider the hypersurfaces $`y=const.`$, which have the induced metric $`\widehat{g}_{\mu \nu }`$ and the tangent and normal vectors
$$_\mu x^a=\{\begin{array}{cc}\delta _\mu ^a\hfill & \text{for }a=0,1,2,3\text{,}\hfill \\ 0\hfill & \text{for }a=5\text{,}\hfill \end{array}$$
(8)
and
$$N_a=(0,0,0,0,n),N^a=\frac{1}{n}(n^\mu ,1),$$
(9)
respectively. Then, the second fundamental form measuring the extrinsic curvature of the hypersurfaces has the form
$$H_{\mu \nu }=\frac{1}{2n}\left(_y\widehat{g}_{\mu \nu }\widehat{}_\mu n_\nu \widehat{}_\nu n_\mu \right).$$
(10)
Einstein’s equation,
$$R^{ab}\frac{1}{2}g^{ab}R=g^{ab}\mathrm{\Lambda }+8\pi T^{ab},$$
(11)
can now be rewritten in terms of $`\widehat{g}_{\mu \nu }`$, $`n_\mu `$ and $`n^2`$. First, using the Gauss-Codazzi equations, the normal and mixed components of eqn. (11) with respect to the hypersurfaces become
$`\widehat{R}+H_\nu ^\mu H_\mu ^\nu H^2`$ $`=2\mathrm{\Lambda },`$ (12)
$`_\mu H\widehat{}_\nu H_\mu ^\nu `$ $`=0,`$ (13)
where we have used the fact that $`N_aT^{ab}=0`$ for the energy momentum tensor derived from the action (4). For the tangential components of eqn. (11) we prefer the form
$$R_{\mu \nu }=\frac{2}{3}\widehat{g}_{\mu \nu }\mathrm{\Lambda }+8\pi \left(T_{\mu \nu }\frac{1}{3}\widehat{g}_{\mu \nu }T\right),$$
(14)
because it saves us from calculating the scalar curvature $`R`$. In the standard approach, one fixes $`n_\mu `$ and $`n^2`$ to some convenient value. Then, eqn. (14) is the equation of motion for $`\widehat{g}_{\mu \nu }`$, whereas eqns. (12) and (13) are constraints. Notice that until now all expressions have been exact.
In a previous paper , some of us followed the standard approach choosing globally $`n_\mu =n^21=0`$ and found that the linearized approximation was inconsistent for the Randall-Sundrum background. Therefore, we would now like to take a different approach. Instead of fixing the lapse function $`n`$ and the shift vector $`n_\mu `$ a priori, we leave them present at first and fix them in the course of our calculations by the condition that the linearization be consistent. For the linearization of the Randall-Sundrum and the alternative background, we write the induced metric as
$$\widehat{g}_{\mu \nu }=\mathrm{e}^{2k|y|}\left(\eta _{\mu \nu }+\gamma _{\mu \nu }\right),$$
(15)
and we consider $`\gamma _{\mu \nu }`$, $`n_\mu `$ and $`n^21`$ as small perturbations. Furthermore, we assume that the induced metric perturbations, $`\gamma _{\mu \nu }`$, are continuous at $`y=0`$. The necessary linearized expressions for the connections and curvatures are given in the appendix.
## 3 Linearized Equations and Gauge Choice
Let us start by linearizing Einstein’s equations. The energy momentum tensor, as found from the action (4), has the form
$$T^{\mu \nu }=\frac{3k}{4\pi }\sqrt{\frac{\widehat{g}}{g}}\delta (y)\widehat{g}^{\mu \nu }+\delta (y)t^{\mu \nu }(x),T^{5\mu }=T^{55}=0,$$
(16)
where the first term of $`T^{\mu \nu }`$ is the background from the brane, and $`t^{\mu \nu }`$ is a small matter perturbation sitting on the brane. The covariant conservation law, $`_aT^{ab}=0`$ is satisfied to first order, if and only if $`t^{\mu \nu }`$ is conserved in the conventional sense, $`_\mu t^{\mu \nu }=0`$.
The constraints, eqns. (12) and (13), take the linearized forms
$$\mathrm{e}^{2k|y|}(\gamma ^{\mu \nu }{}_{,\mu \nu }{}^{}\mathrm{}\gamma 6k\mathrm{sgn}y^\mu n_\mu )=3k\mathrm{sgn}y\gamma _{,y}12k^2\varphi $$
(17)
and
$$\mathrm{e}^{2k|y|}(\mathrm{}n_\mu _\mu ^\nu n_\nu )=3k\mathrm{sgn}y_\mu \varphi +_y\left(\gamma ^\nu {}_{\mu ,\nu }{}^{}\gamma _{,\mu }\right),$$
(18)
respectively, where we have defined
$$\varphi =n^21.$$
(19)
Next, let us linearize the tangential equation, eqn. (14), with the energy momentum tensor (16). We have $`\widehat{g}/g1\varphi `$, so that one finds the linearized form
$$\begin{array}{c}\frac{1}{2}(\gamma ^\rho {}_{\mu ,\rho \nu }{}^{}+\gamma ^\rho {}_{\nu ,\rho \mu }{}^{}\mathrm{}\gamma _{\mu \nu }\gamma _{,\mu \nu })+3k^2\varphi \mathrm{e}^{2k|y|}\eta _{\mu \nu }\frac{1}{2}_y\left(\mathrm{e}^{2k|y|}\gamma _{\mu \nu ,y}\right)\hfill \\ \hfill \frac{1}{2}_\mu _\nu \varphi +\frac{1}{2}_y(n_{\mu ,\nu }+n_{\nu ,\mu })k\mathrm{sgn}y\left(n_{\mu ,\nu }+n_{\nu ,\mu }+\eta _{\mu \nu }^\lambda n_\lambda \right)\\ \hfill +k\mathrm{sgn}y\mathrm{e}^{2k|y|}\left[\gamma _{\mu \nu ,y}+\frac{1}{2}\eta _{\mu \nu }\gamma _{,y}\right]\frac{1}{2}k\eta _{\mu \nu }_y\left(\mathrm{e}^{2k|y|}\mathrm{sgn}y\varphi \right)\\ \hfill =8\pi \delta (y)\left(t_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }t\right),\end{array}$$
(20)
where $`t=\eta ^{\mu \nu }t_{\mu \nu }`$.
In the Randall-Sundrum case, the exponentially increasing terms in eqns. (17) and (18) are potentially problematic, but problems can be avoided by choosing a suitable gauge. Thus, instead of setting $`n_\mu =\varphi =0`$, we impose the gauge conditions
$`4k\varphi `$ $`=\mathrm{sgn}y\gamma _{,y},`$ (21)
$`^\mu \stackrel{~}{\gamma }_{\mu \nu }`$ $`=0,`$ (22)
where
$`\stackrel{~}{\gamma }_{\mu \nu }`$ $`=\gamma _{\mu \nu }{\displaystyle \frac{1}{4}}\eta _{\mu \nu }\gamma `$
is the traceless part of $`\gamma _{\mu \nu }`$. Eqn. (22) says that the traceless part $`\stackrel{~}{\gamma }_{\mu \nu }`$ should also be transversal.
Together, the gauge conditions (21) and (22) imply that the right hand sides of eqns. (17) and (18) are zero, so that these equations reduce to
$`\gamma ^{\mu \nu }{}_{,\mu \nu }{}^{}\mathrm{}\gamma `$ $`=6k\mathrm{sgn}y^\mu n_\mu ,`$ (23)
$`\mathrm{}n_\mu _\mu ^\nu n_\nu `$ $`=0,`$ (24)
respectively. Eqns. (23) and (24) are equations for the shift vector $`n_\mu `$, and their general solution is
$$n_\mu =\frac{\mathrm{sgn}y}{8k}\gamma _{,\mu }+A_\mu .$$
(25)
Here, the vector $`A_\mu `$ satisfies the 4-dimensional equations of a free vector field in Lorentz gauge,
$$\mathrm{}A_\mu =^\mu A_\mu =0,$$
(26)
but it can depend also on $`y`$.
After inserting the gauge conditions (21) and (22) as well as the solution (25) for the shift vectors into eqn. (20), we find
$$\begin{array}{c}_y\left(\mathrm{e}^{2k|y|}_y\stackrel{~}{\gamma }_{\mu \nu }\right)2k\mathrm{sgn}y\mathrm{e}^{2k|y|}_y\stackrel{~}{\gamma }_{\mu \nu }+\mathrm{}\stackrel{~}{\gamma }_{\mu \nu }\hfill \\ \hfill _y(A_{\mu ,\nu }+A_{\nu ,\mu })+2k\mathrm{sgn}y(A_{\mu ,\nu }+A_{\nu ,\mu })\\ \hfill =\delta (y)\left[16\pi \left(t_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }t\right)\frac{1}{2k}\gamma _{,\mu \nu }\right].\end{array}$$
(27)
Notice that by virtue of eqns. (22) and (26) the left hand side of eqn. (27) is traceless and four-divergence-free. Thus, we should expect the same of the right hand side. The former property translates into
$$\mathrm{}\gamma |_{y=0}=\frac{32\pi k}{3}t,$$
(28)
whereas the latter property is expressed as
$$16\pi \left(^\mu t_{\mu \nu }\frac{1}{3}_\nu t\right)\frac{1}{2k}_\nu \mathrm{}\gamma |_{y=0}=0.$$
(29)
Using eqn. (28) we find $`^\mu t_{\mu \nu }=0`$, which is a good check of concistency, since Einstein’s equation should imply the covariant energy-momentum conservation law.
Looking at the equations presented so far we realize that there is no equation of motion for $`\gamma `$, and eqns. (26) and (27) are insufficient to determine $`A_\mu `$ and $`\stackrel{~}{\gamma }_{\mu \nu }`$. Therefore, we suspect that there is residual gauge freedom, the determination of which is the subject of the next section.
## 4 Physical Degrees of Freedom
We would like to discuss the residual gauge freedom left after imposing the gauge conditions (21) and (22) and keeping the brane fixed at $`y=0`$. We shall find that traceless transversal spin-2 excitations are the only physical degrees of freedem. Moreover, we shall show how to remove the step functions in the solutions for $`n_\mu `$ and $`\varphi `$ \[cf. eqns. (21) and (25)\].
To start, consider two coordinate systems with the metrics
$`ds^2`$ $`=\mathrm{e}^{2k|y|}\left(\eta _{\mu \nu }+\gamma _{\mu \nu }\right)dx^\mu dx^\nu +2n_\mu dx^\mu dy+(1+\varphi )(dy)^2`$ (30)
$`=\mathrm{e}^{2k|y^{}|}(\eta _{\mu \nu }+\gamma _{\mu \nu }^{})dx^{}{}_{}{}^{\mu }dx^{}{}_{}{}^{\nu }+2n_\mu ^{}dx^{}{}_{}{}^{\mu }dy^{}+(1+\varphi ^{})(dy^{})^2,`$
related by an infinitesimal coordinate transformation,
$$x^{}{}_{}{}^{\mu }=x^\mu \xi ^\mu (x,y),y^{}=y\xi ^5(x,y).$$
(31)
The location of the brane remains unchanged, i.e. we have the restriction $`\xi ^5(x,0)=0`$. This also ensures that no normal components of the energy momentum tensor, $`T^{a5}`$, are generated.
Under the coordinate transformation (31) the first order elements transform as
$`\gamma _{\mu \nu }`$ $`=\gamma _{\mu \nu }^{}+2k\mathrm{sgn}y\xi ^5\eta _{\mu \nu }\xi _{\mu ,\nu }\xi _{\nu ,\mu },`$ (32)
$`n_\mu `$ $`=n_\mu ^{}\xi _{,\mu }^5\mathrm{e}^{2ky}\xi _{\mu ,y},`$ (33)
$`\varphi `$ $`=\varphi ^{}2\xi _{,y}^5.`$ (34)
Here, indices have been lowered using the Minkowski metric. With some calculations one can check that the linearized Einstein equations, eqns. (17), (18) and (20), are invariant under this transformation, provided that $`\xi ^5|_{y=0}=0`$.
Coordinate transformations can be applied separately for $`y<0`$ and $`y>0`$. For simplicity, we shall restrict our discussion to $`y>0`$. For $`y>0`$, the trace of eqn. (32) yields
$$\gamma =\gamma ^{}+8k\xi ^52^\lambda \xi _\lambda ,$$
(35)
which can be combined with eqn. (32) to obtain the transformation of $`\stackrel{~}{\gamma }_{\mu \nu }`$,
$$\stackrel{~}{\gamma }_{\mu \nu }=\stackrel{~}{\gamma }_{\mu \nu }^{}\xi _{\mu ,\nu }\xi _{\nu ,\mu }+\frac{1}{2}\eta _{\mu \nu }^\lambda \xi _\lambda .$$
(36)
Thus, imposing the gauge condition (22) on both metrics yields
$$\mathrm{}\xi _\mu +\frac{1}{2}_\mu ^\nu \xi _\nu =0.$$
(37)
Next, we substitute the gauge condition (21) into eqn. (34) and obtain $`_y^\lambda \xi _\lambda =0`$, i.e.
$$^\lambda \xi _\lambda =f(x)$$
(38)
is a function of $`x`$ only. Moreover, from eqn. (37) we find that $`f`$ must satisfy $`\mathrm{}f=0`$. Finally, substituting the solution for $`n_\mu `$, eqn. (25), into eqn. (33), one finds that $`A_\mu `$ transforms as
$$A_\mu =A_\mu ^{}\frac{1}{4k}_\mu f(x)\mathrm{e}^{2ky}\xi _{\mu ,y}.$$
(39)
As a check of consistency, we observe that eqn. (26) and the boundary condition (28) remain unchanged.
Let us now fix the remaining gauge freedom and identify the physical degrees of freedom. First, we can use
$$\gamma |_{y=0}=\gamma ^{}|_{y=0}2f(x)$$
to obtain a *unique* solution for the boundary condition $`(\text{28})`$, which we shall formally write as
$$\gamma |_{y=0}=\frac{32\pi k}{3\mathrm{}}t.$$
(40)
Second, we make use of $`\xi ^5`$ in order to choose a convenient function for $`\gamma `$ in the bulk satisfying the boundary conditions (40). For example, one could pick
$$\gamma =\frac{32\pi k}{3\mathrm{}}t\mathrm{e}^{ay^2}$$
(41)
with some positive coefficient $`a`$. Notice that on the brane $`\gamma `$ is determined by the matter content, which cannot be gauged away. Last, we make use of the remaining freedom, $`\xi _\mu `$ satisfying $`^\mu \xi _\mu =\mathrm{}\xi _\mu =0`$, and
$$A_\mu =A_\mu ^{}\mathrm{e}^{2ky}\xi _{\mu ,y}$$
in order to pick a convenient function for $`A_\mu `$. Of course, the most convenient value is $`A_\mu =0`$, which we shall adopt.
After this gauge fixing, we are left with only the physical degrees of freedom $`\stackrel{~}{\gamma }_{\mu \nu }`$, which describe spin-2 gravity excitations.
We would like to point out the following subtle point regarding the bent-brane formulation used by Garriga and Tanaka , Giddings, Katz and Randall and others. If one did not impose the condition $`\xi ^5|_{y=0}=0`$ \[cf. eqn. (31)\], it would seem from eqns. (32)–(34) that one could transform a metric in our gauge into a metric in Randall-Sundrum gauge using $`\xi ^5=\gamma ^{}/8k`$. The obvious effect would be that the brane appears bent to an observer, and non-zero normal components $`T^{5\mu }`$ are generated. However, we would like to emphasize that eqn. (30) is not valid in this case, because one cannot expand $`\mathrm{e}^{2k|y^{}+\xi ^5|}\mathrm{e}^{2k|y^{}|}(12k\mathrm{sgn}y\xi ^5)`$, which would imply that the brane is located again at $`y^{}=0`$ (the brane is where the singularity of the curvature is). Rather, one should write
$$ds^2=\mathrm{e}^{2k|y^{}+\xi ^5|}(\eta _{\mu \nu }+\gamma _{\mu \nu }^{\prime \prime })dx_{}^{}{}_{}{}^{\mu }dx_{}^{}{}_{}{}^{\nu }+dy_{}^{}{}_{}{}^{2}.$$
Thus, eqn. (10) of should contain $`\delta [y\xi ^5(x)]`$ and not $`\delta (y)`$. In our opinion, this seems to be a drawback of the bent-brane formulation and has not been addressed properly.
As a last point in this section, we would like to demonstrate how to remove the step functions in the solutions for $`n_\nu `$ and $`\varphi `$ \[cf. eqn. (2)\] and to obtain a continuous solution. First let us note that it is enough to remove the discontinuity of $`n_\mu `$, because one can choose $`\gamma `$ such that $`\varphi |_{y=0}=0`$ \[cf. eqn. (41)\], i.e. $`\varphi `$ is already continuous. We take a gauge transformation of the form
$$\xi ^5=0,\xi _\mu =\frac{1}{16k^2}\gamma _{,\mu }\mathrm{e}^{2k|y|}\psi (y),$$
(42)
where $`\psi (y)`$ is a smooth function with a compact support such that $`\psi (0)=1`$. Then, the gauge-transformed components of the metric will be
$`\gamma _{\mu \nu }^{}`$ $`=\gamma _{\mu \nu }+{\displaystyle \frac{1}{8k^2}}\mathrm{e}^{2k|y|}\psi (y)\gamma _{,\mu \nu },`$ (43)
$`n_\mu ^{}`$ $`={\displaystyle \frac{\mathrm{sgn}y}{8k}}\gamma _{,\mu }+{\displaystyle \frac{1}{16k^2}}\mathrm{e}^{2k|y|}_y\left[\gamma _{,\mu }\mathrm{e}^{2k|y|}\psi (y)\right],`$ (44)
$`\varphi ^{}`$ $`={\displaystyle \frac{\mathrm{sgn}y}{4k}}\gamma _{,y}.`$ (45)
One can easily see that $`n_\mu ^{}`$ is continuous at $`y=0`$. Thus, we have found an explicit variable transformation for a vicinity of the brane, which transforms our solution into one that is continuous at $`y=0`$. We can, therefore, conclude that the step functions in eqn. (2) correspond only to an apparent singularity.
## 5 Particular Solution and Gravity on the Brane
Let us now fix the gauge as described in the last section and solve eqn. (27). Substituting $`A_\mu =0`$ and eqn. (40), eqn. (27) becomes
$$\begin{array}{c}_y\left(\mathrm{e}^{2k|y|}_y\stackrel{~}{\gamma }_{\mu \nu }\right)2k\mathrm{sgn}y\mathrm{e}^{2k|y|}_y\stackrel{~}{\gamma }_{\mu \nu }+\mathrm{}\stackrel{~}{\gamma }_{\mu \nu }\hfill \\ \hfill =16\pi \delta (y)\left[t_{\mu \nu }\frac{1}{3}\left(\eta _{\mu \nu }\frac{_\nu _\mu }{\mathrm{}}\right)t\right].\end{array}$$
(46)
We note that the inverse of the d’Alembertian is unique after the residual gauge fixing.
The particular solutions for the source $`t_{\mu \nu }`$ are even in $`y`$. Therefore, we shall consider eqn. (46) for $`y>0`$,
$$_y\left(\mathrm{e}^{2ky}_y\stackrel{~}{\gamma }_{\mu \nu }\right)2k\mathrm{e}^{2ky}_y\stackrel{~}{\gamma }_{\mu \nu }+\mathrm{}\stackrel{~}{\gamma }_{\mu \nu }=0$$
(47)
and impose the Neumann boundary condition
$$_y\stackrel{~}{\gamma }_{\mu \nu }|_{y=+0}=8\pi \left[t_{\mu \nu }\frac{1}{3}\left(\eta _{\mu \nu }\frac{_\nu _\mu }{\mathrm{}}\right)t\right],$$
(48)
which arises from integrating eqn. (46) over the singularity.
Fourier transforming eqn. (47) with respect to the brane coordinates and changing variables to $`z=\mathrm{e}^{2ky}`$ leads to the differential equation
$$\left(z^2_z^2z_z\frac{p^2}{4k^2}z\right)\stackrel{~}{\gamma }_{\mu \nu }=0,$$
(49)
whose solutions are given in terms of Bessel functions .
Let us consider the case $`p^2>0`$. The static case ($`p_0=0`$) is included here, and the case $`p^2=0`$ can be obtained as a limit from $`p^2>0`$. The two linearly independent solutions of eqn. (49) are
$$\stackrel{~}{\gamma }_{\mu \nu }(p,y)=c_{\mu \nu }(p)\mathrm{e}^{2ky}\{\begin{array}{cc}\mathrm{K}_2\left(\mathrm{e}^{ky}|p|/k\right)\hfill & \\ \mathrm{I}_2\left(\mathrm{e}^{ky}|p|/k\right)\hfill & \end{array}(p^2>0).$$
(50)
Since blowing-up solutions are inconsistent with the linearization (and clearly unphysical), we are led to choose the solution with the $`\mathrm{K}`$ function for $`k>0`$ (Randall-Sundrum background) and the solution with the $`\mathrm{I}`$ function for $`k<0`$ (alternative background). Then, we can determine the coefficients $`c_{\mu \nu }(p)`$ from the boundary condition (48). After doing so, the final solution for $`\stackrel{~}{\gamma }_{\mu \nu }`$ reads
$$\stackrel{~}{\gamma }_{\mu \nu }(p,y)=\frac{8\pi }{|p|}\left[t_{\mu \nu }\frac{1}{3}\left(\eta _{\mu \nu }\frac{p_\mu p_\nu }{p^2}\right)t\right]\mathrm{e}^{2k|y|}\{\begin{array}{cc}\frac{\mathrm{K}_2\left(\mathrm{e}^{k|y|}|p|/k\right)}{\mathrm{K}_1(|p|/k)}\hfill & \text{(RS),}\hfill \\ \frac{\mathrm{I}_2\left(\mathrm{e}^{k|y|}|p|/k\right)}{\mathrm{I}_1(|p|/k)}.\hfill & \end{array}$$
(51)
We can use the solution (51) to discuss the effective laws of gravity on the brane. First, we would like to know the gravitational potential due to, but far away from a static point source, which we introduce by
$$t_{00}(x)=M/M_{Pl(5)}^3\delta (\stackrel{}{x}),t_{00}(p)=2\pi \delta (p_0)M/M_{Pl(5)}^3,$$
where $`M_{Pl(5)}`$ is the Planck mass in five dimensions.
There are two possibilities for the behaviour of a test particle on the brane. First, one could assume that the particle is free to move in five dimensions, i.e. it will follow a geodesic in five-space, which for small velocities is
$$\frac{d^2x^i}{dt^2}\mathrm{\Gamma }^i_{00}$$
(52)
for $`i=1,2,3`$. Then,
$$\mathrm{\Gamma }^i{}_{00}{}^{}=\widehat{\mathrm{\Gamma }}^i{}_{00}{}^{}+k\mathrm{sgn}yn_i=\frac{1}{2}_i\gamma _{00}\frac{1}{8}_i\gamma =\frac{1}{2}_i\stackrel{~}{\gamma }_{00},$$
(53)
where we have used eqn. (25) and the fact that our particular solution is static. Since we are interested in the long distance behaviour on the brane, we set $`y=0`$ in eqn. (51) and consider small $`|p|`$. Using $`\mathrm{K}_2(z)2/z^2`$, $`\mathrm{K}_1(z)1/z`$, $`\mathrm{I}_2(z)z^2/8`$ and $`\mathrm{I}_1(z)z/2`$, we find
$$\stackrel{~}{\gamma }_{00}(p,0)\{\begin{array}{cc}\frac{32\pi k}{3p^2}t_{00}+\mathrm{}\hfill & \text{(RS)},\hfill \\ \frac{4\pi }{3k}t_{00}+\mathrm{}.\hfill & \end{array}$$
The $`1/p^2`$ term for the Randall-Sundrum background generates a $`1/r`$ potential term about the static point source, whereas there is no $`1/r`$ term in the alternative case. More specifically, we can deduce from eqns. (52) and (53) that
$$V_{\text{unconstr}}(r)=\{\begin{array}{cc}\frac{4kM}{3M_{Pl(5)}^3r}+\mathrm{}\hfill & \text{(RS),}\hfill \\ \text{no }1/r\text{ term}\hfill & \text{(alternative)}.\hfill \end{array}$$
(54)
The second possibility is that the particle is constrained to move along the brane by some non-gravitational mechanism. This would mean that it follows a geodesic on the brane, i.e.
$$\frac{d^2x^i}{dt^2}\widehat{\mathrm{\Gamma }}^i{}_{00}{}^{}.$$
(55)
However,
$$\widehat{\mathrm{\Gamma }}^i{}_{00}{}^{}=\frac{1}{2}_i\gamma _{00}=\frac{1}{2}_i\stackrel{~}{\gamma }_{00}+\frac{1}{8}_i\gamma ,$$
(56)
where $`\gamma `$ is given by the Dirichlet boundary condition (40). Hence, the gravitational potential is
$$V_{\text{constr}}(r)=V_{\text{unconstr}}(r)\frac{4\pi k}{3\mathrm{}}t_{00}\{\begin{array}{cc}\frac{kM}{M_{Pl(5)}^3r}+\mathrm{}\hfill & \text{(RS),}\hfill \\ \frac{kM}{3M_{Pl(5)}^3r}+\mathrm{}\hfill & \text{(alternative)}.\hfill \end{array}$$
(57)
Obviously, in the unconstrained case the dynamics on the brane is affected by the non-zero shift vectors. For the RS background, the freedom of a test-particle to move in five dimensions would imply that the trajectories on the brane are unstable, since the brane is repulsive . Hence, we are lead to favour the situation, in which the test particle is confined to move along the brane by some non-gravitational mechanism. Then, potential shortcuts via the fifth dimension are not allowed, the dynamics is determined by the intrinsic metric only, and eqn. (57) applies. Moreover, in the alternative background, only confinement to the brane leads to a $`1/r`$ potential.
Let us also comment on the solution of eqn. (49) for $`p^2<0`$, i.e. for the case of tachyonic matter sources on the brane. In that case, the linearly independent solutions of eqn. (49) are
$$\stackrel{~}{\gamma }_{\mu \nu }(p,y)=c_{\mu \nu }(p)\mathrm{e}^{2ky}\{\begin{array}{cc}\mathrm{N}_2\left(\mathrm{e}^{ky}|p|/k\right)\hfill & \\ \mathrm{J}_2\left(\mathrm{e}^{ky}|p|/k\right)\hfill & \end{array}(p^2<0).$$
(58)
In the Randall-Sundrum background, both modes diverge for large $`y`$. (Although the Bessel functions both go to zero, the exponential factor in front diverges faster.) Interestingly, these modes are integrable, since the norm integral contains a factor $`\mathrm{e}^{4k|y|}`$ in the invariant integration measure cancelling the diverging factor; and they describe the massive Kaluza Klein modes used to construct the five dimensional Green’s function in . In fact, we have performed the integral over the Kaluza Klein states for $`p^2>0`$ in the Green’s function \[17, eqn. (13)\] and found perfect agreement with our solution (51). However, classical solutions containing the $`p^2<0`$ modes will diverge for $`y\mathrm{}`$. Therefore, in the RS scenario, the existence of tachyonic matter on the brane is inconsistent with linearized gravity in the five-dimensional space-time, as is the existence of free Kaluza Klein gravity excitations. The non-existence of both might be quite desirable in a theory describing the real world.
As a second objective we would like to consider the zero-mode truncation of the solution (51) on the brane. The intrinsic Einstein tensor on the brane can be written as
$$\widehat{R}_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\widehat{R}=\frac{1}{2}\mathrm{}\stackrel{~}{\gamma }_{\mu \nu }\frac{1}{4}\left(\gamma _{,\mu \nu }\eta _{\mu \nu }\mathrm{}\gamma \right),$$
(59)
where we have used the gauge (22). In the Randall-Sundrum case, using only the zero-mode of eqn. (51), we find
$$\mathrm{}\stackrel{~}{\gamma }_{\mu \nu }\stackrel{\text{zero-mode}}{=}16\pi k\left[t_{\mu \nu }\frac{1}{3}\left(\eta _{\mu \nu }\frac{_\mu _\nu }{\mathrm{}}\right)t\right].$$
(60)
Furthermore, we can substitute the Dirichlet boundary condition (40) for $`\gamma `$, so that eqn. (59) becomes
$$\widehat{R}_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\widehat{R}\stackrel{\text{zero-mode}}{=}8\pi kt_{\mu \nu }.$$
(61)
The same equation can be obtained using $`\gamma _{\mu \nu }^{}`$, which is related to $`\gamma _{\mu \nu }`$ by a four-dimensional gauge transformation \[cf. eqn. (43)\]. Thus, the zero-mode truncation of the solution (51) yields Einstein’s equation on the brane for the Randall-Sundrum background, which by now is a well-established result (cf. e.g. ).
It is easy to see that this is not the case for the alternative background, as in that case eqn. (60) would not be valid.
## 6 Conclusions
In this paper, we have used a novel gauge in order to obtain a solution of the linearized Einstein equations in the Randall-Sundrum and an alternative background, where the brane is kept straight in spite of matter perturbations located on it. Our solution is consistent in each of the two half-spaces $`y>0`$ and $`y<0`$, and the two patches can be connected by making the gauge transformation (43)–(45).
The explicit solution was summarized already by the equations (1), (2) and (3), and a particular solution of eqn. (3) is found in eqn. (51). Our analysis of the gauge degrees of freedom showed that the traceless transversal part of $`\gamma _{\mu \nu }`$, $`\stackrel{~}{\gamma }_{\mu \nu }`$, represents all physical degrees of freedom. In particular, we conclude that the unphysical graviscalar mode mentioned in is a gauge mode. This was found independently in .
Based on our solution, we studied the effective laws of gravity on the brane and found, in the Randall-Sundrum background, that the zero-mode truncation yields Einstein’s equation for the intrinsic metric on the brane. This implies the validity of the Newtonian limit, if the dynamics is determined by intrinsic quantities on the brane only, and agrees with our derivation of the Newton potential. Moreover, it emphasizes the importance of non-gravitational confinement of matter to the brane in order to prevent the extrinsic geometry from entering the dynamics of matter on the brane. A non-gravitational confinement is also necessary for the geodesics along the brane to be stable . In this article, we did not discuss the corrections to the Newton potential on the brane, as this has been done elsewhere .
In this work, we have restricted our discussion to thin branes. It would be interesting to study the analogue for thick branes, which would also provide an appropriate regularization.
## Acknowledgments
We would like to thank R. Brandenberger, G. Dvali, S. Giddings, L. Randall, K. Sfetsos, R. Sundrum and T. Tanaka for fruitful discussions and correspondence. I. A. and I. V. are grateful to the Physics Department of Simon Fraser University for its kind hospitality. Partial support of this research came from the following grants: RFFI 99-01-00166 and INTAS 99-0545 for I. A., the grant for leading scientific schools 96-15-96208 for M. I. and I. V., NSERC for W. M. and K. V., and RFFI 99-01-00105 and INTAS 99-0590 for I. V.
## Appendix
Here, we list various linearized expressions necessary for the calculations in the main text. The metric tensor has the form (7), where the induced metric $`\widehat{g}_{\mu \nu }`$ is linearized by eqn. (15), and $`\gamma _{\mu \nu }`$, $`n_\mu `$ and $`n^21`$ are small perturbations. We shall henceforth raise and lower the indices of $`\gamma _{\mu \nu }`$ and of $`_\mu `$ (and only of these) with the Lorentz metric.
First, the connection coefficients intrinsic to the hypersurfaces are
$`\widehat{\mathrm{\Gamma }}^\mu _{\nu \lambda }`$ $`={\displaystyle \frac{1}{2}}(\gamma ^\mu {}_{\nu ,\lambda }{}^{}+\gamma ^\mu {}_{\lambda ,\nu }{}^{}\gamma _{\nu \lambda }{}_{}{}^{,\mu }),`$
and the intrinsic Ricci tensor and curvature scalar are, respectively
$`\widehat{R}_{\nu \rho }`$ $`={\displaystyle \frac{1}{2}}(\gamma ^\mu {}_{\nu ,\mu \rho }{}^{}+\gamma ^\mu {}_{\rho ,\mu \nu }{}^{}\mathrm{}\gamma _{\nu \rho }\gamma _{,\nu \rho }),`$
$`\widehat{R}`$ $`=\mathrm{e}^{2k|y|}\left(\gamma ^{\mu \nu }{}_{,\mu \nu }{}^{}\mathrm{}\gamma \right).`$
The linearized expression for the second fundamental form, (cf. eqn. (10)) is
$$H_\nu ^\mu =\frac{1}{2n}[2k\mathrm{sgn}y\delta _\nu ^\mu +\gamma ^\mu {}_{\nu ,y}{}^{}\mathrm{e}^{2k|y|}\eta ^{\mu \lambda }(n_{\nu ,\lambda }+n_{\lambda ,\nu })].$$
The necessary connection coefficients of the five-space are
$`\mathrm{\Gamma }^\mu _{\nu \lambda }`$ $`=\widehat{\mathrm{\Gamma }}^\mu {}_{\nu \lambda }{}^{}k\mathrm{sgn}y\eta _{\nu \lambda }\eta ^{\mu \rho }n_\rho ,`$
$`\mathrm{\Gamma }^y_{\nu \lambda }`$ $`={\displaystyle \frac{1}{n^2}}\left(k\mathrm{sgn}y\widehat{g}_{\nu \lambda }\right)+{\displaystyle \frac{1}{2}}\left(n_{\nu ,\lambda }+n_{\lambda ,\nu }\mathrm{e}^{2k|y|}\gamma _{\nu \lambda ,y}\right),`$
$`\mathrm{\Gamma }^\mu _{\nu y}`$ $`=k\mathrm{sgn}y\delta _\nu ^\mu +{\displaystyle \frac{1}{2}}\left[\gamma ^\mu {}_{\nu ,y}{}^{}+\mathrm{e}^{2k|y|}\eta ^{\mu \lambda }(n_{\lambda ,\nu }n_{\nu ,\lambda })\right],`$
$`\mathrm{\Gamma }^y_{\nu y}`$ $`=k\mathrm{sgn}yn_\nu +{\displaystyle \frac{1}{2}}(n^21)_{,\nu },`$
$`\mathrm{\Gamma }^y_{yy}`$ $`={\displaystyle \frac{1}{2}}(n^21)_{,y}.`$
Thus, the linearized Ricci tensor becomes
$`R_{\mu \nu }`$ $`=\widehat{R}_{\mu \nu }{\displaystyle \frac{4k^2}{n^2}}\widehat{g}_{\mu \nu }+{\displaystyle \frac{2k}{n^2}}\delta (y)\widehat{g}_{\mu \nu }{\displaystyle \frac{1}{2}}\left(\mathrm{e}^{2k|y|}\gamma _{\mu \nu ,y}\right)_{,y}{\displaystyle \frac{1}{2}}(n^21)_{,\mu \nu }`$
$`+{\displaystyle \frac{1}{2}}(n_{\mu ,\nu }+n_{\nu ,\mu })_{,y}k\mathrm{sgn}y(n_{\mu ,\nu }+n_{\nu ,\mu }+\eta _{\mu \nu }n_\lambda {}_{}{}^{,\lambda })`$
$`+k\mathrm{sgn}y\mathrm{e}^{2k|y|}\left\{\gamma _{\mu \nu ,y}+{\displaystyle \frac{1}{2}}\eta _{\mu \nu }\left[\gamma _{,y}(n^21)_{,y}\right]\right\}.`$
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# The Properties of Matter in White Dwarfs and Neutron Stars
## I. Introduction
Astronomical phenomena provide many examples where matter exists in extreme conditions not found in terrestrial environments. One example is the high density of degenerate matter in “compact objects” - the relics of stars that have ceased burning thermonuclear fuel, and thereby no longer generate thermal pressure to support themselves against gravitational collapse. By contracting appreciably from their original sizes, the interiors of compact objects reach sufficiently high densities to produce nonthermal pressure via degenerate fermion pressure and particle interactions. Compact objects provide cosmic laboratories for studying the properties of matter at high densities.
Firm observational evidence and well-founded theoretical understanding both exist for two classes of compact objects which support themselves against collapse by cold, degenerate fermion pressure: white dwarfs, whose interiors resemble a very dense solid, with an ion lattice surrounded by degenerate electrons, and neutron stars, whose cores resemble a giant atomic nucleus - a mixture of interacting nucleons and electrons, and possibly other elementary particles and condensates. White dwarfs are supported by the pressure of degenerate electrons, while neutron stars are supported by pressure due to a combination of nucleon degeneracy and nuclear interactions. These unique states of matter are achieved by significant compression of stellar material. Table 1 compares the principal physical quantities of a typical white dwarf and neutron star with those of the sun<sup>2</sup><sup>2</sup>2Throughout this review we will be using units which are the standard in astrophysical research: cgs for microscopic properties of matter and solar units (denoted by $``$) for macroscopic properties of astronomical objects..
Condensed matter in compact objects spans an enormous range of densities, which we loosely refer to as “high densities”. These extend from about $`7`$gm$`\text{cm}^3`$ (e.g., the density of terrestrial $`{}_{26}{}^{}{}_{}{}^{56}`$Fe), at the surface of a cold neutron star or white dwarf, to as much as $`\rho 10^{15}`$gm$`\text{cm}^3`$ , several times the density in atomic nuclei, in the cores of neutron stars. Matter at the various densities found in compact objects exhibits a variety of novel properties. Electromagnetic, strong, and weak interactions all play an important role in determining the character of compact objects. Since these objects are bound by gravity, they are a meeting point of all four of the fundamental forces of nature. Correspondingly, the astrophysics of white dwarfs and neutron stars incorporates a wide variety of physics including nuclear, particle, solid state and gravitation physics, to name a few areas.
In this review we briefly survey the theory of condensed matter at high densities in compact objects and illustrate how the basic theory is tested through astronomical observations. Since we must cover fifteen orders of magnitude in density, our presentation is at most introductory in nature, and we encourage the interested reader to pursue the cited references. A detailed introduction to the physics of high density matter and compact objects can be found in the textbook Black Holes, White Dwarfs and Neutron Stars: the Physics of Compact Objects, by Shapiro and Teukolsky .
We begin by considering the fundamental nature of cold ($`T=0`$) high density matter in Section II., and distinguish between different regimes of high density. In Section III. we connect these microscopic properties with the fundamental macroscopic parameters of a compact object through the hydrostatic equilibrium dependence of mass and radius on central density. We summarize the fundamental predictions regarding the structure of white dwarfs and neutron stars. In Section IV.-V. we examine how observations of white dwarfs and neutron stars can be used to probe the properties high density matter. We briefly discuss the perturbative effects of a finite temperature in Section VI..
## II. The Cold, High Density Equation of State
The pressure and energies in white dwarfs and neutron stars are nonthermal; thermal effects due to a finite-temperature can be treated as a perturbation. We may therefore treat high density matter as having a zero temperature to very good approximation. The equation of state (EoS) of the matter then reduces to a single-parameter function, $`P(\rho _0)`$ and $`\rho (\rho _0)`$, or $`P(\rho )`$, where $`P`$ is the pressure, $`\rho _0`$ is the rest-mass density and $`\rho `$ is the total mass-energy density which accounts for internal (possibly relativistic) particle energies as well as rest-mass energy.
There exist two main regimes of high density, distinguished as follows. As long as all nucleons are confined to nuclei, their contribution to the total pressure is negligible compared to that of the degenerate electrons. At some threshold density, $`\rho _{ndrip}`$, it becomes favorable for the nuclei to disintegrate, i.e., neutrons “drip” out of the nuclei and form a “nucleon gas”. The standard EoS of Baym, Pethick and Sutherland suggests that $`\rho _{ndrip}4\times 10^{11}`$gm$`\text{cm}^3`$ . We may distinguish between the EoS for $`\rho \rho _{ndrip}`$ which characterizes matter in white dwarfs and in the outermost layers of neutron stars, and $`\rho >\rho _{ndrip}`$, which describes matter in the interior of neutron stars.
### A. The Equation of State below Neutron Drip Density: $`\rho <4\times 10^{11}`$gm$`\text{cm}^3`$
In matter below the neutron drip density the ions provide a Coulomb lattice of point-like charges, which is (to good approximation) independent of the properties of the surrounding electrons. The EoS of such matter is governed mainly by the electron gas. To lowest approximation, we may treat the electrons as an ideal fermion gas, incorporating some Coulomb corrections at relatively low density ($`\rho 10^4`$gm$`\text{cm}^3`$ ) and corrections due to inverse $`\beta `$decay just below the neutron drip density ($`10^9`$gm$`\text{cm}^3`$ $`\rho \rho _{ndrip}`$). The EoS of condensed matter below neutron drip density is well understood. The standard equations for cold, degenerate matter in white dwarfs (helium, carbon, oxygen and, possibly iron dominated models) have been derived by Chandrasekar and Salpeter , whereas models for equilibrated matter<sup>3</sup><sup>3</sup>3The equilibrium isotope of matter is the nucleus of highest binding energy per nucleon. At low densities this isotope is normal $`{}_{26}{}^{}{}_{}{}^{56}`$Fe, but as density increases so do the atomic mass and neutron to proton ratio ., are based on the works of Dirac and Feynman, Metropolis and Teller for $`\rho 10^4`$gm$`\text{cm}^3`$ , and e.g., of Harrison and Wheeler and Baym, Pethick and Sutherland for higher densities.
#### 1. The Ideal Fermion Gas
For almost the entire range of high densities the electrostatic energy associated with the structure of matter is much smaller than the Fermi energies. Consequently, Coulomb forces are generally negligible in a first order treatment of the high density EoS. The electron component of high density matter can therefore be described by a cold, single species gas of noninteracting fermions. At zero temperature the fermions fill all the states with momentum $`pp_F`$ and none of the states with $`p>p_F`$, where $`p_F`$ is the Fermi momentum. The corresponding Fermi energy of the particle species is
$$E_F\left((p_Fc)^2+(mc^2)^2\right)^{1/2},$$
(1)
where $`c=3\times 10^{10}`$cm$`\text{sec}^1`$ is the speed of light in vacuum and $`m`$ is the fermion rest mass.
For electrons, their number density, $`n_e`$, is directly related to their Fermi momentum, $`p_{F,e}`$ by integrating over all occupied phase space ($`h=6.63\times 10^{27}`$ergs sec is Plank’s constant):
$$n_e=_0^{p_{F,e}}n_e(p)𝑑p2\frac{1}{h^3}_0^{p_{F,e}}4\pi p^2𝑑p=\frac{8\pi p_{F,e}^3}{3h^3}.$$
(2)
The factor of 2 in the second equation arises from the electron spin degeneracy.
The pressure the electrons supply is calculated through the mean momentum flux of the electron gas,
$$P=\frac{1}{3}_0^{p_{F,e}}v_e(p)n_e(p)p𝑑p=\frac{2}{h^3}_0^{p_{F,e}}\frac{p^2c^2}{(p^2c^2+(mc^2)^2)^{1/2}}4\pi p^2𝑑p=\frac{m_ec^2}{\lambda _e^3}\varphi (x),$$
(3)
where $`x_ep_{F,e}/m_ec`$ is the electron “relativity parameter”, $`\lambda _eh/(2\pi m_ec)`$ is the electron Compton wavelength, $`v_e=p_ec^2/E_e`$ is the electron velocity and
$$\varphi (x)=\frac{1}{8\pi ^2}\left\{x(1+x^2)^{1/2}(2x^2/31)+\mathrm{ln}\left[x+(1+x^2)^{1/2}\right]\right\}.$$
(4)
The mass-energy density of the free electrons is also uniquely related to the Fermi momentum as
$$\epsilon _e=_0^{p_{F,e}}E_e(p)n_e(p)𝑑p=\frac{2}{h^3}_0^{p_{F,e}}(p^2c^2+(mc^2)^2)^{1/2}4\pi p^2𝑑p=\frac{mc^2}{\lambda _e^3}\chi (x_e);$$
(5)
where
$$\chi (x)=\frac{1}{8\pi ^2}\left\{x(1+x^2)^{1/2}(1+2x^2)\mathrm{ln}\left[x+(1+x^2)^{1/2}\right]\right\}.$$
(6)
However, even when the degenerate electrons contribute most of the pressure, the mass-energy density is dominated by the rest mass of the ions, which are very nonrelativistic at these densities. Thus, the density of the matter may be simply taken as
$$\rho =\rho _0=\frac{n_em_B}{Y_e},$$
(7)
where $`Y_e`$ is the mean number of electrons per nucleon and $`m_B`$ is the mean nucleon mass. In the case of white dwarfs we may set $`Y_e=Z/A=0.5`$ ($`Z=`$atomic number, $`A=`$atomic weight) which is appropriate for fully ionized helium, carbon or oxygen, the most abundant constituents in a white dwarf, and $`m_B=1.66\times 10^{24}`$gm. Combining Eqs. (2), (3) and (7) provides the basic EoS of electron-pressure dominated high-density condensed matter in this regime.
There exist opposite limits to the cold fermion gas equation of state: the low density, nonrelativistic ($`x1`$) and the high density, extremely relativistic ($`x1`$) limits. From Eqs. (2) and (7) we find that in cold matter with $`Y_e=0.5`$ the electron relativity parameter satisfies $`x_e=1`$ at $`\rho 10^6`$gm$`\text{cm}^3`$ , which therefore marks the transition density between a nonrelativistic (NR) and extremely relativistic (ER) electron gas. It is convenient to write the EoS in the two limits in a polytropic form,
$$P(\rho )=K\rho ^\mathrm{\Gamma }$$
(8)
where
NR $`x_e1,\rho 10^6\text{gm}\text{cm}^3:`$ $`\mathrm{\Gamma }={\displaystyle \frac{5}{3}},K=1.0036\times 10^{13}Y_e^{5/3};`$ (9)
ER $`x_e1,\rho 10^6\text{gm}\text{cm}^3:`$ $`\mathrm{\Gamma }={\displaystyle \frac{4}{3}},K=1.2435\times 10^{15}Y_e^{4/3}.`$ (10)
The constant $`K`$ is in cgs units, yielding a pressure in dyne$`\text{cm}^2`$ for a density in gm$`\text{cm}^3`$ . Note that in this approximation the composition of the matter enters only through $`Y_e`$. Correspondingly, helium, carbon and oxygen, which all have $`Y_e=0.5`$ have identical ideal equations of state, slightly stiffer than that of iron ($`Y_e=26/560.43`$). Fully equilibrated matter (often referred to as “catalyzed”) has a $`Y_e`$ that decreases with density and is therefore softer than matter composed of a single element.
The free, degenerate electron pressure EoS outlined here is a good approximation for the equation of state below neutron-drip density. It was employed by Chandrasekar in his pioneering analysis of equilibrium white dwarfs , for which he received the Nobel prize in 1983. More exact treatments were introduced in later years, which included the two main required corrections - electrostatic effects at low densities, and neutronization (or inverse $`\beta `$decay) at higher densities.
#### 2. Electrostatic Corrections to the Cold Equation of State: $`\rho <10^4`$gm$`\text{cm}^3`$
There exists a net electrostatic correction to the ideal equation of state due to the fact that the local distribution of charge is very nonuniform. The fact that positive charge is concentrated in point-like ions causes the average electron-ion separation to be smaller than the average distance between electrons. The net electrostatic potential felt by the electrons is thus an attractive one, which effectively reduces the pressure for a given density.
Electrostatic corrections to the cold equation state are mostly important at relatively low densities. Electrostatic energies are inversely proportional to the average separation between particles, $`<r>`$ which is naturally proportional to $`n_e^{1/3}`$. The relative importance of electrostatic energy, $`E_C`$, between a degenerate, nonrelativistic electron and an ion of charge $`Z`$, can be estimated through
$$\frac{E_C}{E_F^{}}\frac{Ze^2/<r>}{p_{F,e}^2/2m_e}n_e^{1/3}$$
(11)
where $`E_{F,e}^{}=p_{F,e}^2n_e^{2/3}`$ is the Fermi kinetic energy of the nonrelativistic electrons. Unlike the case of hot matter (where the mean electron kinetic energy is $`k_BT`$), the relative importance of electrostatic corrections decreases with density as $`n_e^{1/3}`$.
A rough estimate of the quantitative electrostatic correction can be performed by using the Wigner-Seitz approximation, which describes the lattice as neutral sphere with a central point-like ion and an ambient uniform electron gas<sup>4</sup><sup>4</sup>4We note that such an approximation is actually better suited for white dwarfs than for laboratory solids, where the electron distribution is much more nonuniform.. For a carbon or oxygen lattice, the correction to the pressure is typically of the order of a few percent. At lower densities the electron distribution deviates from uniformity, and more sophisticated approaches, namely the Thomas-Fermi and Thomas-Fermi-Dirac models, must be invoked .
#### 3. Neutronization Corrections to the Cold Equation of State: $`10^9`$gm$`\text{cm}^3`$ $`\rho <4\times 10^{11}`$gm$`\text{cm}^3`$
Stable high density matter must be in chemical equilibrium to all types of reactions, including the weak interactions which drive $`\beta `$ decay and electron capture (“inverse $`\beta `$decay”):
$$np+e+\overline{\nu }_e,p+en+\nu _e,$$
(12)
where $`n`$ and $`p`$ denote a neutron and a proton, respectively, $`e`$ denotes and electron and $`\nu _e(\overline{\nu }_e)`$ denote an electron neutrino (anti-neutrino). If the matter’s composition is out of $`\beta `$equilibrium, it will adjust through $`\beta `$decays or electron capture. Both types of reactions change the electron per nucleon fraction, $`Y_e`$, and thus affect the EoS (Eqs. (9-10)).
In cold white dwarfs and neutron stars the weakly interacting neutrinos freely escape the system: a zero neutrino abundance implies a zero neutrino chemical potential. The condition of chemical equilibrium is then stated as:
$$\mu _n=\mu _p+\mu _e,$$
(13)
where $`\mu _x`$ denotes the chemical potential of species $`x`$. The condition of chemical equilibrium is the fundamental origin of the stable existence of neutrons in nuclei and in uniform $`npe`$ matter. For free, single, particles the chemical potential is identical to the rest mass. The masses in MeV of the three elementary particles of Eq. (13) are $`m_n=939.6,m_p=938.3`$ and $`m_e=0.511`$ (1 MeV is equivalent to $`1.78\times 10^{27}`$gm). It is energetically allowed to have a free neutron decay in the reaction $`np+e+\overline{\nu }_e`$ ($`m_nm_pm_e0.8\text{MeV}`$). Indeed, the life-time of a free neutron is only about $`1000`$sec before it undergoes a $`\beta `$-decay. By contrast, in a cold, degenerate, noninteracting $`npe`$ gas neutron decay can be “blocked”: since the protons and electrons must obey the Pauli principle, the decays will be suppressed if the energy available to the newly formed electron and proton is insufficient to place them above their respective Fermi levels. At such high densities, the equilibrium state of the gas includes a finite fraction of neutrons.
In bulk matter, the electron chemical potential, which is equal to the electron Fermi energy, rises with electron number density. Hence, maintaining chemical equilibrium (Eq. 13)) may require some protons to capture electrons and covert to neutrons. In matter below the neutron drip density these conversions occur in nuclei, so their neutron fraction increases - they “neutronize”. The result is a net decrease in the electron abundance ($`Y_e`$) at high densities, which lowers the pressure and “softens” the EoS.
The fact that “neutronization” becomes energetically favorable at densities below the neutron drip must be taken into account when formulating an exact EoS for this range. In matter composed of nuclei and free electrons, $`\beta `$decay is limited to relatively high densities ($`\rho 10^9`$gm$`\text{cm}^3`$ ), and the exact threshold for the onset of neutronization depends on the nature of the nucleus, since the nuclear interactions must be taken into account in addition to Fermi kinetic energies. For example, the reaction $`{}_{6}{}^{12}C+e_5^{12}B`$ requires the nucleus to absorb an energy of $`14`$MeV, which an electron can supply if the density of a pure carbon lattice reaches $`3\times 10^{10}`$gm$`\text{cm}^3`$ . The EoS of matter close to the neutron drip density is thus dependent on its chemical composition even for elements of equal $`Y_e`$, such as helium and carbon .
### B. The Equation of State Above Neutron Drip Density: $`\rho >4\times 10^{11}`$gm$`\text{cm}^3`$
As neutronization proceeds, the nuclei become increasingly neutron-rich. The so-called “tensor-force” component of the nuclear interactions causes like nucleons to repel one another, so that the binding energy of a neutron rich nucleus is smaller than in one where $`Z/A=0.5`$. Fully equilibrated matter, reaches the last stable isotope <sup>118</sup>Kr ($`Y_e0.31`$) just before the density of $`\rho \rho _{ndrip}4.3\times 10^{11}`$gm$`\text{cm}^3`$ . At higher densities it becomes favorable for some neutrons to “drip” out of their parent nuclei. The nucleon component of the matter can no longer be confined to point-like objects and the nuclei begin to dissolve. As density is increased, a larger fraction of the nucleons exists as “free” particles, outside the nuclei. This is a gradual transition, until the matter approaches the density of atomic nuclei, where all nuclei have essentially dissolved, and the distribution of nucleons becomes uniform. The EoS above the neutron drip density must be formulated by consistently including the effects of nuclear physics, which become the governing component of the properties of matter as the density of atomic nuclei, $`\rho _{nuc}2.8\times 10^{14}\text{gm/cm}^3`$ is approached.
### C. Subnuclear Densities: $`4\times 10^{11}`$gm$`\text{cm}^3`$ $`<\rho <2.8\times 10^{14}`$gm$`\text{cm}^3`$
An analysis of cold matter in the range $`\rho _{ndrip}\rho \rho _{nuc}`$ is complicated by requiring chemical equilibrium between nucleons inside the nuclei and those that have dripped outside. One must account for the effects of the surrounding gas of free nucleons on the nuclei, as well as other effects such nuclei surface and Coulomb energies. In addition, the nuclei are expected to be very neutron rich, deviating from the $`Z/A=0.40.5`$ found in terrestrial nuclei. Nonetheless, the properties of matter in this range of densities can still be derived by a natural extrapolation from ordinary nuclei, and, indeed, the EoS for such matter is believed to be well understood. The principal studies of matter (e.g., Ref. ) are based on the nuclear liquid-drop model, and are suitable for most applications regarding neutron star structure. For problems where a more accurate description of the neutron star inner crust is required, attention must be given to lattice effects, as the nuclei and free nucleons arrange in a distinct spatial structure (where the nuclei settle into bubbles, slabs or rods, depending on density ). It is expected that at $`\rho \frac{1}{2}\rho _{nuc}`$ all nuclei will have dissolved so that the matter is completely uniform.
### D. Supernuclear Densities: $`\rho >2.8\times 10^{14}`$gm$`\text{cm}^3`$
The upper end of the high density regime is referred to as “supernuclear” where $`\rho \rho _{nuc}`$. As matter is compressed to such densities the EoS becomes gradually dominated by the degenerate nucleons and the nucleon interactions.
It is instructive to begin by considering an ideal uniform mixture of neutrons, protons and electrons. For any given nucleon number density, $`n_N`$, solving the equilibrium composition $`n_n,n_p`$ and $`n_e`$ (the neutron, proton and electron number densities, respectively) requires three equations. The first is baryon number conservation,
$$n_N=n_n+n_p,$$
(14)
and the other two arise from imposing chemical equilibrium (Eq. 13) and charge neutrality. The condition for charge neutrality is
$$n_p=n_e.$$
(15)
For a noninteracting mixture of nucleons and electrons all chemical potentials are simply the Fermi energies (Eq.1)). Electrons are extremely relativistic at nuclear densities so their electron chemical potential is
$$\mu _ep_{F,e}c=\mathrm{}c(3\pi ^2Y_e\rho _B)^{1/3}100\left(\frac{Y_e}{0.03}\right)^{1/3}\left(\frac{n_N}{n_{nuc}}\right)^{1/3}\text{MeV},$$
(16)
where $`n_{nuc}=0.16\text{fm}^3`$ is the number density of nuclei at the saturation density (and $`1\text{fm}=10^{13}`$cm). Even if the electron fraction per baryon is only a few percent, the electron chemical potential exceeds the mass difference between neutrons and protons by two orders of magnitude. The only way to maintain a finite electron fraction (which is required to balance the protons for charge neutrality) and satisfy chemical equilibrium is by having a significantly larger neutron than proton fraction (note that for nonrelativistic fermions $`E_Fn^{2/3}`$ and for extremely relativistic fermions $`E_Fn^{1/3}`$). Equilibrium matter at nuclear densities must be very neutron dominated, and objects composed of such matter are thus “neutron stars”. For an noninteracting gas the ratio of neutrons to protons in equilibrium must be about $`8:1`$, which is also representative of more realistic models of supernuclear densities.
The noninteracting gas approximation is not reliable for deriving the EoS at supernuclear densities. Unlike electrostatic perturbations, nucleon-nucleon interactions are not negligible, and the interaction energies are comparable to the Fermi energies of the degenerate nucleons (electrons do not feel the strong interaction, and may still be treated as noninteracting). Modeling of the nucleon-nucleon interaction is one of the longest-standing problems in nuclear physics, still only partially solved. Profound difficulties exist due to the absence of a comprehensive theory of the interactions and the difficulty of obtaining experimental data for $`\rho >\rho _{nuc}`$. Further complications arise due to the fact the Fermi and interaction energies at $`\rho 23\rho _{nuc}`$ reach a sizable fraction of the rest mass, and relativistic effects must be taken into account as well. It is not yet possible to apply quantum chromodynamics (QCD), the fundamental theory of strong interactions, to the many-body nuclear domain at $`\rho \rho _{nuc}`$. Instead, the most useful approaches are still based on phenomenological potential formalisms, and many-body Schrödinger-like systems of equations .
The EoS of supernuclear matter remains to date a field of active research. Current approaches include variational methods based on deduced two and three nucleon interactions (, see for a recent study), and relativistic mean field approximations (, see for a review). The models are based on fitting parameters to reproduce the empirically determined properties of finite nuclei with $`n=n_{nuc}`$. A rough approximation of the properties of the nucleon component of supernuclear matter can be obtained through the effective, nonrelativistic model of Ref. ,
$$\rho (n,x)=E_{gas}(n,x)+an^2+b(n_nn_p)^2+cn^{\delta +1}.$$
(17)
The first term $`E_{gas}`$ is the mass-energy density of a noninteracting gas of nucleons at density $`n`$ and composition $`n_n,n_p`$. The coefficient $`a`$ is negative, representing the long range attractive component of the inter-nucleon force, while $`c`$ is positive, representing the short-range repulsive component. The power $`\delta `$ is larger than unity, so that short-range repulsion dominates at high density. The symmetry term includes the positive coefficient $`b`$ which describes the “tensor-force” that repels like nucleons, and its contribution is therefore minimized at for symmetric nuclear matter ($`n_n=n_p`$). The values of $`a,b,c`$ and $`\delta `$ are derived by requiring Eq. (17) to reproduce the assumed properties of symmetric matter at $`n_N=n_{nuc}`$. See Refs. for typical (model dependent) values of these parameters.
Finally, we note that it is quite possible that other particles, besides neutrons, protons and electrons, coexist in stable equilibrium at supernuclear matter. The most obvious example is the muon, which is a lepton similar to the electron, but has a rest mass of $`m_\mu 105\text{MeV}`$. The equilibrium condition for the muons will simply be $`\mu _\mu =\mu _e`$, so from Eq. (16) it is evident that at densities $`n>n_{nuc}`$ it is energetically favorable to convert some electrons into muons through the weak interactions. At densities of $`n>2n_{nuc}`$ it possible that other exotic particles will appear, such as hyperons (baryons that are heavier than nucleons), Bose-Einstein condensates of mesons (i.e., pions or kaons), or even conversion of the nucleons into an uniform mixture of quarks (see Refs. for recent reviews of the possible presence and roles of such particles in neutron stars).
### E. Basic Properties of High Density Matter
In Table 2 we present an EoS of cold, fully catalyzed matter ranging from normal iron at zero pressure to supernuclear densities. The EoS in the subnuclear regime is compiled from Refs. and are “standard” for studying catalyzed high density matter. For realistic white dwarf models, which are presumably composed mostly of helium or carbon and oxygen, and never reached sufficiently high temperatures to catalyze their nuclei, a slightly different EoS must be used, based on a single species; see Ref. . We also tabulate one “state-of-the-art” EoS for the supernuclear range
The EoS $`P(\rho )`$ is plotted in Figure 1a. It is especially instructive to examine corresponding values of the adiabatic index, $`\mathrm{\Gamma }d\mathrm{ln}P/d\mathrm{ln}\rho `$, plotted in Figure 1b, along with the sound speed ($`c_s=\sqrt{dP/d\rho }`$). Note that around $`\rho 10^6`$gm$`\text{cm}^3`$ the EoS transforms, as expected, from a nonrelativistic $`\mathrm{\Gamma }=5/3`$ polytrope to a relativistic $`\mathrm{\Gamma }=4/3`$ one as the Fermi energy of the electrons, which dominate the pressure, gradually becomes relativistic. There is a sharp drop in the adiabatic index around the neutron drip density, since at first the dripped neutrons contribute to mass-density but not to pressure. Only as matter approaches nuclear densities does the nucleon pressure become important, pushing the adiabatic index back up to values in the range $`23`$.
Along with the EoS, we also examine some thermodynamic quantities of cold high density matter. In Table 2 we also list the adiabatic index, sound speed, bulk modulus $`Y`$, and incompressibility, $`K`$, defined as
$$Y=n\frac{dP}{dn}\text{and}K=9\frac{dP}{dn}9n\frac{d^2\rho }{dn^2}.$$
(18)
Note that the bulk modulus (which has units of pressure) is the reciprocal of the quantity usually defined as the compressibility ($`\chi Y^1`$), while the “incompressibility” (units of energy) is more commonly used in nuclear physics applications (and is indeed measurable in nuclei). In Table 3 we list typical values of $`\mathrm{\Gamma },c_s,Y,K`$ and the specific heat capacity, $`c_v`$, for condensed matter found in white dwarfs and neutron stars. For the purpose of comparison, we also give values for the sun, where the pressure roughly follows an ideal Maxwell-Boltzmann law $`Pnk_BT`$. Note how the extreme conditions of high densities leads to properties that are very different than those found for terrestrial materials. For example, the speed of sound in a neutron star is expected to be several tens of percent of the speed of light!
## III. White Dwarf and Neutron Star Structure
Compact objects are self-gravitating equilibria with a mass comparable to that of the sun, $`1M_{}1.989\times 10^{33}`$grams. Both white dwarfs and neutron stars are centrally condensed objects - most of their mass is located in a high density core, which is limited to a fraction of the volume; furthermore, the radius of the objects decreases with increasing mass. These traits are characteristic of a configuration supported by degenerate-fermion pressure.
### A. Construction of a Self-gravitating, Equilibrium Star
The equilibrium structure of a self-gravitating object is derived from the equations of hydrostatic equilibrium. The simplest case is that of a spherical, nonrotating, static configuration, where for a given EoS all macroscopic properties are parameterized by a single parameter, for example, the central density. In the case of compact objects, the gravitational fields are strong enough that calculations must be performed in the context of general relativistic (rather than Newtonian) gravity. The fundamental equation of hydrostatic equilibrium in its general relativistic form has been derived by Tolman and Oppenheimer and Volkoff , and is known as the “TOV” equation:
$$\frac{dP(r)}{dr}=\frac{Gm(r)\rho (r)}{r^2}\left(1+\frac{P(r)}{c^2\rho (r)}\right)\left(1+\frac{4\pi r^3P(r)}{c^2m(r)}\right)\left(1\frac{2Gm(r)}{c^2r}\right)^1.$$
(19)
This equation simply states that at any radial distance $`r`$, the gravitational pull by the mass interior to $`r`$, $`m(r)`$, is balanced by the gradient of the pressure $`P(r)`$; $`G=6.67\times 10^8\text{cm}^3\text{gm}^1\text{sec}^2`$ is the gravitational constant. Note that the first term on the right hand side is the only term in the nonrelativistic case, where $`P/(\rho c^2)1`$ and $`2Gm/(rc^2)1`$. The second and third factors arise from the pressure being a form of energy density (“regeneration of pressure” effect), while the last term includes the correction due to the curvature of space in the strong gravitational field of the star.
### B. Stable Configurations of High Density Matter
Given an EoS, Eq. (19) can be integrated simultaneously with the mass equation,
$$\frac{dm(r)}{dr}=4\pi r^2\rho (r),$$
(20)
to determine the entire profile of the object. The actual integration is typically solved by numerical means, assuming a central density $`\rho (r=0)=\rho _c`$, and integrating outward from the center until reaching the surface where $`P(r)=0`$, which identifies $`r=R`$ as the radius of the star and $`M(R)`$ as its mass. We emphasize that the TOV equations provide the gravitational mass (or total mass-energy) of the star, which includes the effect of the gravitational binding energy, $`E_{GB}GM^2/R`$, as well as the internal and rest mass energy of the stellar constituents. Note that in the case of neutron stars, $`(GM^2/R)/(Mc^2)1020\%`$; indeed, the gravitational mass of a neutron star is measurably lower than the total rest mass of its constituents.
Varying $`\rho _c`$ produces a sequence of models for the given EoS, yielding $`M(\rho _c)`$ and $`R(\rho _c)`$ along the sequence. The resulting sequences for stars composed of pure carbon (based on ) and of cold, catalyzed condensed matter (with the EoS tabulated in table 2) are both shown in Figure 2 (note that carbon dominated matter beyond the neutron drip density is unphysical and therefore omitted). The most distinct feature about the equilibrium sequences is the existence of local maxima in the $`M(\rho _c)`$ curves. A hydrostatic equilibrium configuration is dynamically unstable to catastrophic gravitational collapse if $`dM/d\rho _c<0`$, since a radial perturbation would cause it to collapse on itself . Therefore stable carbon white dwarfs exist only if the central density is $`\rho _c6\times 10^9`$gm$`\text{cm}^3`$ , while cold, catalyzed, matter has two distinct regimes of stable configurations: one with central densities below $`10^{10}`$gm$`\text{cm}^3`$ (equivalent to white dwarfs), and one where the central density lies roughly in the range $`10^{13}\rho _c10^{15}`$gm$`\text{cm}^3`$ (neutron stars). Configurations where $`\rho _c>\rho _{ndrip}`$ are unstable, and cannot be found in nature. This situation is a direct consequence of the nature of the EoS as shown in Figure 1. Quantitatively, it can be shown that a hydrostatic configuration is stable only if its mass averaged adiabatic index, $`\overline{\mathrm{\Gamma }}`$, satisfies
$$\overline{\mathrm{\Gamma }}\frac{4}{3}>\frac{2GM}{Rc^2}.$$
(21)
It is evident therefore, that the effect of neutronization and neutron drip, which cause $`\overline{\mathrm{\Gamma }}`$ to drop below $`4/3`$, is to separate astrophysical objects with high density matter into two distinct classes of stable configurations, which are shown schematically in Figure 3.
#### 1. White Dwarfs
White dwarfs are composed of matter below the neutronization density. The astrophysical scenario which creates white dwarfs, i.e., quasi-static contraction of a progenitor star at the end point of thermonuclear burning, does not allow matter to reach high enough temperatures to achieve equilibrium composition through further burning. Accordingly, the composition of white dwarfs is mostly dominated by its nuclear ashes. In the standard scenario white dwarfs are composed mostly of carbon and oxygen, but there is observational evidence that helium white dwarfs (and perhaps even iron-core white dwarfs - see below) exist as well. The pressure throughout the star is due to degenerate electrons, which are nonrelativistic in the outer layers, but are relativistic in the interior of the more massive stars with $`MM_{}`$.
By approximating the EoS as a Newtonian polytrope (Eq. (8)), Chandrasekar derived the basic features of white dwarfs:
$`\rho _c<10^6\text{gm}\text{cm}^3`$ $`\mathrm{\Gamma }={\displaystyle \frac{5}{3}}:`$ $`R1.12\times 10^9\times \left({\displaystyle \frac{\rho _c}{10^6\text{gm}\text{cm}^3}}\right)^{1/6}\left({\displaystyle \frac{Y_e}{0.5}}\right)^{5/6}\text{cm}`$
$`M0.496\times \left({\displaystyle \frac{\rho _c}{10^6\text{gm}\text{cm}^3}}\right)^{1/2}\left({\displaystyle \frac{Y_e}{0.5}}\right)^{5/2}M_{},`$
$`\rho _c>10^6\text{gm}\text{cm}^3`$ $`\mathrm{\Gamma }={\displaystyle \frac{4}{3}}:`$ $`R3.35\times 10^9\times \left({\displaystyle \frac{\rho _c}{10^6\text{gm}\text{cm}^3}}\right)^{1/3}\left({\displaystyle \frac{Y_e}{0.5}}\right)^{2/3}\text{cm}`$
$`M1.46\times \left({\displaystyle \frac{Y_e}{0.5}}\right)^2M_{}.`$
White dwarfs are expected to have radii of the order of $`10^4`$km - roughly the size of the Earth, or about one percent of that of the sun (which has a radius $`R_{}7\times 10^5`$km).
As long as the EoS is approximated as a pure polytrope, so that electrostatic and neutronization corrections are ignored, helium, carbon and oxygen white dwarfs ($`Y_e=0.5`$) are all identical, while an iron-dominated white dwarf is only slightly different ($`Y_e0.43`$). Cold, catalyzed matter, which is used for the mass sequence in Figure 2, has a significantly softer EoS (due to $`Y_e`$ decreasing with increasing density), and its sequence lies, therefore, lower than those of stars composed of a single species. The most significant result in Eq. (1.) is that there is no dependence of $`M`$ on the central density. The mass $`M_{Ch}1.46M_{}`$, known as the Chandrasekar mass, is the asymptotic value a white dwarf can reach if it achieves sufficiently high density, so that its entire structure is governed by relativistic fermions (in practice this mass cannot be reached, since the outermost layer have nonrelativistic electrons). Nonrotating white dwarfs cannot have a mass exceeding $`M_{Ch}`$, and the precise limit on their mass is lower by several percent, due to neutronization at high densities (Section III.B.3.) and general-relativistic effects. In spite of the negligible effect on the hydrostatic equilibrium profile of white dwarfs, general relativity is required for a full analysis of white dwarf stability, since an exact $`\mathrm{\Gamma }=\frac{4}{3}`$ star is unstable to gravitational collapse (Eq. 21).
#### 2. Neutron Stars
Neutron stars are composed mostly of matter at nuclear densities, including a core with supernuclear densities, and is topped by a thin crust at subnuclear densities. The crust, composed of cold, catalyzed matter, may also be divided into an inner part with $`\rho _{ndrip}\rho \rho _{nuc}`$ and an outer part where $`\rho \rho _{ndrip}`$.
Immediately after the neutron was discovered (1932) Landau modeled a neutron star as a gas of noninteracting, degenerate neutrons . He found that neutron stars would have a maximum mass of about $`1M_{}`$ and a radius of several kilometers. Although this was a very crude approximation, it does suggest the correct orders of magnitude regarding the structure of these objects: a neutron star has a mass comparable to that of the sun compressed to the size of a medium city! If supported by degenerate fermion pressure, a neutron star cannot have a mass that exceeds the Chandrasekar limit by much, while its radius is inversely proportional to the mass of the pressure providing fermion. A neutron star should indeed have radius about $`m_n/m_e10^3`$ times smaller that of a white dwarf.
A realistic treatment of neutron star structure requires general relativity , whose effects are appreciable in this case. Furthermore, quantitative estimates must be based on a realistic EoS at all densities and especially a realistic model for the supernuclear regime where particle interactions are most important. In particular, these interactions are repulsive at short distances and oppose compression, thereby stiffening the EoS in comparison to a free particle gas. Indeed, plausible equations of state predict that $`M_{max}(NS)2M_{}`$, ($`2.2M_{}`$ for Ref. ), whereas the star’s radius will lie in the range $`R1015`$km. As an example, we show in Figure 4 the $`M`$ vs. $`R`$ relation found for a few representative equations of state of supernuclear densities. The exact value of $`M_{max}(NS)`$ depends on the assumed EoS and therefore provides an integral measure of the properties of matter at supernuclear densities . We also note that the existence of a maximum mass for neutron stars has important astrophysical implications, since it suggests that a larger mass cannot be sustained by cold pressure and must inevitably collapse to a black hole.
## IV. Observations of Condensed Matter in White Dwarfs
Over 2000 white dwarfs have been discovered to date, based on the spectroscopic properties of observed stars . While the thermal energy in white dwarfs is small compared to the Fermi energies of the electrons in the interiors, it is still sufficient to generate an observable surface luminosity for several billion years. The luminosity of a star is dependent on its radius, $`R`$, and its effective surface temperature, $`T_{eff}`$, according to
$$L=4\pi \sigma R^2T_{eff}^4,$$
(24)
where $`\sigma =5.67\times 10^5\text{erg}\text{cm}^2\text{sec}^1\text{deg}^4`$ is the Stephan-Boltzmann constant. The luminosities of white dwarfs are typically $`10^3`$ to $`10^2`$ of that of the sun, $`L_{}4\times 10^{33}`$ergs sec<sup>-1</sup>, but their radii are also smaller than those of ordinary stars by a factor of $`100`$. As a result, white dwarfs have apparent surface temperature of several times $`10^4`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ , unusually high in comparison with most observable stellar objects, so they do appear very “white”. The existence of white dwarfs was established spectroscopically (by determining the surface temperature of sources) as early as 1910.
White dwarfs are believed to be remnants of stars with initial masses in the range $`0.18M_{}`$, which are not massive enough to complete the thermonuclear burning process all the way to iron. Mostly they are mostly composed of carbon and oxygen, but in some cases, thermonuclear burning ceased before these elements were produced, and such stars are dominated be helium.
### A. Radii
Observational determination of a white dwarf radius is straightforward if its flux, $`F`$, is measured and its distance from the Earth, $`D`$, is known:
$$F(D)=\frac{L}{4\pi D^2}R^2=\frac{FD^2}{\sigma T_{eff}^4}.$$
(25)
The stellar radius is derived only after the effective surface temperature is obtained spectroscopically . Estimated radii of white dwarfs reside in the range of $`0.0070.013R_{}`$, which is consistent with an object supported by degenerate electron pressure (Section III.B.) and thereby confirms the basic nature of white dwarfs.
### B. Mass-Radius Relations
The most significant test of the nature of matter in a white dwarf is obtained by comparing observed mass-radius relations with the theoretical predictions. Once the radius of a white dwarf is determined, details of the surface emission, namely effects of gravitational acceleration on line emission (which scales as $`M/R^2`$), or gravitational redshift (which scales as $`M/R`$), are then used to estimate the surface gravity. An independent (and usually more accurate) estimate of the mass can be obtained directly for white dwarfs in binary-star systems, through Kepler’s third law . Through these methods, masses of a few white dwarfs have been estimated with reasonable accuracy, and most seem to cluster around $`0.6M_{}`$ . Some larger mass white dwarfs are also known, including the most famous Sirius B , with $`M1M_{}`$. Note that these masses are smaller than $`M_{Ch}`$: it is believed that stellar evolution, and especially periods of mass loss, limit the masses of most white dwarfs to $`M1M_{}`$, even though the EoS could sustain somewhat higher masses.
Some examples of observationally determined radii and mass for white dwarfs are presented in Table 4. Note that larger mass stars have smaller radii, as is expected (Eqs. (1.-1.)). It is interesting to note that for some time, there was a nonnegligible discrepancy between the observations and theory, where observed radii seemed to be $`1020\%`$ smaller than estimated by theory for carbon white dwarfs . Only improved estimates of distances to several white dwarfs with the Hipparcos satellite, and better modeling of white dwarf atmospheres have allowed this discrepancy to be mostly resolved . A comparison of current observed mass-radii determinations with theoretical curves is presented in Figure 5, and in general, the fit is indeed very good: these new results seem to confirm that the composition of most white dwarfs is indeed dominated by medium weight elements (carbon and oxygen). However, they also imply that a small minority of white dwarfs do have relatively small radii, and may therefore contain iron cores, which presents an intriguing puzzle from the point of view of stellar evolution.
## V. Observations of Condensed Matter in Neutron Stars
Neutron stars are often identified observationally as very accurately pulsating sources, where the pulsation is attributed to the star’s rotation. The period of pulsation limits the size of the source to no more than a few tens of kilometers, thus indicating the presence of a very compact object. Since the first discovery in 1967, close to 800 pulsars have been identified, over 700 of those in radio waves . About thirty X-ray pulsars have also been observed, and a smaller number of nonpulsating X-ray sources of various types are also most likely to contain a neutron star .
Theory suggests that neutron stars are formed when a massive star is disrupted in a “supernova” explosion. As the massive star evolves, matter in its core undergoes thermonuclear burning all the way to iron, which is the most tightly bound nucleus and therefore cannot burn further. When the iron core is massive enough it collapses under its own weight, until the collapse is halted when nuclear densities are reached. Most of the gravitational binding energy is released in the form of neutrinos, but some is transformed into an outgoing shock wave that expels the star’s envelope and leaves behind the newly formed neutron star. Indeed, several dozen young neutron stars have been found in sites of past supernovae, confirming this scenario.
### A. Neutron Star Masses
The masses of over twenty neutron stars have been determined observationally through their gravitational pull on a binary companion star. Within the errors of measurement, the masses of these stars appear to cluster in a narrow range of $`1.35\pm 0.1M_{}`$ . This is especially true of the masses of neutron stars in the four known binary-pulsar systems (where the companion object is also a neutron star), listed in Table 5, where mass determinations are most accurate. Another example is the neutron star - white dwarf binary (once believed to be a double binary neutron star system) B2303+46. The properties of the orbit allow to limit the mass of the neutron star to $`1.2<M_{NS}<1.4M_{}`$ . In particular, the Hulse-Taylor binary pulsar 1913+16 (the coordinates of its location in the sky) has been very accurately determined to have $`M_{1913+16}=1.4411\pm 0007M_{}`$. Clearly, a high density EoS must satisfy $`M_{max}1.44M_{}`$ to be consistent with this observation.
This limit does not impose a serious constraint on realistic models for the supernuclear EoS and is satisfied by all but the softest equations. There has been much debate whether the narrow range of observed neutron star masses is evidence for the maximum mass being rather low ($`M_{max}(NS)1.5M_{}`$) , or whether it is due to astrophysical effects, which could restrict the range of masses of observable stars. Recently, there is growing evidence that heavier neutron stars exist: one in the Vela X-1 binary is now estimated to have a mass of $`1.9M_{}`$ , and the oscillations observed in some galactic X-ray sources seem to indicate that they include neutron stars with masses larger than $`2M_{}`$ . If these estimates are confirmed with improved observations in the future, they will provide a much more stringent test that would be consistent only with the stiffer models for the supernuclear EoS.
### B. Radii
The mass and radius relation of a nonrotating neutron star is uniquely defined for any given EoS through the solution of the TOV equations (Figure 4). However, to date, observational methods for estimating neutron star radii lack the accuracy required to critically differentiate between realistic equations of state. Hopefully, data from a new generation of X-ray satellites will provide substantial constraints on the compactness ($`M/R`$) of neutron stars (either through their thermal emission or from emission of material accreting upon them ).
### C. Rotation periods
The observed pulsation period of a neutron star is attributed to rotation. Most measured periods are in the range $`0.251`$sec, but a subclass of millisecond pulsars are also known to exist. The fastest spinning neutron star observed to date has a period of $`P=1.56`$msec (so that it spins faster than an egg-beater!).
Rotation provides an additional centrifugal barrier which assists the internal pressure in supporting the gravitational load. The resulting configuration is then dependent on both central density and rotation period, and must be solved self-consistently including the effects of general relativity. The maximum mass a given equation of state can support increases with respect to the static (TOV) value as the star is allowed to rotate faster . There must exist a lower limit on the rotation period (upper limit on the rotation rate), since too rapid rotation would cause the star to shed mass at the equator. The general trend is that a softer EoS predicts a smaller lower limit on the rotation period, since a more compact star is more tightly bound gravitationally and so is better suited to resist the centrifugal force. At present, the limit of $`1.56`$msec is not very restrictive, and is consistent with basically all realistic models of high density matter. Discovery of a sub-millisecond ($`P>0.5`$msec) pulsar would allow to distinguish more directly between competing supernuclear equations of state .
### D. Pulsar Glitches
Pulsar periods are observed to increase gradually with time, implying that rotational energy of the star is being lost (primarily by electromagnetic dipole radiation). However, several pulsars have been observed to undergo sudden decreases in the rotation period, which are believed to originate from the transfer of angular momentum between different parts of the star. The current leading model suggests that the source of excess angular momentum is the inner curst, and thus the magnitude of glitch phenomena can be used to set a lower limit on the fraction of the total moment of inertia the inner crust must carry. Recent results suggest the crust of the Vela pulsar must carry at least $`1.4\%`$ of the total moment of inertia of the star, which is sufficient to rule out equations of state which are very soft in the range $`1\rho /\rho _{nuc}2`$. Possible progress in theoretical modeling of glitch phenomena could provide further limits on the properties of high density matter near the nuclear saturation density.
## VI. Thermal Properties of Matter in White Dwarfs and Neutron Stars
Fermions in both white dwarfs and neutron stars are expected to have degenerate energies which are much larger than their thermal energies. However, compact objects do have finite temperatures that are a relic of their progenitors. Since degenerate matter has very long scattering lengths for individual particles, the cores of both white dwarfs and neutron stars are expected to be excellent thermal conductors, and practically isothermal.
Finite temperature effects on the properties of condensed matter in compact objects has become a very active theoretical and observational field in recent years, especially in the case of white dwarfs. A finite temperature leads to cooling through thermal emission processes which are potentially observable. Examining the cooling history (i.e., temperature vs. age) of compact objects thus provides empirical evidence regarding the thermal properties of high density matter. Here we can only briefly mention some of the topics of current interest.
### A. Thermal Effects in White Dwarf Structure
The region most affected by a finite temperature is the low density outer-layers of a white dwarf, which may include a thin atmosphere composed of helium and possibly also hydrogen. For lower densities the electrons are not degenerate and a surface temperature of $`10^4`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ will severely alter the equation of state, especially in the regime where electrostatic corrections are important. Modern equations of state for a hydrogen/helium atmosphere have been incorporated in several studies of white dwarf structure (see for helium stars, for carbon-oxygen, and for a recent compilation of various compositions). The thermal pressure in the outer layers is generally found to be very effective in inflating the stellar radius. As seen from Figure 5, a surface temperature of $`T_{eff}=10^4`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ increases the radius by a significant fraction with respect to the cold models.
There has been a recent revival of studies of white dwarf thermal evolution, and its dependence on the properties of matter at the relevant high densities. We briefly review here the two main aspects of these studies, namely cooling and pulsations.
#### 1. White Dwarf Cooling
The basic theory of white dwarf cooling was established by Mestel in 1952 . The heat is originally stored in the nondegenerate ion-lattice, and the dominant cooling mechanism is photon diffusion to the surface from the isothermal core through the nondegenerate outer layers. The simplest results are based on Kramer’s approximation for the photon opacities in the nondegenerate regime:
$$\kappa =\kappa _0\rho T^{3.5},$$
(26)
where $`\kappa `$ is the mean opacity (in cm$`{}_{}{}^{2}\text{gm}_{}^{1}`$ ) and $`\kappa _0=4.34\times 10^{24}Z(1+X_H)`$cm$`{}_{}{}^{2}\text{gm}_{}^{1}`$ is a composition dependent constant, $`X_H`$ and $`Z`$ being the mass-fraction of hydrogen and heavy (not hydrogen or helium) elements, respectively. The luminosity through the outer layers, $`L(r)`$, is dependent on the temperature gradient according to the radiative diffusion equation
$$L(r)=4\pi r^2\frac{c}{3\kappa \rho }\frac{d}{dr}(aT^4)$$
(27)
(where $`a=7.7565\times 10^{15}\text{erg}\text{deg}^4`$ is the radiation constant). Combining the above thermal temperature gradient with the pressure gradient required for hydrostatic stability (Eq. (19), see ), one finds that the surface luminosity satisfies
$$L=5.7\times 10^6\mu Y_e^2\frac{1}{Z(1+X)}\left(\frac{M}{M_{}}\right)T_c^{3.5},$$
(28)
where $`\mu `$ is the mean molecular weight in the nondegenerate atmosphere and $`T_c`$ is the core temperature in the star. For typical white dwarfs ($`M1M_{}`$) with no hydrogen in the atmosphere, $`X=0,Z0.1`$, $`\mu =1.4,Y_e=0.5`$, an observed luminosity in the range $`10^510^2L_{}`$ corresponds to a central temperature of $`T_c10^610^7`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ .
The key elements of Eq. (28) are that the surface luminosity (which is observable) is clearly related to the mass of the star, the composition and the internal temperature. If the mass is also determined separately (as discussed above), luminosity can then provide a direct test of the star’s thermal and composition profiles. Furthermore, the age of the white dwarf can also be estimated based on its cooling time scale, which is dependent on the luminosity function (Eq. 28), the stellar mass and the specific heat capacity of a white dwarf material. The latter (in ergs$`\text{deg}^1\text{gm}^1`$) if roughly that of the ion lattice,
$$c_v(\text{ions})=2\frac{3}{2}\frac{k_B}{\mu },$$
(29)
where the factor of 2 comes from the three degrees of collective vibration and rotation, in addition to the three single ion degrees of motion. Due to their degeneracy, the specific heat capacity of the electrons is suppressed by a factor of $`k_BT/E_F`$ , which for a density of $`10^8`$gm$`\text{cm}^3`$ and an internal temperature of $`10^7`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ is $`10^3`$.
While models which include more detailed microphysics do reproduce the general results of the Mestel model (see, for example, in Refs. ), accurate work that is compatible with the quality of current observations must include several important perturbations. Most notably, accurate low temperature opacities, semidegenerate electron thermal conductivity and pressure-induced ionization must be accounted for in modeling the white dwarf atmosphere. Crystallization of the ion gas may also have a significant effect on the thermal history of the star, both as a transient source of (latent) heat and by effectively reducing the heat capacity. In particular, below the Debye temperature, $`\mathrm{\Theta }_D4\times 10^3\rho ^{1/2}`$ $`{}_{}{}^{\mathrm{o}}\text{K}`$ , the specific heat capacity becomes temperature dependent ,
$$c_v=\frac{12\pi ^4}{5}k_B\left(\frac{T}{\mathrm{\Theta }_D}\right)^3.$$
(30)
The specific effects of the finite temperature on the thermodynamic properties of matter in white dwarf atmosphere are most notable through studies of pulsations (see below).
It is noteworthy that for several years there seemed to be a paucity of low-luminosity (and therefore, old) observed white dwarfs, which posed a nontrivial puzzle in terms of models for the stellar evolution. Only rather recently it was realized that the properties of the finite temperature atmospheres would make old white dwarfs dimmer than previously anticipated (especially if their atmosphere included only a helium layer and no hydrogen).
#### 2. Pulsation of White Dwarfs
The basic theory of stellar pulsation is reviewed in detail by Sarbani Basu in this volume. We note that studies of pulsations in general allow one to estimate the sound speed profile of the star, its temperature gradient, and even the mean molecular weight in its core. It has been well established that stellar pulsation provides an effective probe for the properties of matter in a star. Commonly referred to as “asteroseismology”, this field is rapidly emerging due significant advances in observational instrumentation and theoretical modeling.
White dwarfs, like many other stars, may undergo stable pulsations which leave a detectable imprint on the time dependence of the star’s luminosity . Since pulsations tend to damp over time (or the star evolves out of the instability strip, if it is overstable to the pulsations), they are most easily observable in younger, and therefore hotter, white dwarfs. Modeling observed pulsations of young white dwarfs is particularly important for probing the EoS of semidegenerate matter . With the aid of theoretical models, studies of white dwarfs pulsations have been used for independently restricting the mass-radius relationships and dimensions of the hydrogen envelope , and in general seems to hold much potential for future studies of the details of white dwarf physics. For example, some pulsation modes are theoretically predicted to be more sensitive to the extent of crystallization of the atmosphere material. Recent analysis of the periods and relative magnitudes various pulsation modes in the white dwarf BPM 37093 helps determine the extent of crystallization in the atmosphere (about 50% in mass), which compares favorably with theory. In another recent example , it is possible to identify, for the first time, the contraction rate of the pre-white dwarf star PG 1159-035, by measuring the secular changes over time of the periods of various modes.
### B. Thermal Effects in Neutron Star Structure
Although finite temperature effects are practically negligible when considering the overall structure of evolved neutron stars, they are important in assessing their thermal history and as probes of matter in the interior of these objects. Some thermal effects on structure do exist in newly-born neutron stars (often called “proto-neutron star”).
#### 1. Proto-Neutron Stars
Neutron stars are born in supernovae with initial internal temperatures of several times $`10^{11}`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ and initial entropies of $`2k_B`$ per nucleon ($`k_B`$ is the Boltzmann constant). These values are large enough to impose nontrivial effects on the composition and structure of the neutron star compared to the $`T=0`$ case . At such high temperatures the matter is not transparent to the thermal neutrinos , and the proto-neutron star must cool through neutrino diffusion to the surface, which occurs over a time scale of several seconds (much longer than a dynamical time of milliseconds) . This fundamental prediction was dramatically confirmed by the observed neutrino pulse from the supernova 1987A (the closest supernova observed from the Earth in 400 years), which lasted for $`12`$ seconds. (see for a review). Further studies of the properties of high temperature, supernuclear-density matter are required to model this brief early cooling epoch and its possible observational signatures .
#### 2. Cooling of Neutron Stars
After the initial rapid cooling stage which lasts several days, a neutron star settles on a slower cooling curve . Once neutrinos escape the system freely, they must be continuously produced in order to fuel the emission process. Subsequent neutrinos are mostly produced by the so-called URCA ($`\beta `$ and inverse $`\beta `$ decays as in Eq. (12)) and other processes. The isothermal core at this stage is expected to have temperature of several times $`10^8`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ , while at the surface the temperature is down to several times $`10^6`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ \- which provides for continuous surface thermal luminosity in soft X-rays. This soft X-ray luminosity is much more difficult to detect than the optical ultraviolet luminosity white dwarfs. Current observations are generally limited to young (age $`<10^6`$yrs) neutron stars, where the main cooling mechanism is still the neutrino emission from the core. An important consequence of this situation is that estimates of thermal emission from neutron stars serve as a probe of the properties of matter in the core .
Over 20 compact, isolated soft X-ray sources observed with X-ray imaging telescopes have been identified as neutron stars . In most cases the source is either too faint to allow for a spectroscopic determination of the surface temperature, or the emission is dominated by other phenomena (most likely magnetospheric emission). However, in a hand-full of cases surface thermal emission has been strong enough to be detected and the surface temperature has been deduced to within a factor of two . The measurements are still not accurate enough to determine whether the observed layer is the actual surface or a possible hydrogen atmosphere, which would have different emission properties, but they are sufficient to place some constraints on the rates of neutrino emitting processes in the core .
Perhaps the most notable conclusion to date of observed neutron star cooling rates is that they support the theoretical prediction that the nucleons in the core couple to a superfluid state. Theoretical models suggest that the strong interactions will pair the neutrons in the core in a $`{}_{}{}^{3}P_{1}^{}`$ superfluid and the protons to a $`{}_{}{}^{1}S_{0}^{}`$ superconductor with critical temperatures of the order of $`10^9`$$`{}_{}{}^{\mathrm{o}}\text{K}`$ . Neutrino emitting processes, such as $`\beta `$-decays must break a coupled Cooper pair before its constituents can participate in the decay. The main effect of nucleon superfluidity is therefore damping the efficiency of neutrino emission from the core (there is also a secondary effect due to the modulation of the heat capacity, which first increases discontinuously as the star cools to the critical temperature, and then decreases exponentially at lower temperatures). Observed neutron star surface temperatures are apparently too high to be consistent with the cooling rates predicted for normal core, but they agree with the suppressed cooling rates when the nucleons couple to a superfluid state . Further progress in analyses of the matter in the interiors of neutron stars is expected through measurements by NASA’s recently launched Chandra X-ray satellite.
## VII. Concluding Remarks
The theoretical and observational study of compact objects remains one of the most exciting fields in modern astronomy. In essence, this research is also an exploration of the properties of condensed matter at extreme densities. Predictions regarding the properties of white dwarfs and neutron stars serve to test our understanding of matter at these high densities, while theories of high density matter serve as a basis for interpreting observational results regarding these objects. Most exciting, these objects bring together all four of the fundamental forces of nature and probe regimes not accessible in the terrestrial laboratory. They provide the most numerous and accessible sample of objects where relativistic gravitation - general relativity - plays a role in determining their physical properties.
In this review we have described the tight interconnection between the microscopic (local) properties of condensed matter at high densities and the macroscopic (global) properties of white dwarfs and neutron stars. While the fundamental principles of cold, high density matter are believed to be well understood, and are generally consistent with observations, key questions still remain, and new observations may give rise to new puzzles. The current boom in capabilities of Earth-bound telescopes and satellite instrumentation promises that many more puzzles - and hopefully, answers - are in store regarding the nature of cosmic matter at high densities.
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# Molecular dynamics of a dense fluid of polydisperse hard spheres
## 1 Introduction
The study of hard spheres has taught us much about the equilibrium and dynamic behaviour of dense fluids and crystals. We have learnt that crystallisation can be a result of the more efficient packing in the crystal of atoms or colloids, and that diffusion is inhibited in dense fluids by repeated collisions of particles with their neighbours. One outstanding problem is that of understanding and perhaps trying to predict, the dramatic slowdown of the dynamics of fluids and the formation of a glass; see Refs. 1, 2, 3 for introductions to slow dynamics and glasses. An obvious thing to do is to study this phenomenon in hard spheres. This has been done and the results are presented here.
Unfortunately, monodisperse hard spheres crystallise readily, preventing the study of hard spheres at densities high enough to observe very slow dynamics. Crystallisation is inhibited if the spheres are not all the same size. Speedy studied a binary mixture of hard spheres of two diameters , which did not crystallise. Here mixtures of spheres with continuous ranges of diameters are studied, these are called polydisperse mixtures. This will enable us to compare with experiments on hard-sphere colloids, which are inevitably on polydisperse particles, . Monte Carlo simulations of polydisperse hard spheres have been performed by Doliwa and Heuer .
In this contribution we will quantify the slow down in the dynamics, characterise these slow dynamics, and assess the effect of having spheres of different diameters present. The slow down is measured by calculating a relaxation time. The dynamics are characterised by calculating their deviation from what we would expect if each particle was diffusing independently. When the deviations are large we call the dynamics heterogeneous. When the dynamics are heterogeneous, the distribution of speeds and relaxation times of individual particles is much broader than the Gaussian function obtained when particles are diffusing independently. There are many more particles which have traveled much farther than the mean than we would expect for independent diffusion of each particle. These fast-moving particles are not distributed randomly in space, they are concentrated in clusters. In other words, fast particles tend to be surrounded by other fast particles .
Three different widths of the distribution of diameters are studied. The polydispersity width is characterised by the standard deviation of the distribution in diameters. For the three distributions this is 9, 14 and 20%. As in experiments on hard-sphere colloids if the distribution is narrower the spheres can crystallise. Surprisingly, although increasing the polydispersity at fixed volume fraction speeds up relaxation, the relaxation time is, within simulation resolution, independent of the width of the distribution at fixed compressibility factor. The relaxation time appears to depend only on the compressibility factor. We also study the effect of the heterogeneity in diameters on the heterogeneity in the dynamics.
In the next section we briefly describe our, standard, methodology. Section 3 outlines when our simulations crystallised and quantifies how much polydispersity in the diameters is required to avoid crystallisation. Section 4 discusses a thermodynamic property, the pressure. The last but one section contains our results for the dynamics, and the last section is a discussion. All results are for fluids at equilibrium not glasses. The focus is on the relaxation of the structure of the fluids, diffusion coefficients and other properties of motion at times large with respect to the relaxation time are not calculated.
## 2 Simulation methodology
Conventional NpT Monte Carlo and NVE molecular dynamics techniques are used . $`N`$, $`p`$, $`V`$, $`E`$, and $`T`$ are the number of particles, the pressure, the volume, the energy and the temperature, respectively. The only nonstandard part concerns the polydispersity of the system, the presence of spheres of different diameters. In theoretical treatments of polydispersity first the thermodynamic limit $`N\mathrm{}`$, is taken and then the number of components present is taken to $`\mathrm{}`$. This gives rise to a continuous distribution of sizes of spheres with a number density $`\rho x(\sigma )\mathrm{d}\sigma `$ of spheres in the range $`\sigma `$ to $`\sigma +\mathrm{d}\sigma `$. $`\rho =N/V`$ is the total number density of spheres. The width of the distribution of diameters can be characterised by a (dimensionless) width parameter $`w`$. The larger $`w`$ is the broader the distribution of sizes present in the mixture. In the limit $`w0`$ we recover a monodisperse system. A simple functional form for the distribution function $`x`$ is the hat function:
$$x(\sigma )=\{\begin{array}{cc}0\hfill & \sigma <\overline{\sigma }(1w/2)\hfill \\ (\overline{\sigma }w)^1\hfill & \overline{\sigma }(1w/2)\sigma \overline{\sigma }(1+w/2)\hfill \\ 0\hfill & \sigma >\overline{\sigma }(1+w/2)\hfill \end{array},$$
(1)
where $`\overline{\sigma }`$ is the mean diameter. The standard deviation divided by its mean, of the distribution $`x(\sigma )`$, $`s=w/\sqrt{12}`$.
Simulations are of course of a finite system — we cannot take the thermodynamic limit in a simulation. For finite $`N`$, $`x(\sigma )`$ cannot be a continuous function. For our simulations we generate a set of $`N`$ diameters by generating a set of $`N`$ pseudo-random numbers in the range $`\overline{\sigma }(1w/2)`$ to $`\overline{\sigma }(1+w/2)`$. Ideally we would average over many different realisations of the polydispersity by performing computer simulations for many different sets of the $`N`$ diameters. However, as our systems have long relaxation times and so require long runs this is not feasible. Thus all the results we will present for a given polydispersity width will be for a single system of 1372 particles with one set of the 1372 diameters. We have performed some simulations with different sets of the diameters but the same number of particles and the results did not change significantly. We have also performed simulations with 256 and 500 particles and obtained results which did not differ significantly with the exception of some results for the dynamics. We discuss these differences below.
NpT Monte Carlo simulations are used to start with and to compress the fluid phase up to the required density. NVE molecular dynamics are used to equilibrate at a particular density and to obtain all the averages shown, including the pressure. All simulation results are for fluids at (possibly metastable) equilibrium, not for glasses. Our results are in reduced units. We use the mean diameter $`\overline{\sigma }`$ as our length scale. The reduced pressure $`p_r`$ is then $`p_r=Z\rho \overline{\sigma }^3`$, where $`\rho `$ is the number density of particles and $`Z=pV/(NkT)`$ is the compressibility factor. The results are all averages over 8 runs. We use the variation between the results of the individual runs to estimate statistical errors. In no case are they significantly larger than the plot symbols used in the figures.
A time scale can obtained by the ratio of a distance to a velocity. The velocity we choose is the root mean square velocity along one of the 3 Cartesian axes, of a particle with a diameter equal to the mean diameter. This velocity is equal to $`(kT/\overline{m})^{1/2}`$, where $`\overline{m}`$ is the mass of particle with diameter $`\overline{\sigma }`$. The time scale is thus $`\overline{\sigma }(\overline{m}/kT)^{1/2}`$ and all times will be made dimensionless by dividing by this. For molecular dynamics we have to specify the masses of the particles. We assume that the mass of a sphere scales with the cube of its radius, correct if all spheres are made from the same substance with the same uniform density. The density we use is the volume fraction $`\eta =(\pi /6)\rho \overline{\sigma ^3}`$, where $`\overline{\sigma ^3}`$ is the third moment of $`x(\sigma )`$.
## 3 Crystallisation
Although we will present results only for systems which showed no sign of crystallising, some of our systems did start to crystallise. Speedy has undertaken a study of the freezing of monodisperse hard spheres within molecular dynamics simulations. He finds crystallisation only above a volume fraction of 0.54. We simulated 1372 monodisperse hard spheres. They were stable for a time much longer than the relaxation time of the fluid at a volume fraction $`\eta =0.54`$ but crystallised at $`\eta =0.56`$. Thus we agree with the findings of Speedy. Simulations of 1372 spheres with a polydispersity width of $`w=0.1`$ also crystallised but when the width was increased to $`w=0.2`$ ($`s=5.8\%`$) the fluid was stable, it never crystallised. Simulations at polydispersity widths $`w=0.3`$, 0.5 and 0.7, of both 256 and 1372 spheres, never showed any sign of crystallisation. They remained amorphous up to the highest densities simulated. Widths $`w=0.3`$, 0.5 and 0.7 correspond to values of $`s`$ of 8.7%, 14.4% and 20.2%, respectively.
Experiments on colloids find no crystallisation for polydispersities beyond a standard deviation $`s`$ close to 8% . Thus our finding of crystallisation when the standard deviation is 3% ($`w=0.1)`$ but not when it is 6% ($`w=0.2`$) is consistent with experiment. However, our simulations are for systems much much smaller than experiment and the time elapsed until nucleation occurs will depend on system size. Theoretical and computer simulation studies of the equilibrium phase behaviour of polydisperse hard spheres show that polydispersity destabilises the crystalline phase, pushing the fluid-crystal transition to higher volume fractions.
Our systems with polydispersity widths of $`w=0.2`$, 0.3, 0.5 and 0.7 do not crystallise even at volume fractions up to 0.58 because either the fluid phase is still the equilibrium phase even at these high densities or it is metastable with respect to fluid-crystal coexistence but the fluid is not sufficiently deep into the coexistence region for the barrier to nucleation of the crystalline phase to be small . The nucleation rate varies as the inverse of the product of the relaxation time and the exponential of the free energy barrier . If the barrier is large up to densities at which the relaxation time is very large, say $`\eta =0.57`$-0.58, then the rate will always be very small. Perhaps too small to be observed.
Theoretical studies suggest that for polydispersity widths of around 0.3 and greater, the crystalline phase cannot form with the full width of the distribution. It can only form if it accepts only a narrow range of diameters. Neither in experiments nor here is the formation of a crystalline phase with a narrow range of diameters observed to form from a polydisperse mixture of spheres. It is not clear whether or not the fact that the crystalline phase can only form from a fraction of the distribution of diameters is a cause of the fact that crystallisation is not observed.
We checked for segregation of the large and small spheres, i.e., for a tendency to phase separate into a phase of large spheres coexisting with a phase of small spheres. There was no sign of such segregation. Segregation is not expected at the polydispersity widths studied here in the fluid phase. The Boublik-Mansoori-Carnahan-Starling-Leland (BMCSL) theory as generalised to polydisperse hard spheres by Salacuse and Stell does however predict phase separation at much broader distributions of diameters . This fluid-fluid separation is metastable with respect to phase separation plus crystallisation of the large spheres .
## 4 Pressure
The reduced pressure is plotted as a function of volume fraction in Fig. 1. Results from simulation are shown along with the predictions of the BMCSL expression for the pressure . The agreement between the BMCSL equation and simulation is excellent. The only noticeable deviation is that the BMCSL is a little too high at the highest densities for $`w=0.3`$ and 0.5.
The pressure tends to decrease as the polydispersity width increases. The random-close-packing density of polydisperse hard spheres increases as the polydispersity width increases . The pressure diverges at random-close-packing. Therefore, at constant $`\eta `$ we move farther from the density where the pressure diverges as the polydispersity width increases and so it is unsurprising that the pressure decreases. The distribution used in Refs. 33, 34 was different to the hat function used here but the trend should be relatively insensitive to the exact form of the distribution when the polydispersity is not too great.
## 5 Dynamics
As we are considering dense states the motion of a particle is strongly restricted by its neighbours. Over some timescale a particle will rattle back and fore in a cage formed by its neighbours; then as we are considering a fluid not a glass, the particle will ‘break out’ of the cage and eventually its motion will be diffusive. We can assess whether or not particles are just rattling back and fore or are moving significant distances, i.e., distances comparable to their own diameter, by examining the intermediate scattering function, sometimes called the self incoherent intermediate scattering function, $`F_s(q,t)`$. As indicated it is a function of wavevector $`q`$ and time $`t`$, it is also a function of volume fraction and polydispersity. It is defined by
$$F_s(q,t)=\frac{1}{N}\underset{i=1}{\overset{N}{}}\mathrm{cos}(𝐪.[𝐫_i(t)𝐫_i(0)]),$$
(2)
where $`𝐫_i(t)`$ is the position vector for the $`i`$th particle at time $`t`$. For the fluids we consider, it is not a function of the initial time, taken to be $`t=0`$ above, only of the time elapsed, $`t`$. For our isotropic systems it is also only a function of the magnitude $`q`$ of the wavevector $`𝐪`$, not of its orientation. To improve statistics we average over wavevectors along the $`x`$, $`y`$ and $`z`$ axes. Clearly $`F_s(q,0)=1`$ and at times $`t`$ sufficiently long that the positions of the particles are no longer correlated with their positions at time $`t=0`$, then $`F_s(q,t)=0`$. As the particles move, their positions become decorrelated with their positions at $`t=0`$ and $`F_s`$ decays to 0. For a fixed $`q`$, $`F_s(q,t)`$ will become small when the particles have moved a distance of about $`\pi /(2q)`$. All our results are for $`q=2\pi /\overline{\sigma }`$, and so $`F_s`$ will become small when most of the particles have moved a distance of around a quarter of their diameter. $`F_s`$ enables us to define a relaxation or correlation time for the dynamics, $`\tau `$. We define $`\tau `$ by
$$F_s(2\pi /\overline{\sigma },\tau )=1/\mathrm{e}.$$
(3)
Values of $`\tau `$ for $`w=0.3`$, 0.5 and 0.7 are shown in Fig. 2. As expected $`\tau `$ increases rapidly as the volume fraction increases; by about 2 orders of magnitude between $`\eta =0.5`$ and 0.58. It is also clear that at constant volume fraction, increasing the polydispersity decreases the relaxation time $`\tau `$. However, if the relaxation time is plotted as a function of the compressibility factor $`Z`$, Fig. 3, then the data for the 3 polydispersity widths essentially follow the same curve. Note that in Fig. 3 we have included data for monodisperse spheres (at densities below where they crystallise) and shown results down to $`\eta =0.3`$. Spheres with a polydispersity width of $`w=0.7`$ at $`\eta =0.57`$ have almost the same $`Z`$ and almost the same relaxation time as spheres with $`w=0.3`$ at $`\eta =0.56`$. These are the two points which are almost superposed in Fig. 3. To a good approximation the relaxation time as defined by Eq. (3) is a function only of the compressibility factor $`Z=pV/(NkT)`$. It is not obvious to the author why $`\tau `$ should depend only on $`Z`$ although there are theories which attempt to relate dynamic quantities such as $`\tau `$ to thermodynamic quantities . If we plot $`\tau `$ as a function of the reduced pressure $`p_r`$ (not shown) then the results for the 3 polydispersity widths do not fall on the same curve, in particular the results for $`w=0.7`$ are above those for the narrower distributions.
We should note that for $`\eta =0.56`$ and above the relaxation times calculated using 256 spheres are consistently above those calculated with 1372 spheres. For $`w=0.3`$ and $`\eta =0.57`$ the relaxation times are 77 and 110 for 1372 and 256 spheres, respectively. This suggests that the result for 1372 spheres is 10-30% higher than that of an infinite system. When $`\tau `$ is large, we observe finite size effects for the dynamics but not for the statics. This seems to be a quite general finding . It suggests that there is some length scale associated with the dynamics but not the statics which grows as the relaxation time $`\tau `$ increases. Recently this possibility has been extensively studied by Glotzer and coworkers . They define a purely dynamical correlation length and observe that it increases as the dynamics slow down. Perhaps the easiest way to see the correlations in the dynamics is to plot out the fastest and slowest moving particles. We simulate the particles with $`w=0.3`$ at the highest volume fraction studied for this polydispersity, $`\eta =0.57`$. Fig. 4 shows the fastest 5% of the 1372 particles, while Fig. 5 shows the slowest 5%. Fastest and slowest are defined as having the largest or smallest displacement over a time interval of length 29.4. This is a little less than half the relaxation time $`\tau `$. Clearly, neither the fastest nor the slowest particles are randomly distributed, both show strong clustering — the dynamics are highly correlated. This has been observed in previous simulation studies .
So over times of order $`\tau `$ the positions of the fastest particles are correlated. The movement of the particles is very far from being independent diffusion of each particle. If the particles were diffusing independently then for each particle, the probability of finding the particle at a point would be a Gaussian function of the distance between the point and the position of the particle at $`t=0`$. As they are not diffusing independently it will be not be a Gaussian. We can measure deviations of the distribution of particles from a Gaussian by using the non-Gaussian parameter introduced by Rahman and defined by
$$\alpha (t)=\frac{3N^1_{i=1}^N[𝐫_i(t)𝐫_i(0)]^4}{5\left(N^1_{i=1}^N[𝐫_i(t)𝐫_i(0)]^2\right)^2}1.$$
(4)
For a Gaussian distribution the ratio of the fourth moment to the square of the second is 5/3 and $`\alpha =0`$. In Fig. 6 we have plotted $`\alpha `$ as a function of time for a number of volume fractions. For monodisperse particles the Maxwell-Boltzmann distribution of velocities enforces $`\alpha =0`$ in the short time, ballistic, regime, and at times much longer than $`\tau `$ we have diffusion and again $`\alpha =0`$. For polydisperse particles $`\alpha `$ is nonzero in the ballistic regime due to different particles having different masses, and nonzero at very long times due to larger particles having smaller diffusion coefficients than smaller particles. So, in Fig. 6 $`\alpha `$ is always nonzero. The value of $`\alpha `$ in the ballistic limit depends only on the polydispersity width $`w`$, its value in the other limit comes from the spread in diffusion coefficients with particle size and so will also depend on density. Positive values of $`\alpha `$ come from a large fourth moment, due to more particles with large displacements than would be found for a Gaussian distribution with the same second moment. It is clear if we compare Fig. 6 with the relaxation times of Fig. 2 that the non-Gaussian parameter $`\alpha `$ goes through a maximum at a time comparable to but (at least at the highest densities) less than the relaxation time $`\tau `$. This maximum value increases rapidly as $`\tau `$ increases. The deviations from a Gaussian are much larger at a volume fraction of 0.57 than at 0.5.
Recently, experiments have imaged colloidal suspensions of hard-sphere-like particles at the high volume fractions considered here . The suspensions are of particles with somewhat smaller polydispersities than considered here, 8% , 6% and 5% . The dynamics of colloidal particles , are rather different from those of free particles such as our hard spheres. Colloidal particles are immersed in fluid and so their motion is diffusive on all time scales on which the particles move a significant distance. In addition, the motion of nearby particles is coupled by hydrodynamic interactions through the fluid. If we compare the $`\alpha `$ values measured by Weeks et al. with those of Fig. 6 we find that the simulation and experimental results are rather different. The experimental values of $`\alpha `$ are consistently larger than those in Fig. 6, Weeks et al. measure a maximum $`\alpha `$ at $`\eta =0.52`$ of about 1.5 and at $`\eta =0.56`$ of a little under 2.5. The dynamics on a time scale of the relaxation time $`\tau `$ must be rather different for colloidal particles and free particles, although in both cases the relaxation time increases very rapidly and in both cases glasses can form. The results of Kegel and van Blaaderen for $`\alpha `$ are also larger than our simulation results; those of Kasper et al. are much larger, at $`\eta =0.56`$ the maximum value of $`\alpha `$ is over 5.
Two obvious sources of the difference between molecular dynamics simulation and experiment are: i) the motion of colloidal particles is diffusive even on time scales much less than $`\tau `$, and ii) hydrodynamic interactions between particles. i) and ii) could be distinguished by performing both molecular dynamics and Brownian dynamics of the same potential . Brownian dynamics is an approximation to the dynamics of particles in a fluid which accounts for the fact that particles diffuse even over very small length scales but neglects hydrodynamic interactions between particles .
Another way to look at the large spread in the amount individual particles are moving is to define $`F_s`$’s for each particle, and then use this to obtain a relaxation time for each particle. Averaging over wavevectors along the $`x`$, $`y`$ and $`z`$ axes we have for the $`i`$th particle
$`{\displaystyle \frac{1}{3}}\{\mathrm{cos}(q[x_i(\zeta _i)x_i(0)])+\mathrm{cos}(q[y_i(\zeta _i)y_i(0)])+`$
$`\mathrm{cos}(q[z_i(\zeta _i)z_i(0)])\}=1/\mathrm{e},`$ (5)
which defines the relaxation time for the $`i`$th particle, $`\zeta _i`$. As always $`q=2\pi /\overline{\sigma }`$. In Fig. 7 we have plotted the probability function $`P(\zeta )`$, where $`P(\zeta )\mathrm{d}\zeta `$ is the probability that a particle will have a relaxation time between $`\zeta `$ and $`\zeta +\mathrm{d}\zeta `$. For both $`\eta =0.5`$ and 0.56 there is a peak at short times and then $`P`$ decays. The decay is clearly close to exponential, which makes sense as (at not-too-short times) the probability of a particle leaving its cage in some short time interval should be independent of time; at equilibrium the environment of a particle or cluster of particles is on average not changing with time. For $`\eta =0.5`$ the peak is a little below $`\tau `$ at that density but for $`\eta =0.56`$ the peak is at about one tenth of $`\tau `$. At the higher density the probability distribution of $`\zeta `$ is very broad.
### 5.1 Polydispersity
The spheres are polydisperse, different spheres have different diameters and so their static and dynamic properties will differ. We expect the smaller spheres to move larger distances than larger spheres in the same time. At very short times, in the ballistic regime, this is trivially true due to the smaller mass and hence higher velocity of the smaller spheres. We can look at motion over a time scale of order $`\tau `$ by examining the sizes of the fastest and slowest 5% of the spheres. Figs. 4 and 5 show these spheres for $`w=0.3`$. We average over 8 runs with an average time interval close to that in Figs. 4 and 5. The average diameter of the fastest 5% (69) of the spheres, over a time interval of 29.2, is $`0.95\overline{\sigma }\pm 0.01\overline{\sigma }`$, and for the slowest 5% it is $`1.01\overline{\sigma }\pm 0.01\overline{\sigma }`$. The effect is small but as we would expect the fastest spheres are smaller than average. For a polydispersity width of 0.7, over a time interval of 5.88, the average diameter of the fastest 5% is $`0.79\overline{\sigma }\pm 0.02\overline{\sigma }`$, and for the slowest 5% it is $`1.08\pm 0.02\overline{\sigma }`$. A time of 5.88 is, as with the less polydisperse spheres, somewhat less than half $`\tau `$ and so near the maximum in $`\alpha `$. Of course the effect is now larger but still the slowest spheres are not much larger than the mean, their mean diameter would have been above 1.3 if all the slowest spheres were also the largest. The difference between the mean diameter of the fastest spheres and the overall mean diameter is twice the difference between the mean diameter of the slowest spheres and the overall mean diameter.
We have defined an intermediate scattering function, $`F_s`$, for all spheres regardless of diameter in Eq. (2). We can also define an intermediate scattering function for subsets of the particles with diameters in some range $`\sigma _{min}`$ to $`\sigma _{max}`$, $`F_s(q,t;\sigma _{min},\sigma _{max})`$. The definition is completely analogous to that of the total $`F_s(q,t)`$ in Eq. (2),
$`F_s(q,t;\sigma _{min},\sigma _{max})=`$
$`{\displaystyle \frac{1}{N_{mm}}}{\displaystyle \underset{i=1}{\overset{N_{mm}}{}}}\mathrm{cos}(𝐪.[𝐫_i(t)𝐫_i(0)]),`$
(6)
where the sum is over all spheres with diameters in the range $`\sigma _{min}`$ to $`\sigma _{max}`$, and $`N_{mm}`$ is the number of spheres in this range. As before $`q=2\pi /\overline{\sigma }`$ and we average over wavevectors along the $`x`$, $`y`$ and $`z`$ axes. We will study intermediate scattering functions for the smallest, $`F_s^{(s)}(q,t)`$, and largest spheres, $`F_s^{(l)}(q,t)`$, defined by
$`F_s^{(s)}(q,t)`$ $`=`$ $`F_s(q,t;\overline{\sigma }(1w/2),\overline{\sigma }(1w/2+0.05))`$
$`F_s^{(l)}(q,t)`$ $`=`$ $`F_s(q,t;\overline{\sigma }(1+w/20.05),\overline{\sigma }(1+w/2)).`$
These are for the particles within $`0.05\overline{\sigma }`$ of the minimum or maximum diameter. For $`w=0.7`$ this corresponds to spheres with diameters in the ranges 0.65 to $`0.7\overline{\sigma }`$ and 1.3 to $`1.35\overline{\sigma }`$. We also define relaxation times $`\tau ^{(s)}`$ and $`\tau ^{(l)}`$ from $`F_s^{(s)}(q,t)`$ and $`F_s^{(l)}(q,t)`$, respectively, using the analogues of Eq. (3). $`\tau ^{(s)}`$ and $`\tau ^{(l)}`$ for $`w=0.7`$ are plotted as filled triangles in Fig. 2. The ratio $`\tau ^{(l)}/\tau ^{(s)}`$ increases with increasing density but slower than $`\tau `$ does. At $`\eta =0.57`$ the ratio is approximately 10.
A feature of the dynamics when $`\tau `$ is large is that they are heterogeneous, some particles travel much farther than others in a time $`\tau `$. In a fluid of monodisperse particles these heterogeneities are dynamic, a particle may be fast at one time but later on may be slow. For polydisperse particles in addition to this source of heterogeneity there is the spread in diameters which means that some particles, the smaller ones, are on average faster than others over all time scales. In Fig. 6 we can compare the $`\alpha `$ parameter for $`w=0.3`$ and 0.7 at the same volume fraction $`\eta =0.57`$. $`\alpha `$ is larger for the more polydisperse spheres; this is the case despite the fact that the relaxation time for the more polydisperse particles is smaller. The system with $`w=0.7`$ and $`\eta =0.57`$ is at almost the same compressibility factor as the one with $`w=0.3`$ and $`\eta =0.56`$ so we see that at constant compressibility factor $`\alpha `$ increases sharply with polydispersity.
For weakly polydisperse spheres the relaxation times for all sizes of sphere increase together and so if we are at high enough density then $`\tau `$ will be very large and on a timescale much less than $`\tau `$ all spheres will be localised, none or at least very few will have moved more than a small fraction of their diameter. However, if spheres with a wide range of diameters are present then the smallest particles will have a relaxation time scale which is much less than $`\tau `$. Then over some time scale much less than $`\tau `$ but much larger than the relaxation time scale for the smallest particles not all the spheres will be localised. The smallest spheres will be diffusing. This has been observed in binary mixtures of hard spheres in which the smaller spheres are much smaller than the larger spheres, by Jackson et al. . The same effect is seen in crystals of binary mixtures . This is an extreme example of polydispersity making the dynamics even more heterogeneous than they are when the particles are all the same size. Note that polydispersity and the heterogeneous nature of the dynamics when $`\tau `$ is large have the same effect, they produce a wide spread in speeds/relaxation times of the individual particles. Even for monodisperse particles, for times much less than $`\tau `$ the heterogeneous dynamics mean that some particles have left their cages. However, if we have some particles much smaller than others then the small ones may have a relaxation time not much larger and a diffusion constant not much smaller than at lower densities.
## 6 Discussion
We have studied the dynamics of dense polydisperse hard spheres using molecular dynamics. The relaxation time $`\tau `$ at a given volume fraction was found to depend on the polydispersity, it decreased as the polydispersity increased. So, the glass transition is pushed to higher volume fractions as the spheres become more polydisperse. This is consistent with the increase in the density of random-close-packing with increasing polydispersity . Although we have found this using molecular dynamics, we expect that experiments on polydisperse colloids would show that as polydispersity increased the kinetic glass transition observed in experiment would move to higher volume fractions.
The relaxation time $`\tau `$ is, to a good approximation, a function only of the compressibility factor $`Z`$, changing the polydispersity at fixed $`Z`$ has no effect on $`\tau `$. This is despite the fact that a characteristic of the dynamics, such as $`\alpha `$, changes a great deal at fixed $`Z`$ as $`w`$ increases; compare the $`w=0.3`$, $`\eta =0.56`$ and $`w=0.7`$, $`\eta =0.57`$ curves in Fig. 6, which have almost the same $`Z`$ but very different $`\alpha `$’s. Apparently, increasing the polydispersity at constant volume fraction reduces the relaxation time by reducing $`Z`$. The virial equation ,
$$Z=1\frac{1}{6NkT}\underset{ij}{}𝐫_{ij}.𝐟_{ij}$$
(8)
relates $`Z`$ to the forces between the particles. $``$ denotes an ensemble average. In Eq. (8) $`𝐫_{ij}`$ and $`𝐟_{ij}`$ are the vector between the centres of the $`i`$th and $`j`$th particles and the force on the $`i`$th particle due to the $`j`$th particle, respectively. The sum is over all pairs of particles. So, $`Z`$ is a constant, 1, plus a term proportional to the sum over the product of the interparticle forces and the interparticle separations. Why the relaxation time should be a function only of this product is not clear to the author.
Although $`Z`$ can be expressed in terms of forces it is also of course a thermodynamic quantity. At constant temperature $`Z`$ varies as $`p/\rho `$ so $`\tau `$ being a function of $`Z`$ is equivalent to it being a function of the ratio of the pressure to the number density. As the energy of hard spheres is purely kinetic it is independent of volume, thus the pressure is simply the volume derivative of the entropy times the temperature. So, the relaxation time is a function only of $`Z`$ and so is a function only of the ratio of volume derivative of the entropy to the number density. But it is not a simple function of it; $`\tau `$ does not vary as the exponential of the $`Z`$, it increases more rapidly. The pressure shows no sign of a discontinuity in slope so we conclude that there is no phase transition in the density range studied here. It is an open question whether or not there is a phase transition to an ‘ideal’ glass phase at higher densities, as assumed by Speedy , see also Refs. 1, 3.
We have characterised the dynamics using the non-Gaussian parameter $`\alpha `$, Fig. 6, and the distribution of relaxation times $`P(\zeta )`$, Fig. 7. The results for $`\alpha `$ have been compared to experimental results for colloidal suspensions . The dynamics are heterogeneous by which we mean that there is a broad distribution of speeds and of relaxation times of the particles. This is true not only at times of order the relaxation time $`\tau `$ but also, at higher density and hence $`\tau `$, for times at least an order of magnitude larger, see Fig. 7. Indeed given the very wide spread of individual relaxation times $`\zeta `$ it is clear that the overall relaxation time $`\tau `$ is inadequate to characterise the time dependence of the dynamics. In particular, relaxation is far from complete even at times much longer than $`\tau `$. Comparison of our non-Gaussian parameter with that obtained in experiments on colloidal hard spheres showed significant differences. One possible explanation of this is that hydrodynamic interactions between colloidal particles are acting to make the relaxation more cooperative, i.e., acting to increase the clustering shown in Fig. 4. It is known that hydrodynamic interactions tend to favour collective over relative motion .
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# Magnetic Coherence in Cuprate Superconductors
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## Abstract
Recent inelastic neutron scattering (INS) experiments on La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> observed a magnetic coherence effect, i.e., strong frequency and momentum dependent changes of the spin susceptibility, $`\chi ^{\prime \prime }`$, in the superconducting phase. We show that this effect is a direct consequence of changes in the damping of incommensurate antiferromagnetic spin fluctuations due to the appearance of a d-wave gap in the fermionic spectrum. Our theoretical results provide a quantitative explanation for the weak momentum dependence of the observed spin-gap. Moreover, we predict (a) a Fermi surface in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> which is closed around $`(\pi ,\pi )`$ up to optimal doping, and (b) similar changes in $`\chi ^{\prime \prime }`$ for all cuprates with an incommensurate magnetic response.
\]
The spin excitation spectrum in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> in the normal and superconducting state has been intensively studied during the last few years in inelastic neutron scattering (INS) experiments. The normal state spectrum for compounds with $`x>0.04`$ is characterized by peaks in $`\chi ^{\prime \prime }(𝐪,\omega )`$ at incommensurate wave-vectors $`𝐐_i=(1\pm \delta ,1)\pi `$ and $`𝐐_i=(1,1\pm \delta )\pi `$ , where $`\delta `$ increases with increasing doping. Recent INS experiments in the superconducting state of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> (LSCO) by Mason et al. (x=0.14) and Lake et al. (x=0.16) show striking momentum and frequency dependent changes in $`\chi ^{\prime \prime }`$ upon entering the superconducting state, which the authors called the magnetic coherence effect. For both compounds, $`\chi ^{\prime \prime }(𝐐_i)`$ in the superconducting state is considerably decreased from its normal state value below $`\omega 7`$ meV, while it increases above this frequency. For frequencies in the vicinity of 7 meV, the incommensurate peaks sharpen in the superconducting state, while at higher frequencies the peak widths in the normal and superconducting state are approximately equal. Moreover, by employing a Kramers-Kronig transformation, the authors found that the static susceptibility, $`\chi ^{}`$, at $`𝐐_i`$ decreases in the superconducting state .
In this communication, we show that the magnetic coherence effect is a direct consequence of changes in the damping of incommensurate antiferromagnetic spin fluctuations due to the appearance of a d-wave gap in the fermionic spectrum. We obtain results for the frequency and momentum dependence of $`\chi ^{\prime \prime }`$ that are in good qualitative, and to a large extent quantitative agreement with the experimental data, and also explain the weak momentum dependence of the spin-gap. We show that INS data in the superconducting state provide information on the symmetry of the order parameter and the topology of the Fermi surface (FS) and that for La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, INS experiments suggest a FS closed around $`(\pi ,\pi )`$. We predict that the magnetic coherence effect is to be expected for any cuprate superconductor with an incommensurate spin spectrum, and thus in particular for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> (YBCO), in which an incommensurate spin structure at low frequencies has been observed .
The starting point for our calculations is a spin-fermion model in which the damping of incommensurate spin-excitations arises from their interaction with fermionic quasi-particles. In this model, the spin propagator, $`\chi `$, is given by
$$\chi ^1=\chi _0^1\mathrm{\Pi },$$
(1)
where $`\chi _0`$ is the bare propagator, and $`\mathrm{\Pi }`$ is the bosonic self-energy given by the irreducible particle-hole bubble. $`\chi _0`$ is in general obtained by integrating out the high-energy fermionic degrees of freedom. However, since the form of fermionic excitations at high frequencies is so far not well understood, a microscopic calculation of $`\chi _0`$ is not yet feasible. We therefore make the experimentally motivated ansatz
$$\chi _0^1=\frac{\xi _0^2+(𝐪𝐐_i)^2}{\alpha },$$
(2)
where $`\xi _0`$ is defined as the “bare” magnetic correlation length (unrenormalized by the coupling to low-frequency particle-hole excitations) and $`\alpha `$ is a temperature independent constant. In general one would expect a frequency term in Eq.(2) which is omitted here because experimentally there is no observed dispersion in the spin excitation spectrum below $`\omega 25`$ meV , well above the frequency range we consider here. The above form of $`\chi _0`$ thus only determines the position of the incommensurate peaks in momentum space; it does not affect the frequency dependence of $`\chi ^{\prime \prime }`$, which arises solely from $`\mathrm{\Pi }`$. In the following we define the renormalized magnetic correlation length as $`\xi ^2=\xi _0^2\alpha \mathrm{Re}\mathrm{\Pi }`$.
We first consider $`\chi ^{\prime \prime }`$ at $`𝐐_i`$ in the normal state. Calculating $`\mathrm{\Pi }_N`$ to lowest order in the spin-fermion coupling $`g`$ yields Im$`\mathrm{\Pi }_N\omega `$, while Re$`\mathrm{\Pi }_N`$ const. , yielding $`\xi _N^2(\omega )=const`$, and a frequency dependent dynamic susceptibility, $`\chi ^{\prime \prime }(𝐐_i,\omega )`$, of the MMP form which quantitatively describes the results of INS experiments in the normal state of LSCO and YBCO .
In the superconducting state, $`\mathrm{\Pi }_{SC}`$ is given by (to lowest order in $`g`$)
$`\mathrm{\Pi }_{SC}(𝐪,i\omega _n)`$ $`=`$ $`g^2T{\displaystyle \underset{𝐤,m}{}}\{G(𝐤,i\mathrm{\Omega }_m)G(𝐤+𝐪,i\mathrm{\Omega }_m+i\omega _n)`$ (4)
$`+F(𝐤,i\mathrm{\Omega }_m)F(𝐤+𝐪,i\mathrm{\Omega }_m+i\omega _n)\},`$
where $`G`$ and $`F`$ are the normal and anomalous Green’s functions
$`G`$ $`=`$ $`{\displaystyle \frac{i\omega _n+ϵ_k}{(i\omega _n)^2ϵ_k^2\mathrm{\Delta }_k^2}},F={\displaystyle \frac{\mathrm{\Delta }_k}{(i\omega _n)^2ϵ_k^2\mathrm{\Delta }_k^2}}.`$ (5)
$`E_𝐤=\sqrt{ϵ_𝐤^2+|\mathrm{\Delta }_𝐤|^2}`$ is the fermionic dispersion in the superconducting state, $`\mathrm{\Delta }_𝐤=\mathrm{\Delta }_0\left(\mathrm{cos}(k_x)\mathrm{cos}(k_y)\right)/2`$ is the d-wave gap and
$`ϵ_𝐤`$ $`=`$ $`2t\left(\mathrm{cos}(k_x)+\mathrm{cos}(k_y)\right)`$ (7)
$`4t^{}\mathrm{cos}(k_x)\mathrm{cos}(k_y)\mu ,`$
is the electronic tight-binding dispersion where $`t,t^{}`$ are the hopping elements between nearest and next-nearest neighbors, respectively, and $`\mu `$ is the chemical potential. Since our theoretical results for both doping levels of LSCO $`x=0.14(0.16)`$ are quantitatively similar, we consider for definiteness $`x=0.16`$ and choose $`t^{}/t=0.22`$ and $`\mu /t=0.84`$, a choice which will be seen to yield agreement with the experimental data. The superconducting gap, $`\mathrm{\Delta }_010`$ meV, is taken to be that extracted from Raman scattering experiments by Chen et al. and the incommensurate wave-vector $`𝐐_i`$ is at $`\delta 0.25`$ .
Our theoretical results for the spin-damping, Im$`\mathrm{\Pi }`$ in Eq.(4) at $`𝐐_i`$ are presented in Fig. 1. Since $`𝐐_i`$ is incommensurate, we obtain four decay channels for spin excitations; the two channels in the first Brillouin zone are shown in the inset of Fig. 1. In the normal state all four channels for particle-hole excitations are excited in the low frequency limit, which yields Im$`\mathrm{\Pi }_N\omega `$, as noted above. In the superconducting state, the four channels split into two pairs with degenerate non-zero threshold energies, $`\omega _c^{(1,2)}`$, that are determined by the momentum dependence of the order parameter and the shape of the Fermi surface. In particular, we find $`\omega _c^{(1,2)}=|\mathrm{\Delta }_𝐤|+|\mathrm{\Delta }_{𝐤+𝐐_𝐢}|`$, where $`𝐤`$ and $`𝐤+𝐐_i`$ both lie on the Fermi surface, as shown in the inset of Fig. 1. For the band parameters chosen, the threshold energies are $`\omega _c^{(1)}=0.70\mathrm{\Delta }_{SC}`$ for quasiparticle excitations close to the nodes of the superconducting gap (excitation 1), and $`\omega _c^{(2)}=1.86\mathrm{\Delta }_{SC}`$ for excitations which connect momenta around $`(0,\pi )`$ and $`(\pi ,0)`$ (excitation 2). Note that due to the superconducting coherence factors in Eq.(4), Im$`\mathrm{\Pi }_{SC}`$ exhibits sharp jumps at the threshold frequencies. Since for $`T=0`$ and $`\omega <\omega _c^{(1)}`$, Im$`\mathrm{\Pi }_{SC}0`$, and thus $`\chi _{SC}^{\prime \prime }0`$, $`\omega _c^{(1)}`$ is often referred to as the spin-gap in the superconducting state.
We now turn to the calculation of Re$`\mathrm{\Pi }_{SC}`$. The gap in Im$`\mathrm{\Pi }_{SC}`$ gives rise to a $`\omega ^2`$-term in Re$`\mathrm{\Pi }_{SC}`$ for $`\omega \omega _c^{(1)}`$, while Re$`\mathrm{\Pi }_{SC}const.`$ for $`\omega \omega _c^{(2)}`$. The steps in Im$`\mathrm{\Pi }_{SC}`$ at $`\omega _c^{(1,2)}`$ create logarithmic divergences in Re$`\mathrm{\Pi }_{SC}`$; as has recently been demonstrated these are an artifact of our restriction to the second order bosonic self-energy correction. When fermionic lifetimes are calculated within a self-consistent strong-coupling approach, the authors of Ref. found that the steps in Im$`\mathrm{\Pi }_{SC}`$ are smoothed out, while the gap below a frequency $`\omega _c^{(1)}`$ still persists. The weak logarithmic divergences in Re$`\mathrm{\Pi }_{SC}`$ become a smooth function of frequency. Since the spin-gap survives the inclusion of realistic fermionic lifetimes, we expect the conclusions we draw in the following to be valid beyond the current level of approximation.
In Fig. 2 we present a fit of our theoretical results for $`\chi ^{\prime \prime }`$ to the experimental data of Ref. in the normal and superconducting state.
The fit to the experimental data in the superconducting state for $`\omega _c^{(1)}<\omega <16`$ meV was obtained by making the ansatz that at these frequencies, $`\xi _{SC}(\omega )`$ is frequency independent and given by
$$\xi _{SC}^2(\omega )\frac{3}{2}\xi _N^2=const.$$
(8)
To account for the experimental energy resolution we convolute our theoretical results with a Gaussian distribution of width $`\sigma 2`$ meV.
The “high frequency” result, Eq.(8), is a consequence of the redistribution in the spectral weight of $`\chi ^{\prime \prime }(𝐐_i,\omega )`$ in the superconducting state. By using the Kramers-Kronig relation, we find $`\xi _{SC}(\omega =0)<\xi _N(\omega =0)`$ while, as seen in Fig. 2 quite generally, for $`\omega _c^{(1)}<\omega <\omega _c^{(2)}`$, $`\chi _{SC}^{\prime \prime }(\omega )`$ exceeds $`\chi _N^{\prime \prime }(\omega )`$, so that $`\xi _{SC}(\omega )>\xi _N(\omega )`$ in this frequency range. As may be seen in Fig. 2, the simple ansatz, Eq.(8), yields good agreement with experiment. Because the form of both $`\chi _{SC}^{\prime \prime }(\omega )`$ and $`\chi _N^{\prime \prime }(\omega )`$ above $`\omega _c^{(2)}`$ is not well known, one cannot at present arrive at a self-consistent description of the frequency dependence of $`\xi _{SC}`$ using the Kramers-Kronig relation. We note, however, that upon restricting the frequency integration in the Kramers-Kronig relation to $`\omega <\omega _c^{(2)}`$, we find $`\chi _{SC}^{}(\omega =0)0.65\chi _N^{}(\omega =0)`$, in agreement with the results of Ref. .
To understand the momentum dependent changes in $`\chi ^{\prime \prime }`$ between the normal and superconducting state, we consider the momentum dependence of the spin-gap, $`\omega _c^{(1)}(𝐪)`$. In Fig. 3 we plot the experimental intensity in the $`(\omega ,𝐪)`$-plane for the momentum space path shown in the inset of Fig. 4b, together with our theoretical results for $`\omega _c^{(1)}(𝐪)`$ (red line).
We also included $`\omega _c^{(1)}(𝐪)`$ as a dashed red line into the normal state data so as to demonstrate the transfer of spectral weight between the normal and the superconducting state from frequencies below $`\omega _c^{(1)}(𝐪)`$ to frequencies above the spin-gap. A comparison of Fig. 3a and b clearly shows that the experimental intensity that exists in the normal state for $`\omega <\omega _c^{(1)}(𝐪)`$ vanishes as expected in the superconducting state. Note that the momentum dependence of the spin-gap is rather weak; it only changes from $`\mathrm{\Delta }_{sg}^{min}5.5`$ meV at its minimum (close to $`𝐐_i`$) to $`\mathrm{\Delta }_{sg}^{max}8`$ meV at its local maximum midway between the incommensurate positions. We thus conclude that our theoretical result for the momentum dependence of the spin gap provides a good quantitative description of the area in the $`(\omega ,𝐪)`$-plane where the spectral weight vanishes in the superconducting state.
We now turn to momentum dependence of $`\chi ^{\prime \prime }`$ in the normal and superconducting state. Our theoretical results, which we present in Fig. 4, correspond to horizontal cuts in the $`(\omega ,𝐪)`$-plane of Fig. 3 at $`\omega =7`$ and 10 meV.
For $`\omega =7`$ meV (Fig. 4a), the peak intensity in the superconducting state is anisotropically reduced, with a stronger suppression of $`\chi _{SC}^{\prime \prime }`$ towards the center of the scan. The anisotropic suppression of $`\chi _{SC}^{\prime \prime }`$ is a direct consequence of the momentum dependence of the spin-gap which increases when moving from $`𝐐_i`$ towards the center of the scan, but decreases in the opposite direction. As a result, $`\chi _{SC}^{\prime \prime }`$ is rapidly cut off by the spin-gap when moving towards the center of the scan, but is scarcely reduced in the opposite direction. This peak anisotropy should be observable for all frequencies between $`\mathrm{\Delta }_{sg}^{min}`$ and $`\mathrm{\Delta }_{sg}^{max}`$. For $`\omega =10`$ meV $`>\mathrm{\Delta }_{sg}^{max}`$ (Fig. 4b), the anisotropy vanishes, and the peak intensity increases in the superconducting state, as expected from Fig. 2. Since the anisotropy of $`\chi _{SC}^{\prime \prime }`$ around $`𝐐_i`$ is reduced with increasing frequency, the peak maximum seems to slightly shift towards the center of the scan. All these results, i.e., narrow peaks around 7-8 meV, a simultaneous increase in peak width and height and a shift of the peak maximum towards the center of the momentum scan with increasing frequency agree qualitatively with the experimental findings of Ref. as may be seen by comparing Fig. 4a (Fig. 4b) with the bottom (top) part of Fig. 2 in Ref. or with Fig. 2b (Fig. 2c) in Ref. . Mason et al. considered the sharpening of the incommensurate peaks in the superconducting state as an effect of magnetic coherence which is shown here to arise solely from the momentum dependence of the spin-gap.
Due to the symmetry of the Fermi surface, $`\chi _{SC}^{\prime \prime }`$ exhibits four different threshold frequencies for momenta away from $`𝐐_i`$. While the two upper thresholds remain close, the energy separation between the two lower threshold increases rapidly with distance from $`𝐐_i`$ . We thus predict that $`\chi _{SC}^{\prime \prime }`$ will acquire additional frequency structure for $`𝐪𝐐_i`$.
The INS data also provide insight into the form of the FS in LSCO and thus complement the results of angle-resolved photoemission (ARPES) experiments. Since excitation (1) is located in the vicinity of the nodes where the superconducting gap changes rapidly with momentum, $`\omega _c^{(1)}`$ sensitively depends on the form of the FS and the symmetry of the order parameter. The frequency location of $`\omega _c^{(1)}`$ can therefore be used to extract information on the form of the Fermi surface, and, within the framework of Eq.(7), on the value of $`t^{}/t`$. In particular, we find that the INS data provide a lower bound for $`t^{}/t`$. Within our scenario, excitation (1) across the FS (see inset of Fig. 1) becomes impossible in the superconducting state for $`|t^{}|<0.2t`$. Since this implies $`\chi ^{\prime \prime }(𝐐_i)=0`$ for frequencies below $`\omega _c^{(2)}`$, in contradiction to the experimental results, we conclude $`|t^{}|0.2t`$. Assuming a weak doping dependence of $`t^{}/t`$, this constraint for $`t^{}`$ yields a FS of LSCO which is closed around $`(\pi ,\pi )`$ up to optimal doping. The FS thus possesses the same topology as that in YBCO; this explains the occurrence of incommensurate peaks in the spin spectrum along the same direction in momentum space in the latter materials .
Though the details of the magnetic coherence effect are sensitive to material specific parameters, e.g., Fermi surface topology, the extent of the incommensuration $`\delta `$, their experimental observation only depends on two criteria: the existence of an incommensurate spin structure and the d-wave symmetry of the superconducting gap. We thus predict a similar effect for all cuprate superconductors in which these criteria are met, and are currently studying its form in YBCO .
Finally, the theoretical scenario for the magnetic coherence effect presented here is conceptually different from that recently proposed for the resonance peak . Not only do the effects take place in different wave-vector and energy regions of $`\chi ^{\prime \prime }`$ , but the origin of the resonance peak is ascribed to a dispersing spin mode, while our scenario for the coherence effect is solely based on the existence of a relaxational spin mode.
In summary, we find that the frequency and momentum dependent changes of $`\chi ^{\prime \prime }`$ in the superconducting state are a direct consequence of changes in the quasiparticle spectrum due to the appearance of a d-wave gap. We show that the available INS data constrain the Fermi surface topology, and suggest a Fermi surface in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> which is closed around $`(\pi ,\pi )`$ up to optimal doping. We make several predictions for the frequency dependence of $`\chi ^{\prime \prime }(\omega )`$ at and around $`𝐐_i`$ which await further experimental testing. Finally, we predict the presence of comparable changes in $`\chi ^{\prime \prime }`$ in all cuprate superconductors with an incommensurate spin-structure.
We would like to thank G. Aeppli, A.V. Chubukov, P. Dai, A. Millis, H. Mook, and J. Schmalian for valuable discussions and particularly B. Lake and T. Mason for very stimulating discussions and for providing us with their experimental data prior to publication. This work has been supported in part by the Science and Technology Center for Superconductivity through NSF-grant DMR91-20000, and by DOE at Los Alamos.
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# Rigidly rotating dust in general relativity
## 1 Introduction
Recently, Bonnor studied axially symmetric stationary solutions of Einstein field equations coupled with dust and showed the reasons why a density gradient parallel to the axis is allowed for General Relativity but not for Newtonian mechanics . He also considered an analytic solution which represents a rigidly rotating dust, and found that the solution is asymptotically flat in all the three spatial directions but with a total mass equal to zero. He attributed this to the center singularity that was believed to have an infinitely large negative mass, which just balances the positive infinite mass of the dust that fills all the spacetime. For the details, we refer the readers to . Motivated by this particular solution, Bonnor wondered whether or not a non-singular rotating dust exists. In the present article, we obtain a solution that generates an axially symmetric rigidly rotating dust, which is accompanied by a rigidly rotating thin disk, namely, a singular hypersurface perpendicular to the axis of the rigidly rotating dust.
It is worth mentioning that there is a paper by Georgiou , who studied rotating Einstein-Maxwell fields<sup>1</sup><sup>1</sup>1It has been one of the referees who brought to our attention this paper by Georgiou.. Georgiou obtains exact exterior and matching interior stationary axially symmetric solutions for a rigidly rotating charged dust. His solution generates an infinitely long cylinder and a thin singular disk perpendicular to the axis of the cylinder. Later on we consider the differences and the similarities of our solution and that of Georgiou’s.
The present paper has also been motivated by a study by Opher, Santos and Wang (OSW) concerning the origin of extragalactic jets. These authors argued that, under certain circumstances, the spacetime given by what they refer to as the van Stockum metric, which is associated with a dust cylinder (studied by van Stockum and extensively analyzed by Bonnor ), can account for the collimating effect present in extragalactic jets. OSW showed that this dust cylinder produces confinement for the geodesic motion of test particles for certain values of the radial energy and angular momentum. In fact, it was one of our aims to improve the OSW’s model by looking for a spacetime that could more realistically describe a jet.
It is worth mentioning, however, that van Stockum in fact rediscovered a solution which was first obtained by Lanczos <sup>2</sup><sup>2</sup>2The referees brought to our attention that Lanczos’ paper preceded that by van Stockum. . Hereafter, therefore, we refer to the metric related to the dust cylinder as Lanczos metric.
The extragalactic jets are ubiquitous in active galaxies, they are highly collimated and the matter which forms them is highly relativistic . It is worth mentioning that there is no consensus in the literature to explain why they are the way they are. Many authors assume that the jets propagate along a direction provided by, most probably, rapidly rotating Kerr black holes present at the centers of active galaxies (see, e.g., Begelman, Blandford and Rees and also ). This fact also suggest that putative general relativistic effects could be important (see, e.g., , among others).
The central engine that gives rise to the jets could be more complex than a simple super massive black hole, it could well occur that jets be driven by an axially symmetric structure present at the center of active galaxies. The Lanczos solution (referred to by OSW as van Stockum solution), however, has its weakness when applied to an actual physical situation: it represents an infinitely long cylinder. As a result, the spacetime is not asymptotically flat and has infinite mass.
The solution here studied generates an axially symmetric rigidly rotating dust accompanied by a surface layer, which does not satisfy the energy conditions (i.e., weak, dominant and strong ) in part of the hypersurface. We point out however, that other authors (see Refs. , and references cited therein) have investigated such structures and some of them, as in our case, do not satisfy any of the energy conditions. In some of these cases, the energy conditions may be satisfied by a suitable choice of parameters.
We argue that the present study may be by itself of interest, because it represents a new axially symmetric dust solution, which could also motivate other authors to find solutions physically satisfactory.
In section 2 a closed form of the solutions are given and the main properties of them are studied, while in section 3 our main conclusions and some discussions are presented.
## 2 The Rotating Dust Metric
Our starting point is the Lanczos metric given by
$$ds^2=dt^22kdtd\phi ld\phi ^2e^\mu (dr^2+dz^2),$$
(1)
where
$`k`$ $`=`$ $`\alpha \eta ,\eta =r\xi _{,r},l=r^2\alpha ^2\eta ^2,`$
$`\mu _{,r}`$ $`=`$ $`{\displaystyle \frac{\alpha ^2}{2r}}(\eta _{,z}^2\eta _{,r}^2),\mu _{,z}={\displaystyle \frac{\alpha ^2}{2r}}\eta _{,r}\eta _{,z},`$ (2)
and the function $`\xi (r,z)`$ satisfies the Laplacian equation $`^2\xi =0`$, with $`^2`$ being the Laplacian operator in Euclidean three-space. The symbol $`()_{,x}`$ denotes partial derivative with respect to the argument $`x`$, and $`\{x^\mu \}\{t,r,z,\phi \},(\mu =0,1,2,3)`$ are the usual axisymmetric coordinates. One can show that the above solutions satisfy the Einstein field equations<sup>3</sup><sup>3</sup>3In this paper we choose units such that $`G=1=c`$, where $`G`$ is the gravitational constant, and $`c`$ the speed of light. $`R_{\mu \nu }g_{\mu \nu }R/2=8\pi \rho u_\mu u_\nu `$ with the energy density and four-velocity of the dust being given, respectively, by
$$\rho =\frac{e^\mu }{8\pi r^2}(\eta _{,z}^2+\eta _{,r}^2),u^\mu =\delta _0^\mu .$$
(3)
The above solutions represent rigidly rotating dust. This can be seen, for example, by calculating the shear in this non-expanding spacetime, $`q_{\mu \nu }(u_{\mu ;\nu }+u_{\nu ;\mu })/2`$, which is identically zero for the solutions given by Eqs.(1) and (2). However, the angular velocity of the dust, which is given by $`w_{\mu \nu }(u_{\mu ;\nu }u_{\nu ;\mu })/2`$, does not vanish.
The specific solution considered by Bonnor , for example, is $`\xi =2h/\sqrt{r^2+z^2}`$ with $`h`$ being a constant. As shown in this spacetime is free of any spacetime singularities, except for that located at the origin of the coordinate system, namely, $`r=z=0`$. This singularity is a curvature singularity with an infinitely large negative mass.
In this paper, we consider the solution with $`\xi =J_0(r)e^z`$, where $`J_0(r)`$ denotes the zero-order Bessel function. Then, substituting it into Eqs.(2) and (3) we find
$`k`$ $`=`$ $`\alpha rJ_1(r)e^z,l=r^2\left[1\alpha ^2J_1^2(r)e^{2z}\right],`$
$`\mu `$ $`=`$ $`{\displaystyle \frac{\alpha ^2r}{2}}J_0(r)J_1(r)e^{2z},`$ (4)
$`\rho `$ $`=`$ $`{\displaystyle \frac{\alpha ^2e^{2z}}{8\pi }}\left[J_0^2(r)+J_1^2(r)\right]e^\mu .`$ (5)
From the above equations we can see that the spacetime is singular when $`z\mathrm{}`$. To remedy this undesirable feature we can replace $`z`$ by $`|z|`$ in Eq.(4), i.e.,
$`k`$ $`=`$ $`\alpha rJ_1(r)e^{|z|},l=r^2\left[1\alpha ^2J_1^2(r)e^{2|z|}\right],`$
$`\mu `$ $`=`$ $`{\displaystyle \frac{\alpha ^2r}{2}}J_0(r)J_1(r)e^{2|z|}.`$ (6)
Before proceeding it is worth mentioning that in a paper by Georgiou exact exterior and matching interior solutions are found, where the interior solution is similar to the solution present here. As shown in the spacetime refers to solutions of the Einstein-Maxwell field equations for a rigidly rotating charged dust with vanishing Lorentz force. The solutions generate an infinitely long cylinder of charged dust rigidly rotating about its axis and a 4-current located on a singular hypersurface perpendicular to the axis of the cylinder at the origin of the coordinate system.
Due to the fact that Georgiou’s interior solutions present a non null electromagnetic 4-potencial, namely, $`A_\mu =(0,0,0,A_3)`$, his equations related to the $`\mu `$ function involves the $`A_3`$ function (see Eqs.(4.2) and (4.3) in ). In such a way he could set $`F=e^\mu =1`$. On the other hand, we have $`F=1`$ and $`\mu `$ is given by Eq.(4), as a result our solutions are different. Apart from the constants our mass density contains in addition the term $`e^\mu `$, as a result a stronger dependence on the z coordinate occurs as compared to the Georgiou’s mass density. Our solution describes rigidly rotating neutral dust and Georgiou’s solution, rigidly rotating charged dust. Consequently the resulting spacetimes are different.
One can show that such resulted spacetime present here is asymptotically flat in the z direction, since as $`|z|\mathrm{}`$ for any particular finite value of $`r`$, $`k0`$, $`lr^2`$ and $`e^\mu 1`$. On the other hand, the behaviour of the solution as $`r\mathrm{}`$, for any particular finite value of z, shows that, for example, $`k`$ oscillates infinitely between $`\mathrm{}`$ and $`+\mathrm{}`$. This would indicate that the spacetime is not asymptotically flat.
On the other hand, the Kretschmann scalar is given by
$``$ $``$ $`R_{\alpha \beta \gamma \delta }R^{\alpha \beta \gamma \delta }={\displaystyle \frac{\alpha ^2e^{2(|z|+\mu )}}{4r^2}}\{16[J_1^2(r)rJ_0(r)J_1(r)`$ (7)
$`+r^2(J_0^2(r)+J_1^2(r))]4\alpha ^2r^2e^{2|z|}[2J_0^4(r)+J_1^4(r)`$
$`+7J_0^2(r)J_1^2(r)4rJ_0(r)J_1(r)(J_0^2(r)+J_1^2(r))]`$
$`+\alpha ^4r^4e^{4|z|}[J_0^6(r)+J_1^6(r)+3J_0^2(r)J_1^2(r)(J_0^2(r)+J_1^2(r))]\}`$
$`+_0\delta (z),`$
where $`_0(r)`$ is a bounded function of $`r`$ (see the discussions following Eq.(14) below), and $`\delta (z)`$ denotes the Dirac delta function. Using the relations
$`J_n(x)\{\begin{array}{cc}\frac{x^n}{2^nn!},\hfill & x0,\hfill \\ & \\ \sqrt{\frac{2}{\pi x}}\mathrm{cos}\left(x\frac{2n+1}{4}\pi \right),\hfill & x+\mathrm{},\hfill \end{array}`$ (11)
we can see from Eq.(7) that $``$ finite, as $`r0`$, and that $`0`$, as $`|z|`$ or $`r+\mathrm{}`$ indicating that the spacetime is asymptotically flat. Although we have considered here some discussion on flatness, it is worth bearing in mind that some authors argue that this concept is not well defined (see,e.g., ).
Also, from Eqs.(6) and (8) we have
$`X_\phi ^2=|g_{\phi \phi }|`$ $``$ $`O(r^2),`$
$`{\displaystyle \frac{X_{,\alpha }X^{,\alpha }}{4X}}`$ $``$ $`1,`$ (12)
as $`r0`$. Hence, the axis ($`r=0`$) of the spacetime is well defined and locally flat. From the above equations we can also see that, by properly choosing the constant $`\alpha `$, we may have $`g_{\phi \phi }<0`$ for any $`r`$. That is, the spacetime may be free of any closed time-like curves. Therefore, the solution given by Eq.(6) represents an axially symmetric and rigidly rotating dust spacetime.
It should be noted that the replacement of $`z`$ by $`|z|`$ gives rise to a thin matter shell. As a matter of fact, this replacement mathematically is equivalent first to cut the original spacetime given by Eq.(4) into two parts, $`z>0`$ and $`z<0`$, and then join the part $`z>0`$ with a copy of it along the hypersurface $`z=0`$, so that the resulted spacetime has a reflection symmetry with respect to the surface. After this cut-paste operation, the spacetime is no longer analytic across the surface $`z=0`$. Actually, the metric coefficients are continuous, but their first derivatives with respect to $`z`$ are not. Then, according to Taub’s theory , a thin matter shell appears on the hypersurface. Introducing the quantity $`b_{\mu \nu }`$ via the relation
$$b_{\mu \nu }g_{\mu \nu ,z}^+|_{z=0^+}g_{\mu \nu ,z}^{}|_{z=0^{}},$$
(13)
where $`g_{\mu \nu }^+(g_{\mu \nu }^{})`$ are quantities defined in the region $`z>0(z<0)`$, we find that the non-vanishing components of $`b_{\mu \nu }`$ are given by
$`b_{11}`$ $`=`$ $`b_{22}=2\alpha ^2rJ_0(r)J_1(r)e^{\mu _0},`$
$`b_{03}`$ $`=`$ $`b_{30}=2\alpha rJ_1(r),b_{33}=4\alpha ^2r^2J_1^2(r),`$ (14)
where $`\mu _0\mu (r,z)|_{z=0}`$. Then, the surface energy-momentum tensor $`\tau _{\mu \nu }`$ is given by
$$\tau _{\mu \nu }=\frac{1}{16\pi }\left\{b(ng_{\mu \nu }n_\mu n_\nu )+n_\lambda (n_\mu b_\nu ^\lambda +n_\nu b_\mu ^\lambda )(nb_{\mu \nu }+n_\lambda n_\delta b^{\lambda \delta }g_{\mu \nu })\right\},$$
(15)
where $`n_\mu `$ is the normal vector to the hypersurface $`z=0`$, given by $`n_\mu =\delta _\mu ^2`$, with $`nn_\lambda n^\lambda `$ and $`bb_\lambda ^\lambda `$. Substituting Eq.(11) into Eq.(12), we find that $`\tau _{\mu \nu }`$ can be written in the form
$$8\pi \tau _{\mu \nu }=\sigma u_\mu u_\nu +px_\mu x_\nu +q(u_\mu x_\nu +u_\nu x_\mu ),$$
(16)
with
$`\sigma `$ $`=`$ $`p=\alpha ^2rJ_0(r)J_1(r)e^{\mu _0},`$
$`q`$ $`=`$ $`\alpha J_1(r)e^{\mu _0},`$ (17)
where $`u_\mu `$ is the four-velocity of the dust restricted to the surface $`z=0`$, and $`x_\mu `$ is a space-like unit vector on the surface, given by $`x_\mu =r\delta _\mu ^3`$, and has the properties: $`x_\lambda x^\lambda =1`$ and $`x_\lambda u^\lambda =0`$. The non null components of $`\tau _\nu ^\mu `$ are explicit shown in the appendix, where we also obtain them using the alternative technique derived by Israel .
Eq.(13) shows that $`\sigma `$ represents the surface energy density of the shell measured by observers who are comoving with the dust fluid, and $`p`$ is the pressure in the direction perpendicular to the observers’ world lines, while $`q`$ represents the heat flow. On the other hand, Eqs.(14) show that all these quantities are finite for any $`r`$. One can also see that $`\sigma `$ oscillates between positive and negative values; when $`\sigma <0`$ the energy conditions are violated. Note that the function $`_0(r)`$ appearing in Eq.(7) is a combination of them and is finite, too.
Note that the quantities $`q_{\mu \nu }`$ and $`w_{\mu \nu }`$ defined above contain only first derivatives of the metric coefficients, as a result they can be written generally in the form
$$Y_{\mu \nu }=Y_{\mu \nu }^+H(z)+Y_{\mu \nu }^{}H(z),$$
where $`Y_{\mu \nu }^\pm (\{q_{\mu \nu }^\pm ,w_{\mu \nu }^\pm \})`$ are quantities calculated respectively in the regions $`z>0`$ and $`z<0`$, and $`H(z)`$ is the Heavside function, which is $`1`$ for $`z>0`$, $`1/2`$ for $`z=0`$, and $`0`$ for $`z<0`$. Since $`q_{\mu \nu }^\pm `$ are all equal zero, we can see that $`q_{\mu \nu }`$ is zero even on the surface $`z=0`$. That is, the matter shell is also rigidly rotating.
On the other hand, in terms of $`\tau _{\mu \nu }`$, the Einstein field equations for the solution (6) takes the form
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R=8\pi [\rho u_\mu u_\nu +\tau _{\mu \nu }\delta (z)],`$ (18)
where $`\rho `$ is given by Eq.(5) but with $`z`$ being replaced by $`|z|`$. Thus, the solution given by Eq.(6) represents a rigidly rotating dust accompanied by a matter shell of the same type.
To further study the spacetime of this solution, let us consider the total mass of the spacetime. Using the formula of Tolman , we find that the mass inside the cylinder $`r=R`$ is given by
$`M(R)`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _0^R}\left\{t_{\mathrm{\hspace{0.33em}0}}^0+\sqrt{g}\left[\rho u^0u_0+\tau _{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^0\delta (z)\right]\right\}𝑑\phi 𝑑z𝑑r`$ (19)
$`=`$ $`{\displaystyle \frac{\alpha ^2R}{16}}\left\{2J_0(R)J_1(R)+R\left[J_0^2(R)J_1^2(R)\right]\right\},`$
where $`R`$ is a constant, and $`t_\mu ^\nu `$ is the so-called gravitational energy-momentum pseudo-tensor . The combination of Eq.(8) and Eq.(16) shows that $`M(R)`$ oscillates infinitely between $`\mathrm{}`$ and $`+\mathrm{}`$ as $`R+\mathrm{}`$. Thus, in the present case the total mass of the spacetime is not well defined. Were $`\sigma `$ always negative it could occur that the total mass of the dust structure were zero. This behaviour of $`M(R)`$ as $`R\mathrm{}`$ is related to the fact that $`\sigma `$ oscillates between positive and negative values. To remedy this problem, one might cut the spacetime and smoothly match the dust structure where $`\sigma >0`$ to vacuum spacetimes. As result one would have a hollow rigidly rotating dust structure with two vacuums, one inside and other outside it. It could occur however that this procedure generated unphysical surface layers. Due to the complexity of these issues, we have not yet been successful in this direction.
## 3 Discussion and Conclusions
Motivated by the article by Bonnor we have studied a particular axially symmetric rigidly rotating dust solution, and found that it is accompanied by a thin disk located on an hypersurface perpendicular to the symmetry axis. The undesirable feature is that the thin disk has negative energy density in part of the hypersurface.
The solution we have found is in some sense similar to a solution found by Georgiou . This author obtains exact interior and matching exterior axially symmetric solutions of the Einstein-Maxwell fields equations, for rigidly rotating charged dust. The fact that our metric function $`\mu `$ be given by Eq.(4), instead of being $`\mu =1`$ as in , turns the spacetimes different. Therefore, the fact that the dust be charged modify significantly the spacetime generated.
A completely satisfactory solution for rigidly rotating dust fluids is yet to be derived and deserves to be investigated. Such a solution must be asymptotically flat, have finite mass, have non-singularities, and whether a thin shell appears it must satisfy the energy conditions.
To see whether or not our solution allow confinement, one needs to follow a similar procedure as in . However, due to the fact that the metric functions now depend on both of the $`r`$ and $`z`$ coordinates, the study of geodesic motions become very complicated. We therefore leave such an issue for a future study.
Acknowledgments
We would like to thank Drs R. Opher and N.O. Santos for helpful discussions, and O.D. Miranda for reading the paper. This work was partially developed while J.C.N.A. had a post-doc position at Departamento de Astronomia - Instituto Astronômico e Geofísico (Universidade de São Paulo). The financial assistance from CNPq and FAPESP is gratefully acknowledged. Finally, we would like to thank the referees for the suggestions and criticisms that helped us to improve our paper.
## Appendix A Alternative Calculation of $`\tau _{\mu \nu }`$
Here we show that Taub’s and Israel’s techniques to calculate the surface energy-momentum tensor, $`\tau _{\mu \nu }`$, are equivalent.
We start from
$`8\pi \tau _\nu ^\mu `$ $``$ $`8\pi \underset{\epsilon 0}{lim}{\displaystyle _\epsilon ^{+\epsilon }}T_\nu ^\mu 𝑑z=\underset{\epsilon 0}{lim}{\displaystyle _\epsilon ^{+\epsilon }}\left(R_\nu ^\mu {\displaystyle \frac{1}{2}}\delta _\nu ^\mu R\right)𝑑z`$ (20)
$`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle _\epsilon ^{+\epsilon }}\left(_\nu ^\mu {\displaystyle \frac{1}{2}}\delta _\nu ^\mu \right)𝑑z`$
(see, e.g., ).
Due to the presence of the thin disk, $`R_\nu ^\mu `$ and $`R`$ contain the terms $`_\nu ^\mu `$ and $``$, respectively. These terms contain delta functions, and appear only with components of $`R_\nu ^\mu `$ whose expressions contain second order partial derivatives with respect to z. Following , it is easy to show that the non null components of $`_\nu ^\mu `$ read
$`_{\mathrm{\hspace{0.33em}\hspace{0.33em}1}}^1`$ $`=`$ $`_{\mathrm{\hspace{0.33em}\hspace{0.33em}2}}^2={\displaystyle \frac{1}{2}}e^\mu (\mu _{,z}^{}\mu _{,z}^+)\delta (z)`$
$`_{\mathrm{\hspace{0.33em}\hspace{0.33em}3}}^3`$ $`=`$ $`_{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^0={\displaystyle \frac{1}{2}}e^\mu {\displaystyle \frac{k}{r^2}}(k_{,z}^{}k_{,z}^+)\delta (z)`$
$`_{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^3`$ $`=`$ $`{\displaystyle \frac{e^\mu }{2r^2}}(k_{,z}^{}k_{,z}^+)\delta (z)`$
$`_{\mathrm{\hspace{0.33em}\hspace{0.33em}3}}^0`$ $`=`$ $`2k_{\mathrm{\hspace{0.33em}\hspace{0.33em}3}}^3l_{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^3,`$ (21)
where the functions with superscript $`+`$ ($``$) stand for the functions defined in the region $`z>0(z<0)`$.
Substituting Eq.(18) into (17) it is straightforward to show that the non null components of $`\tau _\nu ^\mu `$ are given by:
$`8\pi \tau _{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^0`$ $`=`$ $`\alpha ^2[J_1^2(r)+rJ_0(r)J_1(r)]e^{\mu _0}`$
$`8\pi \tau _{\mathrm{\hspace{0.33em}\hspace{0.33em}3}}^0`$ $`=`$ $`\alpha rJ_1(r)[1+\alpha ^2J_1^2(r)]e^{\mu _0}`$
$`8\pi \tau _{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^3`$ $`=`$ $`{\displaystyle \frac{\alpha }{r}}J_1(r)e^{\mu _0}`$
$`8\pi \tau _{\mathrm{\hspace{0.33em}\hspace{0.33em}3}}^3`$ $`=`$ $`\alpha ^2[J_1^2(r)rJ_0(r)J_1(r)]e^{\mu _0},`$ (22)
which agree with Eqs.(13), showing, therefore, that Taub’s and Israel’s techniques are equivalent.
Finally, it is worth mentioning that Georgiou’s definition for $`\tau _\nu ^\mu `$ is a little bit different from Israel’s. The former integrates $`T_\nu ^\mu `$ with respect to the proper distance measured perpendicularly through the thin disk. On the other hand, Eq.(17) is an integral of $`T_\nu ^\mu `$ with respect to the $`z`$ coordinate. If one followed Georgiou’s definition, instead of having the term $`e^{\mu _0}`$ in Eq.(19) one would have $`e^{\mu _0/2}`$.
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# limit-from𝐽-holomorphic Curves, Legendre Submanifolds and Reeb Chords Project 19871044 Supported by NSF
## 1 Introduction and results
Let $`\mathrm{\Sigma }`$ be a smooth closed oriented manifold of dimension $`2n1`$. A contact form on $`\mathrm{\Sigma }`$ is a $`1`$form such that $`\lambda (d\lambda )^{n1}`$ is a volume form on $`\mathrm{\Sigma }`$. Associated to $`\lambda `$ there are two important structures. First of all the so-called Reed vectorfield $`\dot{x}=X`$ defined by
$$i_X\lambda 1,i_Xd\lambda 0;$$
and secondly the contact structure $`\xi =\xi _\lambda \mathrm{\Sigma }`$ given by
$$\xi _\lambda =\mathrm{ker}(\lambda )T\mathrm{\Sigma }.$$
By a result of Gray, , the contact structure is very stable. In fact, if $`(\lambda _t)_{t[0,1]}`$ is a smooth arc of contact forms inducing the arc of contact structures $`(\xi _t)_{t[0,1]}`$, there exists a smooth arc $`(\psi _t)_{t[0,1]}`$ of diffeomorphisms with $`\psi _0=Id`$, such that
$$T\mathrm{\Psi }_t(\xi _0)=\xi _t$$
(1.1)
here it is importent that $`\mathrm{\Sigma }`$ is compact. From (1.1) and the fact that $`\mathrm{\Psi }_0=Id`$ it follows immediately that there exists a smooth family of maps $`[0,1]\times \mathrm{\Sigma }(0,\mathrm{}):(t,m)f_t(m)`$ such that
$$\mathrm{\Psi }_t^{}\lambda _t=f_t\lambda _0$$
(1.2)
In contrast to the contact structure the dynamics of the Reeb vectorfield changes drastically under small perturbation and in general the flows associated to $`X_t`$ and $`X_s`$ for $`ts`$ will not be conjugated.
Concerning the dynamics of Reeb flow, there is a well-known conjecture raised by Arnold in which concerned the Reeb orbit and Legendre submanifold in a contact manifold. If $`(\mathrm{\Sigma },\lambda )`$ is a contact manifold with contact form $`\lambda `$ of dimension $`2n1`$, then a Legendre submanifold is a submanifold $``$ of $`\mathrm{\Sigma }`$, which is $`(n1)`$dimensional and everywhere tangent to the contact structure $`\mathrm{ker}\lambda `$. Then a characteristic chord for $`(\lambda ,)`$ is a smooth path
$$x:[0,T]M,T>0$$
with
$$\dot{x}(t)=X_\lambda (x(t))fort(0,T),$$
$$x(0),x(T)$$
Arnold raised the following conjecture:
Conjecture(see). Let $`\lambda _0`$ be the standard tight contact form
$$\lambda _0=\frac{1}{2}(x_1dy_1y_1dx_1+x_2dy_2y_2dx_2)$$
on the three sphere
$$S^3=\{(x_1,y_1,x_2,y_2)R^4|x_1^2+y_1^2+x_2^2+y_2^2=1\}.$$
If $`f:S^3(0,\mathrm{})`$ is a smooth function and $``$ is a Legendre knot in $`S^3`$, then there is a characteristic chord for $`(f\lambda _0,)`$.
The main results of this paper as following:
###### Theorem 1.1
Let $`(\mathrm{\Sigma },\lambda )`$ be a contact manifold with contact form $`\lambda `$, $`X_\lambda `$ its Reeb vector field, $``$ a closed Legendre submanifold, then there exists at least one characteristic chord for $`(X_\lambda ,)`$.
###### Corollary 1.1
() Let $`(S^3,f\lambda _0)`$ be a tight contact manifold with contact form $`f\lambda _0`$, $`X_{f\lambda _0}`$ its Reeb vector field, $``$ a closed Legendre submanifold, then there exists at least one characteristic chord for $`(X_{f\lambda _0},)`$.
Sketch of proofs: We work in the framework as in . In Section 2, we study the linear Cauchy-Riemann operator and sketch some basic properties. In section 3, first we construct a Lagrangian submanifold $`W`$ under the assumption that there does not exists Reeb chord conneting the Legendre submanifold $``$; second, we study the space $`𝒟(V,W)`$ of contractible disks in manifold $`V`$ with boundary in Lagrangian submanifold $`W`$ and construct a Fredholm section of tangent bundle of $`𝒟(V,W)`$. In section 4, following , we prove that the Fredholm section is not proper by using an special anti-holomorphic section as in . In section 5-6, we use a geometric argument to prove the boundaries of $`J`$holomorphic curves remain in a finite part of Lagrangian submanifold $`W`$. In the final section, we use nonlinear Fredholm trick in to complete our proof.
## 2 Linear Fredholm Theory
For $`100<k<\mathrm{}`$ consider the Hilbert space $`V_k`$ consisting of all maps $`uH^{k,2}(D,C\times C^n)`$, such that $`u(z)\{izR\}\times R^nC\times C^n`$ for almost all $`zD`$. $`L_{k1}`$ denotes the usual Sobolev space $`H_{k1}(D,C\times C^n)`$. We define an operator $`\overline{}:V_kL_{k1}`$ by
$$\overline{}u=u_s+iu_t$$
(2.1)
where the coordinates on $`D`$ are $`(s,t)=s+it`$, $`D=\{z||z|1\}`$. The following result is well known(see).
###### Proposition 2.1
$`\overline{}:V_kL_{k1}`$ is a surjective real linear Fredholm operator of index $`n+3`$. The kernel consists of $`(a_0+isz\overline{a}_0z^2,s_1,\mathrm{},s_n)`$, $`a_0C`$, $`s,s_1,\mathrm{},s_nR`$.
Let $`(C^n,\sigma =Im(,))`$ be the standard symplectic space. We consider a real $`n`$dimensional plane $`R^nC^n`$. It is called Lagrangian if the skew-scalar product of any two vectors of $`R^n`$ equals zero. For example, the plane $`\{(p,q)|p=0\}`$ and $`\{(p,q)|q=0\}`$ are two transversal Lagrangian subspaces. The manifold of all (nonoriented) Lagrangian subspaces of $`R^{2n}`$ is called the Lagrangian-Grassmanian $`\mathrm{\Lambda }(n)`$. One can prove that the fundamental group of $`\mathrm{\Lambda }(n)`$ is free cyclic, i.e. $`\pi _1(\mathrm{\Lambda }(n))=Z`$. Next assume $`(\mathrm{\Gamma }(z))_{zD}`$ is a smooth map associating to a point $`zD`$ a Lagrangian subspace $`\mathrm{\Gamma }(z)`$ of $`C^n`$, i.e. $`(\mathrm{\Gamma }(z))_{zD}`$ defines a smooth curve $`\alpha `$ in the Lagrangian-Grassmanian manifold $`\mathrm{\Lambda }(n)`$. Since $`\pi _1(\mathrm{\Lambda }(n))=Z`$, one have $`[\alpha ]=ke`$, we call integer $`k`$ the Maslov index of curve $`\alpha `$ and denote it by $`m(\mathrm{\Gamma })`$, see().
Now let $`z:S^1\{R\times R^nC\times C^n\}\mathrm{\Lambda }(n+1)`$ be a constant curve. Then it defines a constant loop $`\alpha `$ in Lagrangian-Grassmanian manifold $`\mathrm{\Lambda }(n+1)`$. This loop defines the Maslov index $`m(\alpha )`$ of the map $`z`$ which is easily seen to be zero.
Now Let $`(V,\omega )`$ be a symplectic manifold, $`WV`$ a closed Lagrangian submanifold. Let $`(\overline{V},\overline{\omega })=(D\times V,\omega _0+\omega )`$ and $`\overline{W}=D\times W`$. Let $`\overline{u}=(id,u):(D,D)(D\times V,D\times W)`$ be a smooth map homotopic to the map $`\overline{u}_0=(id,u_0)`$, here $`u_0:(D,D)pWV`$. Then $`\overline{u}^{}TV`$ is a symplectic vector bundle on $`D`$ and $`(\overline{u}|_D)^{}T\overline{W}`$ be a Lagrangian subbundle in $`\overline{u}^{}T\overline{V}|_D`$. Since $`\overline{u}:(D,D)(\overline{V},\overline{W})`$ is homotopic to $`\overline{u}_0`$, i.e., there exists a homotopy $`h:[0,1]\times (D,D)(\overline{V},\overline{W})`$ such that $`h(0,z)=(z,p),h(1,z)=\overline{u}(z)`$, we can take a trivialization of the symplectic vector bundle $`h^{}T\overline{V}`$ on $`[0,1]\times (D,D)`$ as
$$\mathrm{\Phi }(h^{}T\overline{V})=[0,1]\times D\times C\times C^n$$
and
$$\mathrm{\Phi }((h|_{[0,1]\times D})^{}T\overline{W})[0,1]\times S^1\times C\times C^n$$
Let
$$\pi _2:[0,1]\times D\times C\times C^nC\times C^n$$
then
$$\stackrel{~}{h}:(s,z)[0,1]\times S^1\pi _2\mathrm{\Phi }(h|_{[0,1]\times D})^{}T\overline{W}|(s,z)\mathrm{\Lambda }(n+1).$$
###### Lemma 2.1
Let $`\overline{u}:(D,D)(\overline{V},\overline{W})`$ be a $`C^k`$map $`(k1)`$ as above. Then,
$$m(\stackrel{~}{u})=2.$$
Proof. Since $`\overline{u}`$ is homotopic to $`\overline{u}_0`$ in $`\overline{V}`$ relative to $`\overline{W}`$, by the above argument we have a homotopy $`\mathrm{\Phi }_s`$ of trivializations such that
$$\mathrm{\Phi }_s(\overline{u}^{}TV)=D\times C\times C^n$$
and
$$\mathrm{\Phi }_s((\overline{u}|_D)^{}T\overline{W})S^1\times C\times C^n$$
Moreover
$$\mathrm{\Phi }_0(\overline{u}|_D)^{}T\overline{W}=S^1\times izR\times R^n$$
So, the homotopy induces a homotopy $`\stackrel{~}{h}`$ in Lagrangian-Grassmanian manifold. Note that $`m(\stackrel{~}{h}(0,))=0`$. By the homotopy invariance of Maslov index, we know that $`m(\stackrel{~}{u}|_D)=2`$.
Consider the partial differential equation
$`\overline{}\overline{u}+A(z)\overline{u}=0onD`$ (2.2)
$`\overline{u}(z)\mathrm{\Gamma }(z)(izR\times R^n)forzD`$ (2.3)
$`\mathrm{\Gamma }(z)GL(2(n+1),R)Sp(2(n+1))`$ (2.4)
$`m(\mathrm{\Gamma })=2`$ (2.5)
For $`100<k<\mathrm{}`$ consider the Banach space $`\overline{V}_k`$ consisting of all maps $`uH^{k,2}(D,C^n)`$ such that $`u(z)\mathrm{\Gamma }(z)`$ for almost all $`zD`$. Let $`L_{k1}`$ the usual Sobolev space $`H_{k1}(D,C\times C^n)`$
###### Proposition 2.2
$`\overline{}:\overline{V}_kL_{k1}`$ is a real linear Fredholm operator of index n+3.
## 3 Nonlinear Fredholm Theory
### 3.1 Constructions of Lagrangian submanifolds
Let $`(\mathrm{\Sigma },\lambda )`$ be a contact manifolds with contact form $`\lambda `$ and $`X`$ its Reeb vector field, then $`X`$ integrates to a Reeb flow $`\eta _t`$ for $`tR^1`$. Consider the form $`d(e^a\lambda )`$ on the manifold $`(R\times \mathrm{\Sigma })`$, then one can check that $`d(e^a\lambda )`$ is a symplectic form on $`R\times \mathrm{\Sigma }`$. Moreover One can check that
$`i_X(e^a\lambda )=e^a`$ (3.1)
$`i_X(d(e^a\lambda ))=de^a`$ (3.2)
So, the symplectization of Reeb vector field $`X`$ is the Hamilton vector field of $`e^a`$ with respect to the symplectic form $`d(e^a\lambda )`$. Therefore the Reeb flow lifts to the Hamilton flow $`h_s`$ on $`R\times \mathrm{\Sigma }`$(see).
Let $``$ be a closed Legendre submanifold in $`(\mathrm{\Sigma },\lambda )`$, i.e., there exists a smooth embedding $`Q:\mathrm{\Sigma }`$ such that $`Q^{}\lambda |_{}=0`$, $`\lambda |Q(L)=0`$. We also write $`=Q()`$. Let
$$(V^{},\omega ^{})=(R\times \mathrm{\Sigma },d(e^a\lambda ))$$
and
$`W^{}=\times R,W_s^{}=\times \{s\};`$ (3.3)
$`L^{}=(0,_s\eta _s(Q())),L_s^{}=(0,\eta _s(Q()))`$ (3.4)
define
$`G^{}:W^{}V^{}`$ (3.5)
$`G^{}(w^{})=G^{}(l,s)=(0,\eta _s(Q(l)))`$ (3.6)
###### Lemma 3.1
There does not exist any Reeb chord connecting Legendre submanifold $``$ in $`(\mathrm{\Sigma },\lambda )`$ if and only if $`G^{}(W_s^{})G^{}(W_s^{}^{})`$ is empty for $`ss^{}`$.
Proof. Obvious.
###### Lemma 3.2
If there does not exist any Reeb chord for $`(X_\lambda ,)`$ in $`(\mathrm{\Sigma },\lambda )`$ then there exists a smooth embedding $`G^{}:W^{}V^{}`$ with $`G^{}(l,s)=(0,\eta _s(Q(l)))`$ such that
$$G_K^{}:\times (K,K)V^{}$$
(3.7)
is a regular open Lagrangian embedding for any finite positive $`K`$. We denote $`W^{}(K,K)=G_K^{}(\times (K,K))`$
Proof. One check
$$G_{}^{}{}_{}{}^{}(d(e^a\lambda ))=\eta (,)^{}d\lambda =(\eta _s^{}d\lambda +i_Xd\lambda ds)=0$$
(3.8)
This implies that $`G^{}`$ is a Lagrangian embedding, this proves Lemma3.2.
In fact the above proof checks that
$$G_{}^{}{}_{}{}^{}(\lambda )=\eta (,)^{}\lambda =\eta _s^{}\lambda +i_X\lambda ds=ds.$$
(3.9)
i.e., $`W^{}`$ is an exact Lagrangian submanifold.
Now we construct an isotopy of Lagrangian embeddings as follows:
$`F^{}:\times R\times [0,1]R\times \mathrm{\Sigma }`$ (3.10)
$`F^{}(l,s,t)=(a(s,t),G^{}(l,s))=(a(s,t),\eta _s(Q(l)))`$ (3.11)
$`F_t^{}(l,s)=F^{}(l,s,t)`$ (3.12)
###### Lemma 3.3
If there does not exist any Reeb chord for $`(X_\lambda ,)`$ in $`(\mathrm{\Sigma },\lambda )`$ and we choose the smooth $`a(s,t)`$ such that $`_0^sa(\tau ,t)𝑑\tau `$ and $`_s^0a(\tau ,t)𝑑\tau `$ exists, then $`F^{}`$ is an exact isotopy of Lagrangian embeddings(not regular). Moreover if $`a(s,0)a(s,1)`$, then $`F_0^{}(\times R)F_1^{}(\times R)=\mathrm{}`$.
Proof. Let $`F_t^{}=F^{}(,t):\times RR\times \mathrm{\Sigma }`$. It is obvious that $`F_t^{}`$ is an embedding. We check
$`F^{}(d(e^a\lambda ))`$ $`=`$ $`d(F^{}(e^a\lambda ))`$ (3.13)
$`=`$ $`d(e^{a(s,t)}G^{}\lambda )`$ (3.14)
$`=`$ $`d(e^{a(s,t)}ds)`$ (3.15)
$`=`$ $`e^{a(s,t)}(a_sds+a_tdt)ds`$ (3.16)
$`=`$ $`e^{a(s,t)}a_tdtds`$ (3.17)
which shows that $`F_t^{}`$ is a Lagrangian embedding for fixed $`t`$. Moreover for fixed $`t`$,
$`F_t^{}(e^a\lambda )`$ $`=`$ $`e^{a(s,t)}ds`$ (3.18)
$`=`$ $`\{\begin{array}{cc}d(_0^se^{a(\tau ,t)}𝑑\tau )\hfill & \text{for }s0\hfill \\ d(_s^0e^{a(\tau ,t)}𝑑\tau )\hfill & \text{for }s0\hfill \end{array}`$ (3.21)
which shows that $`F_t^{}`$ is an exact Lagrangian embedding, this proves Lemma3.3.
Now we take $`a(s,t)=\frac{a_0\epsilon t}{8}e^{s^2}`$ which satisfies the assumption in Lemma 3.3, then
$$F^{}:\times R\times [0,1]R\times \mathrm{\Sigma }$$
$$F^{}(l,s,t)=(a(s,t),\eta _s(G(l))).$$
Let
$$\psi _0(s,t)=se^{a(s,t)}a_s=2\frac{a_0\epsilon t}{8}e^{(\frac{a_0\epsilon t}{8}e^{s^2})s^2}s^2$$
(3.22)
$$\psi _1(s,t)=_{\mathrm{}}^s\psi _0(\tau ,t)𝑑\tau $$
(3.23)
$$\psi =\frac{\psi _1}{t}se^{a(s,t)}a_t$$
(3.24)
and compute
$`F_{}^{}{}_{}{}^{}(e^a\lambda )`$ $`=`$ $`e^{a(s,t)}ds`$ (3.25)
$`=`$ $`d(se^{a(s,t)})sde^{a(s,t)}`$ (3.26)
$`=`$ $`d(se^{a(s,t)})se^{a(s,t)}a_sdsse^{a(s,t)}a_tdt`$ (3.27)
$`=`$ $`d(se^{a(s,t)})d_s\psi _1se^{a(s,t)}a_tdt`$ (3.28)
$`=`$ $`d((se^{a(s,t)})\psi _1)+{\displaystyle \frac{\psi _1}{t}}dtse^{a(s,t)}a_tdt`$ (3.29)
$`=`$ $`d\mathrm{\Psi }^{}+{\displaystyle \frac{\psi _1}{t}}dtse^{a(s,t)}a_tdt`$ (3.30)
$`=`$ $`d\mathrm{\Psi }^{}\psi (s,t)dt`$ (3.31)
$`=`$ $`d\mathrm{\Psi }^{}\stackrel{~}{l}^{}`$ (3.32)
Let $`(V^{},\omega ^{})=(R\times \mathrm{\Sigma },d(e^a\lambda ))`$, $`W^{}=\times R`$, and $`(V,\omega )=(V^{}\times C,\omega ^{}\omega _0)`$. As in , we use figure eight trick invented by Gromov to construct a Lagrangian submanifold in $`V`$ through the Lagrange isotopy $`F^{}`$ in $`V^{}`$. Fix a positive $`\delta <1`$ and take a $`C^{\mathrm{}}`$-map $`\rho :S^1[0,1]`$, where the circle $`S^1`$ via parametrized by $`\mathrm{\Theta }[1,1]`$, such that the $`\delta `$neighborhood $`I_0`$ of $`0S^1`$ goes to $`0[0,1]`$ and $`\delta `$neighbourhood $`I_1`$ of $`\pm 1S^1`$ goes $`1[0,1]`$. Let
$`\stackrel{~}{l}`$ $`=`$ $`\psi (s,\rho (\mathrm{\Theta }))\rho ^{}(\mathrm{\Theta })d\mathrm{\Theta }`$ (3.33)
$`=`$ $`\mathrm{\Phi }d\mathrm{\Theta }`$ (3.34)
be the pull-back of the form $`\stackrel{~}{l}^{}=\psi (s,t)dt`$ to $`W^{}\times S^1`$ under the map $`(w^{},\mathrm{\Theta })(w^{},\rho (\mathrm{\Theta }))`$ and assume without loss of generality $`\mathrm{\Phi }`$ vanishes on $`W^{}\times (I_0I_1)`$.
Next, consider a map $`\alpha `$ of the annulus $`S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+]`$ into $`R^2`$, where $`\mathrm{\Phi }_{}`$ and $`\mathrm{\Phi }_+`$ are the lower and the upper bound of the fuction $`\mathrm{\Phi }`$ correspondingly, such that
$`(i)`$ The pull-back under $`\alpha `$ of the form $`dxdy`$ on $`R^2`$ equals $`d\mathrm{\Phi }d\mathrm{\Theta }`$.
$`(ii)`$ The map $`\alpha `$ is bijective on $`I\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+]`$ where $`IS^1`$ is some closed subset, such that $`II_0I_1=S^1`$; furthermore, the origin $`0R^2`$ is a unique double point of the map $`\alpha `$ on $`S^1\times 0`$, that is
$$0=\alpha (0,0)=\alpha (\pm 1,0),$$
and $`\alpha `$ is injective on $`S^1=S^1\times 0`$ minus $`\{0,\pm 1\}`$.
$`(iii)`$ The curve $`S_0^1=\alpha (S^1\times 0)R^2`$ “bounds” zero area in $`R^2`$, that is $`_{S_0^1}x𝑑y=0`$, for the $`1`$form $`xdy`$ on $`R^2`$.
###### Proposition 3.1
Let $`V^{}`$, $`W^{}`$ and $`F^{}`$ as above. Then there exists an exact Lagrangian embedding $`F:W^{}\times S^1V^{}\times R^2`$ given by $`F(w^{},\mathrm{\Theta })=(F^{}(w^{},\rho (\mathrm{\Theta })),\alpha (\mathrm{\Theta },\mathrm{\Phi }))`$.
Proof. We follow as in \[8, 2.3$`B_3^{}`$\]. Now let $`F^{}:W^{}\times S^1V^{}\times R^2`$ be given by $`(w^{},\mathrm{\Theta })(F^{}(w,\rho (\mathrm{\Theta })),\alpha (\mathrm{\Theta },\mathrm{\Phi }))`$. Then
$`(i)^{}`$ The pull-back under $`F^{}`$ of the form $`\omega =\omega ^{}+dxdy`$ equals $`d\stackrel{~}{l}^{}d\mathrm{\Phi }d\mathrm{\Theta }=0`$ on $`W^{}\times S^1`$.
$`(ii)^{}`$ The set of double points of $`F^{}`$ is $`W_0^{}W_1^{}V^{}=V^{}\times 0V^{}\times R^2`$.
$`(iii)^{}`$ If $`F^{}`$ has no double point then the Lagrangian submanifold $`W=F^{}(W^{}\times S^1)(V^{}\times R^2,\omega ^{}+dxdy)`$ is exact if and only if $`W_0^{}V^{}`$ is such.
This completes the proof of Proposition 3.1.
### 3.2 Formulation of Hilbert bundles
Let $`(\mathrm{\Sigma },\lambda )`$ be a closed $`(2n1)`$ dimensional manifold with a contact form $`\lambda `$. Let $`S\mathrm{\Sigma }=R\times \mathrm{\Sigma }`$ and put $`\xi =\mathrm{ker}(\lambda )`$. Let $`J_\lambda ^{}`$ be an almost complex structure on $`S\mathrm{\Sigma }`$ tamed by the symplectic form $`d(e^a\lambda )`$.
We define a metric $`g_\lambda `$ on $`S\mathrm{\Sigma }=R\times \mathrm{\Sigma }`$ by
$$g_\lambda =d(e^a\lambda )(,J_\lambda )$$
(3.35)
which is adapted to $`J_\lambda `$ and $`d(e^a\lambda )`$ but not complete.
In the following we denote by $`(V^{},\omega ^{})=((R\times \mathrm{\Sigma }),d(e^a\lambda ))`$ and $`(V,\omega )=(V^{}\times R^2,\omega ^{}+dxdy)`$ with the metric $`g=g^{}g_0`$ induced by $`\omega (,J)`$($`J=J^{}i`$ and $`WV`$ a Lagrangian submanifold which was constructed in section 3.1.
Let $`\overline{V}=D\times V`$, then $`\pi _1:\overline{V}D`$ be a symplectic vector bundle. Let $`\overline{J}`$ be an almost complex structure on $`\overline{V}`$ such that $`\pi _1:\overline{V}D`$ is a holomorphic map and each fibre $`\overline{V}_z=\pi _1(z)`$ is a $`\overline{J}`$ complex submanifold. Let $`H^k(D)`$ be the space of $`H^k`$maps from $`D`$ to $`\overline{V}`$, here $`H^k`$ represents Sobolev derivatives up to order $`k`$. Let $`\overline{W}=D\times W`$, $`\overline{p}=\{1\}\times p`$, $`W^\pm =\{\pm i\}\times W`$ and
$$𝒟^k=\{\overline{u}H^k(D)|\overline{u}(x)\overline{W}a.eforxDand\overline{u}(1)=\overline{p},\overline{u}(\pm i)\{\pm i\}\times W\}$$
for $`k100`$.
###### Lemma 3.4
Let $`W`$ be a closed Lagrangian submanifold in $`V`$. Then, $`𝒟^k`$ is a pseudo-Hilbert manifold with the tangent bundle
$$T𝒟^k=\underset{\overline{u}𝒟^k}{}\mathrm{\Lambda }^{k1}$$
(3.36)
here
$$\mathrm{\Lambda }^{k1}=\{\overline{w}H^{k1}(\overline{u}^{}(T\overline{V})|\overline{w}(1)=0,and\overline{w}(\pm i)T\overline{W}\}$$
###### Note 3.1
Since $`W`$ is not regular we know that $`𝒟^k`$ is in general complete, however it is enough for our purpose.
Proof: See .
Now we consider a section from $`𝒟^k`$ to $`T𝒟^k`$ follows as in , i.e., let $`\overline{}:𝒟^kT𝒟^k`$ be the Cauchy-Riemmann section
$$\overline{}\overline{u}=\frac{\overline{u}}{s}+J\frac{\overline{u}}{t}$$
(3.37)
for $`\overline{u}𝒟^k`$.
###### Theorem 3.1
The Cauchy-Riemann section $`\overline{}`$ defined in (3.37) is a Fredholm section of Index zero.
Proof. According to the definition of the Fredholm section, we need to prove that $`\overline{u}𝒟^k`$, the linearization $`D\overline{}(\overline{u})`$ of $`\overline{}`$ at $`\overline{u}`$ is a linear Fredholm operator. Note that
$$D\overline{}(\overline{u})=D\overline{}_{[\overline{u}]}$$
(3.38)
where
$$(D\overline{}_{[\overline{u}]})v=\frac{\overline{v}}{s}+J\frac{\overline{v}}{t}+A(\overline{u})\overline{v}$$
(3.39)
with
$$\overline{v}|_D(\overline{u}|_D)^{}T\overline{W}$$
here $`A(\overline{u})`$ is $`2n\times 2n`$ matrix induced by the torsion of almost complex structure, see for the computation.
Observe that the linearization $`D\overline{}(\overline{u})`$ of $`\overline{}`$ at $`\overline{u}`$ is equivalent to the following Lagrangian boundary value problem
$`{\displaystyle \frac{\overline{v}}{s}}+\overline{J}{\displaystyle \frac{\overline{v}}{t}}+A(\overline{u})\overline{v}=\overline{f},\overline{v}\mathrm{\Lambda }^k(\overline{u}^{}T\overline{V})`$ (3.40)
$`\overline{v}(t)T_{\overline{u}(t)}W,tD`$ (3.41)
One can check that (3.41) defines a linear Fredholm operator. In fact, by proposition 2.2 and Lemma 2.1, since the operator $`A(\overline{u})`$ is a compact, we know that the operator $`\overline{}`$ is a nonlinear Fredholm operator of the index zero.
###### Definition 3.1
Let $`X`$ be a Banach manifold and $`P:YX`$ the Banach vector bundle. A Fredholm section $`F:XY`$ is proper if $`F^1(0)`$ is a compact set and is called generic if $`F`$ intersects the zero section transversally, see .
###### Definition 3.2
$`deg(F,y)=\mathrm{}\{F^1(0)\}mod2`$ is called the Fredholm degree of a Fredholm section (see).
###### Theorem 3.2
Assum that $`\overline{J}=iJ`$ on $`\overline{V}`$ and $`i`$ is complex structure on $`D`$ and $`J`$ the almost complex structure on $`V`$ which is integrable near point $`p`$. Then the Fredholm section $`F=\overline{}:𝒟^kT𝒟^k`$ constructed in (3.37) has degree one, i.e.,
$$deg(F,0)=1$$
Proof: We assume that $`\overline{u}:D\overline{V}`$ be a $`\overline{J}`$holomorphic disk with boundary $`\overline{u}(D)\overline{W}`$ and by the assumption that $`\overline{u}`$ is homotopic to the map $`\overline{u}_1=(id,\overline{p})`$. Since almost complex structure $`\overline{J}`$ splits and is tamed by the symplectic form $`\overline{\omega }`$, by stokes formula, we conclude the second component $`u:DV`$ is a constant map. Because $`u(1)=p`$, We know that $`F^1(0)=(id,p)`$. Next we show that the linearizatioon $`DF_{(id,p)}`$ of $`F`$ at $`(id,p)`$ is an isomorphism from $`T_{(id,p)}𝒟^k`$ to $`E`$. This is equivalent to solve the equations
$`{\displaystyle \frac{\overline{v}}{s}}+J{\displaystyle \frac{\overline{v}}{t}}=f`$ (3.42)
$`\overline{v}|_DT_{(id,p)}\overline{W}`$ (3.43)
here $`\overline{J}=i+J(p)`$. By Lemma 2.1, we know that $`DF((id,p))`$ is an isomorphism. Therefore $`deg(F,0)=1`$.
## 4 Anti-holomorphic sections
In this section we construct a Fredholm section which is not proper as in .
Let $`(V^{},\omega ^{})=(S\mathrm{\Sigma },d(e^a\lambda ))`$ and $`(V,\omega )=(V^{}\times C,\omega ^{}\omega _0)`$, $`W`$ as in section3 and $`J=J^{}i`$, $`g=g^{}g_0`$, $`g_0`$ the standard metric on $`C`$.
Now let $`cC`$ be a non-zero vector. We consider $`c`$ as an anti-holomorphic homomorphism $`c:TDTV^{}TC`$, i.e., $`c(\frac{}{\overline{z}})=(0,c\frac{}{z}).`$ Since the constant section $`c`$ is not a section of the Hilbert bundle in section 3 due to $`c`$ is not tangent to the Lagrangian submanifold $`W`$, we must modify it as follows:
Let $`c`$ as above, we define
$`c_{\chi ,\delta }(z,v)=\{\begin{array}{cc}c\hfill & \text{if }|z|12\delta \text{,}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ (4.3)
Then by using the cut off function $`\phi _h(z)`$ and its convolution with section $`c_{\chi ,\delta }`$, we obtain a smooth section $`c_\delta `$ satisfying
$`c_\delta (z,v)=\{\begin{array}{cc}c\hfill & \text{if }|z|13\delta \text{,}\hfill \\ 0\hfill & \text{if }|z|1\delta \text{.}\hfill \end{array}`$ (4.6)
$`|c_\delta ||c|`$ (4.7)
for $`h`$ small enough, for the convolution theory see \[12, ch1,p16-17,Th1.3.1\]. Then one can easily check that $`\overline{c}_\delta =(0,0,c_\delta )`$ is an anti-holomorphic section tangent to $`\overline{W}`$.
Now we put an almost complex structure $`\overline{J}=iJ`$ on the symplectic fibration $`D\times VD`$ such that $`\pi _1:D\times VD`$ is a holomorphic fibration and $`\pi _1^1(z)`$ is an almost complex submanifold. Let $`g=\overline{\omega }(,\overline{J})`$ be the metric on $`D\times V`$.
Now we consider the equations
$`\overline{v}=(id,v)=(id,v^{},f):DD\times V^{}\times C`$ (4.8)
$`\overline{}_Jv=c_\delta or`$ (4.9)
$`\overline{}_J^{}v^{}=0,\overline{}f=c_\delta onD`$ (4.10)
$`v|_D:DW`$ (4.11)
here $`v`$ homotopic to constant map $`\{p\}`$ relative to $`W`$. Note that $`WV\times B_R(0)`$ for $`\pi R^2=2\pi R(\epsilon )^2`$, here $`R(\epsilon )0`$ as $`\epsilon 0`$ and $`\epsilon `$ as in section 3.1.
###### Lemma 4.1
Let $`\overline{v}=(id,v)`$ be the solutions of (4.11), then one has the following estimates
$`E(v)=\{{\displaystyle _D}(g^{}({\displaystyle \frac{v^{}}{x}},J^{}{\displaystyle \frac{v^{}}{x}})+g^{}({\displaystyle \frac{v^{}}{y}},J^{}{\displaystyle \frac{v^{}}{y}})`$
$`+g_0({\displaystyle \frac{f}{x}},i{\displaystyle \frac{f}{x}})+g_0({\displaystyle \frac{f}{y}},i{\displaystyle \frac{f}{y}}))d\sigma \}4\pi R(\epsilon )^2.`$ (4.12)
Proof: Since $`v(z)=(v^{}(z),f(z))`$ satisfy (4.11) and $`v(z)=(v^{}(z),f(z))V^{}\times C`$ is homotopic to constant map $`v_0:D\{p\}W`$ in $`(V,W)`$, by the Stokes formula
$$_Dv^{}(\omega ^{}\omega _0)=0$$
(4.13)
Note that the metric $`g`$ is adapted to the symplectic form $`\omega `$ and $`J`$, i.e.,
$$g=\omega (,J)$$
(4.14)
By the simple algebraic computation, we have
$$_Dv^{}\omega =\frac{1}{4}_{D^2}(|v|^2|\overline{}v|^2)=0$$
(4.15)
and
$$|v|=\frac{1}{2}(|v|^2+|\overline{}v|^2$$
(4.16)
Then
$`E(v)`$ $`=`$ $`{\displaystyle _D}|v|`$ (4.17)
$`=`$ $`{\displaystyle _D}\{{\displaystyle \frac{1}{2}}(|v|^2+|\overline{}v|^2)\})d\sigma `$
$`=`$ $`{\displaystyle _D}|c_\delta |_{\overline{g}}^2𝑑\sigma `$
By Cauchy integral formula,
$$f(z)=\frac{1}{2\pi i}_D\frac{f(\xi )}{\xi z}𝑑\xi +\frac{1}{2\pi i}_D\frac{\overline{}f(\xi )}{\xi z}𝑑\xi d\overline{\xi }$$
(4.18)
Since $`f`$ is smooth up to the boundary, we integrate the two sides on $`D_r`$ for $`r<1`$, one get
$`{\displaystyle _{D_r}}f(z)𝑑z`$ $`=`$ $`{\displaystyle _{D_r}}{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}{\displaystyle \frac{f(\xi )}{\xi z}}𝑑\xi 𝑑z+{\displaystyle _{D_r}}{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}{\displaystyle \frac{\overline{}f(\xi )}{\xi z}}𝑑\xi d\overline{\xi }`$ (4.19)
$`=`$ $`0+{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}{\displaystyle _{D_r}}{\displaystyle \frac{\overline{}f(\xi )}{\xi z}}𝑑z𝑑\xi d\overline{\xi }`$ (4.20)
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _D}2\pi i\overline{}f(\xi )𝑑\xi d\overline{\xi }`$ (4.21)
Let $`r1`$, we get
$`{\displaystyle _D}f(z)𝑑z={\displaystyle _D}\overline{}f(\xi )𝑑\xi d\overline{\xi }`$ (4.22)
By the equations (4.11), one get
$$\overline{}f=conD_{12\delta }$$
(4.23)
So, we have
$$2\pi i(12\delta )c=_Df(z)𝑑z_{DD_{12\delta }}\overline{}f(\xi )𝑑\xi d\overline{\xi }$$
(4.24)
So,
$`|c|`$ $``$ $`{\displaystyle \frac{1}{2\pi (12\delta )}}|{\displaystyle _D}f(z)𝑑z|+|{\displaystyle _{DD_{12\delta }}}\overline{}f(\xi )𝑑\xi d\overline{\xi }|`$ (4.25)
$``$ $`{\displaystyle \frac{1}{2\pi (12\delta )}}2\pi |diam(pr_2(W))+c_1c_2|c|(\pi \pi (12\delta )^2))`$ (4.26)
Therefore, one has
$`|c|`$ $``$ $`c(\delta )R(\epsilon )`$ (4.27)
and
$`E(v)`$ $`=`$ $`\pi {\displaystyle _D}|c_\delta |_{\overline{g}}^2`$ (4.28)
$`=`$ $`\pi c(\delta )^2R(\epsilon )^2.`$ (4.29)
This finishes the proof of Lemma.
###### Proposition 4.1
For $`|c|2c(\delta )R(\epsilon )`$, then the equations (4.11) has no solutions.
Proof. By 4.27, it is obvious.
###### Theorem 4.1
The Fredholm section $`F_1=\overline{}_{\overline{J}}+\overline{c}_\delta :𝒟^kE`$ is not proper.
Proof. By the Proposition 4.1 and Theorem 3.2, it is obvious(see).
## 5 $`J`$holomorphic section
Recall that $`W(K,K)WV^{}\times R^2`$ as in section 3. The Riemann metric $`g`$ on $`V^{}\times R^2`$ induces a metric $`g|W`$.
Now let $`cC`$ be a non-zero vector and $`c_\delta `$ the induced anti-holomorphic section. We consider the nonlinear inhomogeneous equations (4.11) and transform it into $`\overline{J}`$holomorphic map by considering its graph as in \[8, p319,1.4.C\] or \[4, p312,Lemma5.2.3\].
Denote by $`Y^{(1)}D\times V`$ the bundle of homomorphisms $`T_s(D)T_v(V)`$. If $`D`$ and $`V`$ are given the disk and the almost Kähler manifold, then we distinguish the subbundle $`X^{(1)}Y^{(1)}`$ which consists of complex linear homomorphisms and we denote $`\overline{X}^{(1)}D\times V`$ the quotient bundle $`Y^{(1)}/X^{(1)}`$. Now, we assign to each $`C^1`$-map $`v:DV`$ the section $`\overline{}v`$ of the bundle $`\overline{X}^{(1)}`$ over the graph $`\mathrm{\Gamma }_vD\times V`$ by composing the differential of $`v`$ with the quotient homomorphism $`Y^{(1)}\overline{X}^{(1)}`$. If $`c_\delta :D\times V\overline{X}`$ is a $`H^k`$ section we write $`\overline{}v=c_\delta `$ for the equation $`\overline{}v=c_\delta |\mathrm{\Gamma }_v`$.
###### Lemma 5.1
(Gromov\[8, 1.4.$`C^{}`$\])There exists a unique almost complex structure $`J_g`$ on $`D\times V`$(which also depends on the given structures in $`D`$ and in $`V`$), such that the (germs of) $`J_\delta `$holomorphic sections $`v:DD\times V`$ are exactly and only the solutions of the equations $`\overline{}v=c_\delta `$. Furthermore, the fibres $`z\times VD\times V`$ are $`J_\delta `$holomorphic( i.e. the subbundles $`T(z\times V)T(D\times V)`$ are $`J_\delta `$complex) and the structure $`J_\delta |z\times V`$ equals the original structure on $`V=z\times V`$. Moreover $`J_\delta `$ is tamed by $`k\omega _0\omega `$ for $`k`$ large enough which is independent of $`\delta `$.
## 6 Gromov’s $`C^0`$convergence theorem
### 6.1 Analysis of Gromov’s figure eight
Since $`W^{}S\mathrm{\Sigma }`$ is an exact Lagrangian submanifold and $`F_t^{}`$ is an exact Lagrangian isotopy(see section 3.1). Now we carefully check the Gromov’s construction of Lagrangian submanifold $`WV^{}\times R^2`$ from the exact Lagrangian isotopy of $`W^{}`$ in section 3.
Let $`S^1T^{}S^1`$ be a zero section and $`S^1=_{i=1}^4S_i`$ be a partition of the zero section $`S^1`$ such that $`S_1=I_0`$, $`S_3=I_1`$. Write $`S^1\{I_0I_1\}=I_2I_3`$ and $`I_0=(\delta ,\frac{5}{6}\delta ](\frac{5\delta }{6},+\frac{5\delta }{6})[\frac{5\delta }{6},\delta )=I_0^{}I_0^{}I_0^+`$, similarly $`I_1=(1\delta ,1\frac{5}{6}\delta ](1\frac{5\delta }{6},1+\frac{5\delta }{6})[1+\frac{5\delta }{6},1+\delta )=I_1^{}I_1^{}I_1^+`$. Let $`S_2=I_0^+I_2I_1^{}`$, $`S_4=I_1^+I_3I_0^+`$. Moreover, we can assume that the double points of map $`\alpha `$ in Gromov’s figure eight is contained in $`(\overline{I}_0^{}\overline{I}_1^{})\times [\mathrm{\Phi }_{},\mathrm{\Phi }_+]`$, here $`\overline{I}_0^{}=(\frac{5\delta }{12},+\frac{5\delta }{12})`$ and $`\overline{I}_1^{}=(1\frac{5\delta }{12},1+\frac{5\delta }{12})`$. Recall that $`\alpha :(S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$ is an exact symplectic immersion, i.e., $`\alpha ^{}(ydx)\mathrm{\Psi }d\mathrm{\Theta }=dh`$, $`h:T^{}S^1R`$. By the construction of figure eight, we can assume that $`\alpha _i^{}=\alpha |((S^1I_i^{})\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])`$ is an embedding for $`i=0,1`$. Let $`Y=\alpha (S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$ and $`Y_i=\alpha (S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`\alpha _i=\alpha |Y_i(S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])`$. So, $`\alpha _i`$ puts the function $`h`$ to the function $`h_{i0}=\alpha _i^1h`$ on $`Y_i`$. We extend the function $`h_{i0}`$ to whole plane $`R^2`$. In the following we take the liouville form $`\beta _{i0}=ydxdh_{i0}`$ on $`R^2`$. This does not change the symplectic form $`dxdy`$ on $`R^2`$. But we have $`\alpha _i^{}\beta =\mathrm{\Phi }d\mathrm{\Theta }`$ on $`(S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])`$ for $`i=1,2,3,4`$. Finally, note that
$`F:W^{}\times S^1V^{}\times R^2;`$ (6.1)
$`F(w^{},\mathrm{\Theta })=(F_{\rho (\mathrm{\Theta })}^{}(w^{}),\alpha (\mathrm{\Theta },\mathrm{\Phi }(w^{},\rho (\mathrm{\Theta })).`$ (6.2)
Since $`\rho (\mathrm{\Theta })=0`$ for $`\mathrm{\Theta }I_0`$ and $`\rho (\mathrm{\Theta })=1`$ for $`\mathrm{\Theta }I_1`$, we know that $`\mathrm{\Phi }(w^{},\rho (\mathrm{\Theta }))=0`$ for $`\mathrm{\Theta }I_0I_1`$. Therefore,
$`F(W^{}\times I_0)=W^{}\times \alpha (I_0);F(W^{}\times I_1)=W^{}\times \alpha (I_1).`$ (6.3)
### 6.2 Gromov’s Schwartz lemma
In our proof we need a crucial tools, i.e., Gromov’s Schwartz Lemma as in . We first consider the case without boundary.
###### Proposition 6.1
Let $`(V,J,\mu )`$ be as in section 4 and $`V_K`$ the compact part of $`V`$. There exist constants $`\epsilon _0`$ and $`C`$(depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^\alpha `$ norm of $`J`$ and $`A_0`$) such that every $`J`$holomorphic map of the unit disc to an $`\epsilon _0`$-ball of $`V`$ with center in $`V_K`$ and area less than $`A_0`$ has its derivatives up to order $`k+1+\alpha `$ on $`D_{\frac{1}{2}}(0)`$ bounded by $`C`$.
For a proof, see.
Now we consider the Gromov’s Schwartz Lemma for $`J`$holomorphic map with boundary in a closed Lagrangian submanifold as in .
###### Proposition 6.2
Let $`(V,J,\mu )`$ as above and $`LV`$ be a closed Lagrangian submanifold and $`V_K`$ one compact part of $`V`$. There exist constants $`\epsilon _0`$ and $`C`$(depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^\alpha `$ norm of $`J`$ and $`K,A_0`$) such that every $`J`$holomorphic map of the half unit disc $`D^+`$ to a $`\epsilon _0`$-ball of $`V`$ with boundary in $`L`$ and area less than $`A_0`$ has its derivatives up to order $`k+1+\alpha `$ on $`D_{\frac{1}{2}}^+(0)`$ bounded by $`C`$.
For a proof see .
Since in our case $`W`$ is a non-compact Lagrangian submanifold, Proposition 6.2 can not be used directly but the proofs of Proposition 6.1-2 still holds in our case.
###### Lemma 6.1
Recall that $`V=V^{}\times R^2`$. Let $`(V,J,\mu )`$ as above and $`WV`$ be as above and $`V_c`$ the compact set in $`V`$. Let $`\overline{V}=D\times V`$, $`\overline{W}=D\times W`$, and $`\overline{V}_c=D\times V_c`$. Let $`Y=\alpha (S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`Y_i=\alpha (S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`\{X_j\}_{j=1}^q`$ be a Darboux covering of $`\mathrm{\Sigma }`$ and $`V_j^{}=R\times X_j`$. Let $`D=S^{1+}S^1`$. There exist constant $`c_0`$ such that every $`J`$holomorphic map $`v`$ of the half unit disc $`D^+`$ to the $`D\times V_j^{}\times R^2`$ with its boundary $`v((1,1))(S^{1\pm })\times F(\times R\times S_i)\overline{W},i=1,..,4`$ has
$$area(v(D^+))c_0l^2(v(^{}D^+)).$$
(6.4)
here $`^{}D^+=D[1,1]`$ and $`l(v(^{}D^+))=length(v(^{}D^+))`$.
Proof. Let $`\overline{W}_{i\pm }=S^{1\pm }\times F(W^{}\times S_i)`$. Let $`v=(v_1,v_2):D^+\overline{V}=D\times V`$ be the $`J`$holomorphic map with $`v(D^+)\overline{W}_{i\pm }D\times W`$, then
$`area(v)`$ $`=`$ $`{\displaystyle _{D^+}}v^{}d(\alpha _0\alpha )`$ (6.5)
$`=`$ $`{\displaystyle _{D^+}}𝑑v^{}(\alpha _0\alpha )`$ (6.6)
$`=`$ $`{\displaystyle _{D^+}}v^{}(\alpha _0\alpha )`$ (6.7)
$`=`$ $`{\displaystyle _{D^+}}v_1^{}\alpha _0+{\displaystyle _{D^+}}v_2^{}\alpha `$ (6.8)
$`=`$ $`{\displaystyle _{^{}D^+[1,+1]}}v_1^{}\alpha _0+{\displaystyle _{^{}D^+[1,+1]}}v_2^{}(e^a\lambda ydxdh_{i0})`$ (6.9)
$`=`$ $`{\displaystyle _{^{}D^+[1,+1]}}v_1^{}\alpha _0+{\displaystyle _{^{}D^+}}v_2^{}(e^a\lambda ydxdh_{i0})+B_1,`$ (6.10)
here $`B_1=_{[1,+1]}v_{2}^{}{}_{}{}^{}(d\mathrm{\Psi }^{})`$. Now take a zig-zag curve $`C`$ in $`V_j^{}\times Y_i`$ connecting $`v_2(1)`$ and $`v_2(+1)`$ such that
$`{\displaystyle _C}(e^a\lambda +ydx)`$ $`=`$ $`B_1`$ (6.11)
$`length(C)`$ $``$ $`k_1length(v_2(^{}D^+))`$ (6.12)
Now take a minimal surface $`M`$ in $`V_j^{}\times R^2`$ bounded by $`v_2(^{}D^+)C`$, then by the isoperimetric ineqality(see\[\[9, p283\]), we get
$`area(M)`$ $``$ $`m_1length(C+v_2(^{}D^+))^2`$ (6.13)
$``$ $`m_2length(v_2(^{}D^+))^2,`$ (6.14)
here we use the (6.12).
Since $`area(M)_M\omega `$ and $`_M\omega =_{D^+}v_2^{}\omega =area(v)`$, this proves the lemma.
###### Lemma 6.2
Let $`v`$ as in Lemma 6.1, then we have
$$area(v(D^+)c_0(dist(v(0),v(^{}D^+)))^2,$$
(6.15)
here $`c_0`$ depends only on $`\mathrm{\Sigma },J,\omega ,\mathrm{},`$etc, not on $`v`$.
Proof. By the standard argument as in \[4, p79\].
The following estimates is a crucial step in our proof.
###### Lemma 6.3
Recall that $`V=V^{}\times R^2`$. Let $`(V,J,\mu )`$ as above and $`WV`$ be as above and $`V_c`$ the compact set in $`V`$. Let $`\overline{V}=D\times V`$, $`\overline{W}=D\times W`$, and $`\overline{V}_c=D\times V_c`$. Let $`Y=\alpha (S^1\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`Y_i=\alpha (S_i\times [5\mathrm{\Phi }_{},5\mathrm{\Phi }_+])R^2`$. Let $`D=S^{1+}S^1`$. There exist constant $`c_0`$ such that every $`J`$holomorphic map $`v`$ of the half unit disc $`D^+`$ to the $`D\times V^{}\times R^2`$ with its boundary $`v((1,1))(S^{1\pm })\times F(\times R\times S_i)\overline{W},i=1,..,4`$ has
$$area(v(D^+))c_0l^2(v(^{}D^+)).$$
(6.16)
here $`^{}D^+=D[1,1]`$ and $`l(v(^{}D^+))=length(v(^{}D^+))`$.
Proof. We first assume that $`\epsilon `$ in section 3.1 is small enough. Let $`l_0`$ is a constant small enough. If $`length(^{}D^+)l_0`$, then Lemma 6.3 holds. If $`length(^{}D^+)l_0`$ and $`v(D^+)D\times V_j^{}\times R^2`$, then Lemma6.3 reduces to Lemma6.1. If $`length(^{}D^+)l_0`$ and $`v(D^+)\overline{}D\times V_j^{}\times R^2`$, then Lemma6.2 imples $`area(v)\tau _0>100\pi R(\epsilon )^2`$, this is a contradiction. Therefore we proved the lemma.
###### Proposition 6.3
Let $`(V,J,\mu )`$ and $`WV`$ be as in section 4 and $`V_K`$ the compact part of $`V`$. Let $`\overline{V}`$, $`\overline{V}_K`$ and $`\overline{W}`$ as section 5.1. There exist constants $`\epsilon _0`$ (depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^\alpha `$ norm of $`J`$) and $`C`$(depending only on the $`C^0`$ norm of $`\mu `$ and on the $`C^{k+\alpha }`$ norm of $`J`$) such that every $`J`$holomorphic map of the half unit disc $`D^+`$ to the $`D\times V^{}\times R^2`$ with its boundary $`v((1,1))(S^{1\pm })\times F(\times R\times S_i)\overline{W},i=1,..,4`$ has its derivatives up to order $`k+1+\alpha `$ on $`D_{\frac{1}{2}}^+(0)`$ bounded by $`C`$.
Proof. One uses Lemma 6.3 and Gromov’s proof on Schwartz lemma to yield proposition 6.3.
### 6.3 Removal singularity of $`J`$curves
In our proof we need another crucial tools, i.e., Gromov’s removal singularity theorem. We first consider the case without boundary.
###### Proposition 6.4
Let $`(V,J,\mu )`$ be as in section 4 and $`V_K`$ the compact part of $`V`$. If $`v:D\{0\}V_K`$ be a $`J`$holomorphic disk with bounded energy and bounded image, then $`v`$ extends to a $`J`$holomorphic map from the unit disc $`D`$ to $`V_K`$.
For a proof, see.
Now we consider the Gromov’s removal singularity theorem for $`J`$holomorphic map with boundary in a closed Lagrangian submanifold as in .
###### Proposition 6.5
Let $`(V,J,\mu )`$ as above and $`LV`$ be a closed Lagrangian submanifold and $`V_K`$ one compact part of $`V`$. If $`v:(D^+\{0\},^{\prime \prime }D^+\{0\})(V_K,L)`$ be a $`J`$holomorphic half-disk with bounded energy and bounded image, then $`v`$ extends to a $`J`$holomorphic map from the half unit disc $`(D^+,^{\prime \prime }D^+)`$ to $`(V_K,L)`$.
For a proof see .
###### Proposition 6.6
Let $`(V,J,\mu )`$ and $`WV`$ be as in section 4 and $`V_c`$ the compact set in $`V`$. Let $`\overline{V}=D\times V`$, $`\overline{W}=D\times W`$, and $`\overline{V}_c=D\times V_c`$. Then every $`J`$holomorphic map $`v`$ of the half unit disc $`D^+\{0\}`$ to the $`\overline{V}`$ with center in $`\overline{V}_c`$ and its boundary $`v((1,1)\{0\})(S^{1\pm })\times F(\times [K,K]\times S_i)\overline{W}`$ and
$$area(v(D^+\{0\}))E$$
(6.17)
extends to a $`J`$holomorphic map $`\stackrel{~}{v}:(D^+,^{\prime \prime }D)(\overline{V}_c,\overline{W})`$.
Proof. This is ordinary Gromov’s removal singularity theorem by $`K`$assumption.
### 6.4 $`C^0`$Convergence Theorem
We now recall that the well-known Gromov’s compactness theorem for cusp’s curves for the compact symplectic manifolds with closed Lagrangian submanifolds in it. For reader’s convenience, we first recall the “weak-convergence” for closed curves.
Cusp-curves. Take a system of disjoint simple closed curves $`\gamma _i`$ in a closed surface $`S`$ for $`i=1,\mathrm{},k`$, and denote by $`S^0`$ the surface obtained from $`S_{i=1}^k\gamma _i`$. Denote by $`\overline{S}`$ the space obtained from $`S`$ by shrinking every $`\gamma _i`$ to a single point and observe the obvious map $`\alpha :S^0\overline{S}`$ gluing pairs of points $`s_i^{}`$ and $`s_i^{\prime \prime }`$ in $`S^0`$, such that $`\overline{s}_i=\alpha (s_i^{})=\alpha (s_i^{\prime \prime })\overline{S}`$ are singular (or cuspidal) points in $`\overline{S}`$(see).
An almost complex structure in $`\overline{S}`$ by definition is that in $`S^0`$. A continuous map $`\beta :\overline{S}V`$ is called a (parametrized $`J`$holomorphic) cusp-curve in $`V`$ if the composed map $`\beta \alpha :S^0V`$ is holomorphic.
Weak convergence. A sequence of closed $`J`$curves $`C_jV`$ is said to weakly converge to a cusp-curve $`\overline{C}V`$ if the following four conditions are satisfied
(i) all curves $`C_j`$ are parametrized by a fixed surface $`S`$ whose almost complex structure depends on $`j`$, say $`C_j=f_j(S)`$ for some holomorphic maps
$$f_j:(S,J_j)(V,J)$$
(ii) There are disjoint simple closed curves $`\gamma _iS`$, $`i=1,\mathrm{},k`$, such that $`\overline{C}=\overline{f}(\overline{S})`$ for a map $`\overline{f}:\overline{S}V`$ which is holomorphic for some almost complex structure $`\overline{J}`$ on $`\overline{S}`$.
(iii) The structures $`J_j`$ uniformly $`C^{\mathrm{}}`$converge to $`\overline{J}`$ on compact subsets in $`S_{i=1}^k\gamma _i`$.
(iv) The maps $`f_j`$ uniformly $`C^{\mathrm{}}`$converge to $`\overline{f}`$ on compact subsets in $`S_{i=1}^k\gamma _i`$. Moreover, $`f_j`$ uniformly $`C^0`$converge on entire $`S`$ to the composed map $`S\overline{S}\stackrel{\overline{f}}{}V`$. Furthermore,
$$Area_\mu f_j(S)Area_\mu \overline{f}(\overline{S})forj\mathrm{},$$
where $`\mu `$ is a Riemannian metric in $`V`$ and where the area is counted with the geometric multiplicity(see).
Gromov’s Compactness theorem for closed curves. Let $`C_j`$ be a sequence of closed $`J`$curves of a fixed genus in a compact manifold $`(V,J,\mu ).`$ If the areas of $`C_j`$ are uniformly bounded,
$$Area_\mu A,j=1,..,$$
then some subsequence weakly converges to a cusp-curve $`\overline{C}`$ in $`V`$.
Cusp-curves with boundary. Let $`T`$ be a compact complex manifold with boundary of dimension $`1`$(i.e., it has an atlas of holomorphic charts onto open subsets of $`C`$ or of a closed half plane). Its double is a compact Riemann surface $`S`$ with a natureal anti-holomorphic involution $`\tau `$ which exchanges $`T`$ and $`ST`$ while fixing the boundary $`T`$. IF$`f:TV`$ is a continous map, holomorphic in the interior of $`T`$, it is convenient to extend $`f`$ to $`S`$ by
$$f=f\tau $$
Take a totally real submanifold $`W(V,J)`$ and consider compact holomorphic curves $`CV`$ with boundaries, $`(\overline{C},\overline{C})(V,W)`$, which are, topologically speaking, obtained by shrinking to points some (short) closed loops in $`C`$ and also some (short) segments in $`C`$ between boundary points. This is seen by looking on the double $`C_CC`$.
Gromov’s Compactness theorem for curves with boundary. Let $`V`$ be a closed Riemannian manifold, $`W`$ a totally real closed submanifold of $`V`$. Let $`C_j`$ be a sequence of $`J`$curves with boundary in $`W`$ of a fixed genus in a compact manifold $`(V,J,\mu )`$. If the areas of $`C_j`$ are uniformly bounded,
$$Area_\mu A,j=1,..,$$
then some subsequence weakly converges to a cusp-curve $`\overline{C}`$ in $`V`$.
The proofs of Gromov’s compactness theorem can found in . In our case the Lagrangian submanifold $`W`$ is not compact, Gromov’s compactness theorem can not be applied directly but its proof is still effective since the $`W`$ has the special geometry. In the following we modify Gromov’s proof to prove the $`C^0`$compactness theorem in our case.
Now we state the $`C^0`$convergence theorem in our case.
###### Theorem 6.1
Let $`(V,J,\omega ,\mu )`$ and $`W`$ as in section4. Let $`C_j`$ be a sequence of $`\overline{J}_\delta `$holomorphic section $`v_j=(id,((a_j,u_j),f_j)):DD\times V`$ with $`v_j:DD\times W`$ and $`v_j(1)=(1,p)D\times W`$ constructed from section 4. Then the areas of $`C_j`$ are uniformly bounded,i.e.,
$$Area_\mu (C_j)A,j=1,..,$$
and some subsequence weakly converges to a cusp-section $`\overline{C}`$ in $`V`$(see).
Proof. We follow the proofs in . Write $`v_j=(id,(a_j,u_j),f_j))`$ then $`|a_{ij}|a_0`$ by the ordinary Monotone inequality of minimal surface without boundary, see following Proposition 7.1. Similarly $`|f_j|R_1`$ by using the fact $`f_j(D)`$ is bounded in $`B_{R_1}(0)`$ and $`_D|f_j|4\pi R^2`$ via monotone inequality for minimal surfaces. So, we assume that $`v_j(D)V_c`$ for a compact set $`V_c`$.
1. Removal of a net.
1a. Let $`\overline{V}=D\times V`$ and $`v_j`$ be regular curves. First we study induced metrics $`\mu _j`$ in $`v_j`$. We apply the ordinary monotone inequality for minimal surfaces without boundary to small concentric balls $`B_\epsilon (A_j,\mu _j)`$ for $`0<\epsilon \epsilon _0`$ and conclude by the standard argument to the inequality
$$Area(B_\epsilon )\epsilon ^2,for\epsilon \epsilon _0;$$
Using this we easily find a interior $`\epsilon `$net $`F_j(v_j,\mu _j)`$ containing $`N`$ points for a fixed integer $`N=(\overline{V},\overline{J},\mu )`$, such that every topological annulus $`Av_jF_j`$ satisfies
$$Diam_\mu A10length_\mu A.$$
(6.18)
Furthermore, let $`A`$ be conformally equivalent to the cylinder $`S^1\times [0,l]`$ where $`S^1`$ is the circle of the unit length, and let $`S_t^1A`$ be the curve in $`A`$ corresponding to the circle $`S^1\times t`$for $`t[0,l]`$. Then obviously
$$_a^b(lengthS_t^1)^2𝑑tArea(A)C_5.$$
(6.19)
for all $`[a,b][0,l]`$. Hence, the annulus $`A_tA`$ between the curves $`S_t^1`$ and $`S_{lt}^1`$ satisfies
$$diam_\mu A_t20(\frac{C_5}{t})$$
(6.20)
for all $`t[0,l]`$.
$`1b`$. We consider the sets $`v_j((S^{1\pm })\times F(W^{}\times I_i^\pm )),i=0,1`$. By the construction of Gromov’s figure eight, there exists a finite components, denote it by
$$v_j((S^{1\pm })\times F(\times R\times I_i^\pm ))=\{\overline{\gamma }_{ij}^k\},i=0,1.$$
(6.21)
Let $`m_i^\pm `$ be the middle point of $`I_i^\pm `$. If
$$\overline{\gamma }_{ij}^k((S^{1\pm })\times F(\times R\times m_i^\pm ))\mathrm{},i=0,1,$$
(6.22)
we choose one point in $`\overline{\gamma }_{ij}^k`$ as a boundary puncture point in $`v_j`$. Consider the concentric $`\epsilon `$ half-disks or quadrature $`B_\epsilon (p)`$ with center $`p`$ on $`\overline{\gamma }_{ij}^k`$, then
$$Area(B_\epsilon (p))\tau _0.$$
(6.23)
Since $`Area(v_j)E_0`$, there exists a uniform finite puncture points.
Consider the concentric $`\epsilon `$ half-disks or quadrature $`B_\epsilon (p)`$ with center $`p`$ on $`v_j`$ and
$$Area(B_\epsilon (p))\tau _0,$$
(6.24)
we puncture one point on such half-disk or quadrature. Since $`Area(v_j)E_0`$, there exists a uniform finite puncture points.
So, we find a boundary net $`G_jv_j`$ containing $`N_1`$ points for a fixed integer $`N_1(\overline{V},\overline{J},\mu )`$, such that every topological quadrature or half annulus $`Bv_j\{F_j,G_j\}`$ satisfies
$`^{\prime \prime }B=B\overline{W}(S^{1\pm })\times F(\times R\times S_i),i=1,2,3,4.`$ (6.25)
2. Poincare’s metrics. 2a. Now, let $`\mu _j^{}`$ be a metric of constant curvature $`1`$ in $`v_j(D)F_jG_j`$ conformally equivalent to $`\mu _j`$. Then for every $`\mu _j^{}`$ball $`B_\rho `$ in $`v_jF_jG_j`$ of radius $`\rho 0.1`$, there exists an annulus $`A`$ contained in $`v_jF_jG_j`$ such that $`B_\rho A_t`$ for $`t=0.01|log|`$(see Lemma 3.2.2in \[4, chVIII\]). This implies with $`(6.3)`$ the uniform continuity of the (inclusion) maps $`(v_jF_j,\mu _j^{})(\overline{V},\overline{\mu })`$, and hence a uniform bound on the $`r^{th}`$ order differentials for every $`r=0,1,2,\mathrm{}`$.
2b. Similarly, for every $`\mu _j^{}`$half ball $`B_\rho ^+`$ in $`v_jF_jG_j`$ of radius $`\rho 0.1`$, there exists a half annulus or quadrature $`B`$ contained in $`v_jF_jG_j`$ such that $`B_\rho ^+B`$ with
$`^{\prime \prime }B=B\overline{W}(S^{1\pm })\times F(\times R\times S_i),i=1,2,3,4.`$ (6.26)
Then, by Gromov’s Schwartz Lemma, i.e., Proposition 6.1-6.3 implies the uniform bound on the $`r^{th}`$ order differentials for every $`r=0,1,2,\mathrm{}`$.
3. Convergence of metrics. Next, by the standard (and obvious ) properties of hyperbolic surfaces there is a subsequence(see), which is still denoted by $`v_j`$, such that
$`(a)`$. There exist $`k`$ closed geodesics or geodesic arcs with boundaries in $`v_jF_j`$, say
$$\gamma _i^j(v_jF_j,\mu _j^{}),i=1,\mathrm{},k,j=1,2,\mathrm{},$$
whose $`\mu _j^{}`$length converges to zero as $`j\mathrm{}`$, where $`k`$ is a fixed integer.
$`(b)`$. There exist $`k`$ closed curves or geodesic arcs with boundaries in $`S`$ of a fixed surface, say $`\gamma _j`$ in $`S`$, and an almost complex structure $`\overline{J}`$ on the corresponding (singular) surface $`\overline{S}`$, such that the almost complex structure $`J_j`$ on $`v_jF_j`$ induced from $`(V,J)`$ $`C^{\mathrm{}}`$converge to $`\overline{J}`$ outside $`_{j=1}^k\gamma _j`$. Namely, there exist continuous maps $`g^j:v_j\overline{S}`$ which are homeomorphisms outside the geodesics $`\gamma _i^j`$, which pinch these geodesics to the corresponding singular points of $`\overline{S}`$(that are the images of $`\gamma _i`$) and which send $`F_j`$ to a fixed subset $`F`$ in the nonsingular locus of $`\overline{S}`$. Now, the convergence $`J_j\overline{J}`$ is understood as the uniform $`C^{\mathrm{}}`$convergence $`g_{}^j(J_j)\overline{J}`$ on the compact subsets in the non-singular locus $`\overline{S}^{}`$ of $`\overline{S}`$ which is identified with $`S_{i=1}^k\gamma _i`$.
4. $`C^0`$interior convergence. The limit cusp-curve $`\overline{v}:\overline{S}^{}\overline{V}`$, that is a holomorphic map which is constructed by first taking the maps
$$\overline{v}_j=(g_j)^1:S_{i=1}^k\gamma _i\overline{V}$$
Near the nodes of $`\overline{S}`$ including interior nodes and boundary nodes, by the properties of hyperbolic metric $`\mu ^{}`$ on $`\overline{S}`$, the neighbourhoods of interior nodes are corresponding to the annulis of the geodesic cycles. By the reparametrization of $`v_j`$, called $`\overline{v}_j`$ which is defined on $`S`$ and extends the maps $`\overline{v}_j:SS_jV`$(see). Now let $`\{z_i|i=1,\mathrm{},n\}`$ be the interior nodes of $`\overline{S}`$. Then the arguments in yield the $`C^0`$interir convergece near $`z_i`$.
5. $`C^0`$boundary convergence. Now it is possible that the boundary of the cusp curve $`\overline{v}`$ does not remain in $`\overline{W}`$. Write $`\overline{v}(z)=(h(z),(a(z),u(z)),f(z))`$, here $`h(z)=z`$ or $`h(z)z_i`$, $`i=1,\mathrm{},n`$, $`z_i`$ is cusp-point or bubble point. We can assume that $`\overline{p}=(1,p)\overline{v}_n`$ is a puncture boundary point. Let $`\overline{v}_1`$ be the component of $`\overline{v}`$ which through the point $`\overline{p}`$. Let $`D=\{z|z=re^{i\theta },0r1,0\theta 2\pi \}`$. We assume that $`\overline{v}_1:D\{e^{i\theta _i}\}_{i=1}^kV_c`$, here $`e^{i\theta _i}`$ is node or puncture point. Near $`e^{i\theta _i}`$, we take a small disk $`D_i`$ in $`D`$ containing only one puncture or node point $`e^{i\theta _i}`$. By the reparametrization and the convergence procedure, we can assume that $`\overline{v}_{1i}=(\overline{v}_1|D_i)`$ as a map from $`D^+\{0\}V_c`$ with $`\overline{v}_1([1,1]\{0\})S^1\times F(W^{}\times S^1)`$ and $`area(\overline{v}_{1i})a_0`$, $`a_0`$ small enough. Since $`Area(\overline{v}_{1i})a_0`$, there exist curves $`c_k`$ near $`0`$ such that $`l(\overline{v}_{1i}(c_k))\delta _1`$. By the construction of convergence, we can assume that $`l(\overline{v}_n(c_k))2\delta _1`$. If $`\overline{v}_{1i}(c_k)(S^1)\times F(\times [N_0,N_0]\times S^1)`$, we have $`\overline{v}_n(c_k)(S^1)\times F(\times [2N_0,2N_0]\times S^1)`$ for $`n`$ large enough. Now $`\overline{v}_n(c_k)`$ cuts $`\overline{v}_n(D)`$ as two parts, one part corresponds to $`\overline{v}_{1i}`$, say $`\overline{u}_n(D)`$. Then $`area(\overline{u}_n(D))=area(h_{n1})+|\mathrm{\Psi }^{}(u_{n2}(c_k^1))\mathrm{\Psi }^{}(u_{n2}(c_k^2))|`$, here $`c_k=\{c_k^1,c_k^2\}`$. Then by the proof of Lemma6.1-6.3, we know that $`\overline{u}_n(Dc_k)(S^1)\times F(\times [100N_0,100N_0]\times S^1)`$. So, $`\overline{v}_{1i}([1,1]\{0\})S^1\times F(\times [100N_0,100N_0]\times S^1)`$. By proposition 6.6, one singularity of $`\overline{v}_1`$ is deleted. We repeat this procedure, we proved that $`\overline{v}_1`$ is extended to whole $`D`$. So, the boundary node or puncture points of $`\overline{v}`$ are removed. Then by choosing the sub-sub-sequences of $`\mu _j^{}`$ and $`\overline{v}_j`$, we know that $`\overline{v}_j`$ converges to $`\overline{v}`$ in $`C^0`$ near the boundary node or puncture point. This proved the $`C^0`$boundary convergence. Since $`\overline{v}_j(1)=\overline{p}`$, $`\overline{p}\overline{v}(D)`$, $`\overline{v}(D)\overline{W}`$.
6. Convergence of area. Finally by the $`C^0`$convergence and $`area(v_j)=_Dv_j^{}\overline{\omega }`$, one easily deduces
$$area(v(S))=\underset{j\mathrm{}}{lim}(v_j(S_j)).$$
### 6.5 Bounded image of $`J`$holomorphic curves in $`W`$
###### Proposition 6.7
Let $`v`$ be the solutions of equations (4.16), then
$$d_W(p,v(D^2))=max\{d_W(p,q)|qf(D^2)\}d_0<+\mathrm{}$$
Proof. It follows directly from Gromov’s $`C^0`$convergence theorem.
## 7 Proof of Theorem 1.1
###### Proposition 7.1
If $`J`$holomorphic curves $`C\overline{V}`$ with boundary
$$CD^2\times ([0,\epsilon ]\times \mathrm{\Sigma })\times R^2$$
and
$$C(D^2\times (\{3\}\times \mathrm{\Sigma })\times R^2)\mathrm{}$$
Then
$$area(C)2l_0.$$
Proof. It is obvious by monotone inequality argument for minimal surfaces.
###### Note 7.1
we first observe that any $`J`$holomorphic curves with boundary in $`R^+\times \mathrm{\Sigma }`$ meet the hypersurface $`\{3\}\times \mathrm{\Sigma }`$ has energy at least $`2l_0`$, so we take $`\epsilon `$ small enough such that the Gromov’s figure eight contained in $`B_{R(\epsilon )}C`$ for $`\epsilon `$ small enough and the energy of solutions in section 4 is smaller than $`l_0`$. we specify the constant $`a_0`$, $`\epsilon `$ in section 3.1-3 such that the above conditions satisfied.
###### Theorem 7.1
There exists a non-constant $`J`$holomorphic map $`u:(D,D)(V^{}\times C,W)`$ with $`E(u)4\pi R(\epsilon )^2`$ for $`\epsilon `$ small enough such that $`4\pi R(\epsilon )^2l_0`$.
Proof. By Proposition 5.1, we know that the image $`\overline{v}(D)`$ of solutions of equations (4.11) remains a bounded or compact part of the non-compact Lagrangian submanifold $`W`$. Then, all arguments in for the case $`W`$ is closed in $`S\mathrm{\Sigma }\times R^2`$ can be extended to our case, especially Gromov’s $`C^0`$converngence theorem applies. But the results in section 4 shows the solutions of equations (4.11) must denegerate to a cusp curves, i.e., we obtain a Sacks-Uhlenbeck-Gromov’s bubble, i.e., $`J`$holomorphic sphere or disk with boundary in $`W`$, the exactness of $`\omega `$ rules out the possibility of $`J`$holomorphic sphere. For the more detail, see the proof of Theorem 2.3.B in .
Proof of Theorem 1.1. If $`(\mathrm{\Sigma },\lambda )`$ has no Reeb chord, then we can construct a Lagrangian submanifold $`W`$ in $`V=V^{}\times C`$, see section 3. Then as in , we construct an anti-holomorphic section $`c`$ and for large vector $`cC`$ we know that the nonlinear Fredholm section or Cauchy-Riemann section has no solution, this implies that the section is non-proper, see section 4. The non-properness of the section and the Gromov’s compactness theorem in section 6 implies the existences of the cusp-curves. So, we must have the $`J`$holomorphic sphere or $`J`$holomorphic disk with bounadry in $`W`$. Since the symplectic manifold $`V`$ is an exact symplectic mainifold and $`W`$ is an exact Lagrangian submanifold in $`V`$, by Stokes formula, we know that the possibility of $`J`$holomorphic sphere or disk elimitated. So our priori assumption does not holds which implies the contact maifold $`(\mathrm{\Sigma },\lambda )`$ has at least Reeb chord. This finishes the proof of Theorem 1.1.
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# 𝑁^∗ RESONANCES IN 𝑒⁺𝑒⁻ COLLISIONS AT BEPC
## 1 Introduction
To understand the internal quark-gluon structure of nucleon and its excited states $`N^{}`$’s is one of the most important tasks in nowadays particle and nuclear physics. The main source of information for the baryon internal structure is their mass spectrum, various production and decay rates. Our present knowledge of this aspect came almost entirely from the old generation of $`\pi N`$ experiments of more than twenty years ago. Considering its importance for the understanding of the nonperturbative QCD, a new generation of experiments on $`N^{}`$ physics with electromagnetic probes (real photon and space-like virtual photon) has recently been started at new facilities such as CEBAF at JLAB, ELSA at Bonn, GRAAL at Grenoble.
The $`J/\mathrm{\Psi }`$ experiment at the Beijing Electron-Positron Collider (BEPC) has long been known as the best place for looking for glueballs. But in fact it is also an excellent place for studying $`N^{}`$ resonances. The corresponding Feynman graphs for the $`N^{}`$ and $`\overline{N}^{}`$ production are shown in Fig. 1.
These graphs are almost identical to those describing the $`N^{}`$ electro-production process if the direction of the time axis is rotated by $`90^o`$. The only difference is that the virtual photon here is time-like instead of space-like and couples to $`NN^{}`$ through a real vector charmonium meson $`J/\mathrm{\Psi }`$. This fact leads to a few advantages for studying $`N^{}`$ resonances at the BEPC as the following.
(1) $`J/\mathrm{\Psi }N\overline{N}\pi `$ and $`N\overline{N}\pi \pi `$ provide a natural isospin $`I=1/2`$ filter for the $`\pi N`$ and $`\pi \pi N`$ systems due to isospin conservation.
Although the existence of the $`N^{}(1440)`$ and $`N^{}(1535)`$ is well-established, their properties, such as mass, width and decay branching ratios etc., still suffer large experimental uncertainties$`^\mathrm{?}`$. A big problem in extracting information on these $`N^{}`$ resonances from $`\pi N`$ and $`\gamma N`$ experiments is the isospin decomposition of 1/2 and 3/2 for $`\pi N`$ and $`\pi \pi N`$ systems$`^\mathrm{?}`$. We expect that the results from $`J/\mathrm{\Psi }`$ decays will provide better determination of the properties of these $`N^{}`$ resonances.
(2) Interference between $`N^{}`$ and $`\overline{N}^{}`$ bands in $`J/\mathrm{\Psi }N\overline{N}\pi `$ Dalitz plots may help to distinguish some ambiguities in the partial wave analysis of $`\pi N`$ two-body channel alone from $`\pi N`$ and $`\gamma N`$ experiments.
(3) The annihilation cross section of $`e^+e^{}`$ through $`J/\mathrm{\Psi }`$ is about two order of magnitude larger than that without going through $`J/\mathrm{\Psi }`$ and the branching ratios for our interested channels from $`J/\mathrm{\Psi }`$ decay are quite large$`^\mathrm{?}`$, cf. Table 1.
With present available 7.8 million $`J/\mathrm{\Psi }`$ events at BES-I and forthcoming 50 millions more at BES-II in near future, we expect to have enough statistics to get the best determination of properties for $`N^{}\pi N`$ and $`\pi \pi N`$ from $`p\overline{n}\pi ^{}`$, $`p\overline{p}\pi ^0`$ and $`p\overline{p}\pi ^+\pi ^{}`$ channels; We can also search for the “missing” $`N^{}`$ states and study known $`N^{}`$ states decaying into $`\eta N`$, $`\eta ^{}N`$, $`\omega N`$ and $`K\mathrm{\Lambda }`$ etc.
(4) On theoretical side, the process $`J/\mathrm{\Psi }\overline{p}N^{}`$ or $`p\overline{N}^{}`$ provides a new way to probe the internal quark-gluon structure of the $`N^{}`$ resonances$`^\mathrm{?}`$.
In the simple three-quark picture of baryons, the process can be described by Fig. 2(a). In this picture, three quark-antiquark pairs are created independently via a symmetric three-gluon intermediate state; the quarks and antiquarks have momenta of very similar magnitude. This is quite different from the $`\gamma pN^{}`$ process, cf. Fig. 2(b), where the photon couples to only one quark and unsymmetric configuration of quarks is favored with $`q_1^{}=q_1+Q`$. Therefore the processes $`J/\mathrm{\Psi }\overline{p}N^{}`$ and $`\gamma pN^{}`$ should probe different aspects of the quark distributions inside baryons. This may help us to distinguish various quark models$`^\mathrm{?}`$
For a hybrid $`N^{}`$, since the $`J/\mathrm{\Psi }`$ decay is a glue-rich process, it can be produced via diagram Fig. 2(c) and is expected to have larger production rate than a pure three-quark $`N^{}`$ resonance$`^\mathrm{?}`$.
Considering these advantages, a $`N^{}`$ program at BEPC has been proposed$`^\mathrm{?}`$ and started$`^\mathrm{?}`$.
## 2 Status of $`N^{}`$ data at BES
Based on 7.8 million $`J/\mathrm{\Psi }`$ events collected at BEPC before 1996, the events for $`J/\mathrm{\Psi }\overline{p}p\pi ^0`$ and $`\overline{p}p\eta `$ have been selected and reconstructed$`^\mathrm{?}`$.
The $`\pi ^0`$ and $`\eta `$ are detected in their $`\gamma \gamma `$ decay mode. For selected $`J/\mathrm{\Psi }\overline{p}p\gamma \gamma `$ events, the invariant mass spectrum of the $`2\gamma `$ is shown in Fig.4. The $`\pi ^0`$ and $`\eta `$ signals are clearly there.
Fig.4 and Fig.6 show Dalitz plot and $`p\pi ^0`$ invariant mass spectrum for the $`J/\mathrm{\Psi }\overline{p}p\pi ^0`$. There are clear peaks around 1480 and 1650 MeV of $`p\pi ^0`$ invariant mass.
Fig.6 shows $`p\eta `$ invariant mass spectrum for the $`J/\mathrm{\Psi }\overline{p}p\eta `$. There are clear enhancement around the $`p\eta `$ threshold, peaks at 1540 and 1650 MeV.
We have also selected and reconstructed events for $`J/\mathrm{\Psi }\overline{p}n\pi ^+`$ and $`p\overline{n}\pi ^{}`$ channels. There are similar $`p\pi `$ structures as in $`J/\mathrm{\Psi }\overline{p}p\pi ^0`$ process. The data processing for other channels, such as $`\overline{p}\mathrm{\Lambda }K`$, $`\overline{\mathrm{\Lambda }}pK`$, $`\overline{p}p\pi ^+\pi ^{}`$ and $`\overline{p}p\omega `$, are in progress.
## 3 Partial wave analyses of $`J/\mathrm{\Psi }\overline{p}p\eta `$
Because the $`J/\mathrm{\Psi }\overline{p}p\eta `$ has the simplest possible resonance contributions, a partial wave analysis is performed for this channel first. We are mainly interested in the structures at 1540 and 1650 MeV of the $`p\eta `$ invariant mass. Only $`J^P=\frac{1}{2}^\pm `$ $`N^{}`$ resonances are included in the analysis, since according to the information$`^{\mathrm{?},\mathrm{?}}`$ from $`\pi N\eta N`$ and $`\gamma N\eta N`$ experiments resonances with higher spins have much smaller couplings to $`p\eta `$ in our interested mass range.
We use the effective Lagrangian approach$`^{\mathrm{?},\mathrm{?}}`$ for the partial wave analysis. The relavent spin-1/2 interaction Lagrangians are
$`_{\eta PR}`$ $`=`$ $`ig_{\eta PR}\overline{P}\mathrm{\Gamma }R\eta +H.c.,`$ (1)
$`_{\mathrm{\Psi }PR}^{(1)}`$ $`=`$ $`{\displaystyle \frac{ig_{T_R}}{M_R+M_P}}\overline{R}\mathrm{\Gamma }_{\mu \nu }q^\nu P\mathrm{\Psi }^\mu +H.c.,`$ (2)
$`_{\mathrm{\Psi }PR}^{(2)}`$ $`=`$ $`g_{V_R}\overline{R}\mathrm{\Gamma }_\mu P\mathrm{\Psi }^\mu +H.c.`$ (3)
where $`R`$ is the generic notation for the resonance with mass $`M_R`$, $`P`$ for proton with mass $`M_P`$ and $`\mathrm{\Psi }`$ for $`J/\mathrm{\Psi }`$ with four-momentum $`q`$. The operator structures for the $`\mathrm{\Gamma }`$, $`\mathrm{\Gamma }_\mu `$ and $`\mathrm{\Gamma }_{\mu \nu }`$ are
$$\mathrm{\Gamma }=1,\mathrm{\Gamma }_\mu =\gamma _\mu ,\mathrm{\Gamma }_{\mu \nu }=\gamma _5\sigma _{\mu \nu },$$
(4)
$$\mathrm{\Gamma }=\gamma _5,\mathrm{\Gamma }_\mu =\gamma _5\gamma _\mu ,\mathrm{\Gamma }_{\mu \nu }=\sigma _{\mu \nu },$$
(5)
where (4) and (5) correspond to nucleon resonances of odd and even parities, respectively. The relative magnitudes and phases of the amplitudes are determined by a maximum likelihood fit to the data. A fit with three $`N^{}`$ resonances is shown in Figs.6-8 for the $`p\eta `$, $`\overline{p}p`$ invariant mass spectra and the angular distribution of the proton relative to the beam direction, respectively.
The peak near the $`p\eta `$ threshold is fitted with a $`N^{}`$ resonance of mass and width optimised at $`M=1540_{17}^{+15}`$ and $`\mathrm{\Gamma }=178_{22}^{+20}`$ MeV, respectively. The data favor $`J^P=\frac{1}{2}^{}`$ over $`\frac{1}{2}^+`$. A fit with $`J^P=\frac{1}{2}^+`$ instead gives likelihood value $`lnL`$ worse by 16.0 than for $`\frac{1}{2}^{}`$ assignment (Our definition of $`lnL`$ is such that it increases by 0.5 for a one standard deviation change in one parameter). The statistical significance of the peak is above 6.0$`\sigma `$. It is obviously the $`S11`$ $`N^{}(1535)`$ resonance. It makes the largest contribution $`(84\pm 5)\%`$ to the $`p\overline{p}\eta `$ final states. The errors here and later include both statistics and systematic errors from the fit.
The second peak around 1650 MeV is also fitted with a $`J^P=\frac{1}{2}^{}`$ resonance $`N^{}(1650)`$. Its mass optimises at $`M=1648_{16}^{+18}`$ MeV with $`\mathrm{\Gamma }=150`$ MeV fixed to PDG value. Its width cannot be well determined by our data due to a correlation with the parameters used for the third resonance. It contributes $`(11\pm 3)\%`$ to the $`p\overline{p}\eta `$ final states.
A small improvement to the fit is given by including an addtional $`J^P=\frac{1}{2}^+`$ resonance, which optimises at $`M=1834_{55}^{+46}`$ MeV with $`\mathrm{\Gamma }=200`$ MeV fixed. The statistical significance of the peak is only $`2.0\sigma `$. We have tried $`J^P=\frac{1}{2}^{}`$ instead $`\frac{1}{2}^+`$, the fit is worse.
An interesting result is that the $`_{\mathrm{\Psi }PR}^{(2)}`$ term given by Eq.(3) makes insignificant contribution for both $`N^{}(1535)`$ and $`N^{}(1650)`$. If we drop this kind of couplings for both resonances, the likelihood value for the fit is only worse by 0.8 for 4 less free parameters. This kind of couplings should vanish for the real photon coupling to $`NN^{}`$. Why it also vanishes for the $`\mathrm{\Psi }NN^{}`$ coupling needs to be understood. A theoretical calculation$`^\mathrm{?}`$ assuming pure $`_{\mathrm{\Psi }PR}^{(2)}`$ coupling without $`_{\mathrm{\Psi }PR}^{(1)}`$ coupling failed to reproduce the basic feature of the $`J/\mathrm{\Psi }\overline{p}p\eta `$ data. This is consistent with our observation that the $`_{\mathrm{\Psi }PR}^{(1)}`$ coupling dominates for both $`N^{}(1535)`$ and $`N^{}(1650)`$.
In the Vector Meson Dominance (VMD) picture, the virtual photon couples to the $`NN^{}`$ through vetor mesons, and the electro-magnetic $`NN^{}`$ transition form factors $`g_{\gamma ^{}NN^{}}`$ can be expressed in terms of photon-meson coupling strengths $`C_{\gamma V}`$ and meson-$`NN^{}`$ vertex form factors $`g_{_{VNN^{}}}`$:
$$g_{\gamma ^{}NN^{}}(q^2)=\underset{j}{}\frac{m_j^2C_{\gamma V_j}}{m_j^2q^2im_j\mathrm{\Gamma }_j}g_{_{V_jNN^{}}}(q^2)$$
(6)
with
$$C_{\gamma V}=\sqrt{\frac{3\mathrm{\Gamma }_{Ve^+e^{}}}{\alpha m_V}}.$$
(7)
At $`q^2=M_\mathrm{\Psi }^2`$, the $`J/\mathrm{\Psi }`$ meson dominates. The terms from other vector mesons are negligible. From our PWA results here and other relavent information from PDG$`^\mathrm{?}`$, we can deduce the transition form factor for the time-like virtual photon to $`PN^{}(1535)`$ as
$$|g_{\gamma ^{}pN^{}}(q^2=M_\mathrm{\Psi }^2)|=2.8\pm 0.5,$$
(8)
which is related to the more familiar helicity amplitude $`A_{1/2}^P`$ for $`N^{}\gamma P`$ by
$$|A_{1/2}^P|^2=\left(\frac{g_{\gamma pN^{}}(q^2=0)}{M_N^{}+M_P}\right)^2\frac{(M_N^{}^2M_P^2)}{2M_P}.$$
(9)
## 4 Summary and outlook
In summary, the $`J/\mathrm{\Psi }`$ decay at BEPC provides a new excellent laboratory for studying the $`N^{}`$ resonances. On experimental side, it provides a natural isospin 1/2 filter for $`\pi N`$ and $`\pi \pi N`$ systems and many interesting channels for studying $`N^{}`$ and hyperon resonances; on theoretical side, it provides a new way to explore the internal structure of baryons and may help us to pin down hybrid baryon(s). Almost all subjects on the $`N^{}`$ resonances at the CEBAF$`^\mathrm{?}`$ and other $`\gamma p(ep)`$ facilities can be studied here complementally with the virtual time-like photon.
Based on 7.8 million $`J/\psi `$ events collected at BEPC, a partial wave analysis of $`J/\psi p\overline{p}\eta `$ data has been performed. Two $`J^P=\frac{1}{2}^{}`$ resonances, $`N^{}(1535,S_{11})`$ and $`N^{}(1650,S_{11})`$, are observed. Now we are collecting more $`J/\mathrm{\Psi }`$ events with the improved BES detector. With the forthcoming 50 millions more $`J/\mathrm{\Psi }`$ events in near future, more precise partial wave analyses can be carried out on many channels involving $`N^{}`$ resonances and should offer some best determinations of $`N^{}`$ properties. A systematic theoretical and experimental study of the $`N^{}`$ and hyperon production from the $`J/\psi `$ decays is underway.
There is also a plan to upgrade the BEPC to BEPC2 which will increase the luminosity by an order of magnitude. This will provide us more precise data for the study of the $`N^{}`$ resonances from $`J/\mathrm{\Psi }`$ decays and also enable us to extend the $`N^{}`$ program to a higher energy at the $`\mathrm{\Psi }^{}`$ resonance which now suffers low statistics for the $`N^{}`$ study.
With the new generation of $`\gamma p(ep)`$, $`J/\psi `$ and $`\mathrm{\Psi }^{}`$ experiments, a new era for the baryon spectroscopy is coming.
## References
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# Introduction
## Introduction
Feynman diagrams computations are often a lengthy and hard task. The complexity of typical computations grows as a factorial with the number of external legs and/or loops involved in the processes. In the Standard Model the main difficulties arise from QCD corrections and the production of jets in the final states; but it is not difficult to imagine extensions of the Standard Model, in which this problem could seriously limit our capabilities of studying and understanding new physical phenomena; in particular if they are involved (in addition to $`\alpha _s`$ ) by new rather strong couplings, which give rise to multiparticle/multiloop amplitudes. This issue motivates a rich and increasing activity. Considerable progresses have been achieved in the past : new methods inspired by string theory \[3-5\], helicity amplitudes \[6-9\], different recursive algorithms to calculate QCD dual amplitudes, or the scattering matrix elements of generic processes with arbitrary initial/final states. Very useful approximations have been proposed \[12-16\], when the exact matrix elements are unknown. In loop computations interesting simplifications \[17-26\] can be used.
While analytical methods can provide us in some specific processes with very simple results, a general approach for generic lagrangian and final states cannot avoid the use of the computer power. A possible strategy is to implement the Feynman rules into some automatic code with the help of some familiar packages for symbolic manipulation of algebraic formulae. However this way to proceed becomes ineffective, when the computer has to manage very complex formulae (as often it is the case).
This suggests a second strategy which tries exploiting the computer power with pure numerical algorithms and thus avoiding lengthy symbolic manipulation. In the recent past these algorithms have been successfully used for the computation of several tree level amplitudes both in QCD and in electroweak processes \[29-31\]. In this paper we discuss the possibility to apply pure numerical techniques in loop calculations.
The problem of performing loop computations can be divided in two parts. The first part is similar to tree level computations: from a given lagrangian one has to produce an explicit function of virtual and real momenta which is equivalent to the sum of all the Feynman diagrams contributing to the process. Once this task is accomplished, loops calculations have an additional problem with respect to tree level ones: integrals over virtual loop momenta are affected by ultraviolet and/or infrared divergences. These divergences require the introduction of a regularization procedure which usually obliges us to perform all virtual loop integrations in an analytical and essentially symbolic fashion (while in tree level calculations one usually performs the integration over the real momenta by numerical montecarlo methods). This is a serious difficulty in any numerical approach. Here we address this last issue. We will define a numerical procedure to perform virtual loop integrations even when these include divergences.
Let us consider a function $`f[l_\mu ,p_\mu ]`$ of the external momenta $`p_\mu `$ and of the virtual momenta $`l_\mu `$. The $`n`$ dimensional integration is performed in a region of complex values of $`n`$ where the integral is convergent; then one considers the analytical continuation to values of $`n`$ where the integral is not well defined. In general for $`n4`$ the analytical continuation can have a pole $`1/\epsilon `$ ($`\epsilon =n4`$). We can use the Laurent expansion around $`\epsilon 0`$ to define some functionals $`I_k[f]`$ such that<sup>1</sup><sup>1</sup>1At one loop.
$$\mathrm{d}^nlf[l_\mu ]=I_1[f]\epsilon ^1+I_0[f]+I_1[f]\epsilon ^1+I_2[f]\epsilon ^2+\mathrm{}$$
(1)
Since the $`n`$ dimensional integration is well defined for any function $`f`$, also the functionals $`I_k[f]`$ are unambiguously defined from equation (1). For instance, $`I_0[f]`$ coincides with the ordinary four dimensional integral
$$\mathrm{d}^4\overline{l}f[l_\mu ]$$
(2)
if $`f`$ is a convergent function. While the left-hand side of the (1) is essentially a formal definition which does not lead us to an obvious numerical integration procedure, each $`I_k[f]`$ can be written in terms of well defined and convergent integrals. In the following we show how to build explicit integral representations for the functionals $`I_k[f]`$. In the next sections, we explain two different definitions (even if equivalent) of the $`I_k[f]`$; in the final part some very simple numerical examples will be given in order to emphasize the practical purposes of the method.
## Method of Integration A
Let us assume that we have to perform an integral in $`n>4`$ dimensions. The virtual momentum $`l_\mu `$ has $`n`$ components. We call $`\overline{l}_\mu `$ the components with $`\mu 4`$ and $`\stackrel{~}{l}_\mu `$ the components with $`\mu >4`$. Then we can rewrite the integral
$$\mathrm{d}^nlf[l_\mu ,p_\mu ]=d^4\overline{l}f[l_\mu ,p_\mu ]\mathrm{d}^\epsilon \stackrel{~}{l}=d^4\overline{l}f[l_\mu ,p_\mu ]\stackrel{~}{l}^{\epsilon 1}d\stackrel{~}{l}d\mathrm{\Omega }_\epsilon $$
(3)
$`\mathrm{d}\mathrm{\Omega }_\epsilon `$ is the solid angle in the subspace ($`\mu >4`$) of dimension $`\epsilon `$ orthogonal to the four-vector $`\overline{l}_\mu `$. Note that $`f[l_\mu ,p_\mu ]`$ is invariant under rotations in this subspace, since the external momenta do not have components $`\mu >4`$. We can exploit this invariance as follows. First we rotate the vector $`\stackrel{~}{l}_\mu `$ such that $`\stackrel{~}{l}_50`$ and $`\stackrel{~}{l}_n=0`$ for any $`n>5`$. Thus we can write $`l_\mu =(\overline{l}_1,\overline{l}_2,\overline{l_3},\overline{l_4},\sqrt{t})`$ with $`t=(l_5)^2+(l_6)^2+\mathrm{}+(l_n)^2`$ and omit all the components with $`\mu >5`$. We have reduced our $`n`$ dimensional space into a five dimensional space. Second, we can perform the integral in $`\mathrm{d}\mathrm{\Omega }_\epsilon `$ and we obtain from eq.(3)
$$\frac{\pi ^{\epsilon /2}}{\mathrm{\Gamma }\left[\frac{\epsilon }{2}\right]}_0^{\mathrm{}}\left(\mathrm{d}^4\overline{l}f[l_\mu ,p_\mu ]\right)t^{\epsilon /21}dt.$$
(4)
Now the integral in $`\mathrm{d}t`$ can be integrated by parts
$$\frac{\pi ^{\epsilon /2}}{\mathrm{\Gamma }\left[\frac{\epsilon }{2}\right]}_0^{\mathrm{}}\frac{4t^{1+\frac{\epsilon }{2}}}{\epsilon (2+\epsilon )}\left(\mathrm{d}^4\overline{l}\frac{\mathrm{d}^2}{\mathrm{d}t^2}f[l_\mu ,p_\mu ]\right)dt.$$
(5)
The function $`f`$ depends on $`t`$ through the fifth component of $`l_\mu `$. The function $`t^{1+\frac{\epsilon }{2}}`$ has a cut in the positive real axis of $`t`$ ( $`Re(t)>0`$ and $`Im(t)=0`$ ) , in fact the rotation $`tte^{i2\pi }`$ gives $`t^{1+\frac{\epsilon }{2}}t^{1+\frac{\epsilon }{2}+i\pi \epsilon }`$. Therefore we can apply the following identity
$$_0^{\mathrm{}}t^{1+\frac{\epsilon }{2}}f[t]dt=\frac{1}{\left(1e^{i\pi \epsilon }\right)}t^{1+\frac{\epsilon }{2}}f[t]dt$$
(6)
where the contour of the last integration in $`t`$ must contain all the poles and singularities lying in the complex $`t`$ plane and must avoid the cut on the positive real axis (see figure 1). Then the integral (5) above becomes
$$\frac{1}{(1e^{i\pi \epsilon })}\frac{\pi ^{\epsilon /2}}{\mathrm{\Gamma }\left[\frac{\epsilon }{2}\right]}\frac{4t^{1+\frac{\epsilon }{2}}}{\epsilon (2+\epsilon )}\left(\mathrm{d}^4\overline{l}\frac{\mathrm{d}^2}{\mathrm{d}t^2}f[l_\mu ,p_\mu ]\right)dt.$$
(7)
Now we can make the expansion in powers of $`\epsilon `$ (the Euler constant and $`\mathrm{log}(\pi )`$ have been re-absorbed into a redefinition of the scale $`\mu `$ as in the $`\overline{MS}`$) and we obtain the final result
$$\frac{it}{2\pi }\left(\frac{2}{\epsilon }+(1+\mathrm{log}[t])\right)\left(\mathrm{d}^4\overline{l}\frac{\mathrm{d}^2}{\mathrm{d}t^2}f[l_\mu ,p_\mu ]\right)dt.$$
(8)
In the above expression each integral is now convergent, due to the action of the derivative $`\mathrm{d}^2/\mathrm{d}t^2`$ appearing inside the integral in $`\mathrm{d}^4\overline{l}`$. Thus one can separately calculate the singular and the non singular part, simply evaluating well defined and convergent integrals. We can check the result above in a specific example and see how the above formula can be used for practical calculations.
Suppose that we have to evaluate the following divergent integral
$$\mathrm{d}^nl\frac{l^2}{l^2+m^2}.$$
(9)
In our approach $`l_\mu `$ is not a $`n`$ dimensional object but it is a more concrete five dimensional vector with components $`(\overline{l}_1,\overline{l}_2,\overline{l_3},\overline{l}_4,\sqrt{t})`$; thus $`l^2=l_\mu l^\mu =\overline{l}^2+t`$ with $`\overline{l}^2`$ equal to the squared length of the real four-vector $`(\overline{l}_1,\overline{l}_2,\overline{l}_3,\overline{l_4})`$ and $`t`$ a complex number. If we are interested in the finite part, we can apply the non singular term in equation (8) and we have
$$\frac{it}{2\pi }(1+\mathrm{log}[t])\left(\frac{\mathrm{d}^2}{\mathrm{d}t^2}\frac{\overline{l}^2+t}{\overline{l}^2+m^2+t}\right)\mathrm{d}^4\overline{l}\mathrm{d}t.$$
(10)
After the derivative $`\frac{\mathrm{d}^2}{\mathrm{d}t^2}`$ the integral in $`\mathrm{d}^4\overline{l}`$ is convergent<sup>2</sup><sup>2</sup>2We consider only value of $`t`$ with non zero imaginary part. See the discussion at the end of this section., and this yields
$$\frac{it}{2\pi }(1+\mathrm{log}[t])\left(\frac{\pi ^2m^2}{t+m^2}\right)dt$$
(11)
The integral over the contour in the complex plane is equal to the residue at the pole $`tm^2`$
$$\pi ^2\left(1+\mathrm{log}[m^2]\right)m^4.$$
(12)
Our final result is correct. The advantage of the procedure above is clear: eqs.(10-11) are well defined and concrete expressions, any step can be done in a pure numerical way. Note that the integral in $`\mathrm{d}^4\overline{l}`$ must be done after the derivative $`\mathrm{d}^2/\mathrm{d}t^2`$, and before the integral in $`\mathrm{d}t`$; otherwise the integral in $`\mathrm{d}^4\overline{l}`$ would be non convergent. In practice, if the above procedure is done numerically, one should replace the integral in $`\mathrm{d}t`$ with a sum over a finite set of points $`t_i`$; for the convergence of all integrals, it is enough to choose the $`t_i`$ in this set with a non zero imaginary part.
## Method of Integration B
The formula (8) is a transparent and compact expression, which makes manifest some properties of the dimensional regularization; for instance, the invariance under the translations $`\overline{l}_\mu \overline{l}_\mu +p_\mu `$ is obvious, since it comes directly from the translational invariance of the integral in $`\mathrm{d}^4\overline{l}`$, which is convergent (after the derivative in $`\mathrm{d}^2/\mathrm{d}t^2`$). However in certain numerical calculations the (8) could be not efficient. We discuss a different method which is more powerful in practical calculations, even if it requires a slightly more involved formula.
There are several different ways of rewriting the (8); each of them depending on the way we choose to parameterize the virtual momentum $`l_\mu `$. Here we do not intend to make an exhaustive study, we will simply describe a quite general procedure to obtain well defined integral representations for the $`I_k`$ in (1).
First we observe that each $`I_k[f]`$ is a linear operator
$$I_k[f_1+f_2]=I_k[f_1]+I_k[f_2]$$
(13)
$$I_k[\lambda f]=\lambda I_k[f].$$
(14)
It is quite natural to think that the linearity (13-14) implies that the $`I_k`$ are authentic integrals or sum of integrals. In fact we can introduce a simple trick to build concrete integral representations for generic linear functionals.
Consider the space of functions which admits a Laurent expansion around a point $`t_0`$. They are defined by a set of $`a_n`$ through the expansion<sup>3</sup><sup>3</sup>3We also assume that the series converges strongly enough to justify the steps below.
$$f[t]=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_n(tt_0)^n.$$
(15)
A linear functional $`I[f]`$ acting on this set of functions can be written
$$I\left[f\right]=I\left[\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_n(tt_0)^n\right]=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_nI[(tt_0)^n]=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_nb_n$$
(16)
where $`b_n=`$ $`I[(tt_0)^n]`$ is a real (or complex) number. Then we can use the identity
$$_C\frac{1}{2\pi i}\frac{(tt_0)^n}{(tt_0)^{k+1}}dt=\delta _{nk}$$
(17)
for any $`n,k=0,\pm 1,\pm 2,\mathrm{}`$ if the integral contour C is a small circle, in the complex plane of $`t`$, around the point $`t_0`$. Using the (17), the (16) can be written
$`I[f]`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \left(\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}a_n(tt_0)^n\right)\left(\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\frac{b_k}{(tt_0)^{k+1}}\right)dt}=`$
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle f[t]w[t]dt}`$
with
$$w[t]=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\left(\frac{b_k}{(tt_0)^{k+1}}\right).$$
(19)
This is an integral representation of the linear functional $`I`$. This simple trick can be used to build explicit and compact integral representations for any linear functional, whose action on each term of a Laurent expansion is known.
Let us apply it to our problem. We start defining the angular integration in $`n=4+\epsilon `$ dimensions. The analytical integration of any tensor with $`k`$ (even) indices and constructed with the $`n`$ dimensional vector $`l_\mu `$ yields<sup>4</sup><sup>4</sup>4Here we assume that the overall factor $`\pi ^{\frac{\epsilon }{2}}/\mathrm{\Gamma }[\frac{n}{2}]`$ has been re-absorbed into a redefinition of the scale $`\mu `$, as prescribed by the $`\overline{MS}`$ scheme.
$$l_{\overline{\mu }}l_{\overline{\nu }}l_{\overline{\alpha }}\mathrm{}l_{\overline{\beta }}l_{\overline{\rho }}l_{\overline{\sigma }}d\mathrm{\Omega }_{4+\epsilon }=2\pi ^2\frac{(2+\epsilon )!!}{(2+k+\epsilon )!!}\left(g_{\overline{\mu }\overline{\nu }}g_{\overline{\rho }\overline{\sigma }}\mathrm{}g_{\overline{\alpha }\overline{\beta }}+\mathrm{}+g_{\overline{\mu }\overline{\alpha }}g_{\overline{\rho }\overline{\nu }}\mathrm{}g_{\overline{\beta }\overline{\sigma }}\right)l^{2(k/2)}$$
(20)
$`\mathrm{d}\mathrm{\Omega }_{4+\epsilon }=\mathrm{sin}^{2+\epsilon }[\gamma ]\mathrm{}\mathrm{sin}^2[\theta ]\mathrm{sin}[\varphi ]\mathrm{d}\gamma \mathrm{d}\theta \mathrm{d}\varphi \mathrm{d}\delta (0<\gamma ,\theta ,\mathrm{},\varphi <\pi ;0<\delta <2\pi )`$ is the solid angle. It is understood that the indices ($`\overline{\mu }`$, $`\overline{\nu }\mathrm{}`$) above are contracted with some external momenta and/or other tensorial indices (like the spin of the external particles), since the scattering matrix is lorentz invariant. Thus<sup>5</sup><sup>5</sup>5We remind that only $`l_\mu `$ has components $`>4`$, and factors $`l_\mu l^\mu =l^2`$ are not relevant for the angular integration. only indices with $`\overline{\mu }4`$ are relevant in our problem. This has been emphasized by the small bar on the Greek indices.
We look for an integral definition which reproduces the (20). To begin the procedure we parameterize $`l_\mu `$ with an array of five components<sup>6</sup><sup>6</sup>6In practice this means that any vector (including external momenta etc.), must be written as an array of five components, but we keep in mind that only $`l_\mu `$ has a non zero fifth component.
$$l_\mu t(x\mathrm{cos}[\theta ],x\mathrm{sin}[\theta ]\mathrm{cos}[\varphi ],x\mathrm{sin}[\theta ]\mathrm{sin}[\varphi ]\mathrm{cos}[\delta ],x\mathrm{sin}[\theta ]\mathrm{sin}[\varphi ]\mathrm{sin}[\delta ],\sqrt{1x^2})$$
(21)
$`\theta `$ $`\varphi `$ and $`\delta `$ are the three angles in the space of four dimensions. $`t`$ and $`x`$ are two complex variables which will be integrated over a complex contour as explained below. The need of the auxiliary variable $`x`$ will be clear in a moment. Suppose that we integrate over the four dimensional solid angle $`\mathrm{d}\mathrm{\Omega }_4=\mathrm{sin}^2[\theta ]\mathrm{sin}[\varphi ]\mathrm{d}\theta \mathrm{d}\varphi \mathrm{d}\delta `$ all the tensors built with the array (21) in place of an $`n`$ dimensional $`l_\mu `$. We get
$$l_{\overline{\mu }}l_{\overline{\nu }}l_{\overline{\alpha }}\mathrm{}l_{\overline{\beta }}l_{\overline{\rho }}l_{\overline{\sigma }}d\mathrm{\Omega }_4=2\pi ^2\frac{(2)!!}{(2+k)!!}(g_{\overline{\mu }\overline{\nu }}g_{\overline{\rho }\overline{\sigma }}\mathrm{}g_{\overline{\alpha }\overline{\beta }}+\mathrm{}+g_{\overline{\mu }\overline{\alpha }}g_{\overline{\rho }\overline{\nu }}\mathrm{}g_{\overline{\beta }\overline{\sigma }})t^kx^k.$$
(22)
Comparing this result with the (20) we see that tensorial structure is identical, but the overall factor is not correct. The term of order $`\epsilon ^0`$ can be obtained simply setting $`x=1`$ and $`t^2=l^2`$. Instead for the term of order $`\epsilon `$, it would be enough to replace $`x^k`$ with a suitable factor<sup>7</sup><sup>7</sup>7 The function $`\psi `$ is the well known derivative $`\mathrm{d}/\mathrm{d}x\mathrm{log}[\mathrm{\Gamma }[x]]`$; it comes out when we expand the factorial in the (20) and take only the first order in $`\epsilon `$.
$$x^kb_k=\frac{\epsilon }{2}(1\gamma _E\psi [2+k/2])$$
(23)
for any even $`k`$.
This can be achieved by means of a linear functional $`I`$ such that $`I[x^k]=b_k`$. In fact, following the recipe above (eqs.(15)-(19)) we can build the integral below
$$I[f]=\frac{\epsilon }{2\pi i}\underset{k=even}{}\frac{b_k}{x^{k+1}}f[x]\mathrm{d}x=\frac{\epsilon }{2\pi i}\frac{x+x^3\mathrm{log}\left[1\frac{1}{x^2}\right]}{2(1+x^2)}f[x]dx$$
(24)
where the integral contour is a circle around the singularity $`x0`$. The integrand above has a cut in the segment of the real axis between -1 and 1. It is easy to see that in the limit of a path of integration very close to this cut, the logarithm can be approximated by its imaginary part. The real part does not contribute to the full integral, and the logarithm can be replaced by $`\pm i\pi `$. Taking also into account that $`f[x]=f[x]`$ we can rewrite the (24)
$$I[f]=\epsilon _0^1\left(\frac{x^3}{1x^2}\right)_+f[x]dx.$$
(25)
Here the notation $`(\mathrm{})_+`$ means that
$$_0^1(h[x])_+f[x]dx=_0^1h[x](f[x]f[1])dx,$$
(26)
and this makes the integral convergent near $`x1`$.
Finally we are able to write a compact formula for the angular integration in $`4+\epsilon `$ dimensions
$$f[l_\mu ]\mathrm{d}\mathrm{\Omega }_{4+\epsilon }=(f[l_\mu ]\mathrm{d}\mathrm{\Omega }_4)_{x=1}+\epsilon (_0^1f[l_\mu ]\left(\frac{x^3}{1x^2}\right)_+\mathrm{d}x\mathrm{d}\mathrm{\Omega }_4)+\epsilon ^2(\mathrm{}).$$
(27)
The last step is to define the integration in $`\mathrm{d}l^2=\mathrm{d}t^2`$. Again we follow the trick (15)-(19): we have to compute the $`b_\sigma `$ in order to get the function $`w[t^2]`$. By definition, we know that for any $`\sigma >2`$
$`b_\sigma `$ $`=`$ $`I\left[{\displaystyle \frac{1}{(t^2+m^2)^\sigma }}\right]={\displaystyle t^{2(1+\epsilon /2)}dt^2\frac{1}{(t^2+m^2)^\sigma }}=`$ (28)
$`=`$ $`{\displaystyle \frac{1}{(2+\sigma )(1+\sigma )}}{\displaystyle \frac{1}{(m^2)^{\sigma 2}}}+\epsilon (\mathrm{}),`$
for $`\sigma `$=1,2
$`b_1`$ $`=`$ $`2{\displaystyle \frac{m^2}{\epsilon }}+m^2\mathrm{log}[m^2]+\epsilon (\mathrm{})`$
$`b_2`$ $`=`$ $`={\displaystyle \frac{2}{\epsilon }}\left(\mathrm{log}[m^2]+1\right)+\epsilon (\mathrm{})`$ (29)
and for $`\sigma `$=0,-1,-2,… the integral is zero ($`b_\sigma =0`$). If the integration contour in the $`t^2`$ complex plane is a small circle around the pole singularity $`t^2=m^2`$ and replacing the $`b_n`$ in the (19) with the $`b_\sigma `$ above we get
$$w[t^2]=\underset{\sigma =0}{\overset{\mathrm{}}{}}b_{\sigma +1}\left(t^2+m^2\right)^\sigma =t^2\left(\frac{2}{\epsilon }+\mathrm{log}\left[t^2\right]\right)$$
(30)
and
$$t^{2(1+\epsilon /2)}dt^2f[t^2]=\frac{2}{\epsilon }\left(\frac{1}{2\pi i}f[t^2]t^2dt^2\right)\left(\frac{1}{2\pi i}f[t^2]t^2\mathrm{log}\left[t^2\right]dt^2\right)+\epsilon (\mathrm{}).$$
(31)
The integral contour is closed, it must contain all the singularities in the complex plane of the function $`f[t^2]`$, except for the singularity in $`t^2=0`$ and the cut for real and positive $`t^2`$ in all integrals where the logarithmic function appears (see for example figure 1). We will comment on this contour of integration later on.
Combing the angular integration (27) and the (31) we get the full formula
$$f[l_\mu ]\mathrm{d}^nl=f[l_\mu ]t^{2(1+\epsilon /2)}\frac{\mathrm{d}t^2}{2}d\mathrm{\Omega }_{4+\epsilon }=\frac{I_1}{\epsilon }+I_0+I_1\epsilon +\mathrm{}$$
(32)
where
$`I_1\left[f[l_\mu ]\right]`$ $`=`$ $`\left({\displaystyle \frac{1}{2\pi i}}{\displaystyle f[l_\mu ]t^2dt^2d\mathrm{\Omega }_4}\right)_{x=1}`$ (33)
$`I_0\left[f[l_\mu ]\right]`$ $`=`$ $`\left({\displaystyle \frac{1}{4\pi i}}{\displaystyle f[l_\mu ]t^2\mathrm{log}\left[t^2\right]dt^2d\mathrm{\Omega }_4}\right)_{x=1}+{\displaystyle \frac{1}{2\pi i}}{\displaystyle _0^1}{\displaystyle f[l_\mu ]t^2dt^2\left(\frac{x^3}{1x^2}\right)_+dxd\mathrm{\Omega }_4}.`$
Clearly the integrals above are well defined and can be evaluated numerically.<sup>8</sup><sup>8</sup>8 Here all integrals are understood to be in Euclidean space. In some kinematical regions, the integration (even if convergent) needs to be regularized by the prescription $`m^2m^2+i\epsilon `$. Also for some infrared singularities the (Method of Integration B) must be rearranged differently or one should use the method A. The practical use of these methods in various realistic situations certainly demands a much more extensive discussion, in this letter we simply present the general idea, a more complete and systematic study will be done elsewhere (see ref.). For example to evaluate the first integral of $`I_0`$ in the expression (Method of Integration B), the contour of integration in $`\mathrm{d}t^2`$ can be chosen as in figure 1. Namely we can divide the contour in three paths. The path A is very close to the cut of the logarithm. This function can be replaced with $`\pm i\pi `$ (the sign depends if we are above or below the cut). This yields the following integral
$$P_1=\left(\frac{1}{2}_{\mathrm{\Lambda }_{ir}^2}^{\mathrm{\Lambda }_{uv}^2}f[l_\mu ]t^2dt^2d\mathrm{\Omega }_4\right)_{x=1}.$$
(35)
Note that this is an ordinary four dimensional integral with an infrared and an ultraviolet cut-off: in fact after setting $`x=1`$, the array (21) may be regarded as a genuine four-vector. Then we have the paths B and C, that are two circles with radius $`\mathrm{\Lambda }_{uv}^2`$ and $`\mathrm{\Lambda }_{ir}^2`$ respectively
$$P_2=\left(\frac{1}{4\pi i}_Bf[l_\mu ]t^2\mathrm{log}\left[t^2\right]dt^2d\mathrm{\Omega }_4\right)_{x=1}+\left(\frac{1}{4\pi i}_Cf[l_\mu ]t^2\mathrm{log}\left[t^2\right]dt^2d\mathrm{\Omega }_4\right)_{x=1}.$$
(36)
These integrals can be seen as some counterterms which cancel the $`\mathrm{\Lambda }_{uv}`$ and $`\mathrm{\Lambda }_{ir}`$ dependence of the integral (35). The sum is equivalent to the finite part ( in the limit $`\epsilon `$$``$0) of the dimensionally regularized integral.
Through the linearity (13-14), this method can be generalized to two (or more) loop calculations. If $`q_\mu `$ and $`l_\mu `$ are the two loop virtual momenta, then we can rewrite the integration variables $`\mathrm{d}^nl=l^{2(1+\epsilon /2)}\mathrm{d}l^2\mathrm{d}\mathrm{\Omega }_{4+\epsilon }/2`$ and $`\mathrm{d}^nq=\mathrm{d}q_{_{}}q_{_{}}^{2(1+\epsilon /2)}\mathrm{d}q_{_{}}\mathrm{d}\mathrm{\Omega }_{3+\epsilon }`$. $`q_{_{}}`$ is the component of $`q_\mu `$ parallel to $`l_\mu `$, $`q_{_{}}`$ is the total length of the components of $`q_\mu `$ orthogonal to $`l_\mu `$. One then defines two six dimensional vectors $`l_\mu `$ and $`q_\mu `$, analogously to (21), as functions of five angles $`\theta ,\varphi ,\delta ,\varphi ^{},\delta ^{}`$, two auxiliary variables $`x`$ and $`y`$ (needed for the integrations $`\mathrm{d}\mathrm{\Omega }_{4+\epsilon }`$ and $`\mathrm{d}\mathrm{\Omega }_{3+\epsilon }`$) and the complex variables $`l^2,q_{_{}},q_{_{}}`$. Applying the trick (15)-(19) one obtains the function $`w[l^2,q_{_{}},q_{_{}}]`$. Needless to say, this contains dilogarithms in addition to some logarithms. From this guidelines, one gets the analogue of the (Method of Integration B), for two loop calculations .
### Some Numerical Examples
In order to make clear the practical purposes of these methods, we discuss here some very simple examples. Suppose that we want to integrate
$$\mathrm{d}^nl\frac{1}{(l+p)^2+m^2}$$
(37)
with $`m^2=2`$ and $`p^2=1`$. Our numerical approach must reproduce the analytical result
$$2\frac{\pi ^2}{\epsilon }m^2+\pi ^2m^2\mathrm{log}[m^2]+\mathrm{}$$
(38)
Suppose that we want to extract the finite part<sup>9</sup><sup>9</sup>9 It is more instructive to discuss the finite part, since the integral contour in (33) is trivial, and extracting the singular part is straightforward. $`\pi ^2m^2\mathrm{log}[m^2]`$. We apply the second method B described above. We use the five dimensional array (21) for $`l_\mu `$. Then the integrand is a function of $`x,t^2,\theta ,\varphi ,\delta `$. The finite part of the (37) is obtained computing the two integrals in $`I_0`$ (eq.(Method of Integration B)). For the first one we choose the contour of integration as in figure 1: the numerical values of $`\mathrm{\Lambda }_{uv}`$ and $`\mathrm{\Lambda }_{ir}`$ have to be chosen in such a way that the contour A+B+C contains all the physical singularities in the complex $`t^2`$ plane. Since the integral is infrared convergent we can take the limit $`\mathrm{\Lambda }_{ir}0`$. Instead, in the ultraviolet region, we can cut the integral at $`\mathrm{\Lambda }_{uv}^2=9`$. We also must set $`x=1`$. Then the integral (35) becomes
$$P_1=_0^9dt^2d\mathrm{\Omega }_4\frac{1}{2}\frac{t^2}{t^2+2t\mathrm{cos}[\theta ]+3}=51.929.$$
(39)
One can recognize in the integral above the ordinary four dimensional integration with a cut-off $`\mathrm{\Lambda }_{uv}^2`$. Then the integration in $`\mathrm{d}t^2`$ follows the path B
$$P_2=\frac{1}{2\pi i}_Bdt^2d\mathrm{\Omega }_4\mathrm{log}[t^2]\frac{1}{2}\frac{t^2}{t^2+2t\mathrm{cos}[\theta ]+3}=43.1815.$$
(40)
This completes the first integral in $`I_0`$. The second one in eq.(Method of Integration B) contains a non trivial integration in the variable $`x`$, which must be performed using the prescription (26) . The contour of integration in $`\mathrm{d}t^2`$ is rather simple: there is no logarithm, and no cut in the real axis; then we do not need to follow the path A. There is no infrared singularity and also the path C vanishes. The only non trivial contribution comes from the path B, a circle of radius $`\mathrm{\Lambda }_{uv}^2`$. This gives
$$P_3=\frac{1}{i2\pi }_0^1_B\frac{1}{t^2+2xt\mathrm{cos}[\theta ]+3}t^2dt^2\left(\frac{x^3}{1x^2}\right)_+dxd\mathrm{\Omega }_4=4.935.$$
(41)
Note in the denominator the appearance of the variable $`x`$, which comes out when we take the scalar product $`lp`$ with $`l_\mu `$ from eq.(21) and $`p^\mu `$ in the direction $`\mu =1`$. The sum of the three contributions yields
$$I_0\left[\frac{1}{(l+p)^2+m^2}\right]=P_1+P_2+P_3=13.6825\pi ^2m^2\mathrm{log}[m^2],$$
(42)
in perfect agreement with the analytical expression. The method can also be applied for infrared divergent integrals. The integral
$$\mathrm{d}^nl\frac{1}{l^6}\frac{1}{(l+p)^2+m^2}$$
(43)
is ultraviolet convergent and infrared divergent. Thus we choose $`\mathrm{\Lambda }_{uv}=\mathrm{}`$ and $`\mathrm{\Lambda }_{ir}=1/4`$. We get
$$P_1=_{1/4}^{\mathrm{}}t^2dt^2\frac{1}{2}d\mathrm{\Omega }\frac{1}{t^2+2t\mathrm{cos}[\theta ]+3}\frac{1}{t^6}=10.843,$$
(44)
$`P_2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _C}dt^2d\mathrm{\Omega }_4\mathrm{log}[t^2]{\displaystyle \frac{1}{2t^4}}{\displaystyle \frac{1}{t^2+2t\mathrm{cos}[\theta ]+3}}=12.1255`$
$`P_3`$ $`=`$ $`{\displaystyle \frac{1}{i2\pi }}{\displaystyle _0^1}{\displaystyle _B}{\displaystyle \frac{1}{t^2+2xt\mathrm{cos}[\theta ]+3}}t^2dt^2\left({\displaystyle \frac{x^3}{1x^2}}\right)_+dxd\mathrm{\Omega }_4=0.1825.`$ (45)
The sum of the three integrals gives
$$I_0\left[\frac{1}{l^6}\frac{1}{(l+p)^2+m^2}\right]=P_1+P_2+P_3=1.465,$$
(46)
very close to the exact result
$$\frac{1}{27}\pi ^2\left(1+\mathrm{log}\left[\frac{81}{4}\right]\right).$$
(47)
## Conclusions
In the dimensional regularization of loop integrals, usually one defines the integral as a function of the number of dimensions in a region of $`n`$ where the integral is convergent. Then one makes an analytic continuation to obtain a definition of the integral also in regions of $`n`$ where the integral is divergent. This procedure is clear and unambiguous, however it can only be applied in pure analytic (and symbolic) calculations. It cannot be converted into an obvious numerical procedure. One is obliged to perform symbolic calculations, which sometimes becomes lengthy and/or unaffordable. In this paper we have shown how to set up a different approach, where each coefficient the Laurent expansion in powers of $`\epsilon `$ in eq. (1) can be written in terms of concrete and convergent integrals (well) defined in a five dimensional space (see eq.(8) and eqs.(32),(33) and (Method of Integration B) . Instead of abstract $`n`$ dimensional vectors we have to deal with “concrete” five dimensional vectors, and we have to perform integrations over a compact space. This formulation has the advantage to allow us the evaluation of all integrals through pure numerical methods. Clearly further work is needed, to generalize the above result to more loops and to prove the efficiency of this technique in realistic physical problems. However we believe that the simplicity of the numerical approach makes it very promising.
## Acknowledgements
I would like to thank P. Nason for very helpful discussions, and Prof. R. Ferrari and the Theoretical Physics Department of Milano University for its kind hospitality.
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# NORDITA-2000/43 HE LAPTH-793/00 hep-ph/0004234 April 26, 2000 Quarkonium Production through Hard Comover Scattering. II [1]
## I Introduction
Quarkonium production is a rich domain of QCD which involves several hardness scales, namely the transverse momentum $`p_{}`$ of the heavy quark pair, the quark mass $`m`$ and the Bohr momentum of the heavy quarks in quarkonium $`\alpha _sm`$. Furthermore, bremsstrahlung gluons are produced along with the pair at an intermediate hardness scale. The proximity of the open flavour threshold makes the production of a quarkonium state highly sensitive to its production environment. Relatively soft rescatterings of the quark pair with surrounding fields can make or break the bound state.
Quarkonium production may be contrasted with open heavy flavour production, for which the total cross section is unaffected by late rescattering. This insensitivity is technically expressed through the QCD factorization theorem , which allows the cross section to be written in terms of process-independent parton distributions. An analogous factorization of hard and soft physics does not apply to quarkonium production rates, which constitute a small fraction of the total open flavour cross section. This makes quarkonium physics a challenging and potentially rewarding field which can teach us new aspects of the dynamics of hard processes.
Attempts to ignore rescattering effects in quarkonium production have met with success in photoproduction but failed in hadroproduction. The PQCD prediction for quarkonium production without rescattering, at lowest order in the quark pair relative velocity $`v`$ (the ‘Color Singlet Model’, CSM ), accounts well for the HERA data on $`\gamma pJ/\psi X`$ in the photon fragmentation region . On the contrary, the CSM underestimates the Tevatron $`\overline{p}p`$ data on direct $`J/\psi `$ and $`\psi ^{}`$ production at large $`p_{}`$ by a factor $`50`$ . Similar discrepancies are also observed in hadroproduction at low $`p_{}`$ .
The CSM amplitudes contain infrared divergencies, which for $`P`$-wave states appear already at lowest order in $`\alpha _s`$. The divergencies cancel systematically between terms of different orders in $`v`$, as shown by the ‘NRQCD’ expansion of the amplitude. In the absence of rescattering effects NRQCD is thus a theoretically viable framework for calculating quarkonium production in QCD.
The NRQCD terms of higher order in $`v`$ correspond to relativistic effects in the quarkonium wave function, and are related to its higher Fock components such as $`|Q\overline{Q}g`$. The magnitude of the relativistic corrections is poorly known, and higher powers of $`v`$ may be compensated by lower powers of $`\alpha _s`$ at the hard vertex. Hence one may consider the possibility that the observed discrepancies between the CSM and hadroproduction data are due to higher order contributions in the NRQCD expansion (the ‘Color Octet Model’, COM ). This implies, however, new contributions also in photoproduction which tends to lead to an overestimate of the data . The phenomenological difficulties of a COM approach have recently been made more acute by the (preliminary) CDF data on $`\psi ^{}`$ polarization at large $`p_{}`$ . The COM prediction of a transverse $`\psi ^{}`$ polarization, considered as a crucial test of the COM approach, fails by more than 3 standard deviations .
It would thus appear that rescattering is important in quarkonium hadroproduction. Since gluons, unlike photons, carry a comoving color field one expects more rescattering in hadroproduction than in photoproduction. In (referred to as I in the following) a scenario involving a perturbative rescattering of $`Q\overline{Q}`$ pairs produced at low $`p_{}m`$ was studied. A simple form of the comoving field allowed several of the puzzling features of the data to be understood.
Here we shall extend the rescattering picture of I to quarkonium production at high $`p_{}`$. We refer to I for a detailed discussion of the rescattering mechanism and its phenomenological justification. In section II we present the low $`p_{}`$ calculation in a more compact form than in I and recall the results which were obtained. In section III we propose a generalization of the low $`p_{}`$ approach to high $`p_{}`$, and evaluate $`S`$-wave and $`P`$-wave quarkonium production amplitudes. All results obtained within the rescattering picture (at low and high $`p_{}`$) are summarized in section IV, together with a discussion of our assumptions.
## II Low $`p_{}`$ Quarkonium Production
### A Physical picture
The rescattering mechanism of quarkonium hadroproduction proposed in I is based on an interaction of the heavy quark pair with a comoving color field. Let us recall why a color field comoving with the $`Q\overline{Q}`$ pair may be created in hadroproduction, but not in photoproduction.
In the QED process $`e^+e^{}\mu ^+\mu ^{}`$ near threshold (Fig. 1a), the incoming electrons are surrounded by an electromagnetic (photon) field. As the electrons annihilate, these fields pass through each other without interacting (except via $`𝒪(\alpha _{em})`$ processes). The $`\mu ^+\mu ^{}`$ pair is thus created in a field-free environment.
In the analogous QCD process $`ggQ\overline{Q}`$ (Fig. 1b), the color fields associated with the color charge of the colliding gluons on the contrary interact strongly. Such multiple interactions may create a remnant color field with low rapidity components in the $`Q\overline{Q}`$ rest frame. Rescattering may occur between this remnant field and the $`Q\overline{Q}`$ pair.
In photoproduction, $`\gamma gQ\overline{Q}`$ (Fig. 1c), the (unresolved) photon carries no radiation field and the situation is analogous to QED: no remnant field is comoving with the $`Q\overline{Q}`$ pair. Rescattering thus occurs only in hadroproduction, and may explain some of the observed ‘anomalies’ of quarkonium production. It is absent in photoproduction (except for resolved photon contributions), where the CSM predictions in fact are quite satisfactory.
### B The model
The effects on low $`p_{}`$ quarkonium hadroproduction of a rescattering of the $`Q\overline{Q}`$ pair off a comoving color field were studied in I in the framework of a simple model. Before extending this scenario to quarkonium production at high $`p_{}`$ we summarize the main results of I.
The model applies to high energy production, where the quarkonium carries a moderate fraction $`x_F`$ of the hadron beam momentum and has $`p_{}m`$, where $`m`$ is the heavy quark mass. The amplitude for quarkonium production in the bound state rest frame shown in Fig. 2,
$``$ $`=`$ $`{\displaystyle \underset{L_z,S_z}{}}LL_z;SS_z|JJ_z{\displaystyle \underset{\lambda \overline{\lambda },\sigma \overline{\sigma }}{}}{\displaystyle \frac{d^3𝒑}{(2\pi )^3}\frac{d^3𝒑^{}}{(2\pi )^3}d^3𝒒}`$ (1)
$`\times `$ $`\delta ^3(𝒑+𝒑^{}+\mathbf{})\mathrm{\Phi }_{\lambda \overline{\lambda }}^{[8]}(𝒑,𝒑^{})_{\lambda \overline{\lambda },\sigma \overline{\sigma }}(\mathbf{},𝒑,𝒒)\mathrm{\Psi }_{\sigma \overline{\sigma }}^{L_zS_z}(𝒒)^{}`$ (2)
is a convolution of the $`ggQ\overline{Q}`$ color octet wave function $`\mathrm{\Phi }^{[8]}`$, the kernel $``$ describing the rescattering between the pair and the comoving color field, and the quarkonium wave function $`\mathrm{\Psi }`$.
The $`ggQ\overline{Q}`$ process occurs on a proper time scale of order $`m^1`$, whereas the time scale for quarkonium bound state formation is of order $`(\alpha _sm)^1`$. The rescattering is assumed to be characterized by the ‘semi-hard’ scale $`\mu `$ of DGLAP gluon radiation
$$\alpha _sm\mu m$$
(4)
As a consequence, the rescattering is well separated in time from both the $`Q\overline{Q}`$ pair creation and the bound state formation. The heavy quarks thus propagate nearly on-shell both before and after the rescattering. In Eq. (LABEL:prodampl) the helicities and momenta of the quarks (in the quarkonium rest frame) before the rescattering are denoted $`\lambda ,\overline{\lambda }`$ and $`𝒑,𝒑^{}`$. After the rescattering of momentum transfer $`\mathbf{}`$ the same quantities are denoted by $`\sigma ,\overline{\sigma }`$ and $`𝒒,𝒒`$. The calculation of I can be summarized by expressing every factor of (LABEL:prodampl) as a contraction between the spinors $`\chi _\lambda `$ and $`\chi _\lambda ^{}`$, with $`\chi _+^{}=\left(\mathrm{1\; 0}\right)`$, $`\chi _{}^{}=\left(\mathrm{0\; 1}\right)`$,
$`\mathrm{\Phi }_{\lambda \overline{\lambda }}^{[8]}(𝒑,𝒑^{})`$ $`=`$ $`\chi _\lambda ^{}\stackrel{~}{\mathrm{\Phi }}^{[8]}(𝒑,𝒑^{})\chi _{\overline{\lambda }}`$ (6)
$`_{\lambda \overline{\lambda },\sigma \overline{\sigma }}(\mathbf{},𝒑,𝒒)`$ $`=`$ $`(2\pi )^3\delta ^3(𝒑𝒒+\mathbf{})\delta _{\overline{\lambda }}^{\overline{\sigma }}\chi _\sigma ^{}\stackrel{~}{}^Q(\mathbf{},𝒒)\chi _\lambda `$ (7)
$`+`$ $`(2\pi )^3`$ $`\delta ^3(𝒑𝒒)\delta _\lambda ^\sigma \chi _{\overline{\lambda }}^{}\stackrel{~}{}^{\overline{Q}}(\mathbf{},𝒒)\chi _{\overline{\sigma }}`$ (8)
$`\mathrm{\Psi }_{\sigma \overline{\sigma }}^{L_zS_z}(𝒒)^{}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{LL_z}^{}(𝒒)}{\sqrt{2}}}\chi _{\overline{\sigma }}^{}\widehat{e}(S_z)^{}\chi _\sigma `$ (9)
In $`\stackrel{~}{}^Q(\mathbf{},𝒒)`$ and $`\stackrel{~}{}^{\overline{Q}}(\mathbf{},𝒒)`$, $`𝒒`$ stands for the momentum of the scattered quark or antiquark after the scattering. Eq. (9) holds for $`S=1`$ quarkonia, $`𝒆(S_z)`$ being the bound state spin polarization vector defined in the rest frame as
$`𝒆(\pm 1)`$ $`=`$ $`(1,i,0)/\sqrt{2}`$ (11)
$`𝒆(0)`$ $`=`$ $`(0,0,1)`$ (12)
We use the notation $`\widehat{a}=𝒂𝝈`$ for a vector $`𝒂`$, with $`\sigma _k`$ the Pauli matrices. The amplitude (LABEL:prodampl) takes a form which can be readily inferred from Fig. 2,
$`(S=1)`$ $`=`$ $`{\displaystyle \underset{L_z,S_z}{}}LL_z;SS_z|JJ_z{\displaystyle \frac{d^3𝒒}{(2\pi )^3}\frac{\mathrm{\Psi }_{LL_z}^{}(𝒒)}{\sqrt{2}}}`$ (13)
$`{\displaystyle \frac{1}{\sqrt{3}}}T_{ji}^b`$ $`\{\mathrm{Tr}\left[\widehat{e}(S_z)^{}\stackrel{~}{}^Q(\mathbf{},𝒒)\stackrel{~}{\mathrm{\Phi }}^{[8]}(𝒒\mathbf{},𝒒)\right]`$ (14)
$`+`$ $`\mathrm{Tr}\left[\stackrel{~}{}^{\overline{Q}}(\mathbf{},𝒒)\widehat{e}(S_z)^{}\stackrel{~}{\mathrm{\Phi }}^{[8]}(𝒒,𝒒\mathbf{})\right]\}`$ (15)
where $`\mathrm{Tr}`$ denotes the trace of $`2\times 2`$ matrices. For convenience we separate the rescattering color factor from the kernel $``$. The color indices of the rescattering gluon and of the quark and antiquark before the rescattering are denoted by $`b`$, $`i`$, $`j`$. The general form (13) will also be used in the case of high $`p_{}`$ quarkonium production via gluon fragmentation in the next section.
In Eq. (13) the zeroth and first order terms in $`𝒒`$ correspond to $`S`$ and $`P`$-wave production, respectively. Recalling that $`|\mathbf{}|\mu m`$ we work only to first order in the small quantities $`\mathbf{}/m`$, $`𝒒/m`$. The rescattering kernel then reads
$`\stackrel{~}{}^{Q,\overline{Q}}(\mathbf{},𝒒)`$ $`=`$ $`{\displaystyle \frac{g}{2m}}\left[iK+\widehat{G}\right]`$ (17)
$`K(\mathbf{},𝒒)`$ $`=`$ $`2m\mathrm{\Gamma }^0(\mathrm{}^0,\mathbf{})+𝚪(\mathrm{}^0,\mathbf{})(\mathbf{}2𝒒)`$ (18)
$`𝑮(\mathbf{})`$ $`=`$ $`𝚪(\mathrm{}^0,\mathbf{})\times \mathbf{}`$ (19)
The color field comoving with the $`Q\overline{Q}`$ pair is modelled as an external source $`\mathrm{\Gamma }^\mu (\mathrm{})`$, the form of which is not known a priori. The hardness scale $`\mu `$ of $`\mathrm{\Gamma }`$ is given by Eq. (4).
The wave function $`\stackrel{~}{\mathrm{\Phi }}^{[8]}`$ of the quark pair produced in $`ggQ\overline{Q}`$ is
$`\stackrel{~}{\mathrm{\Phi }}^{[8]}(`$ $`𝒑`$ $`,𝒑^{})=g^2T^c_{ij}\{i(d_{a_1a_2c}+{\displaystyle \frac{\delta 𝒑^z}{2m}}if_{a_1a_2c})(𝒆_1\times 𝒆_2)^z`$ (20)
$`+`$ $`{\displaystyle \frac{d_{a_1a_2c}}{2m}}[𝒆_1𝒆_2\delta 𝒑^z\sigma _3+𝒆_1\delta 𝒑\widehat{e}_2+𝒆_2\delta 𝒑\widehat{e}_1]\}`$ (21)
with $`\delta 𝒑=𝒑𝒑^{}`$, $`a_1`$, $`a_2`$ the color indices of the incoming gluons and $`𝒆_i=𝒆(\lambda _i)`$, $`i=1,2`$ their polarization vectors. Using (13), (2) and (20) we recover Eq. (18) of I for $`{}_{}{}^{3}S_{1}^{}`$ production<sup>*</sup><sup>*</sup>*In Eq. (18) of I, the gluon propagator $`i/\mathbf{}^2`$ should be included in $`\mathrm{\Gamma }^\mu `$ and the color factor $`D^d`$ should be multiplied by $`\delta _i^i^{}\delta _j^j^{}`$. These details are not important for the conclusions of I.
$`(`$ $`{}_{}{}^{3}S_{1}^{}`$ $`,S_z)={\displaystyle \frac{d_{a_1a_2b}}{4\sqrt{3}}}{\displaystyle \frac{2g^3R_0}{\sqrt{2\pi m^3}}}\{i\lambda _1\delta _{\lambda _1}^{\lambda _2}𝑮𝒆(S_z)^{}`$ (23)
$``$ $`\mathrm{\Gamma }^0`$ $`(\mathrm{})[\delta _{\lambda _1}^{\lambda _2}\mathbf{}^z\delta _{S_z}^0𝒆_1\mathbf{}\delta _{S_z}^{\lambda _2}𝒆_2\mathbf{}\delta _{S_z}^{\lambda _1}]\}`$ (24)
Here we used for $`S`$-wave bound states
$$\frac{d^3𝒒}{(2\pi )^3}\mathrm{\Psi }_{00}^{}(𝒒)=\frac{R_0}{\sqrt{4\pi m}}$$
(25)
where $`R_0`$ is the wave function at the origin. Summing over $`a_1`$, $`a_2`$, $`b`$ and $`\lambda _1`$, $`\lambda _2`$ yields
$`{\displaystyle }`$ $`|(^3S_1,S_z)|^2={\displaystyle \frac{(N_c^24)(N_c^21)}{16N_c}}{\displaystyle \frac{4g^6R_0^2}{3\pi m^3}}`$ (26)
$`\times \{\begin{array}{cc}\left|𝑮^z\right|^2+(\mathbf{}^z)^2\left|\mathrm{\Gamma }^0(\mathrm{})\right|^2& (S_z=0)\\ & \\ \begin{array}{c}\frac{1}{2}\left|𝑮_{}\right|^2+\frac{3}{2}(\mathbf{}_{})^2\left|\mathrm{\Gamma }^0(\mathrm{})\right|^2\hfill \\ \mathrm{Re}\left[i𝑮^y𝑮^x\right]\hfill \end{array}& (S_z=+1)\end{array}`$ (32)
where no assumption has so far been made on $`\mathrm{\Gamma }^\mu `$.
### C Results and predictions
$`S`$-wave states are observed to be unpolarized at low $`p_{}`$ . Comparing to Eq. (32), this suggests that the gluons in $`\mathrm{\Gamma }`$ are dominantly transverse ($`|\mathrm{\Gamma }^0||𝚪|`$) and isotropically distributed. Hence we use the ansatz
$$\mathrm{\Gamma }^\mu (\mathrm{})e_\lambda ^\mu (\mathrm{})\mathrm{with}\lambda =\pm 1$$
(33)
where $`e_{\pm 1}^\mu (\mathrm{})`$ denotes a transverse polarization vector for the gluon $`\mathrm{}`$. The isotropy assumption enters in the integral over $`\mathbf{}`$ after the sum over $`\lambda `$ is performed. In (32) we obtain
$`{\displaystyle \underset{\lambda }{}\left|𝑮^z\right|^2}`$ $``$ $`{\displaystyle \left[(\mathbf{}^x)^2+(\mathbf{}^y)^2\right]}{\displaystyle \frac{2}{3}}{\displaystyle \mathbf{}^2}`$ (34)
$`{\displaystyle \underset{\lambda }{}\frac{1}{2}\left|𝑮_{}\right|^2}`$ $``$ $`{\displaystyle \left[\frac{(\mathbf{}^x)^2+(\mathbf{}^y)^2}{2}+(\mathbf{}^z)^2\right]}{\displaystyle \frac{2}{3}}{\displaystyle \mathbf{}^2}`$ (36)
$`{\displaystyle \underset{\lambda }{}𝑮^y𝑮^x}0`$
The treatment of $`P`$-wave production is similar. For low $`p_{}`$ production the results can be summarized as follows:
* Direct $`S`$-wave production is unpolarized.
* P-wave states are produced in the ratio $`\sigma _{dir}(\chi _1)/\sigma _{dir}(\chi _2)=3/5`$.
* Direct $`\chi _1`$ production is transversely polarized.
* $`\chi _2`$ is produced only with $`J_z=0`$ and $`J_z=\pm 1`$ in the ratio $`\sigma _{dir}(\chi _2,J_z=0)/\sigma _{dir}(\chi _2,J_z=\pm 1)=2/3`$.
As noted in I, the decays of $`{}_{}{}^{3}P_{J}^{}`$ states produced via the rescattering mechanism induce a longitudinal $`J/\psi `$ polarization. Since the $`J/\psi `$ is observed to be nearly unpolarized , this indicates that also the CSM mechanism contributes to $`\chi _2`$ production. CSM produced $`\chi _2`$ states have $`J_z=\pm 2`$ only and decay to transversely polarized $`J/\psi `$’s. For $`\sigma _{CSM}(\chi _2)\sigma _{rescatt}(\chi _2)`$, an overall $`J/\psi `$ polarization consistent with the data is obtained. Hence we also get a lower $`\chi _1/\chi _2`$ ratio, $`\sigma _{dir}(\chi _1)/\sigma _{dir}(\chi _2)0.3`$, compatible with the data .
It is interesting to observe that (low $`p_{}`$) $`P`$-wave production via rescattering is actually independent of $`\mathrm{\Gamma }^0`$. This is because only the $`S=0`$, $`L=1`$ part of the $`Q\overline{Q}`$ wave function (20) (the term proportional to $`f_{a_1a_2c}`$) contributes in this case. Thus a spin-flip rescattering (transverse gluon exchange) is required to form a $`{}_{}{}^{3}P_{J}^{}`$ state.
## III Quarkonium production at high $`p_{}`$
### A Scenario for rescattering
Quarkonium production at high $`p_{}m`$ proceeds dominantly through the fragmentation of quasi-transverse gluons . DGLAP radiation from the fragmenting gluon gives rise to a color field comoving with the quark pair, and we assume that rescatterings occur between this field and the pair. The field comoving with a high $`p_{}`$ $`Q\overline{Q}`$ pair is a priori different from the one at low $`p_{}`$.
A schematic view of our model is shown in Fig. 3 for $`{}_{}{}^{3}P_{J}^{}`$ quarkonia (which require a single rescattering). A high $`p_{}`$ transverse gluon (typically produced in a process like $`gggg`$, which does not concern us here) fragments into a heavy $`Q\overline{Q}`$ pair. The pair rescatters from an external field $`\mathrm{\Gamma }_f`$ of hardness $`\mu _f`$ similar to the scale $`\mu `$ of Eq. (4) and forms a quarkonium bound state. The amplitude of the fragmentation process is a convolution, analogous to Eq. (LABEL:prodampl), of the quark pair wave function $`\mathrm{\Phi }_{f,\lambda \overline{\lambda }}^{[8]}`$, the rescattering kernel $`_{\lambda \overline{\lambda },\sigma \overline{\sigma }}(\mathbf{},𝒑,𝒒)`$ given by Eqs. (8) and (2), and the quarkonium wave function $`\mathrm{\Psi }_{\sigma \overline{\sigma }}^{L_zS_z}(𝒒)`$ of Eq. (9). As in the low $`p_{}`$ situation, the heavy quarks are assumed to be on-shell before and after the rescattering.
The wave function $`\mathrm{\Phi }_{f,\lambda \overline{\lambda }}^{[8]}`$ of the pair in the quarkonium rest frame is
Φf,λλ¯[8]=χλΦ~f[8]χλ¯=gTijau¯λ(𝒑)
e
/
(λg)vλ¯(𝒑)superscriptsubscriptΦ𝑓𝜆¯𝜆delimited-[]8superscriptsubscript𝜒𝜆superscriptsubscript~Φ𝑓delimited-[]8subscript𝜒¯𝜆𝑔subscriptsuperscript𝑇𝑎𝑖𝑗subscript¯𝑢𝜆𝒑
e
/
subscript𝜆𝑔subscript𝑣¯𝜆superscript𝒑\displaystyle\Phi_{f,{\lambda\bar{\lambda}}}^{[8]}=\chi_{\lambda}^{\dagger}\,\tilde{\Phi}_{f}^{[8]}\,\chi_{-\bar{\lambda}}=g\,T^{a}_{ij}\,{\bar{u}}_{\lambda}(\mathchoice{{\hbox{\boldmath$\displaystyle p$}}}{{\hbox{\boldmath$\textstyle p$}}}{{\hbox{\boldmath$\scriptstyle p$}}}{{\hbox{\boldmath$\scriptscriptstyle p$}}})\parbox[b]{6.00006pt}{$e$}\parbox[b]{5.50003pt}{ \raisebox{-0.86108pt}{$/$}}(\lambda_{g}){v}_{\bar{\lambda}}(\mathchoice{{\hbox{\boldmath$\displaystyle p$}}}{{\hbox{\boldmath$\textstyle p$}}}{{\hbox{\boldmath$\scriptstyle p$}}}{{\hbox{\boldmath$\scriptscriptstyle p$}}}^{\prime}) (38)
$`\stackrel{~}{\mathrm{\Phi }}_f^{[8]}(𝒑,𝒑^{})=2mgT_{ij}^a\left\{\widehat{e}(\lambda _g)+{\displaystyle \frac{\widehat{p}}{2m}}\widehat{e}(\lambda _g){\displaystyle \frac{\widehat{p}^{}}{2m}}\right\}`$ (39)
Here $`a`$, $`i`$, $`j`$ are the color indices of the gluon and quarks and $`𝒆(\lambda _g)`$ is the polarization vector of the fragmenting gluon. For the range $`p_{}m`$ that we are considering the gluon is quasi-real and hence transversely polarized ($`\lambda _g=\pm 1`$) relative to its direction in the laboratory frame. An expansion in the small ratios $`𝒑/m`$, $`𝒑^{}/m`$ has been made in Eq. (39).
### B Production of $`{}_{}{}^{3}S_{1}^{}`$ states
Two rescatterings are needed to produce $`{}_{}{}^{3}S_{1}^{}`$ states via gluon fragmentation, see Fig. 4. The fragmentation amplitude $`_f`$ can thus be expressed as a convolution similar to (13), with two rescattering factors $``$. Only the first term in the $`Q\overline{Q}`$ wave function (39), corresponding to a gluon fragmenting into a $`Q\overline{Q}`$ pair with $`L=0`$, actually contributes to $`S`$-wave production.
We find for the $`{}_{}{}^{3}S_{1}^{}`$ production amplitude
$`M_f(^3S_1,S_z)={\displaystyle \frac{mgR_0}{\sqrt{6\pi m}}}`$ (40)
$`\times \{`$ $`\mathrm{Tr}`$ $`\left[ab_2b_1\right]\mathrm{Tr}\left[\widehat{e}(S_z)^{}\stackrel{~}{}^Q(\mathbf{}_2,\mathrm{𝟎})\stackrel{~}{}^Q(\mathbf{}_1,\mathbf{}_2)\widehat{e}(\lambda _g)\right]`$ (41)
$`+`$ $`\mathrm{Tr}`$ $`\left[ab_2b_1\right]\mathrm{Tr}\left[\stackrel{~}{}^{\overline{Q}}(\mathbf{}_2,\mathrm{𝟎})\widehat{e}(S_z)^{}\stackrel{~}{}^Q(\mathbf{}_1,\mathrm{𝟎})\widehat{e}(\lambda _g)\right]`$ (42)
$`+`$ $`\mathrm{Tr}`$ $`\left[ab_1b_2\right]\mathrm{Tr}\left[\widehat{e}(S_z)^{}\stackrel{~}{}^Q(\mathbf{}_2,\mathrm{𝟎})\widehat{e}(\lambda _g)\stackrel{~}{}^{\overline{Q}}(\mathbf{}_1,\mathrm{𝟎})\right]`$ (43)
$`+`$ $`\mathrm{Tr}`$ $`\left[ab_1b_2\right]\mathrm{Tr}\left[\stackrel{~}{}^{\overline{Q}}(\mathbf{}_2,\mathrm{𝟎})\widehat{e}(S_z)^{}\widehat{e}(\lambda _g)\stackrel{~}{}^{\overline{Q}}(\mathbf{}_1,\mathbf{}_2)\right]\}`$ (44)
where $`b_1`$, $`b_2`$ are the color indices of the two rescattering gluons and $`\mathrm{Tr}\left[ab_1b_2\right]`$ is a shorthand notation for $`\mathrm{Tr}\left[T^aT^{b_1}T^{b_2}\right]`$. The four terms in (40) correspond to the four Feynman diagrams implied in Fig. 4. Using (2) leads to the surprisingly simple result
$$_f(^3S_1,S_z)=\frac{R_0g^3}{2\sqrt{6\pi m^3}}d_{ab_1b_2}𝑮_1𝒆(\lambda _g)𝑮_2𝒆(S_z)^{}$$
(46)
where $`𝑮_i=𝚪_f(\mathbf{}_i)\times \mathbf{}_i`$. Note that (46) is independent of $`\mathrm{\Gamma }_f^0`$. Only transverse gluons in the comoving gluon field contribute to $`{}_{}{}^{3}S_{1}^{}`$ quarkonium production at high $`p_{}`$.
The expression (46) trivially satisfies the gauge invariance requirement $`_f(\mathrm{\Gamma }_f^\mu (\mathrm{}_i)\mathrm{}_i^\mu )=0`$. Since the quarks are treated as being on-shell, each Feynman diagram contributing to the production amplitude separately satisfies current conservation. However, in order to keep track of gauge invariance when combining the different kernels $``$, one must pay attention to their $`\mathrm{}^0`$-dependence. For instance the value of $`\mathrm{}_1^0`$ in $`\stackrel{~}{}^{Q,\overline{Q}}(\mathbf{}_1,\mathbf{}_2)`$ differs from that in $`\stackrel{~}{}^{Q,\overline{Q}}(\mathbf{}_1,\mathrm{𝟎})`$ by an amount $`\delta \mathrm{}_1^0=\mathbf{}_1\mathbf{}_2/m`$. Now $`\mathrm{\Gamma }_f^\mu (\mathrm{}_1^0,\mathbf{}_1)`$ is related to $`\mathrm{\Gamma }_f^\mu (\mathrm{}_1^0+\delta \mathrm{}_1^0,\mathbf{}_1)`$ through a boost of velocity $`𝜷=\mathbf{}_2/m`$, $`|𝜷|1`$. One realizes that the vector $`𝚪_f(\mathrm{}^0,\mathbf{})`$ (and thus $`𝑮`$) is insensitive to this boost (at the level of our approximation) whereas the difference between $`\mathrm{\Gamma }_f^0(\mathrm{}_1^0,\mathbf{}_1)`$ and $`\mathrm{\Gamma }_f^0(\mathrm{}_1^0+\delta \mathrm{}_1^0,\mathbf{}_1)`$ is such that we get
$$\stackrel{~}{}^{Q,\overline{Q}}(\mathbf{}_1,\mathbf{}_2)\stackrel{~}{}^{Q,\overline{Q}}(\mathbf{}_1,\mathrm{𝟎})$$
(47)
which gives the simple form (46) for $`_f`$.
In the $`S`$-wave production amplitude (46) the dependence in $`\lambda _g`$ and $`S_z`$ is factorized. No correlation appears between the fragmenting gluon and quarkonium polarizations. Consequently, assuming only that $`\mathrm{\Gamma }_f`$ is isotropically distributed in the (comoving) $`Q\overline{Q}`$ pair rest frame we predict that directly produced $`{}_{}{}^{3}S_{1}^{}`$ quarkonia are unpolarized at high $`p_{}`$. Contrary to the low $`p_{}`$ situation, we do not now need to assume the dominance of transverse gluons in $`\mathrm{\Gamma }_f`$. The preliminary CDF data on high $`p_{}`$ $`\psi ^{}`$ production seems to prefer a longitudinal polarization, but the statistics is insufficient for a definite conclusion.
The high $`p_{}`$ cross section ratio $`\sigma (\psi ^{})/\sigma (J/\psi )`$ appears to be larger than the nearly universal ratio measured in low $`p_{}`$ hadroproduction ,
$$\frac{\sigma (\psi ^{})}{\sigma _{dir}(J/\psi )})_{p_{}m}0.4>\frac{\sigma (\psi ^{})}{\sigma _{dir}(J/\psi )})_{p_{}m}0.24$$
(48)
This is in qualitative agreement with our result that the rescattering gluons are transverse and thus couple to the spin of the individual quarks. In low $`p_{}`$ production the target gluon is longitudinal (in the target rest frame) and its coupling is proportional to the $`Q\overline{Q}`$ color dipole size $`r_{}`$. The convolution (LABEL:prodampl) then probes the bound state wave function in relatively larger configurations, where the node in the $`\psi ^{}`$ wave function tends to decrease its contribution . This systematics carries over to diffractive photoproduction, where two longitudinal gluon exchanges give a factor $`r_{}^2`$, and the corresponding ratio (48) is measured to be $`0.15`$ .
### C Production of $`{}_{}{}^{3}P_{J}^{}`$ states
In the case of high $`p_{}`$ $`{}_{}{}^{3}P_{J}^{}`$ production, only one scattering is needed (see Fig. 3). The production amplitude is given by the expression (13) with $`\stackrel{~}{\mathrm{\Phi }}^{[8]}`$ replaced by $`\stackrel{~}{\mathrm{\Phi }}_f^{[8]}`$ of Eq. (3). We expand the curly bracket of (13) at $`𝒒\mathrm{𝟎}`$ and use for $`P`$-wave states
$`{\displaystyle \frac{d^3𝒒}{(2\pi )^3}\mathrm{\Psi }_{1L_z}^{}(𝒒)}`$ $`=`$ $`0`$ (50)
$`{\displaystyle \frac{d^3𝒒}{(2\pi )^3}𝒒\mathrm{\Psi }_{1L_z}^{}(𝒒)}`$ $`=`$ $`i\sqrt{{\displaystyle \frac{3}{4\pi m}}}R_{1}^{}{}_{}{}^{}𝒆(L_z)^{}`$ (51)
where $`R_{1}^{}{}_{}{}^{}`$ is the derivative of the $`P`$-wave function at the origin. In the small $`𝒒`$ expansion the linear terms in $`𝒒`$ cannot arise from the rescattering kernel $`\stackrel{~}{}`$ because of the gauge invariance property mentioned in the preceding subsection. Recalling (47) with $`\mathbf{}_2`$ replaced by $`𝒒`$, one indeed sees that $`\stackrel{~}{}^Q(\mathbf{},𝒒)`$ is actually independent of $`𝒒`$. This means that the final orbital angular momentum of the $`Q\overline{Q}`$ pair $`L=1`$ must be fixed before the scattering, via the $`L=1`$ part of the $`gQ\overline{Q}`$ wave function $`\stackrel{~}{\mathrm{\Phi }}_f^{[8]}`$. This statement remains true for any number of rescatterings. Thus only the second term of (39) contributes to $`P`$-wave production. Inserting (17) and (39) in (13) we arrive at
$`_f(^3P_J,J_z)`$ $`=`$ $`{\displaystyle \frac{g^2R_1^{}/m}{4\sqrt{2\pi m^3}}}\delta _b^aA^{\alpha \beta }B^{\alpha \beta ,ij}𝒆^i(\lambda _g)𝐇^j`$ (53)
$`A^{\alpha \beta }`$ $`=`$ $`{\displaystyle \underset{L_zS_z}{}}LL_z;SS_z|JJ_z𝒆^\alpha (L_z)^{}𝒆^\beta (S_z)^{}`$ (54)
$`B^{\alpha \beta ,ij}`$ $`=`$ $`\delta _i^j\delta _\alpha ^\beta +\delta _i^\alpha \delta _j^\beta \delta _i^\beta \delta _j^\alpha `$ (55)
$`𝐇`$ $`=`$ $`\mathbf{}^2𝚪_f+2m\mathrm{\Gamma }_f^0\mathbf{}`$ (56)
where $`B^{\alpha \beta ,ij}`$ arises from the trace of four Pauli matrices. Using the form of the tensor $`A^{\alpha \beta }`$ given in we get
$`_f(`$ $`{}_{}{}^{3}P_{J}^{}`$ $`,J_z)={\displaystyle \frac{g^2\sqrt{3}R_1^{}/m}{4\sqrt{2\pi m^3}}}\delta ^a_b`$ (57)
$`\times `$ $`\{\begin{array}{cc}𝒆(\lambda _g)𝐇& (J=0)\\ i\sqrt{\frac{2}{3}}𝒆(\lambda _g)(𝐇\times 𝒆(J_z)^{})& (J=1)\\ 0& (J=2)\end{array}`$ (61)
The amplitude vanishes for $`J=2`$ since in this case $`A^{\alpha \beta }`$ is symmetric and $`A_\alpha ^\alpha =0`$. Experimentally $`\sigma (\chi _2)\sigma (\chi _1)`$ at large $`p_{}`$ . The cross section for $`\chi _2`$ production should be generated either by more than one rescattering in our scenario or by the CSM (or COM) mechanism of gluon emission.
To conclude this section we give our prediction for the $`\chi _1`$ polarization at large $`p_{}`$ and discuss its dependence on the rescattering field $`\mathrm{\Gamma }_f`$.
The high $`p_{}`$ fragmenting gluon being quasi transverse we have from (61)
$$|_f^{\chi _1}(J_z)|^2|(𝐇\times 𝒆(J_z)^{})_{}|^2=\{\begin{array}{cc}|𝐇_{}|^2& (J_z=0)\\ |𝐇^z|^2& (J_z=\pm 1)\end{array}$$
(62)
where the direction of the high $`p_{}`$ gluon in the laboratory frame is chosen as the reference axis. For a general field $`\mathrm{\Gamma }_f`$, the $`\chi _1`$ polarization depends on the relative importance of the two terms contributing to the vector $`𝐇`$ given in (56). However, assuming that $`\mathbf{}`$ is isotropically distributed, we find the same result for the $`\chi _1`$ polarization parameter
$$\lambda _{\chi _1}=\frac{\sigma (J_z=+1)\sigma (J_z=0)}{\sigma (J_z=+1)+\sigma (J_z=0)}=\frac{1}{3}$$
(63)
in the two following extreme cases:
* $`\mathrm{\Gamma }_f`$ contains only transverse gluons (the first term of $`𝐇`$ dominates) and is modelled by the ansatz (33) used at low $`p_{}`$.
* $`\mathrm{\Gamma }_f`$ is purely longitudinal (the second term of $`𝐇`$ dominates).
Thus, rather independently of the form of $`\mathrm{\Gamma }_f`$, we predict the $`\chi _1`$ yield to be longitudinally polarized if the rescattering scenario dominates $`\chi _1`$ production at high $`p_{}`$.
## IV Summary and discussion
We presented an extension of the approach of I to high $`p_{}`$ quarkonium production. As in our previous work, we assumed the presence of rescattering between the heavy quark pair and a comoving color field, and found that such a scenario qualitatively agrees with the low and high $`p_{}`$ data. We summarize here our main results and predictions.
At low $`p_{}`$, the rescattering mechanism applies only to hadroproduction, not to photoproduction (in the photon fragmentation region). In order to agree with the observed non-polarization of $`{}_{}{}^{3}S_{1}^{}`$ low $`p_{}`$ hadroproduction, we assume the field $`\mathrm{\Gamma }`$ created at low $`p_{}`$ to contain dominantly transversely polarized gluons. The $`\sigma (\chi _1)/\sigma (\chi _2)`$ ratio is found to be compatible with the measured value and our main predictions at low $`p_{}`$ are:
* Direct $`\chi _1`$ production is transverse.
* $`\sigma _{dir}(\chi _2,J_z=0)/\sigma _{dir}(\chi _2,J_z=\pm 1)=2/3`$. In addition, we expect a significant CSM (no rescattering) contribution to $`\sigma (\chi _2)`$, having $`J_z=\pm 2`$.
The comoving field $`\mathrm{\Gamma }_f`$ arising from high $`p_{}`$ gluon fragmentation is a priori different from the low $`p_{}`$ field $`\mathrm{\Gamma }`$. Our scenario is identical in high $`p_{}`$ hadroproduction and high $`p_{}`$ photoproduction, the $`Q\overline{Q}`$ pair originating from gluon fragmentation in both cases. Our predictions at large $`p_{}`$ are:
* Direct $`{}_{}{}^{3}S_{1}^{}`$ production is unpolarized.
* Direct $`\chi _1`$ production contains a longitudinal component (Eq. (63)).
* $`\chi _2`$ is not produced via (a single) rescattering.
The prediction (iii) is consistent with the present CDF data on $`\psi ^{}`$ polarization , but higher statistics is needed for a definite conclusion. Only transverse gluon rescattering contributes to $`{}_{}{}^{3}S_{1}^{}`$ production. This is in qualitative agreement with the measured hadroproduction ratio $`\sigma (\psi ^{})/\sigma _{dir}(J/\psi )`$ being smaller at low than at high $`p_{}`$ (cf. Eq. (48)), and being even smaller in diffractive photoproduction.
All our results and predictions apply equally to the charmonium and bottomonium families.
We used several simplifying assumptions which we now discuss.
* We supposed the rescattering color field to be isotropically distributed in the $`Q\overline{Q}`$ rest frame. Our predictions on ratios and relative polarization rates depend on this assumption.
* We assumed the rescatterings to be perturbative. Thus only the minimal number of rescatterings required to produce a given bound state was considered. On the other hand the rescattering probability should be large enough for this mechanism to dominate the CSM $`S`$-wave and (low $`p_{}`$) $`\chi _1`$ rates. Whether such a compromise holds is not obvious and will be studied in a future work.
* It is unlikely that the production amplitude is sensitive only to the quarkonium wave function (or its derivative) at the origin, at least for charmonium. This is particularly so for $`P`$-waves (see the discussion of $`\sigma (\chi _2)/\sigma (J/\psi )`$ in I) and when ‘dipole’ factors of $`r_{}`$ enhance larger wave function configurations .
## Acknowledgments
We are grateful for helpful discussions with T. Binoth, S. Brodsky, J. Rathsman and U. Wiedemann. N.M. would like to thank Nordita for its kind hospitality and support.
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# The Lattice Schwinger Model: Confinement, Anomalies, Chiral Fermions and All That Work supported by Department of Energy contract DE-AC03-76SF00515.
## I Introduction
It was argued in an earlier paper that the Contractor Renormalization Group(CORE) method can be used to map a theory of lattice fermions and gauge fields into an equivalent highly frustrated anti-ferromagnet. Although explicit computations were presented only for the free fermion theory, it was argued that a corresponding mapping must exist for the interacting theory because the space of retained states used for the free theory coincides with the set of lowest energy states of the strongly coupled gauge theory. While this argument is true, it is obviously important to have a better understanding of the details of how the mapping works. In order to get some experience with this process for a theory which is well understood we decided to study the lattice Schwinger model (i.e., two-dimensional QED), since the exact continuum solution of this model exists. Before diving into the CORE computation, however, we first needed to understand the degree to which the lattice model exhibits the interesting features of the continuum theory. This paper is devoted to an analytical treatment of the lattice Schwinger model with an eye to clarifying the physics which underlines the continuum solution and identifying those general features of the model which should provide an ultimate check of the correctness of any numerical solution.
The continuum Schwinger model, in addition to being a non-trivial interacting theory of fermions and gauge fields, provides a laboratory for studying a wide range of interesting phenomena. It exhibits: confinement of the fermionic degrees of freedom and the concomitant appearance of a massive boson in the exact spectrum; breaking of chiral symmetry through the axial anomaly; screening of external charges and background electric fields; infinite degeneracy of the vacuum states of the theory (theta parameters); and the ability to produce arbitrary fermionic polarization charge densities by applying an operator of the form $`e^{i{\scriptscriptstyle dx\alpha (x)j_5^0(x)}}`$ to the vacuum state (due to the anomalous commutator of the electric and axial-charge density operators). It is important to ask which of these features can be understood in the lattice theory before taking the continuum limit and how complicated a CORE computation has to be in order to extract this physics. Although the literature contains discussions of various aspects of the model, such as confinement and the axial anomaly, we are not aware of any systematic discussion of the theory which attempts to parallel the derivation of the continuum solution within the lattice framework. This is what we do in this paper.
In order to make the physics as transparent as possible we formulate the Hamiltonian version of the theory in $`A_0=0`$ gauge and only then rewrite it within the super-selected sector of gauge-invariant states. We then study the Hamilton equations of motion for the electric charge density operator, whose form is completely determined by the way in which local gauge invariance is introduced into the lattice theory. Obviously, the form of the operator equations of motion depends upon the specific lattice fermion derivative and so we study this problem for a wide class of different derivatives; in particular, generalizations of the so-called SLAC derivative, which explicitly maintain the lattice chiral symmetry and generalizations of the Wilson derivative, which break the chiral symmetry for non-zero momenta. We find that all of these approaches produce a satisfactory treatment of the continuum theory, however the detailed physical picture of how things work varies greatly.
We show that a key issue for connecting the lattice theory to the continuum theory is which lattice currents go over to the continuum current operators $`j_0(x)`$ and $`j_0^5(x)`$. Obviously the local lattice charge density operator, whose form is fixed by the way in which one introduces gauge invariance, cannot have this property because the normal ordered version of this operator satisfies the identity $`j_0(i)^3=j_0(i)`$ for all values of the lattice spacing (since only the charges $`0,1,1`$ can exist on a single lattice site). On the other hand, as we will show, the Fourier transform $`j_0(k)`$ can be treated as a boson operator and the the dynamics of the theory tells us that the current operators of the continuum theory are obtained by forming an appropriately regulated version of these lattice operators.
In order to make our discussion essentially self-contained we begin by briefly reviewing the $`A_0=0`$ gauge treatment of the Hamiltonian version of the continuum Schwinger model. We discuss: the need for imposing a state condition, such as restricting to gauge-invariant states; why only the total $`Q=0`$ sector of the theory can exist at finite energy; and why different sectors of gauge-invariant states exist and are labelled by a continuous parameter $`1/2ϵ1/2`$, which can be identified as a background electric field. Finally, we review the Hamiltonian derivation of the fact that the electric charge density is a free massive Bose field and the role played by the anomalous commutator of the electric and axial charge density operators in the derivation of this result. After reviewing the continuum theory we set up and discuss the physics of the lattice version of the Schwinger model in $`A_0(x)=0`$ gauge. We then parallel the continuum arguments as closely as possible for a variety of fermion derivatives. A careful treatment of the Hamilton equations of motion for the Fourier transform of the charge density operator leads to an understanding of how regulated versions of these operators go over to the point-split operators of the continuum theory and the sense in which these regulated operators can be treated as Bose fields. The difference between the way in which things work for generalized SLAC-type derivatives and Wilson-type derivatives becomes clear due to this discussion, as does the connection between the strong and weak coupling theory for generalized SLAC-type derivatives.
## II The Continuum Schwinger Model
Hamiltonian formulations of the continuum Schwinger model have been discussed in the literature. Our discussion will parallel these discussions to a degree but will differ in important details. Our goal is to allow the reader to understand the important features of the Schwinger model without unnecessary formalism.
As we have already noted, the Schwinger model is simply QED in $`1+1`$ dimensions, and has a Lagrangian density given by:
$$=\overline{\psi }(i_\mu \gamma _\mu +eA_\mu \gamma _\mu )\psi \frac{1}{4}F_{\mu \nu }F^{\mu \nu }.$$
(1)
In $`1+1`$ dimensions there are only three anti-commuting $`\gamma `$ matrices, $`\gamma _0,\gamma _1,\gamma _5`$, and so they can be realized in terms of the Pauli $`\sigma `$-matrices:
$`\gamma _0=i\sigma _x,`$ (2)
$`\gamma _1=i\sigma _y,`$ (3)
$`\gamma _5=\gamma _0\gamma _1=\sigma _z.`$ (4)
In order to enable us to give the most physical treatment of gauge-invariance of the theory we choose to work in temporal, or $`A_0(x)=0`$, gauge. Making this choice the Lagrangian density becomes:
$$=\overline{\psi }(i_\mu \gamma _\mu eA(x)\gamma _1)\psi +\frac{1}{2}(_0A(x))^2.$$
(5)
Here, for convenience, we have denoted the spatial component of the vector potential as $`A(x)`$ and dropped its subscript. Eq.(5) tells us that the electric field,
$$E(x)=_0A(x),$$
(6)
is the canonical momentum conjugate to $`A(x)`$ and it has the usual equal-time commutation relations with $`A(x)`$:
$$[E(x),A(x^{})]=i\delta (xx^{}).$$
(7)
Similarly, the fermion operators satisfy the anti-commutation relations
$$\{\psi _\alpha ^{}(x),\psi _\beta (x^{})\}=\delta (xx^{})\delta _{\alpha \beta }.$$
(8)
It follows immediately that the Hamiltonian in $`A_0=0`$ gauge is
$$H=dx\left[\frac{E(x)^2}{2}+\psi ^{}(x)\left(i_1+ieA(x)\right)\sigma _z\psi (x)\right].$$
(9)
There is an essential piece of the physics of working in $`A_0=0`$ gauge which requires discussion. Since we begin by setting $`A_0=0`$ in the Lagrangian, we cannot vary $``$ with respect to $`A_0`$ or $`_0A_0`$, and so we do not obtain Gauss’ law
$$G(x)=\left(_xE(x)e\psi ^{}(x)\psi (x)\right)=0$$
(10)
as an operator equation of motion. In fact, using the canonical commutation relation, Eq.(7), we see that
$`e^{i{\scriptscriptstyle dy\alpha (y)A(y)}}G(x)e^{i{\scriptscriptstyle dy\alpha (y)A(y)}}`$ $`=`$ $`\left(_x(E(x)+\alpha (x))e\psi ^{}(x)\psi (x)\right)`$ (11)
$`=`$ $`G(x)+_x\alpha (x).`$ (12)
This means that even if we start with a state $`|\varphi `$ for which
$$G(x)|\varphi =0,$$
(13)
we can generate states of the form
$$|\varphi _\alpha =e^{i{\scriptscriptstyle d\xi \alpha (\xi )A(\xi )}}|\varphi $$
(14)
for which
$$G(x)|\varphi _\alpha =_x\alpha (x)|\varphi _\alpha .$$
(15)
Fortunately, the operators $`G(x)`$ (which we identify with the generators of time-independent gauge transformations) commute with one another and with $`H`$, and so they can all be simultaneously diagonalized. Thus we are free to impose Eq.(13) as a state condition because the Hamiltonian cannot take us out of this sector of the Hilbert space. Actually, we are free to impose the more general condition of Eq.(15) for any arbitrary function $`\alpha (x)`$. What all this means is that the canonical quantization of the Schwinger model in $`A_0=0`$ gauge produces not one, but rather an infinite number of theories distinguished from one another by the fact that they have, in addition to the dynamical fermion fields, different static classical background charge distributions $`\rho (x)_{\mathrm{class}}=_x\alpha (x)`$. This shouldn’t be a surprise because one should be able to formulate QED in the presence of an arbitrary distribution of static classical background charges. By quantizing in $`A_0=0`$ gauge all we are doing is obtaining all of these possibilities at the same time.
Since, on physical grounds, we are not interested in formulating the Schwinger model in the presence of any non-dynamical charge density, it is customary to limit attention to the so-called gauge-invariant states defined by the condition that $`\rho (x)_{\mathrm{class}}=0`$. Note that this doesn’t quite reduce us to a single possibility since all it means is that $`_x\alpha (x)=0`$ or, in other words, $`\alpha (x)`$ can be an arbitrary constant. If we make such a transformation we shift the operators $`E(x)`$ by a constant, which means that we are free to formulate the theory in the presence of a constant background field $`ϵ`$. If we worked in finite volume this would amount to allowing for the possiblity that there are non-vanishing classical charges on the boundaries; i.e. the remaining sectors of the theory differ by a choice of boundary conditions. One key question associated with the Schwinger model is whether or not the physics is different for different values of the background field. In particular, does the ground-state energy density, which is certainly different for the free theory, depend upon the value of $`ϵ`$ when interacting fermions are introduced into the game.
A simple argument given by Coleman shows that values of $`ϵ`$ which differ by an integer must be equivalent to one another. Before giving the details of the argument it is important to note that in one dimension the solution to the equation
$$_xE(x)=\underset{j}{}e_j\delta (xx_j)$$
(16)
for a set of charges $`e_j`$ located at positions $`x_j`$ only has a finite energy solution when the total charge $`_je_j=0`$. This is so because Eq.(16) tells us that in the regions between the two charges $`E(x)`$ is constant and it changes by an amount $`e_j`$ at each point $`x_j`$. If the sum of the $`e_j`$’s is not zero then, assuming the field vanishes to the left of the first charge at $`x_1`$, the field must continue to infinity to the right of the last charge. This means that in order to minimize the field energy $`E^2(x)/2`$ one or more of the charges must move off to infinity leaving behind a totally neutral system. In particular, if we assume no background field then the energy of a pair of particles with charges $`\pm 1`$ separated by a distance $`s`$ is $`s/2`$. In the presence of a background field $`ϵ>0`$ the situation is different. When the field is present, there is a background energy density equal to $`ϵ^2/2`$. If one now separates a pair of charges oriented so as to reduce the field between the charges to $`ϵ1`$, the total change in the energy of the system is given by
$$\delta =\frac{\left(sϵ^2+s(ϵ1)^2\right)}{2}=\frac{s(12ϵ)}{2},$$
(17)
where the term $`sϵ^2`$ occurs because in the region of length $`s`$ we have replaced the original background field $`ϵ`$ by $`ϵ1`$. From Eq.(17) it follows that for $`ϵ<1/2`$ increasing the separation between the charges costs energy, while for $`ϵ>1/2`$ separating the charges will gain energy (i.e., by moving the charges off to infinity one reduces the background field to $`ϵ^{^{}}=ϵ1`$ and gains an infinite amount of energy). Clearly with that kind of energy gain nothing can stop this process from happening and, since the only change in the problem is that now there will be pairs of charges at $`\pm \mathrm{}`$, it will continue until the background field is reduced to the region $`1/2ϵ1/2`$. For historical reasons this reduced range of $`ϵ`$ is usually parametrized by an angle $`\theta =2\pi ϵ`$ and is one of the two angles which label the exact solutions to the continuum Schwinger model.
If we work in the sector of physical states for which
$$Q_{\mathrm{tot}}|\varphi =e\underset{\mathrm{}}{\overset{\mathrm{}}{}}d\xi \rho (\xi )|\varphi =0,$$
(18)
we can solve for $`E(x)`$ in terms of $`\rho (x)`$
$$E=e\stackrel{x}{}d\xi \rho (\xi ),Q_{\mathrm{tot}}=e\underset{\mathrm{}}{\overset{\mathrm{}}{}}d\xi \rho (\xi )=0.$$
(19)
Substituting this into the Hamiltonian we obtain
$$H=dx\stackrel{~}{\psi }^{}(x)i_x\sigma _z\stackrel{~}{\psi }(x)\frac{e^2}{4}dxdy\stackrel{~}{\rho }(x)|xy|\stackrel{~}{\rho }(y),$$
(20)
where $`\stackrel{~}{\psi }(x)=e^{i_{\mathrm{}}^xd\xi A(\xi )}\psi (x)`$. This field transformation enables us to eliminate the term $`A(x)`$ from the Hamiltonian and simultaneously preserve the canonical commutation relations of operators $`\psi (x)`$, $`\psi ^{}(x)`$. It is important to observe that even if we had not been able to eliminate $`E(x)`$ from the Hamiltonian we could have still made this definition but it would not have been particularly useful since in that case $`E(x)`$ would have non-trivial equal time commutators with the fermion fields and we couldn’t use the canonical quantization rules to carry out computations. Note, in what follows we will, by abuse of notation, drop the tilde and simply write $`\stackrel{~}{\psi }(x)`$ as $`\psi (x)`$.
The content of the exact solution of this model is that it is the theory of a free boson of mass $`m^2=e^2/\pi `$ and, moreover, the charge density operator $`\rho (x)`$ can be used as an interpolating field for this particle because it satisfies a free field equation with the same mass. To see how this happens all we need to do is derive the Heisenberg equations of motion for $`\rho (x)`$.
The time derivative of $`\rho (x)`$ is
$$_0\rho (x)=\frac{1}{i}[\rho (x),H].$$
(21)
Since $`\rho (x)`$ commutes with itself, we use canonical equal time anti-commutation relations for the fermionic fields Eq.(8) and obtain:
$$_0\rho (x)=_xj(x),$$
(22)
where $`j(x)=\psi ^{}(x)\sigma _z\psi (x)`$. Eq.(22) simply states that divergence of the vector current vanishes; i.e., the vector current is conserved.
The second derivative of the charge density operator is now given by
$$_0^2\rho (x)=\frac{1}{i}[_xj(x),H],$$
(23)
which evaluates to
$$_0^2\rho (x)=_x^2\rho (x)\frac{e^2}{4}dy_1dy_2|y_1y_2|\left(\rho (y_1)[i_xj(x),\rho (y_2)]+[i_xj(x),\rho (y_1)]\rho (y_2)\right).$$
The key point in the solution of the Schwinger model is the commutator of $`j(x)`$ and $`\rho (x^{})`$. It is known that this commutator acquires a Schwinger term which we will compute by considering Fourier components of the currents:
$$\rho (x)=\frac{\mathrm{d}k}{2\pi }e^{ikx}\rho _k,j(x)=\frac{\mathrm{d}k}{2\pi }e^{ikx}j_k.$$
(24)
By introducing creation and annihilation operators for the upper $`u_k`$ and lower $`d_k`$ components of the fermion fields with standard anticommutation relations:
$$\{u_k^{},u_q\}=2\pi \delta (kq),\{d_k^{},d_q\}=2\pi \delta (kq),$$
(25)
one obtains:
$$[j_k,\rho _q]=\frac{\mathrm{d}l}{2\pi }\left(u_{lk}^{}u_{l+q}u_{lkq}^{}u_l(ud)\right).$$
(26)
At first sight, this is zero, since the integration momenta $`l`$ can be shifted $`llq`$ in the first term of the integrand. This, however, is not true. The problem is that the momenta shifts can be safely done only in the operators that are normal ordered with respect to the vacuum state, otherwise the difference of two infinite $`c`$-numbers appears. Since, in this basis, the $`e=0`$ Hamiltonian $`H_0`$ reads:
$$H_0=\frac{\mathrm{d}k}{2\pi }k\left(u_k^{}u_kd_k^{}d_k\right),$$
(27)
the vacuum (the lowest energy eigenstate of $`H_0`$) is obtained by filling all negative energy states
$$|\mathrm{vac}=\underset{k<0}{}u_k^{}\underset{k>0}{}d_k^{}|0,$$
(28)
where $`|0`$ is the state annihilated by the $`u_k`$’s and $`d_k`$’s. One may see, that for $`qk`$ in Eq.(26), the right hand side annihilates the vacuum and hence momenta shifts are allowed. For $`k=q`$, however, this is not the case, and that can be easily seen by considering $`[j_k,\rho _k]|\mathrm{vac}`$. One finally obtains:
$$[j_k,\rho _q]=\frac{k}{\pi }2\pi \delta (k+q),$$
(29)
which translates to:
$$[j(x),\rho (x^{})]=\frac{i}{\pi }_x\delta (xx^{}).$$
(30)
Consequently
$$[i_xj(x),\rho (x^{})]=\frac{1}{\pi }_x^2\delta (xx^{}),$$
(31)
and we obtain:
$$_0^2\rho (x)=_x^2\rho (x)\frac{e^2}{2\pi }dy_1dy_2|y_1y_2|\rho (y_1)_x^2\delta (y_2x).$$
(32)
Integrating by parts twice and using $`_x^2|xx^{}|=2\delta (xx^{})`$, we obtain:
$$_0^2\rho =_x^2\rho \frac{e^2}{\pi }\rho .$$
(33)
We see therefore, that the charge density operator $`\rho (t,x)`$ satisfies the equation for the free field with the mass $`\mu ^2=e^2/\pi `$.
Let us take another look at the role of the anomalous commutation relation and the gauge invariance in the exact solution of the Schwinger model. First consider the case $`e=0`$. The equations of motion
$$[\rho _k,H_0]=kj_k,[j_k,H_0]=k\rho _k,$$
(34)
allow us to write the free fermion Hamiltonian as a quadratic polynomial in $`\rho _k`$ and $`j_k`$:
$$H_0=\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}dk\left(\rho _k\rho _k+j_kj_k\right),$$
(35)
since, combined with the anomalous commutator Eq.(29), it produces exactly the same Heisenberg equations of motion. Since the gauge invariance of the theory allowed us to eliminate $`A(x)`$ from the Hamiltonian once $`E(x)`$ was replaced by the Coulomb interaction written in terms of the operators $`\rho _k`$ alone, the full Hamiltonian is obtained by adding the operator
$$H_I=e^2\underset{0}{\overset{\mathrm{}}{}}\frac{\mathrm{d}k}{2\pi }\frac{\rho _k\rho _k}{k^2}$$
(36)
to $`H_0`$. Obviously, $`H_I`$ is also a quadratic polynomial in $`\rho _k`$ and therefore, thanks to the equations of motion, the anomalous commutator of the spatial and temporal components of the vector current and the gauge invariance, the total Hamiltonian is quadratic in $`\rho _k`$ and $`j_k`$. This makes the theory completely solvable in the continuum. We will discuss just how much of this picture survives when one moves from the continuum to lattice version of the theory in the next section.
To complete the usual bosonization of the theory we observe that $`\rho _k`$ and $`j_k`$ don’t satisfy canonical commutation relations, however a simple rescaling remedies this problem and at the same time casts the Hamiltonian into a more familiar form. To be precise, since $`\rho _k`$ has no $`k=0`$ term<sup>*</sup><sup>*</sup>*$`\rho _00`$ would imply that the system is not neutral and that would violate the state condition $`G(x)|\varphi =0`$., we can define
$$\sigma _k=\frac{\sqrt{\pi }}{k}\rho _k,\mathrm{\Pi }_k=\sqrt{\pi }j_k.$$
(37)
Then, using Eq.(29), we see that
$$[\mathrm{\Pi }_k,\sigma _q]=2\pi \delta (k+q),$$
(38)
and the Hamiltonian takes the form
$$H=_0^{\mathrm{}}\frac{\mathrm{d}k}{2\pi }\left(\mathrm{\Pi }_k\mathrm{\Pi }_k+(k^2+\frac{e^2}{\pi })\sigma _k\sigma _k\right).$$
(39)
Given the canonical commutation relations for $`\mathrm{\Pi }_k`$ and $`\sigma _k`$ and this form of the Hamiltonian, it is obvious that we are dealing with the theory of a free massive Bose field.
Let us now turn to the question of the dependence of the theory on the background electric field, or rather to the more general question of what happens in the Schwinger model if we introduce static classical charges. The remarkable property of the Schwinger model is that independent of their magnitude these charges are screened completely. Understanding how this occurs will fully answer the question of how the theory depends upon a background electric field, since we already noted that having a background field of magnitude $`1/2ϵ1/2`$ corresponds to having classical charges of magnitude $`\pm ϵ`$ on the boundaries (or equivalently at $`\pm \mathrm{}`$ ).
From the solution of the theory in terms of $`j(x)`$ it is easy to understand the screening phenomena, since it follows immediately from Eq.(33). Let us consider the Schwinger model with two external charges of the opposite sign:
$$\rho _{\mathrm{ext}}(x)=eQ_{\mathrm{ext}}\left(\delta (xx_1)\delta (xx_2)\right).$$
(40)
As we have seen already, Eq.(13) gets modified to include the external charge density. For this reason the part of the Hamiltonian corresponding to the Coulomb interaction acquires an additional term and the new equation of motion becomes:
$$_0^2\rho =_x^2\rho \mu ^2\left(\rho (x)+\rho _{\mathrm{ext}}(x)\right),$$
(41)
where $`\mu ^2=e^2/\pi `$. This equation implies that there is now a classical, time-independent component of the charge density operator induced by the external charge which satisfies:
$$\rho _{\mathrm{ind}}(x)=\mu ^2eQ_{\mathrm{ext}}\frac{\mathrm{d}k}{2\pi }\frac{\mathrm{cos}(k|xx_1|)}{k^2+\mu ^2}(x_1x_2).$$
(42)
Computing the integral, we obtain for the induced charge density:
$$\rho _{\mathrm{ind}}=\frac{eQ_{\mathrm{ext}}\mu }{2}\left(e^{\mu |xx_1|}e^{\mu |xx_2|}\right),$$
(43)
which, as advertised, screens the external charge densities. One interesting feature of the screening is that two external charges get screened independently from each other . Note also that the screening occurs on the scales $`\mathrm{\Delta }x1/\mu `$, which for small coupling constant can be rather large. Nevertheless, if we now move the external charges off to infinity, so as to go over to the sector which in the free theory would have an external background field, we see that this field is totally screened in the groundstate of the interacting theory. Moreover, since all of the screening takes place within a finite distance of the boundary, there is no contribution to the groundstate energy density coming from the background field.
We should point out that while the previous computation makes it clear that there shouldn’t be a change in the energy density of the groundstate, it is not at all obvious that there isn’t a finite change in the energy of the state due to the regions surrounding the screened external charge. In fact, there clearly is such a change when the external charges are located at a finite distance from one another; however, the question of what happens as one moves these charges to plus and minus infinity is a bit subtle. The crux of the issue has to do with a definition of the limiting process. As will become apparent in a moment the conventional treatment of the Schwinger model amounts to a prescription in which one defines the Hamiltonian of the system as a limit
$$H=\underset{\mathrm{\Omega }\mathrm{}}{lim}H_\mathrm{\Omega }=\underset{\mathrm{\Omega }\mathrm{}}{lim}d\xi H(\xi )$$
(44)
where $`\mathrm{\Omega }`$ is the closed finite interval $`\mathrm{\Omega }=[\omega ,\omega ]`$. With this definition in mind the usual prescription is to first take the classical background charges to plus and minus infinity and then to take the limit $`\mathrm{\Omega }\mathrm{}`$. Given this prescription it is clear that the total Hamiltonian defined in this way never sees the classical screened charges and therefore there is no change in the vacuum energy. In order to see that this is the usual prescription which follows from bosonization of the model let us go back to Eq.(39) and modify it to include the possibility of having an arbitrary external classical charge density $`\rho _{\mathrm{ext}}(x)`$. In configuration space we obtain:
$$H=dx\left(\frac{1}{2}\mathrm{\Pi }(x)^2+\frac{1}{2}(_x\sigma (x))^2+\frac{e^2}{\pi }(\sigma (x)+ϵ(x))^2\right),$$
(45)
where $`ϵ(x)`$ is the function which satisfies the equation
$$\rho _{\mathrm{ext}}(x)=\frac{1}{\sqrt{\pi }}_xϵ(x).$$
(46)
Now, if we set $`ϵ(x)`$ equal to a constant $`ϵ`$ we see that all we have to do is define $`\stackrel{~}{\sigma }(x)=\sigma (x)+ϵ`$ and the Hamiltonian becomes identical to the one without a background field:
$$H=dx\left(\frac{1}{2}\stackrel{~}{\mathrm{\Pi }}(x)^2+\frac{1}{2}(_x\stackrel{~}{\sigma }(x))^2+\frac{e^2}{\pi }\stackrel{~}{\sigma }(x)^2\right).$$
(47)
This is the usual way of handling this issue and so we see that this treatment says that the groundstate energy is independent of the external constant background field, which corresponds to the prescription we gave above.
To complete our discussion of the continuum Schwinger model we present another way of seeing the screening of the classical background field which doesn’t require working with the exact solution to the problem, but only the anomalous commutation relation of $`\rho (x)`$ and $`j(x)`$. The key to this discussion is the introduction of the conserved gauge-dependent current
$$\stackrel{~}{j}(x)=j(x)+\frac{e}{\pi }A(x).$$
(48)
Obviously, since $`A(x)`$ doesn’t commute with the gauge-generators $`G(x)`$ defined in Eq.(10), this current mixes states which satisfy different forms of the general state-condition defined in Eq.(14). This means that we should think of $`\stackrel{~}{j}(x)`$ as operating in the full Hilbert space of the theory obtained by canonical quantization in $`A_0=0`$ gauge without imposing any gauge condition. To show that $`\stackrel{~}{j}(x)`$ is conserved we commute it with the Hamiltonian to obtain
$$_0\stackrel{~}{j}(x)=\frac{1}{i}[\stackrel{~}{j},H]=\frac{1}{i}[j(x),H]+\frac{e}{i\pi }[A(x),H].$$
(49)
Now, a slight rewrite of the derivation of Eq.(33) gives
$$\frac{1}{i}[j(x),H]=_0j(x)=_x\rho (x)\frac{e^2}{\pi }_x^1\rho (x)=_x\rho (x)\frac{e}{\pi }E(x)$$
(50)
and since by construction $`_oA(x)=E(x)`$, we obtain
$$_0\stackrel{~}{j}(x)_x\rho (x)=0,$$
(51)
which means the current is conserved. Integrating this equation over all space we obtain, under the usual assumptions about surface terms, that
$$[dx\stackrel{~}{j}(x),H]=0,$$
(52)
a fact we will use in a moment.
To understand the significance of the fact that $`\stackrel{~}{j}(x)`$ is conserved imagine that we start in a sector of the theory whose lowest energy state satisfies $`0|G(x)|0=0`$. Next consider the transformed state
$$U(\alpha )|0=e^{i{\scriptscriptstyle d\xi \alpha (\xi )\left(j(\xi )+{\scriptscriptstyle \frac{e}{\pi }}A(\xi )\right)}}|0.$$
(53)
We already saw in Eq.(14) and Eq.(15) that the effect of the term proportional to $`A(\xi )`$ in the exponent is to shift the field $`E(x)`$ so that
$$0|U^{}(\alpha )G(x)U(\alpha )|0=\frac{e}{\pi }_x\alpha (x).$$
(54)
This equation says that $`U(\alpha )`$ takes us from a state with no background charge density to one with background charge density equal to $`e_x\alpha (x)/\pi `$. Similarly, it follows from the commutations relations of $`\rho (x)`$ and $`j(x)`$ and an integration by parts, that
$$0|U^{}(\alpha )\rho (x)U(\alpha )|0=\frac{e}{\pi }_x\alpha (x).$$
(55)
Thus, the total effect of applying $`U(\alpha )`$ to the vacuum of sector of the theory with no classical charges is to map this state into a sector which has a non-vanishing classical charge density and at the same time to produce a fermionic charge polarization which cancels it exactly. Now imagine that $`\alpha (x)`$ is chosen as in Fig.1. Since $`_x\alpha (x)`$ vanishes except in the two narrow regions around $`x_L`$ and $`x_R`$ we see that the effect of this operator is to map the original state into one which has equal and opposite classical charge densities around $`x_L`$ and $`x_R`$ and induced cancelling fermionic polarization charge densities. As we move $`x_L`$ and $`x_R`$ to minus and plus infinity respectively the function $`\alpha (x)`$ becomes a constant and in the limit, the fact that $`U(\alpha )`$ commutes with $`H`$ implies that
$$0|U^{}HU|0=0|H|0.$$
(56)
Hence, the energy of the vacuum of the sector with an arbitrary background field is the same as the energy of the vacuum of the sector with no background field, which agrees with the previous argument for the bosonized version of the theory.
## III The Lattice Schwinger Model
Let us now discuss the Hamiltonian version of the Schwinger model on a lattice. In the Hamiltonian formalism time is continuous and space is taken to be an infinite lattice whose points are separated by a distance $`a`$. As in the continuum, we work in $`A_0=0`$ gauge. Furthermore, we introduce fermionic variables $`\psi _n^{}`$ and $`\psi _n`$ associated with each site on the spatial lattice and replace the continuous fields $`A(x)`$ and $`E(x)`$ by conjugate variables $`A_n`$ and $`E_n`$ associated with the link $`(n,n+1)`$ joining the sites $`n`$ and $`n+1`$. This leads to a lattice Hamiltonian of the form
$$H=H_E+H_f,$$
(57)
where
$$H_E=\frac{a}{2}\underset{n}{}E_n^2,H_f=\underset{n,n^{^{}}}{}(\psi _n^{})^\alpha K(nn^{^{}})_{\alpha \beta }e^{ie_{j=n}^{n^{^{}}1}A_j}(\psi _n^{^{}})^\beta .$$
(58)
Here the kinetic term $`K(nn^{^{}})_{\alpha \beta }`$ is a two-by-two matrix for each value of $`nn^{^{}}`$, the fermion fields satisfy the anti-commutation relations
$$\{(\psi _n^{})^\alpha ,(\psi _n^{^{}})^\beta \}=\delta _{n,n^{}}\delta _{\alpha ,\beta }$$
(59)
and the link fields satisfy the usual harmonic oscillator commutation relations
$$[A_n,E_n^{^{}}]=i\delta _{n,n^{^{}}}.$$
(60)
Note that the fermion fields are dimensionless and in order to make the connection to continuum fields we will have to rescale them by a factor of $`1/\sqrt{a}`$ to give them dimensions of $`\mathrm{mass}^{1/2}`$. In direct analogy to the continuum theory, the eigenvalue of the operator $`E_n`$ is the electric flux carried by the link $`(n,n+1)`$. Since, as we have seen, the operator $`e^{ieA_n}`$ shifts the flux on the link $`(n,n+1)`$ by $`e`$ it follows that if we define the normal ordered charge density operator to be
$$\rho _n=:(\psi _n^{}\psi _n):,$$
(61)
then the operators
$$G(n)=E_{n+1}E_ne\rho _n$$
(62)
commute with the Hamiltonian. Hence, similar to the continuum, we are free to impose the discrete version of Gauss’ law
$$G(n)|\varphi =\rho _n^{\mathrm{class}}|\varphi $$
(63)
as a general state condition. Therefore we see that the lattice and continuum versions of the Schwinger model are essentially the same, in that canonical quantization in $`A_0=0`$ gauge gives not one version of two-dimensional QED but rather an infinite number of versions of the theory corresponding to quantizing in the presence of an arbitrary classical background charge distribution. Note that the form of Gauss’ law expressed in Eq.(63) requires us to use the local charge density operator $`\rho _n`$ as the lattice analog of the continuum charge density operator.
Paralleling the discussion of the continuum theory as closely as possible, we focus attention on the zero charge sector of the space of gauge-invariant states; i.e., the ones that satisfy the state condition
$$G(n)|\varphi =0.$$
(64)
Once again, in this sector we can explicitly solve for $`E_n`$ in terms of $`\rho _n`$ and eliminate the factors of $`e^{ieA_n}`$ by incorporating them in the definition of $`\psi _n`$. In this way, in the $`Q=0`$ sector of gauge-invariant states, the lattice Hamiltonian can be written as:
$$H=H_f\frac{e^2a}{4}\underset{n,m}{}\rho _n|nm|\rho _m.$$
(65)
Because the kinetic term $`K(nn^{^{}})_{\alpha \beta }`$ is a function of the difference of $`n`$ and $`n^{^{}}`$ we can write the Hamiltonian in momentum space as:
$$H=\underset{\pi /a}{\overset{\pi /a}{}}\frac{\mathrm{d}k}{2\pi }\psi _k^{}\left\{Z_k\sigma _z+X_k\sigma _x\right\}\psi _k+\frac{e^2a^2}{4}\underset{\pi /a}{\overset{\pi /a}{}}\frac{\mathrm{d}k}{2\pi }\frac{\rho _k\rho _k}{1\mathrm{cos}ak}.$$
(66)
Here we have rewritten the Fourier transform of $`K(nn^{^{}})_{\alpha \beta }`$ in terms of two functions $`Z_k`$ and $`X_k`$, allowing for a very general class of fermion derivatives. Note that in Eq.(66) and all the equations to follow we have adopted the convention that all momentum space operators are normalized in a way that the continuum limit is reproduced by taking $`a0`$ without any additional field renormalization. For example:
$$\{(\psi _k^{})^\alpha ,(\psi _q)^\beta \}=2\pi \delta (kq)\delta _{\alpha \beta }.$$
(67)
Taking our clue from the discussion of the continuum theory we now turn to the derivation of the Heisenberg equations of motion for the current $`\rho _n`$. The first step, namely computing
$$_0\rho _n=\frac{1}{i}[\rho _n,H],$$
(68)
leads us to identify the result of this computation with the divergence of the spatial component of the vector current (or, alternatively, the time component of the axial-vector current) $`j_n`$. Since the discussion to follow is necessarily a bit detailed it is helpful to summarize what it will show us in advance. First, we will see that unlike the charge density operator the current $`j_n`$ is intrinsically point-split as a consequence of the equations of motion. Second, as in the continuum, the important part of the computation of $`_0j_n`$, by taking its commutator with $`H`$, involves commuting the $`j_n`$ and $`\rho _n`$. This computation will show that one cannot solve the lattice Schwinger model exactly for any finite value of the lattice spacing $`a`$ because this lattice commutation relation is not the same as its continuum counterpart. Note that this feature is related to the properties of the free lattice Hamiltonian rather than being a consequence of the interaction. The same computation will show that the continuum limit of the naive commutators does not approach the continuum values for the Schwinger model; from this we will see why, on dynamical grounds, one has to study what amounts to a point-split version of $`\rho _n`$ in order to get the correct physics.
For the purpose of illustration, let us consider explicit forms of $`X_k`$ and $`Z_k`$ corresponding to a number of popular fermion derivatives. In the case of the naive fermion derivative $`Z_k=\mathrm{sin}(ka)/a,X_k=0`$; in the case of the Wilson fermion derivative $`Z_k=\mathrm{sin}(ka)/a,X_k=r/a(1\mathrm{cos}(ak))`$; and for the SLAC derivative one has $`Z_k=k,X_k=0`$. Given any one of these derivatives it is easy to find the one-particle energy levels of the non-interacting Hamiltonian $`H_f`$ by rotating the fields:
$$\chi _k=U_k\psi _k,$$
(69)
where
$$U_k=𝒆^{i\frac{\theta }{2}\sigma _y}=\mathrm{cos}\left(\frac{\theta }{2}\right)+i\sigma _y\mathrm{sin}\left(\frac{\theta }{2}\right),$$
(70)
and
$$\mathrm{cos}\theta _k=\frac{Z_k}{E_k},\mathrm{sin}(\theta _k)=\frac{X_k}{E_k},E_k=\sqrt{X_k^2+Z_k^2}.$$
This unitary transformation diagonalizes the Hamiltonian:
$$H_f=\underset{\pi /a}{\overset{\pi /a}{}}\frac{\mathrm{d}k}{2\pi }E_k\chi _k^{}\sigma _z\chi _k,$$
(71)
and if we introduce creation and annihilation operators for the $`\chi `$-fields
$$\chi _k=\left(\begin{array}{c}u_k\\ d_k\end{array}\right),$$
(72)
with $`\{u_k^{},u_q\}=2\pi \delta (kq)`$ and $`\{d_k^{},d_q\}=2\pi \delta (kq)`$, we obtain:
$$H_f=\underset{\pi /a}{\overset{\pi /a}{}}\frac{\mathrm{d}k}{2\pi }E_k\left(u_k^{}u_kd_k^{}d_k\right).$$
(73)
Finally, the vacuum state of the free theory is obtained by filling the negative energy sea; i.e.,
$$|\mathrm{vac}=\underset{\pi /a<k<\pi /a}{}d_k^{}|0.$$
(74)
Given these equations it is a straightforward matter to compute the commutator of $`H`$ with $`\rho _n`$ to obtain $`_0\rho _n`$:
$$_0\rho _n=\frac{1}{i}[\rho _n,H],$$
(75)
which in the continuum theory is equal to $`_xj(x)`$ (where $`j(x)`$ is identified as the spatial component of the vector current, or the time component of the axial-vector current). Computing the commutator of $`H`$ with $`\rho _n`$ is straightforward but we must say a few words about how we identify $`j_n`$. Basically, in order to maintain the parallel to the continuum discussion we define the quantity equal to $`_0\rho _n`$ as the lattice derivative of $`j_n`$; i.e.,
$$_0\rho _n=\frac{1}{a}\left(j_{n+1}j_n\right).$$
(76)
With this identification, the algebra of $`\gamma `$ matrices in two dimensions ensures that the spatial component of the vector current coincides with the temporal component of the axial current, and therefore all the currents we are going to work with appear to be defined. Clearly, different lattice fermion derivatives will produce different definitions of the spatial component of the vector current operator, an inescapable consequence of the Heisenberg equations of motion.
To derive an explicit form for $`j_n`$, we Fourier transform Eq.(76). Defining
$$\rho _k=\underset{n}{}\rho _ne^{ikan},$$
(77)
we obtain:
$$_0\rho _k=\frac{1}{i}[\rho _k,H].$$
(78)
Writing the right hand side of this equation as:
$$\frac{1}{a}\underset{n}{}\left(j_{n+1}j_n\right)e^{ikan}=\frac{2i\mathrm{sin}(ak/2)e^{ika/2}}{a}j_k,$$
(79)
defines the Fourier transform of the spatial component of the vector current. Explicit computation of $`\rho _k`$ yields:
$$\rho _k=\underset{\pi /a}{\overset{\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }\frac{\mathrm{d}k_2}{2\pi }\psi _{k_1}^{}\psi _{k_2}2\pi \delta ^{\mathrm{lat}}(k_1+kk_2),$$
(80)
where $`\delta ^{\mathrm{lat}}(q)`$ is the lattice $`\delta `$-function which implies the momentum conservation modulo $`2\pi /a`$. Focusing, for the sake of definiteness, on momenta $`k>0`$, one finds:
$$\rho _k=\underset{\pi /a}{\overset{\pi /ak}{}}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}\psi _{k_1+k}+\underset{\pi /ak}{\overset{\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}\psi _{k_1+k2\pi /a}.$$
(81)
It is now a straightforward matter to compute the spatial component of the vector current using Eq.(78):
$$j_k=\frac{ae^{ika/2}}{2\mathrm{sin}(ak/2)}[\underset{\pi /a}{\overset{\pi /ak}{}}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}M(k_1,k)\psi _{k_1+k}.+\underset{\pi /ak}{\overset{\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}M(k_1,k)\psi _{k_1+k2\pi /a}.],$$
(82)
where
$$M(k_1,k)=\left\{\left(Z_{k+k_1}Z_{k_1}\right)\sigma _z+\left(X_{k+k_1}X_{k_1}\right)\sigma _x\right\}$$
(83)
and we have used the fact that $`Z_k`$ and $`X_k`$ are periodic functions with the period $`2\pi /a`$.
From the continuum solution of the Schwinger model it is clear that we should focus on the Schwinger term appearing in the commutator $`[j_k^{},\rho _q]`$, since it is the source of the anomalous Heisenberg equation of motion and the reason for the mass of the photon being non-zero. As we saw in the previous section it suffices to take the vacuum expectation value $`\mathrm{vac}|[j_k^{},\rho _q]|\mathrm{vac}`$ in order to compute the Schwinger term. Direct computation yields the following result:
$$\mathrm{vac}|[j_k^{},\rho _q]|\mathrm{vac}=2\pi \delta (kq)W(k),$$
(84)
where the function $`W`$ is:
$$W=\frac{ae^{ika/2}}{2\mathrm{sin}(ak/2)}\left[\underset{\pi /a}{\overset{\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }\left(2Z_{k_1}Z_{k_1k}Z_{k_1+k}\right)\mathrm{cos}\theta _{k_1}+\left(2X_{k_1}X_{k_1k}X_{k_1+k}\right)\mathrm{sin}\theta _{k_1}\right].$$
(85)
To compare the result of this computation with the continuum result we take the limit $`a0`$, in which case Eq.(85) simplifies and one obtains:
$$\underset{ak0}{lim}W=\frac{k}{\pi }\left[\underset{0}{\overset{\pi }{}}d\xi \left(\frac{\mathrm{d}^2Z_\xi }{\mathrm{d}\xi ^2}\mathrm{cos}(\theta _\xi )+\frac{\mathrm{d}^2X_\xi }{\mathrm{d}\xi ^2}\mathrm{sin}(\theta _\xi )\right)\right].$$
This equation gives the $`a0`$ limit of the anomalous commutator for a general lattice fermion Hamiltonian and is therefore useful for the analysis of the continuum limit of the various choices for the fermion derivative. To get a feeling for how things work let us consider several specific examples.
Let us begin with the case of the naive lattice fermion derivative, where $`Z_\xi =\mathrm{sin}\xi `$, $`X_\xi =0`$, $`E_\xi =|\mathrm{sin}\xi |`$. In this case we obtain:
$$\underset{ak0}{lim}W=\frac{k}{\pi }\underset{0}{\overset{\pi }{}}d\xi \mathrm{sin}\xi =\frac{2k}{\pi }.$$
(86)
This shows that the anomalous commutator is two times larger than the continuum result, which implies that in the $`a0`$ limit the mass of the photon is two times larger than in the continuum theory. In principle, this result should have been expected since the lattice theory with the naive fermion derivative has an exact $`SU(2)`$ symmetry for all values of $`a`$ and as a consequence of this symmetry the fermion spectrum is doubled as is evident from the form of $`E_k`$. Thus, it follows that the continuum limit of the naive theory is not the original Schwinger model, but rather an $`SU(2)`$-Schwinger model which is known to have a photon mass which is $`2e^2/\pi `$.
In case of Wilson fermions we have $`Z_\xi =\mathrm{sin}\xi `$, $`X_\xi =r(1\mathrm{cos}\xi )`$ and $`E_\xi =\sqrt{Z_\xi ^2+X_\xi ^2}`$. By explicit calculation one finds:
$$\underset{ak0}{lim}W=\frac{2k}{\pi },$$
(87)
Once again the result is two times larger than the continuum oneThe fact that the anomalous commutator for Wilson fermions is $`r`$-independent in the $`a0`$ limit is a bit of a miracle and we do not quite understand the reason for that. Note however, that a similar situation has been observed in earlier calculations of the chiral anomaly in the lattice Schwinger model with Wilson fermions . It is generally accepted that the continuum limit of the Schwinger model with Wilson fermions gives correct anomaly and correctly reproduces all other continuum results. However, there is no direct contradiction with our statement since, as we explained and as our result seem to illustrate, different currents lead to different results., but in this case the low energy spectrum is clearly undoubled and the reason for the discrepancy between the lattice and continuum results must be different.
In order to clarify the underlying physics, it is instructive to consider somewhat unconventional fermion derivatives. Let us begin by considering a modified SLAC fermion derivative. Consider the free fermion Hamiltonian defined by ($`k>0`$$`Z_k=Z_k`$):
$$Z_k=k\theta \left(\mu \frac{\pi }{a}k\right)+\frac{\mu }{\mu 1}\left(\frac{\pi }{a}k\right)\theta \left(k\mu \frac{\pi }{a}\right),X_k=0,$$
(88)
and $`E_k`$ is equal to $`|Z_k|`$. A plot of $`E_k`$ is shown in Fig.2. Computing the second derivative of $`Z_k`$ one obtains:
$$\frac{\mathrm{d}^2Z_\xi }{\mathrm{d}\xi ^2}=\delta \left(\mu \frac{\pi }{a}\xi \right)+\frac{\mu }{1\mu }\delta \left(\xi \mu \frac{\pi }{a}\right),$$
(89)
and the anomalous commutator becomes:
$$\underset{ak0}{lim}W=\frac{k}{\pi }\left(1+\frac{\mu }{1\mu }\right).$$
(90)
To understand the information encoded in this form of the anomalous commutator let us consider what happens in the continuum limit. It is clear from the plot of $`E_k`$ for this modified SLAC derivative that two species of fermions survive in the limit $`a0`$. Note however that $`\mathrm{d}E_k/\mathrm{d}k`$ is quite different for the two linear regions of the spectrum, which means that the two species propagate with very different speeds. The anomalous commutator is really the sum of two contributions: one coming from $`0<k<\mu \pi /a`$ and the other from $`\mu \pi /a<k<\pi /a`$ and these contributions can be easily identified with the different fermions. Since both fermions are charged, they both contribute to the anomalous commutator and to the mass gap.
Given the simple nature of this fermion derivative it is clear how to separate the contributions of the two fermion species to the total current. The easiest way to do this is to put a sharp momentum cut-off somewhere below and above the turning point $`\mu \pi /a`$. With this prescription we write the charge density operator as a sum of three contributions:
$$\rho _k=\rho _k^{(1)}+\rho _k^{(2)}+\rho _k^{(3)},$$
(91)
where
$$\rho _k^{(1)}=\underset{\mu _1\pi /a}{\overset{\mu _1\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}\psi _{k_1+k},\rho _k^{(3)}=\underset{|k_1|>\mu _2\pi /a}{}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}\psi _{k_1+k}+\underset{\pi /ak}{\overset{\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }\psi _{k_1}^{}\psi _{k_1+k2\pi /a},$$
(92)
and $`\mu _1<\mu <\mu _2`$. The $`\rho _k^{(2)}`$ provides for the remaining contribution to the charge density and it is only sensitive to fermions with momenta $`\mu _1\pi /a<|k|<\mu _2\pi /a`$. Now, following our previous argument, we define the corresponding spatial components of the vector current by explicitly commuting the above charge densities with the Hamiltonian.
A straightforward computation shows that in the limit of vanishingly small lattice spacing the anomalous commutators of the above currents are given by:
$`[\left(j_k^{(1)}\right)^{},\rho _q^{(1)}]={\displaystyle \frac{k}{\pi }}2\pi \delta (kq),`$ (93)
$`[\left(j_k^{(2)}\right)^{},\rho _q^{(2)}]=0,`$ (94)
$`[\left(j_k^{(3)}\right)^{},\rho _q^{(3)}]=c{\displaystyle \frac{k}{\pi }}2\pi \delta (kq),`$ (95)
where we introduced $`c=\mu /(1\mu )`$, which is the velocity of the fermions in the region $`k\pi `$. Note that in all of these formulas there are also non-vanishing normal-ordered operators coming from large momentum excitations which we have not displayed. These operators annihilate the vacuum and one can argue that they are unimportant for small $`k`$ physics. This final point, which is intimately related to the sense in which $`j_k`$ can be treated as a boson operator, merits elaboration and we will return to it immediately after completing our discussion of the equations of motion.
Proceeding with our computation of the equations of motion for $`\rho _k^{(1)}`$ and $`\rho _k^{(3)}`$ we obtain:
$`_0^2\rho _k^{(1)}=k^2\rho _k^{(1)}+{\displaystyle \frac{e^2}{\pi }}\rho _k^{\mathrm{tot}},`$ (96)
$`_0^2\rho _k^{(3)}=c^2k^2\rho _k^{(3)}+c{\displaystyle \frac{e^2}{\pi }}\rho _k^{\mathrm{tot}},`$ (97)
where the total charge density operator appears on the right hand side of these equations and so $`\rho _k^{(2)}`$ is still included. Consistent with the point made above and subject to the discussion to follow we will set it to zero, since for small energies fermions with such high momenta are not excited.
From the equations for the commutators we see that the field $`\rho _k^{(3)}`$ is not canonically normalized and so we introduce the new field $`\stackrel{~}{\rho }_k^3=1/\sqrt{c}\rho _k^{(3)}`$ and its commutation relation with the corresponding current becomes canonical. The equations of motion become:
$`_0^2\rho _k^{(1)}=k^2\rho _k^{(1)}+{\displaystyle \frac{e^2}{\pi }}\left(\rho _k^{(1)}+\sqrt{c}\stackrel{~}{\rho }_k^{(3)}\right),`$ (98)
$`_0^2\stackrel{~}{\rho }_k^{(3)}=c^2k^2\stackrel{~}{\rho }_k^{(3)}+{\displaystyle \frac{e^2}{\pi }}\left(\sqrt{c}\rho _k^{(1)}+c\stackrel{~}{\rho }_k^{(3)}\right).`$ (99)
The energies of elementary excitations can be determined from the eigenvalues of the matrix
$$=\left(\begin{array}{cc}k^2+e^2/\pi & \sqrt{c}e^2/\pi \\ \sqrt{c}e^2/\pi & c^2k^2+ce^2/\pi \end{array}\right).$$
(100)
The matrix is easily analyzed in the limit of large $`c`$ which corresponds to the large slope of the fermion derivative in the region $`k\pi `$. Note that the original SLAC fermion derivative corresponds to $`c\mathrm{}`$ limit.
It is then easy to see that there are two different limits in this equation. For small momenta $`ck^2e^2/\pi `$, there are two eigenvalues:
$$E_1=\sqrt{c}k,E_2=\sqrt{1+c}\frac{e}{\pi }.$$
Hence there is a zero mass eigenstate which is the Goldstone boson of the theory. The other excitation is the massive one. If we consider the limit $`c\mathrm{}`$, the region of momenta sensitive to the Goldstone mode shrinks to zero (see Fig.3) and the mass of the other excitation goes to infinity.
For larger momenta, $`ck^2e^2/\pi `$, the mixing of two states becomes small, they propagate independently and their energies are given by:
$$E_1=\sqrt{k^2+\frac{e^2}{\pi }},E_2=\sqrt{c^2k^2+c\frac{e^2}{\pi }}.$$
(101)
In this momentum region the energy of the lower excitation approaches the result of the continuum theory of a bosonic field with the mass $`e^2/\pi `$, the other excitation becomes infinitely heavy and decouples explicitly (see Ref.). Hence, as we approach the limit $`c\mathrm{}`$ (which is equivalent to the orginal form of the SLAC derivative), the continuum limit of the theory has a massive boson of mass $`m^2=e^2/\pi `$ and an isolated state at $`k=0`$ which can be identified with a seized Goldstone mode.
Since the main purpose of this paper is to provide an analytic framework for CORE computations to follow we should point out that the fact that there is a Goldstone mode when $`e^2/\pi ck^2`$ is quite significant since it relates to the whole question of how the strong-coupling limit of the model connects to the weak coupling theory and whether one can get the correct physics by projecting onto the sector spanned by the $`e\mathrm{}`$ eigenstates. We will have more to say about this point in the conclusions, but first we should complete our discussion of terms we ignored in the commutator of $`[j_k,\rho _q]`$.
While this preceding argument leading to Eq.(99) makes the lattice discussion look remarkably like the continuum Schwinger model, we are not really finished. The issue which still needs discussion relates to the interpretation of $`\rho _k`$ and $`j_k`$ as boson fields. This is more than an academic issue. Although, as we have shown, the vacuum expectation value of the commutator of $`j_k`$ and $`\rho _q`$ gives the required Schwinger term, computation of the full commutator contains an extra piece which, if it has non-vanishing matrix elements between states whose energy remains finite in the $`a0`$ limit, ruins the interpretation of $`j_k`$ and $`\rho _q`$ as boson fields. This is the issue which we will now address.
The commutation relation for $`j_k`$ and $`\rho _q`$, which is valid for arbitrary lattice spacing, reads:
$$[\left(j_k^{(1)}\right)^{},\rho _q^{(1)}]=\frac{e^{ika/2}ak}{2\mathrm{sin}(ka/2)}\left(2\pi \delta (kq)\frac{k}{\pi }+O(q,k)\right),$$
(102)
where the normal ordered operator $`O(k,q)`$ has the form:
$$O(k,q)=\theta (kq)\{\underset{\mu _1\pi /ak}{\overset{\mu _1\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }:\psi _{k_1+k}^{}\sigma _z\psi _{k_1+q}:+\underset{\mu _1\pi /aq}{\overset{\mu _1\pi /a}{}}\frac{\mathrm{d}k_1}{2\pi }:\psi _{k_1+k}^{}\sigma _z\psi _{k_1+q}:\}+(kq).$$
(103)
We will now argue that even though the term $`O(k,q)`$ is not explicitly suppressed by a power of $`a`$, nevertheless this operator does not contribute to the dynamics of any state whose energy remains finite as $`a0`$; in particular, any state which can be created by applying arbitrary powers of $`\rho _k`$ to the groundstate of the theory. As we pointed out in the Introduction, the operator $`\rho _n`$ cannot be considered a boson operator since $`\rho _n^3=\rho _n`$, thus arbitrary powers of $`\rho _n`$ can produce at most three linearly independent states when they are applied to the groundstate. The situation is quite different for $`\rho _k^{(1)}`$ and $`\rho _k^{(3)}`$ which are sums of $`\rho _n`$’s and therefore, for an infinite lattice, will not satisfy an identity of this type.
The argument begins by considering the non-interacting theory and thinking of $`O(k,q)`$ as acting on a Hilbert space constructed by applying polynomials in the current operators to the groundstate of the free theory. For values of $`k`$ and $`q`$ which are small compared to $`\pi /a`$, the operator $`O(k,q)`$ can only act on the part of a state which contains left and right moving fermions with momenta $`k\pm \mu _1\pi /a`$ since it has to first absorb a high momentum fermion and create another one with a momentum which differs by a small amount. This means that in order to have a matrix element of this operator between states generated by polynomials in $`\rho _k^{(1)}`$ and $`\rho _k^{(3)}`$ these states have to have non-vanishing components having fermions with high momenta. Thus, we have to ask how such components can be generated?
Both $`\rho _k^{(1)}`$ and $`\rho _k^{(3)}`$ are bilinears in fermion creation and annihilation operators and, being normal ordered, can only absorb a fermion at one momentum and create a replacement at another momentum. Generically these operators are of the general form:
$$\rho _k^{(i)}=\frac{\mathrm{d}k_1}{2\pi }\left(u_{k_1}^{}u_{k_1+k}+d_{k_1}^{}d_{k_1+k}\right)$$
(104)
and since for the modified SLAC derivative $`X_k=0`$, the vacuum, as in the continuum, is given by Eq.(28). If we now, for the sake of definiteness, consider
$$\rho _{k>0}^{(1)}|\mathrm{vac}=\rho _k^{(1)}\underset{k<0}{}u_k^{}\underset{k>0}{}d_k^{}|0$$
(105)
we see that almost all terms in $`\rho _k^{(i)}`$ annihilate the vacuum state. The only terms which act non-trivially are ones were either $`u_{k_1+k}`$ or $`d_{k_1+k}`$ can absorb a particle and then either $`u_{k_1}^{}`$ or $`d_{k_1}^{}`$ can create a particle. Clearly, for small $`k>0`$ only the $`d_{k_1}`$ terms can act, since if $`k_1<0`$ and $`k_1+k>0`$ then $`d_{k_1+k}`$ can absorb a $`d`$ from the vacuum state and $`d_{k_1}^{}`$ can create a $`d`$. The $`u_k`$ terms cannot act non-trivially because in order for $`u_{k_1+k}`$ to absorb a particle, $`k_1+k`$ has to be less than zero, in which case $`k_1<0`$ and therefore $`u_{k_1}^{}`$ annihilates the resulting state. If, however, $`k<0`$ then it is the $`u_{k_1}^{}u_{k_1+k}`$ term which acts non-trivially and the corresponding $`d`$ term annihilates the state. In either event the important point is that the $`\rho _k^{(i)}`$ only creates and absorbs particles from the vacuum which are within a distance $`|k|`$ of the top of the negative energy sea (i.e., the fermi-surface) thus creating a particle anti-particle pair.
The next step is to see what happens if we apply $`\rho _k^{(i)}`$ to the state we just generated. What we get is
$$\rho _{k}^{(i)}{}_{}{}^{2}|\mathrm{vac}=\frac{1}{(2\pi )^2}dk_1dk_2\left(u_{k_2}^{}u_{k_2+k}+d_{k_2}^{}d_{k_2+k}\right)\left(u_{k_1}^{}u_{k_1+k}+d_{k_1}^{}d_{k_1+k}\right)|\mathrm{vac}.$$
(106)
It should be clear that for almost all $`k_1`$ and $`k_2`$ in the allowed region $`\rho _{k}^{(i)}{}_{}{}^{2}`$ creates two low momentum particle anti-particle pairs and in fact for given allowed $`k_1`$ and $`k_2`$ there are $`2!`$ ways of getting the same two-pair state; however, for a given $`k_1`$ there is exactly one value of $`k_2`$ for which one can create a higher energy one-pair state by absorbing one of the particles in the pair created by the first application of $`\rho _k^{(i)}`$ and promoting it to higher momentum. From this it follows that the factor needed to normalize this state is greater than $`1/\sqrt{2!}`$. Similarly, if one hits this state with another power of $`\rho _k^{(i)}`$ almost all of the terms would create three low energy particle anti-particle pairs and each of these three pair states would be created $`3!`$ times. As in the previous case however there would be a single term which could promote the previous single higher energy one-pair state to yet higher energy. Note, however, that since the normalization of this state would have to be larger than $`1/\sqrt{3!}`$ (which we are beginning to recognize as the normalization factor which goes with a three boson state) the coefficient of this higher energy single-pair state appearing in the normalized version of the state created by $`\rho _{}^{(i)}{}_{k}{}^{3}`$ is getting smaller each time. If one now imagines carrying out this process $`p`$-times the argument generalizes in the obvious way. The state obtained by applying $`p`$ powers of $`\rho _k^{(i)}`$ to $`|\mathrm{vac}`$ is going to be mostly made of $`p`$ different low-energy particle anti-particle states, each of which will be arrived at in $`p!`$ ways. Furthermore, there will be a single particle anti-particle pair state with individual momenta $`p`$ times larger than $`k`$. Since now the normalization factor of this state is bigger than $`p!`$, the coefficient of this single higher energy pair state is getting very small relative to the part of the wavefunction made of $`p`$ low energy pair states. A more careful discussion of this point would also take into account the fact that the same procedure will generate two-pair, three-pair, etc., parts of the wavefunction. However, the point is that if we keep $`k\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }`$ is a maximal energy we wish to consider and $`\mathrm{\Lambda }a0`$ in the continuum limit, then in order to achieve the fermionic level with momenta $`\mu _1\pi /a`$, one should create a state $`j_k^N|\mathrm{vac}`$ with
$$N\frac{\mu _1\pi }{ak}\frac{\mu _1\pi }{a\mathrm{\Lambda }}\mathrm{}.$$
(107)
The energy of this bosonic state is $`E_{\mathrm{typ}}\mu _1\pi /a\mathrm{}`$; and it is easy to see that the probability of finding a single high momentum pair state equals to $`1/N!0`$. The factors of $`p!`$ which appear in the normalization of the $`j^p|\mathrm{vac}`$ states thus produce the explanation of both why we can think of $`\rho _k^{(i)}`$ as a boson operator and why $`O(k,q)`$ has no significant matrix elements between normalized states generated by applying arbitrary powers of $`\rho _k^{(i)}`$ to $`|\mathrm{vac}`$.
The main point of the above discussion is that as we approach the continuum limit the essential physics of the model is taking place near the top of the negative energy sea and so it is useful to limit our attention to modified operators $`\rho _k^{(i)}`$ that only have support in these regions. For the model based upon a modified SLAC derivative we saw that, since the low energy spectrum of the theory was explicitly doubled, the non-split fermion current really was made up of two parts: the first, coming from states near $`k0`$ and the other from $`k\pi `$. Leaving aside the complications related to the existence of the Goldstone mode, we see that the dynamics of the theory tells us that the current constructed out of the fermionic fields with small momenta is essentially the current with the correct continuum limit. The large momentum part of the current decouples from the continuum limit after the limit $`c\mathrm{}`$ is taken. From this point of view, we see that the dynamics of the lattice model tells us that in order to take the continuum limit of the theory we have to restrict attention to only a part of the unregulated lattice current. This is essentially equivalent to adopting a point-splitting procedure for defining the current in the continuum theory.
Though these peculiarities have been made obvious because of the explicit doubling, our calculation of the anomalous commutation relation shows that for the non-point-split currents the large momentum modes do not decouple automatically, even without fermion doubling. To have explicit decoupling one has to construct the currents by explicitly cutting off the region of large momentum. If in the small momentum region the fermion derivative is sufficiently continuum-like (i.e., linear) then we are assured that: the current constructed in this way will have correct anomalous commutation relations modulo corrections suppressed by inverse cut-off; the equations of motion for this current will be free equations of motion with an additional source term given by high momentum fermionic modes; if the cut-off is sufficiently large in physical energy units (as opposed to lattice units), such a current operator will correspond to a continuum bosonic degree of freedom for all low-energy purposes.
To conclude this section let us consider a simple example of what we will refer to as a perfect Wilson model for the fermion derivative. It is defined by
$$Z_k=k\theta \left(\frac{\pi }{2a}k\right)+\frac{\pi }{2a}\mathrm{sin}\left(\pi ka\right)\theta \left(k\frac{\pi }{2a}\right),X_k=\frac{\pi }{2a}\mathrm{cos}\left(\pi ka\right)\theta \left(k\frac{\pi }{2a}\right),$$
where we have defined $`Z_k`$ and $`X_k`$ for $`k>0`$, and assumed that $`Z_k=Z_k`$ and $`X_k=X_k`$. The one-particle energy spectrum is shown in Fig.4.
One may easily check (using our general result for the anomalous commutator) that for this model, in the limit $`a0`$, the commutator for the non-split currents $`[j_k,\rho _q]`$ differs from the continuum limit, even though in this case there is no doubling and the theory remains continuum-like up to momenta $`k\pi /(2a)`$. Note, however, that if we construct a low energy current by restricting to the linear region of the derivative function, as in the case of the modified SLAC derivative, we can guarantee that the high momenta modes do not have an influence on the dynamics of the low energy current and we can verify that this low-energy current and its time derivative satisfy the desired anomalous commutation relations. Once we establish this fact we can proceed to derive the Heisenberg equations of motion for this low-energy (or regulated) version of the current and make the connection to the continuum theory. Conceptually, our example of the perfect Wilson fermion derivative is very close to Wilson’s original proposal. All we have done is to enlarge the region of momentum space in which the lattice derivative looks identical to the continuum derivative so as to make it easier to see why this type of fermion derivative works once the proper current operators have been identified.
## IV Conclusions
As we have said in the Introduction our aim is to use the lattice Schwinger model to test the idea that one can use CORE methods to map gauge theories into highly frustrated spin antiferromagnerst and then use the same methods to study these spin systems. The Schwinger model is a very good place to test this notion since the continuum model exhibits a rich spectrum of physical phenomena, anomalous commutators, background electric fields, charge screening, etc., and so it is important that any numerical treatment of this model be able to see these effects. We had several major goals in this paper: first, to get some analytic control over the physics of the lattice Schwinger model in order to understand which features of the continuum theory we might expect to emerge easily from a numerical computation and which might be difficult to obtain; second, to gain a feeling for how much of this physics we might hope to see if we first use CORE to map the lattice Schwinger model into a highly-frustrated generalized antiferromagnet and then to analyze the physics of that spin system before carrying out detailed numerical computations; third, to get a better understanding of how the low-energy physics of the lattice system depends upon the choice of fermion derivative and why, on dynamical grounds, the lattice currents of interest are those which correspond to continuum point-split currents.
To accomplish our goals we studied Hamiltonian formulations of both the continuum and lattice Schwinger model and then, by paralleling the solution of the continuum version of the theory in the lattice framework, identified those features of the lattice theory which differ from the continuum theory and identified the operators of the lattice theory which go over smoothly to their continuum counterparts. It became apparent from the treatment of the equations of motion for the various forms of the charge density in the lattice theory that getting the right behavior involves showing that one is close enough to the continuum limit so that the appropriately defined currents act as bosons; in other words, it is not sufficient to only show that there is a gap between the vacuum state and the first excited state and that it numerically appears to be of the order of $`e/\sqrt{\pi }`$. At a minimum one should be able to show that the operators $`O(k,q)`$ have negligible matrix elements between the computed low-lying states of the theory.
A surprising outcome of this work was the fact that almost any fermion derivative works for the study of the Schwinger model. As we have seen, the $`c\mathrm{}`$ limit of the chirality conserving modified SLAC derivative and the perfect Wilson derivative had essentially the same low-energy behavior. The interesting fact was that the dynamics of the system, while different for the two cases, managed to automatically eliminate spurious degrees of freedom. Basically this says that we can use any short-range derivative, either chirality preserving or chirality violating, to carry out numerical studies of the lattice Schwinger model and by comparing them get additional control over how well the numerical methods can be expected to converge.
Finally, and most pertinent to our eventual goal, is the fact that the discussion of the modified SLAC derivative shows that the trick of using CORE to map the system into a frustrated generalized antiferromagnet will preserve the relevant low energy physics. The reason for this is that the CORE method is based on defining the set of retained states to be those states which have zero energy in the limit $`e\mathrm{}`$. This, of course, requires that for these states the Coulomb term vanishes. In other words, these are the states for which the normal ordered charge density operator $`\rho _n`$ is zero identically. (This set of states is generated by selecting from the four possible states per site, only the two states having zero charge and then taking tensor products of all of these states.) Note that for large $`e^2`$ these states are all degenerate to order $`e^2`$ and this degeneracy is lifted by the kinetic term which acts on them by creating a pair of separated charges and then acting a second time to bring them back to a neutral state. A second order degenerate perturbation theory calculation shows that the low energy theory in the large $`e^2`$ limit is that of a Heisenberg anti-ferromagnet, which means that in this limit the theory is that of a massless particle. Going back to Eq.(100) we see that perturbing in the kinetic term is the same as taking $`e^2/\pi `$ to be much greater than $`k^2`$ and $`ck^2`$. But this is exactly the situation in which we have one massive and one massless mode in the theory and the low energy physics is that of a massless boson. This matching of the two results at large $`e^2`$ would imply that the space of retained states must have a non-vanishing overlap with the true low lying states of the theory for finite values of $`e^2`$ which is all that is needed to show that the CORE method must work.
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# Shell model description of 16O(p,𝛾)17F and 16O(p,p)16O reactions
## Abstract
We present shell model calculations of both the structure of <sup>17</sup>F and the reactions <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F, <sup>16</sup>O(p,p)<sup>16</sup>O. We use the ZBM interaction which provides a fair description of the properties of $`{}_{}{}^{16}\text{O}`$ and neighbouring nuclei and, in particular it takes account for the complicated correlations in coexisting low-lying states of $`{}_{}{}^{16}\text{O}`$.
A realistic account of the low-lying states properties in exotic nuclei requires taking into account the coupling between discrete and continuum states which is responsible for unusual spatial features of these nuclei. Within the newly developed Shell Model Embedded in the Continuum (SMEC) approach , one may obtain the unified description of the divergent characteristics of these states as well as the reactions involving one-nucleon in the continuum. This provides a stringent test of approximations involved in the SMEC calculations and permits to asses the mutual complementarity of the reaction and structure data for understanding of these nuclei. The quality of the SMEC description depends crucially on the realistic account of the configuration mixing for coexisting low-lying structures and hence on the quality of the Shell Model (SM) effective interactions and the SM space considered. In this work, we shall present the calculation for $`{}_{}{}^{17}\text{F}`$ for which it is believed that configurations of up to four particles and four holes are necessary. Closed shell nuclei are never inert and multiple particle-hole excitations are always observed in their spectra. In <sup>16</sup>O and <sup>17</sup>O, Brown and Green described the low-lying spectra by mixing spherical and deformed states . From the shell model point of view, Zuker-Buck-McGrory (ZBM) set an effective interaction in the basis of $`0p_{1/2}`$, $`1s_{1/2}`$ and $`0d_{5/2}`$ orbitals . This valence space has the advantage to be practically non spurious and most of states at the $`psd`$ interface around <sup>16</sup>O are nicely described through configuration mixing of these three orbitals. In particular, the energy spectra, spectroscopic factors and correlations in the low-lying states of $`A=16`$ and $`A=17`$ nuclei are well reproduced . The wavefunction components for the first three $`0^+`$ states in <sup>16</sup>O are in a fair agreement with the recently developed interactions in the full $`psd`$ shells . The aim of this work is not to provide better new SM wavefunctions for <sup>16</sup>O but to build on them the continuum effects and to investigate the consequences of this coupling both for the structure of <sup>17</sup>F and for the reactions <sup>16</sup>O(p,$`\gamma `$)<sup>17</sup>F, <sup>16</sup>O(p,p)<sup>16</sup>O. For that purpose, the ZBM interaction is satisfactory and we are going to use it in this study.
In the SMEC formalism the subspaces of (quasi-) bound (the $`Q`$ subspace) and scattering (the $`P`$ subspace) states are not separated artificially . (For the review of earlier works see also Ref. ). Using the projection operator technique, we separate the $`P`$ subspace of asymptotic channels from the $`Q`$ subspace of many-body localized states which are build up by the bound single-particle (s.p.) wavefunctions and by the s.p. resonance wavefunctions. $`P`$ subspace contains $`(N1)`$-particle localized states and one nucleon in the scattering state. The s.p. resonance wavefunctions outside of the cutoff radius $`R_{cut}`$ are included in the $`P`$ subspace. The resonance wavefunctions for $`r<R_{cut}`$ are included in the $`Q`$ subspace. The wavefunctions in $`Q`$ and $`P`$ are then properly renormalized in order to ensure the orthogonality of wavefunctions in both subspaces.
In the first step, we calculate the (quasi-) bound many-body states in $`Q`$ subspace. For that we solve the multiconfigurational SM problem : $`H_{QQ}\mathrm{\Phi }_i=E_i\mathrm{\Phi }_i`$, using the code ANTOINE . For $`H_{QQ}QHQ`$ we take the ZBM interaction which yields realistic internal mixing of many-body configurations in $`Q`$ subspace.
To generate the radial s.p. wavefunctions in $`Q`$ subspace and the scattering wavefunctions in $`P`$ subspace we use the average potential of Woods-Saxon (WS) type with the spin-orbit and Coulomb parts included:
$$U(r)=V_0f(r)+V_{SO}\mathrm{¯}\lambda _\pi ^2(2𝐥𝐬)\frac{1}{r}\frac{df(r)}{dr}+V_C,$$
where $`\mathrm{¯}\lambda _\pi ^2=2`$fm<sup>2</sup> is the pion Compton wavelength and $`f(r)`$ is the spherically symmetrical WS formfactor : $`f(r)=\left[1+\mathrm{exp}((rR_0)/a)\right]^1`$. The Coulomb potential $`V_C`$ is calculated for the uniformly charged sphere with radius $`R_0`$. This ’first guess’ potential $`U(r)`$, is then modified by the residual interaction. We shall return to this problem below.
For the continuum part, we solve the coupled channel equations :
$$(E^{(+)}H_{PP})\xi _E^{c(+)}\underset{c^{^{}}}{}(E^{(+)}H_{cc^{^{}}})\xi _E^{c^{^{}}(+)}=0,$$
where index $`c`$ denotes different channels and $`H_{PP}PHP`$ . The superscript $`(+)`$ means that boundary conditions for incoming wave in the channel $`c`$ and outgoing scattering waves in all channels are used. The channel states are defined by coupling of one nucleon in the scattering continuum to the many-body SM state in $`(N1)`$-nucleus. For the coupling between bound and scattering states around $`{}_{}{}^{16}\text{O}`$, we use the density dependent interaction which is close to the Landau - Migdal type of interactions . However, as compared to the original force of Schwesinger and Wambach , the radius parameter $`r_0`$ of the WS density formfactor is somewhat reduced to better fit the experimental matter radius in oxygen ($`r_0=2.64`$ fm). This interaction provides external mixing of SM configurations via the virtual excitations of particles to the continuum states.
The channel - channel coupling potential is :
$`H_{cc^{^{}}}=(T+U)\delta _{cc^{^{}}}+\upsilon _{cc^{^{}}}^J,`$ (1)
where $`T`$ is the kinetic-energy operator and $`\upsilon _{cc^{^{}}}^J`$ is the channel-channel coupling generated by the residual interaction. Reduced matrix elements of the channel - channel coupling, which involve one-body operators of the kind : $`𝒪_{\beta \delta }^K=(a_\beta ^{}\stackrel{~}{a}_\delta )^K`$, depend sensibly on the amount of $`2p2h`$ and $`4p4h`$ correlations in the ground state of $`{}_{}{}^{16}\text{O}`$. The potential for channel $`c`$ in (1) consists of initial WS guess, $`U(r)`$, and of the diagonal part of coupling potential $`\upsilon _{cc}^J`$ which depends on both the s.p. orbit $`\varphi _{l,j}`$ and the considered many-body state $`J^\pi `$. This modification of the initial potential $`U(r)`$ change the generated s.p. wavefunctions $`\varphi _{l,j}`$ defining $`Q`$ subspace which in turn modify the diagonal part of the residual force, etc. In other words, the procedure of solving of the coupled channel equations is accompanied by the self-consistent iterative procedure which yields for each total $`J`$ independently the new self-consistent potential :
$$U^{(sc)}(r)=U(r)+\upsilon _{cc}^{J(sc)}(r),$$
and consistent with it the new renormalized formfactor of the coupling force. $`U^{(sc)}(r)`$ differs significantly from the initial WS potential, especially in the interior of the potential . Parameters of $`U(r)`$ are chosen in such a way that $`U^{(sc)}(r)`$ reproduces energies of experimental s.p. states, whenever their identification is possible.
The third system of equations in SMEC consists of the inhomogeneous coupled channel equations:
$$(E^{(+)}H_{PP})\omega _i^{(+)}=H_{PQ}\mathrm{\Phi }_iw_i$$
with the source term $`w_i`$ which depends on the structure of $`N`$ \- particle SM wavefunction $`\mathrm{\Phi }_i`$. Formfactor of the source term is given by the self-consistently determined s.p. wavefunctions. The solutions : $`\omega _i^{(+)}G_P^{(+)}H_{PQ}\mathrm{\Phi }_i`$, where $`G_P^{(+)}`$ is the Green function for the motion of s.p. in the $`P`$ subspace, describe the decay of quasi-bound state $`\mathrm{\Phi }_i`$ in the continuum. Reduced matrix elements of the source term, which involve products of two annihilation operators and one creation operator of the kind : $`_{\gamma \delta (L)\beta }^{j_\alpha }=(a_\beta ^{}(\stackrel{~}{a}_\gamma \stackrel{~}{a}_\delta )^L)^{j_\alpha }`$, are calculated between different initial state wavefunctions in $`{}_{}{}^{17}\text{F}`$ and a given final state wavefunction in $`{}_{}{}^{16}\text{O}`$. It should be stressed that the matrix elements of the source term depend sensitively on the percentage of the shell closure in <sup>16</sup>O, i.e., on the amount of correlations both in the g.s. of <sup>16</sup>O and in the considered states of $`{}_{}{}^{17}\text{F}`$. Obviously, this kind of couplings are not accounted for by the spectroscopic amplitudes.
The total wavefunction is expressed by three functions: $`\mathrm{\Phi }_i`$ , $`\xi _E^c`$ and $`\omega _i`$ :
$`\mathrm{\Psi }_E^c=\xi _E^c+{\displaystyle \underset{i,j}{}}(\mathrm{\Phi }_i+\omega _i){\displaystyle \frac{1}{EH_{QQ}^{eff}}}<\mathrm{\Phi }_jH_{QP}\xi _E^c>`$ (2)
where : $`H_{QQ}^{eff}(E)=H_{QQ}+H_{QP}G_P^{(+)}H_{PQ},`$ is the new energy dependent effective SM Hamiltonian which contains the coupling to the continuum. Operator $`H_{QQ}^{eff}(E)`$, which is Hermitian for energies below the particle emission threshold, becomes non-Hermitian for energies higher than the threshold. Consequently, the eigenvalues $`\stackrel{~}{E_i}\frac{1}{2}i\stackrel{~}{\mathrm{\Gamma }_i}`$ are complex for decaying states and depend on the energy $`E`$ of the particle in the continuum. The energy and the width of resonance states are determined by the condition: $`\stackrel{~}{E_i}(E)=E`$. The eigenstates corresponding to these eigenvalues can be obtained by the orthogonal but in general non-unitary transformation . Inserting them in (2), one obtains symmetrically the new continuum many-body wavefunctions modified by the discrete states, and the new discrete state wavefunctions modified by the coupling to the continuum states.
The SMEC wavefunctions, can be used to calculate various spectroscopic and reaction quantities. These include for example the proton (neutron) capture data, Coulomb dissociation data, elastic (inelastic) proton (neutron) scattering data, energies and wavefunctions of discrete and resonance states, transition matrix elements between SMEC wavefunctions, static nuclear moments etc. . The application of the SMEC model for the description of structure for mirror nuclei and capture cross sections for mirror reactions in $`p`$-shell has been published in Ref. . The analysis of the structure of $`{}_{}{}^{17}\text{F}`$ and the reactions $`{}_{}{}^{16}\text{O}(\text{p},\gamma )^{17}\text{F}`$, $`{}_{}{}^{16}\text{O}(\text{p,p})^{16}\text{O}`$ in the $`(0p1s0d)`$-space, neglecting the 2p-2h, 4p-4h admixtures in$`{}_{}{}^{16}\text{O}`$ wavefunctions have been reported recently as well . To correct for the missing correlations in the low-energy wavefunctions of <sup>16</sup>O and <sup>17</sup>F, the matrix elements of the $`QP`$ coupling in this study have been quenched by the factors related to the spectroscopic amplitudes for positive parity states in $`{}_{}{}^{17}\text{F}`$ ($`{}_{}{}^{17}\text{O}`$) and to the amount of $`2p2h,4p4h`$ correlations in the g.s. of $`{}_{}{}^{16}\text{O}`$. This quenching correction of the effective operator allowed to obtain a reasonable description of the spectrum of <sup>17</sup>F but failed in solving the problem of ’halo’ of discrete states for positive energies which was observed in the elastic cross section and in the phase shifts . The whole problem results from the non-hermitean corrections to the eigenvalues for positive energies which generate the imaginary part and which are particularly large for pure single-particle (or single-hole) configurations. For this reason, the incorrect internal configuration mixing in SM wavefunctions may lead to an unphysical enhancement of the resonant-like correction from bound states into the scattering data (e.g. the elastic phase-shifts). This aspect of the continuum coupling we want to investigate using the ZBM interaction which includes the most essential for the present studies configuration mixing in low energy wavefunctions.
In Fig. 1 we show SMEC energies and widths for positive parity (l.h.s. of the plot) and negative parity (r.h.s. of the plot) states of $`{}_{}{}^{17}\text{F}`$. The calculations were performed using either the density dependent interaction of Ref. , called DDSM0, or the similar interaction with the overall strength reduced by a factor 0.67 , called DDSM1 . (In both cases, radius of the density formfactor has been reduced as discussed above.) This latter interaction was used before in the SMEC calculations in $`0p1s0d`$ SM space . However, since the external mixing of wavefunctions in the SM space of ZBM is smaller than in the $`0p1s0d`$ SM space, therefore one may consider bigger coupling strength of the residual interaction in SMEC-ZBM calculations. For that reason, we compare results obtained with both DDSM0 and DDSM1 interactions. For the comparison, the experimental data and the SM input in $`Q`$-space is shown in Fig. 1 separately for positive and negative parity states as well. The agreement of the SM-ZBM calculations with the experimental data is even better than obtained in the $`0p1s0d`$ SM space . The agreement of SMEC calculations with experiment is encouraging for both density dependent interactions but the widths of states are better reproduced by SMEC-DDSM0. Only the coupling matrix elements between the $`J^\pi =0_1^+`$ g.s. wavefunction of $`{}_{}{}^{16}\text{O}`$ and all considered states in $`{}_{}{}^{17}\text{F}`$ are included. Zero of the energy scale for $`{}_{}{}^{17}\text{F}`$ is fixed by the experimental position of $`J^\pi =1/2_1^+`$ ($`E=105`$ keV) with respect to the proton emission threshold. In Fig. 1 the theoretical prediction for $`3/2^+`$ state is compared with the experimental $`3/2_2^+`$ state, because the s.p. orbit $`0d_{3/2}`$ is missing in SM space of ZBM. For those s.p. wavefunctions which in a given many body state are not modified by the selfconsistent renormalization, we calculate radial formfactors using the common reference s.p. potential $`U^{(ref)}`$ of the WS type which is adjusted to reproduce experimental binding energies of of $`5/2_1^+`$ and $`1/2_1^+`$ states for protons . This potential has the following parameters : the radius $`R_0=3.214`$fm, the diffuseness $`a=0.58`$fm, the spin-orbit strength $`V_{SO}=3.683`$MeV and the depth $`V_0=52.46`$MeV. The same potential without the Coulomb term is also used to calculate radial formfactors for neutrons. In the case when the self-consistent conditions in a $`J^\pi `$ state determine the radial function of a s.p. wavefunction $`\varphi _{l,j}`$, we readjust the depth of ’first guess’ potential $`U(r)`$ so that the energy of the s.p. state $`\varphi _{l,j}`$ in the converged potential $`U^{(sc)}(r)`$ corresponds to the energy of this s.p. state in $`U^{(ref)}`$. The remaining parameters : $`R_0`$, $`a`$, $`V_{SO}`$, of the initial potential are the same as in the reference potential $`U^{(ref)}`$. This readjustment of $`V_0`$ in the initial potentials $`U(J^\pi )`$ for a s.p. state $`(l,j)`$ guarantees that the binding of self-consistently determined wavefunction $`\varphi _{l,j}`$ is at the experimental value. Moreover, we have the same s.p. energies and, consequently, the same asymptotic behaviours of wavefunctions for large $`r`$ for both residual density dependent interactions. This choice is essential for quantitative description of radiative capture cross-section.
$`B(E2)`$ transition matrix element between $`1/2_1^+`$ and $`5/2_1^+`$ bound states is an interesting test of the wavefunction. In SMEC with ZBM interaction and DDSM0 residual coupling, the value for this transition : $`B(E2)=79.61`$e<sup>2</sup>fm<sup>4</sup>, which is obtained assuming the effective charges : $`e_p1+\delta _p`$, $`e_n\delta _n`$, with the polarization charge $`\delta =\delta _p=\delta _n=0.2`$, agrees with the experimental value $`B(E2)_{exp}=64.92`$e<sup>2</sup>fm<sup>4</sup>. For the range of $`\delta `$ in between 0.1 and 0.3, which is compatible with the theoretical estimates , the SM values with Harmonic Oscillator wavefunctions are much smaller. This difference reflects a more realistic radial dependence of the $`1s_{1/2}`$ s.p. orbit in $`J^\pi =1/2_1^+`$ many body state which is in part a consequence of the external mixing in SMEC wavefunctions due to the coupling to the scattering continuum and provides the halo structure of the $`1s_{1/2}`$ s.p. orbit. Similar effect can be generated using the WS potential or Pöschl-Teller-Ginocchio potential for the appropriately chosen $`1s_{1/2}`$ orbit . The rms radius for this orbit in SMEC is : $`<r^2>^{1/2}=5.24`$fm, as compared to : $`<r^2>^{1/2}=3.03`$fm in SM.
In Fig. 2, the calculated total astrophysical $`S`$-factor as a function of the c.m. energy, as well as its values for the $`{}_{}{}^{16}\text{O}(p,\gamma )^{17}\text{F}(J^\pi =1/2_1^+)`$ and $`{}_{}{}^{16}\text{O}(p,\gamma )^{17}\text{F}(J^\pi =5/2_1^+)`$ branches are compared with the experimental data . Results shown in Fig. 2 have been calculated with DDSM1 interaction. The total $`S`$-factor depends weakly on the strength of the residual coupling and with DDSM0 interactions one obtains very similar results. The energy scale is adjusted to reproduce the experimental position of $`1/2_1^+`$ state with respect to the proton emission threshold. Therefore, the energy scale for excitation energy is the same as c.m. energy in the $`p+^{16}\text{O}`$ system. The photon energy is then given by the difference of c.m. energy of $`[^{16}\text{O}+\text{p}]_{J_i}`$ system and the experimental energy of the final state $`J_f`$ in $`{}_{}{}^{17}\text{F}`$. The dominant contribution to the total capture cross-section for both $`5/2_1^+`$ and $`1/2_1^+`$ final states, is provided by $`E1`$ transitions from the incoming $`p`$ wave to the bound $`0d_{5/2}`$ and $`1s_{1/2}`$ states. We took into account all possible $`E1`$, $`E2`$, and $`M1`$ transitions from incoming $`s`$, $`p`$, $`d`$, $`f`$, and $`g`$ waves but only $`E1`$ from incoming $`p`$ \- waves give important contributions. In the transition to the g.s., the $`E1`$ contribution from incoming $`f_{7/2}`$ wave is by a factor $``$100 smaller than the contribution from $`p_{3/2}`$ wave at $`E_{CM}0`$. This contribution however increases with the energy of incoming proton and becomes $`0.3`$ at 3.5 MeV.
The energy dependence of $`S`$-factor as $`E_{CM}0`$ can be fitted by a second order polynomial to calculated points obtained in the interval from 20 to 50 keV in steps of 1 keV. We have $`S(0)=9.32\times 10^3`$ MeV$``$b, and the logarithmic derivative is $`S^{^{}}(0)/S(0)=4.86`$ MeV<sup>-1</sup>. These results have been obtained for DDSM1 residual interaction but are practically identical for the DDSM0 force. The ratio of $`E2`$ and $`E1`$ contributions in the branch $`{}_{}{}^{16}\text{O}(p,\gamma )^{17}\text{F}(J^\pi =1/2_1^+)`$ for DDSM1 interaction is : $`\sigma ^{E2}/\sigma ^{E1}=1.622\times 10^4`$, $`2.225\times 10^4`$ and $`5.458\times 10^4`$ at 20, 100 and 500 keV, respectively. Also these ratios do not depend on the strength of the residual coupling. On the contrary, the ratio $`\sigma ^{M1}/\sigma ^{E1}`$ depends on the residual coupling and equals 3.90$`\times 10^4`$ and 1.087$`\times 10^3`$ at $`E_{CM}0`$ for DDSM1 and DDSM0 interactions respectively. At 500 keV, this ratio equals 6.11$`\times 10^5`$ and 1.81$`\times 10^4`$ for both interactions, respectively.
For the deexcitation to the g.s. $`5/2_1^+`$, the fit of calculated $`S`$-factor for DDSM1 interaction as $`E_{CM}0`$ yields: $`S(0)=3\times 10^4`$ MeV$``$b and $`S^{}(0)/S(0)=0.649`$ MeV<sup>-1</sup>. Almost identical results are obtained with the DDSM0 interaction. The ratio of $`E2`$ and $`E1`$ contributions for DDSM1 interaction at 20, 100 and 500 keV is: $`\sigma ^{E2}/\sigma ^{E1}=1.336\times 10^3`$, $`1.11\times 10^3`$ and $`9.74\times 10^4`$, respectively. These ratios show some sensitivity on the residual coupling. For DDSM0 one obtains : $`\sigma ^{E2}/\sigma ^{E1}=2.25\times 10^3`$, $`1.458\times 10^3`$ and $`1.043\times 10^3`$ at 20, 100 and 500 keV, respectively. Similarly as in the branch $`{}_{}{}^{16}\text{O}(p,\gamma )^{17}\text{F}(J^\pi =1/2_1^+)`$, the ratio $`\sigma ^{M1}/\sigma ^{E1}`$ depends on the residual coupling and equals 2$`\times 10^3`$ and 2.99$`\times 10^3`$ at $`E_{CM}0`$ for DDSM1 and DDSM0 interactions respectively. In the total cross section, ratio of $`E2`$ and $`E1`$ contributions at 20, 100 and 500 keV equals : $`\sigma ^{E2}/\sigma ^{E1}=2.03\times 10^4`$, $`2.63\times 10^4`$ and $`5.85\times 10^4`$ for DDSM1 interaction , and $`\sigma ^{E2}/\sigma ^{E1}=2.34\times 10^4`$, $`2.8\times 10^4`$ and $`5.92\times 10^4`$ for DDSM0 interaction.
Elastic phase shifts and elastic cross-sections for different proton bombarding energies are shown in Figs. 3 and 4. The elastic phase shifts data are very well reproduced by the SMEC calculations for all considered partial waves including the $`5/2^+`$ for which a significant discrepancy with the data has been reported in SMEC calculations using the $`psd`$ interaction neglecting higher order p-h correlations in the SM wavefunction for positive parity states . In Fig. 3 we show the calculations for the DDSM1 residual interaction. As we have already mentioned, $`0d_{3/2}`$ s.p. orbit is missing in the $`Q`$ subspace. On the other hand, the $`3/2^+`$ partial wave contributes to the elastic cross-section. So we have decided to include this resonance in the $`P`$ subspace, adjusting scattering potential for the $`d_{3/2}`$ proton wave to place it at the experimental position. The encouraging agreement of SMEC results with the data for the spectrum of $`{}_{}{}^{17}\text{F}`$, the proton capture cross-section $`{}_{}{}^{16}\text{O}(p,\gamma )^{17}\text{F}(J^\pi =5/2_1^+)`$ and the $`5/2^+`$ elastic phase shift is not accidental and shows importance of the np-nh excitations across the ’closed core’ $`N=Z=8`$, which are taken into account in the present studies with the ZBM force.
Calculated elastic excitation functions at a laboratory angle of 166 in SMEC with DDSM0 (the solid line) and DDSM1 (the dashed line) are compared with the experimental data in Fig. 4. The agreement is very encouraging for both residual interactions in the almost entire energy range. where the interference pattern depending sensitively on the precise values of the energy and width of resonance states is absent. Earlier SMEC calculations , using the simplified wavefunctions for the g.s. of $`{}_{}{}^{16}\text{O}`$ and $`{}_{}{}^{17}\text{F}`$, failed to reproduce the experimental elastic cross-section in the low-energy domain below the resonances. We have found that an unrealistic account for excitations from $`0p`$ to $`1s0d`$ shells leads to a large resonant contribution from the g.s. $`5/2_1^+`$ wavefunction to the phase shift in the partial wave $`5/2^+`$ and implies a strong decrease of the elastic excitation function at low energies. This effect is caused by the virtual coupling of discrete and continuum states. For energies below the proton emission threshold, coupling to the continuum introduces hermitean modifications of $`H_{QQ}`$ which shift the energy of $`5/2_1^+`$ state with respect to its initial position given by the SM but do not generate any width for this state. For excitation energies above the proton threshold, the $`QP`$ coupling generates non-hermitean corrections which yield the imaginary part of the eigenvalue of $`H_{QQ}^{eff}`$ and, hence, produce the resonant behavior. Internal mixing of configurations in the SM wavefunctions tend to reduce this resonant behavior for positive energies which, in general, is strongest for the states having little internal mixing, i.e. those having a s.p. nature. Actually, these strong resonant-like features associated with certain bound states for positive energies (above the particle-emission threshold) and the large shifts of the real part for certain eigenvalues at negative energies (below the particle-emission threshold), have the same origin in the interplay between external (i.e. via the continuum coupling) and internal (i.e. within $`Q`$ space) configuration mixing in the SMEC wavefunction for this state. This example shows also that in the SMEC approach, one may use different experimental observables to fix those few parameters of the model such as the overall strength of the residual $`QP`$ coupling or the radius and depth of the initial average potential. Moreover, the information about the amount of correlations in the low-lying coexisting states can be extracted not only from the spectroscopic data but also from the elastic excitation function. For nuclei far from the stability line where the amount of experimental data is strongly limited, this feature of the model is very attractive.
We feel that the evidence presented in this work shows that the SM calculation extended to include the coupling to the continuum of the scattering states can go a long way towards providing a detailed explanation not only of the structure of $`{}_{}{}^{17}\text{F}`$, $`{}_{}{}^{16}\text{O}`$ and neighboring nuclei, but also the reaction data involving one nucleon in the continuum. The corner-stone of this model is the effective SM interaction providing a realistic internal mixing of configurations in the $`Q`$ subspace. Several problems remain, such as the missing configurations and/or more realistic asymptotic decay channels which show up in the decay width of resonances. In this latter case, the extension of the SMEC is being investigated.
Acknowledgments
We thank E. Caurier for his help in the early stage of development of SMEC model. This work was partly supported by KBN Grant No. 2 P03B 097 16 and the Grant No. 76044 of the French - Polish Cooperation.
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# Universal features of the order-parameter fluctuations : reversible and irreversible aggregation
## I Introduction
Fluctuations in many physical processes are difficult to analyze because they develop dynamically and often keep the memory of initial conditions. On the other hand, strong fluctuations are ubiquitous, as show examples of hadronization in strong interaction physics, polymerization, colloid aggregation, aerosol coalescence or the formation of large scale structures in the Universe. With the advent of recently developed advanced detection systems, the study of large fluctuations in physical observables became accessible in ’small systems’, such as formed in ultrarelativistic collisions of hadrons, leptons and nuclei, in the heavy-ion collisions at the intermediate energies or in the collisions of atomic aggregates .
In theoretical studies, it is often assumed that fluctuations are irrelevant. In this spirit, many aggregation processes have been studied in the mean-field approximation . This problem has been revisited recently . It was shown that contrary to the usual believe, fluctuations in the size distribution of largest cluster are generally large in the aggregation processes. Large fluctuations in the cluster multiplicity distribution have been also reported in the binary fragmentation process with the inactivation mechanism .
This article deals with the features of fluctuations of physical quantities in a $`N`$-body, $`d`$-dimensional system, with $`N`$ essentially finite. Moreover, the system is not necessarily at the thermodynamic equilibrium. Both these aspects of our approach are important in many areas of physics where, for example, small and strongly fluctuating systems are produced in the violent collision processes. Consequently, these systems live shortly and the typical time-scales are such that standard methods of equilibrium statistical physics may not to applicable.
We shall be particularly interested in self-similar systems, such as the fractal objects or thermodynamic systems at the second-order phase transition. The self-similarity means in particular that one is unable to define the characteristic length $`(N^{})^{1/d}`$, where $`N^{}`$ is the characteristic size, which could be associated with the disappearance of fluctuations. Our aim in this work is the discussion of universal scaling laws of fluctuations of different observables in self-similar systems. In particular, we shall consider the order parameter fluctuations in any system, both equilibrium and non-equilibrium one, in which the second-order critical behaviour can be identified. These considerations will provide an understanding of the relation between the order parameter, the criticality and the scaling law of fluctuations. The notion of relevant variable (relevant observable) for the discussion of critical behaviour in finite systems will appear from this discussion.
The paper is organized as follows. In Sect. II the order parameter fluctuations in statistical systems are analyzed and the relation with the finite-size scaling analysis and Widom’s hypothesis is developed in details (Sect. II.B). The generic features of the tail of the scaling function is addressed in Sect. II.C. In Sects. II.E - II.H, the generalized scaling of the observable quantities , the $`\mathrm{\Delta }`$-scaling, is discussed and the reasons for the deviations from the limiting cases $`\mathrm{\Delta }=1`$ and $`\mathrm{\Delta }=1/2`$ are presented.
Sects. III - VI are devoted to the detailed discussion of several well known generic models, using the results and methods of analysis proposed in Sect. II. In Sect. III, the non-critical fragmentation model is discussed which exhibits the power-law cluster size distribution. Results obtained in Potts model are discussed in Sect. IV. Sect. V is devoted to the discussion of the reversible aggregation, as modelled by the percolation model. Both realistic 3-dimensional bond percolation and the mean-field percolation on the Bethe lattice is being considered. Irreversible aggregation and the fluctuation properties in the Smoluchowski kinetic model are discussed in Sect. VI. Finally, the main conclusions are given in Sect. VII.
## II Order parameter fluctuations
### A Correlation function argument
Let us call $`m`$ the observable under investigation. For the reason of presentation, we shall restrict ourselves to the case where $`m`$ is a scalar quantity and takes real positive values. ( In fact, this restriction is not a true limitation and, generally, one can consider $`|m|^2`$ as well.) Fluctuations of the order parameter in thermodynamic systems are expected to have different properties at the critical point and outside of it. Far from the critical behaviour, the correlations are short-ranged. Fluctuations of the extensive order parameter $`m`$ in this case resemble the ergodic Brownian motion of this variable in its proper configuration space. Consequently $`<(m<m>)^2>/<m>`$ is roughly a constant, meaning independent of the number of constituents in the sample. On the contrary, close to the second-order transition point, the fluctuations are correlated throughout the whole system and the correlation length $`\xi `$ for the infinite system becomes infinite as well. Let us now define in such a system the deviation $`ϵ`$ of the driven parameter from its critical value, $`m`$ the local value of the order parameter and the field $`h`$ conjugated to $`m`$. The following isotropic correlation function is introduced :
$`\sigma _n(ϵ,h,\stackrel{}{r}_1,\mathrm{},\stackrel{}{r}_{n1})`$ $`=`$ $`<m(\stackrel{}{r}_o)m(\stackrel{}{r}_o\stackrel{}{r}_1)\mathrm{}`$ (1)
$`\mathrm{}`$ $`m(\stackrel{}{r}_o\stackrel{}{r}_{n1})>.`$ (2)
Notation $`<\mathrm{}>`$ in the above expression denotes the thermodynamic average at a given $`ϵ`$ and over the position $`\stackrel{}{r}_o`$. For vanishing $`ϵ`$, all the length scales disappear and the correlation length $`\xi `$ diverges algebraically as
$`\xi ϵ^\nu `$ (3)
with a universal exponent $`\nu `$, which depends only on the universality class of the transition. Moreover, the scaling description of second-order critical phenomena leads to the fundamental postulate that the thermodynamic potential $`G`$ verifies
$`G(\lambda ϵ,\lambda ^{2\alpha \beta }h)\lambda ^{2\alpha }G(ϵ,h),`$ (4)
and this defines the two universal exponents $`\alpha `$ and $`\beta `$. Let us now go back to the correlation function $`\sigma _n`$ defined in (1). The integral of $`\sigma _n`$ over the $`n1`$ space-variables : $`\stackrel{}{r}_1,\mathrm{},\stackrel{}{r}_{n1}`$, is equal to the $`n`$-th derivative of $`G`$ with respect to the field $`h`$. Hence, to be consistent with (4), we must get the scaling relation :
$`\sigma _n(\lambda ϵ,\lambda ^{2\alpha \beta }h,\lambda ^\nu \stackrel{}{r}_1,`$ $`\mathrm{}`$ $`,\lambda ^\nu \stackrel{}{r}_{n1})`$ (5)
$``$ $`\lambda ^{n\beta }\sigma _n(ϵ,h,\stackrel{}{r}_1,\mathrm{},\stackrel{}{r}_{n1}).`$ (6)
For $`n=1`$, we recover the well-known scaling behaviour of the averaged order parameter with the critical exponent $`\beta `$. For any integer $`n`$, putting all the space-variables to $`\stackrel{}{0}`$ in formula (5) leads to the scaling of powers of the local order parameter :
$`<m^n>(\lambda ϵ,\lambda ^{2\alpha \beta }h)\lambda ^{n\beta }<m^n>(ϵ,h).`$ (7)
As a consequence, if we let $`h=0`$ and $`\lambda =1/ϵ`$, one finds that the quantities $`<m^n>/<m>^n`$ are independent of $`ϵ`$, when close to the transition. This means, by the finite-size analysis, that this ratio is independent of the size $`N`$ of the system at the transition point.
Let us now introduce the cumulants $`\kappa _q`$ from the general formula for the moment expansion of the order-parameter probability distribution $`P[m]`$ :
$`\mathrm{ln}\left({\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}P[m]\mathrm{exp}(mu)\right)={\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{u^q}{q!}}\kappa _q.`$ (8)
Expanding the l.h.s. of the above expression in the power series in $`u`$, and comparing corresponding powers on l.h.s. and r.h.s., one derives the relations between ordinary moments of $`P[m]`$ and the cumulant moments :
$`\kappa _1`$ $`=`$ $`<m>`$ (9)
$`\kappa _2`$ $`=`$ $`<m^2><m>^2`$ (10)
$`\kappa _3`$ $`=`$ $`<m^3>3<m^2><m>+2<m>^3`$ (11)
$`\kappa _4`$ $`=`$ $`<m^4>4<m^3><m>3<m^2>^2+`$ (12)
$`+`$ $`12<m^2><m>^26<m>^4`$ (13)
$`\mathrm{}`$ (14)
In the case of the second-order phase transition, as a result of scaling relations (7), all cumulant moments scale like
$$\kappa _q<m>^q.$$
Consequently, the generating function of the $`m`$-probability distribution (8) is only a function of the reduced variable $`<m>u`$. This can be written as :
$`P[m]={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}\widehat{G}(iu)\mathrm{exp}(imu)𝑑u`$ (15)
where $`\widehat{G}`$ is the generating function :
$`\widehat{G}(u)={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}P[m]\mathrm{exp}(mu).`$ (16)
In the case when $`<m>`$ tends to $`\mathrm{}`$ but $`m/<m>`$ remains finite, we can rewrite formula (15) in the more compact form :
$`<m>P[m]=\mathrm{\Phi }\left({\displaystyle \frac{m}{<m>}}\right),`$ (17)
which is valid at the critical point. $`\mathrm{\Phi }`$ is the scaling function of the single reduced variable $`m/<m>`$. As stated before, we can express this scaling as :
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}P_N[m]\mathrm{exp}(mu)=\mathrm{\Psi }(<m>u),`$ (18)
which is the necessary and sufficient condition for the applicability of scaling law (17) . This result implies also that if this scaling occurs for the fluctuations of the parameter $`m`$ then it holds also for fluctuations of any power :
$$X=N^am^b$$
of this parameter as well. This is a consequence of (7) and
$$P_N[X]dX=P_N[m]dm.$$
Till now, we did not specify reasons of changing $`<m>`$. Indeed, under the condition that the scaling framework of second-order phase transition holds, the scaling relation (17) is valid independently of the explicit reasons of changing $`<m>`$, and independently of any phenomenological details. In other words, the explicit relation between the size $`N`$ of the system and $`<m>`$ need not to be known at this stage. In the following section, we shall show how to derive supplementary informations about $`\mathrm{\Phi }`$ when the system is at the pseudo-critical point.
### B Widom’s hypothesis and the finite-size scaling argument
The hypothesis of Widom states that in the thermodynamic limit of a system at thermal equilibrium, the free energy density close to the critical point scales as :
$`f(\lambda ^\beta \eta ,\lambda ϵ)\lambda ^{2\alpha }f(\eta ,ϵ),`$ (19)
where $`\alpha ,\beta `$ are the usual critical exponents, $`\eta `$ is the intensive order parameter and $`\lambda `$ is the scale parameter. Even though finite systems do not exhibit the critical behaviour, nevertheless their properties may resemble those of infinite systems if the correlation length $`\xi `$ is larger or comparable to the typical length $`L`$ of the system. In this case, one speaks about the pseudocritical point in a finite system at a distance
$`ϵcN^{1/\nu d}`$ (20)
from the true critical point . The quantity $`N`$ in (20) is the size of the $`d`$-dimensional system and $`c`$ is some dimensionless constant which can be either positive or negative. This constant $`c`$ is negative if the maximum of finite-size susceptibility or of any other divergent macroscopic quantity in the thermodynamic limit lies in the ordered phase, while $`c`$ is positive if this maximum is in the disordered phase. One can then derive the finite-size scaling of the total free energy :
$$F(\eta ,ϵ,N)=Nf(\eta ,ϵ)$$
at the pseudocritical point :
$`F_c(\eta ,N)f(\eta N^{\frac{\beta }{2\alpha }}).`$ (21)
In deriving (21) we used the hyperscaling relation :
$$2\alpha =\nu d.$$
The canonical probability density of the order parameter $`P_N[\eta ]`$ is given by :
$`P_N[\eta ]={\displaystyle \frac{1}{Z_N}}\mathrm{exp}(\beta _TF(\eta ,ϵ,N)),`$ (22)
where the coefficient $`\beta _T(1/T)`$ is independent of $`\eta `$ ($`T`$ is the temperature of the system). Using Eq.(22), one may calculate the most probable value of the order parameter, which is the solution to the equation
$`{\displaystyle \frac{P_N[\eta ]}{\eta }}=0,`$ (23)
as well as the average value of the order parameter and the partition function
$`Z_NN^{\frac{\beta }{2\alpha }}<\eta >\eta ^{}.`$ (24)
$`\eta ^{}`$ in (24) denotes the most probable value of the order parameter. The average value of the order parameter vanishes for large values of $`N`$, since both $`\beta `$ and $`2\alpha =2\beta +\gamma `$ are positive. The probability density $`P_N[\eta ]`$ obeys then the scaling law formally identical to Eq. (17) :
$`<\eta >P_N[\eta ]=\mathrm{\Phi }\left({\displaystyle \frac{\eta }{<\eta >}}\right)\mathrm{\Phi }(z)`$ (25)
where, in addition,
$`\mathrm{\Phi }(z)\mathrm{exp}\left(\beta _Tf(az,c)\right).`$ (26)
In the above formula, we have omitted the temperature-dependent multiplicative factor which can be determined by normalization of $`P_N[\eta ]`$. Coefficients $`a`$ and $`c`$ may both depend on $`\beta _T`$. We can then rewrite the scaling (25) in a standard form for the extensive order parameter $`m=N\eta `$ :
$`<m>P_N[m]=\mathrm{\Phi }(z_{(1)})`$ (27)
with the scaling variable $`z_{(1)}`$ defined by
$`z_{(1)}={\displaystyle \frac{mm^{}}{<m>}}.`$ (28)
$`m^{}`$ denotes the most probable value of the extensive order parameter. We call (27) with (28) the first-scaling law. The scaling domain is defined by this asymptotic behaviour of $`P_N[m]`$ when $`m\mathrm{}`$ and $`<m>\mathrm{}`$, but $`z_{(1)}`$ has a finite value. The normalization of the probability distribution $`P_N[\eta ]`$ and the definition of the average value of $`m`$, provide the two constraints :
$`lim{\displaystyle _{m^{}/<m>}^{\mathrm{}}}\mathrm{\Phi }(z_{(1)})𝑑z_{(1)}`$ $`=`$ $`1`$ (29)
$`lim{\displaystyle _{m^{}/<m>}^{\mathrm{}}}z_{(1)}\mathrm{\Phi }(z_{(1)})𝑑z_{(1)}`$ $`=`$ $`0`$ (31)
The first-scaling law (25) is a consequence of the self-similarity of the statistical system. The self-similarity means here that the fluctuations of the reduced order parameter $`\eta /<\eta >`$ at different scales characterized by different values of the intensive order parameter $`<\eta >`$, have identical properties. This is a qualitative explanation for this scaling.
The logarithm of scaling function (26) corresponds to the non-critical free energy density at the renormalized distance $`ϵ=c`$ from the critical point. If it happens that the order parameter is related to the number of fragments, like in the Fragmentation - Inactivation - Binary (FIB) process , then (25) can be written in an equivalent form to the Koba - Nielsen - Olesen (KNO) scaling , proposed some time ago as an ultimate symmetry of the $`S`$ \- matrix in the relativistic field theory . The multiplicity distribution of produced particles is intensely studied in the strong interaction physics where simple behaviour of much of the data on hadron-multiplicity distribution seems to point to some universality independent of the particular dynamical process.
If, instead of real positive scalar, the parameter under investigation is a vector of dimensionality $`n`$ : $`\stackrel{}{m}=[m_1,\mathrm{},m_n]`$ , then the first-scaling law (27, 28) takes the more general form :
$`<|\stackrel{}{m}|>^nP_N[\stackrel{}{m}]=\mathrm{\Phi }(\stackrel{}{z}_{(1)})`$ (32)
with
$`\stackrel{}{z}_{(1)}={\displaystyle \frac{\stackrel{}{m}\stackrel{}{m}^{}}{<|\stackrel{}{m}|>}}.`$ (33)
The scaling limit in (32) is defined by the asymptotic behaviour of $`P_N[\stackrel{}{m}]`$ when $`m_i\mathrm{}`$ ($`i=1,\mathrm{},N`$) and $`<|\stackrel{}{m}|>\mathrm{}`$, but $`z_{i(1)}`$ ($`i=1,\mathrm{},N`$) have finite values.
### C The tail of the scaling function
The scaling function $`\mathrm{\Phi }`$ introduced in the first-scaling law (27) has some typical features reminiscent of the non-Gaussian critical distribution of the order parameter. In this section, we are interested in the behaviour of scaling function for large values of the reduced parameter $`m/<m>`$, so we have to study the system subject to the small field $`h`$ conjugated to the order parameter. This breaks the symmetry of the distribution by shifting $`m`$ and $`<m>`$ towards larger values. More precisely, let us write the probability to get the value $`\eta `$ of intensive order parameter at the distance $`ϵ`$ from the critical point as :
$`P_N[\eta ,ϵ,h]=P_N[\eta ,ϵ,0]\mathrm{exp}(\beta _T\eta Nh).`$ (34)
Till now, we have studied the behaviour of $`P_N[\eta ,ϵ,0]`$ for which the first-scaling law holds when $`ϵ=0`$ (the critical point) or $`ϵ=cN^{1/\nu d}`$ (the pseudo-critical point). Substituting (25) and using (24), we obtain :
$`P_N[\eta ,0,h]=N^{\frac{\beta }{2\alpha }}\mathrm{exp}\left(\mathrm{ln}\mathrm{\Phi }\left(\eta N^{\frac{\beta }{2\alpha }}\right)+\beta _T\eta Nh\right).`$ (35)
The most probable value
$$\eta ^{}h^{1/\delta }$$
of the order parameter in the limit of small external field $`h`$, is given by the maximum of the term in the exponential. Since
$$\delta =\frac{2\alpha \beta }{\beta }$$
we get :
$`\mathrm{ln}\mathrm{\Phi }(h^{\frac{1}{\delta }}N^{\frac{1}{\delta +1}})\beta _Th^{\frac{1}{\delta }}Nh.`$ (36)
Relation (36) is valid for any value of $`N`$ if and only if
$`\mathrm{\Phi }(z)\mathrm{exp}(az^{\delta +1})\mathrm{exp}(az^{\widehat{\nu }}),`$ (37)
with the coefficient $`a`$ which depends on the temperature regularly.
One can express this relation in a different way. The anomalous dimension for extensive quantity $`m=N\eta `$ can be defined as :
$`g=\underset{N\mathrm{}}{lim}g_N=\underset{N\mathrm{}}{lim}{\displaystyle \frac{d}{d\mathrm{ln}N}}\left(\mathrm{ln}<m>\right).`$ (38)
One can see that due to both (24) and the Rushbrooke relation between critical exponents :
$`\alpha +2\beta +\gamma =2,`$ (39)
the anomalous dimension is :
$`g=1{\displaystyle \frac{\beta }{\gamma +2\beta }}.`$ (40)
Since both $`\alpha `$ and $`\beta `$ are positive, therefore $`g`$ is contained between 1/2 and 1 for equilibrium systems at the critical point of the second-order phase transition . Because of these additional relations between critical exponents, one may note that a behaviour of the tail of the scaling function (37) is governed by the exponent :
$`\widehat{\nu }={\displaystyle \frac{1}{1g}}=\delta +1={\displaystyle \frac{2\alpha }{\beta }},`$ (41)
which is always larger than 2. The limiting case : $`\widehat{\nu }=2`$ , i.e. the Gaussian tail (see (37)), corresponds in this framework to the non-critical system.
Finally, let us mention in passing that whenever cluster-size can be defined in a system exhibiting the second-order phase transition , like e.g. in the case of percolation, Ising model or Fisher droplet model, the exponent $`\tau `$ of the power-law cluster-size distribution
$`n_ss^\tau `$ (42)
satisfies additional relations :
$`\gamma +\beta `$ $`=`$ $`{\displaystyle \frac{1}{\sigma }}`$ (43)
$`\gamma +2\beta `$ $`=`$ $`{\displaystyle \frac{\tau 1}{\sigma }}`$ (45)
what yields :
$`g={\displaystyle \frac{1}{\tau 1}}`$ (46)
and
$`\widehat{\nu }={\displaystyle \frac{\tau 1}{\tau 2}}`$ (47)
Since $`g`$ at the second-order equilibrium phase transition is contained between 1/2 and 1, therefore the allowed values of exponent $`\tau `$ at the critical point are : $`2<\tau <3`$, and the normalized cluster-size distribution :
$$\underset{s=1}{\overset{N}{}}sn_s=N$$
is
$$n_s=Ns^\tau .$$
Consequently, whenever $`\tau `$ is defined, we get an information whether the studied equilibrium system is at the second-order phase transition and whether the considered extensive quantity can be identified with the order parameter of this transition.
Let us define, for example, the multiplicity as the total number of clusters. (The definition of clusters includes also the monomers.) The cluster multiplicity cannot be an order parameter of these kind of equilibrium phase transitions because with $`2<\tau <3`$ :
$`{\displaystyle \underset{s=1}{\overset{N}{}}}n_sN,`$ (48)
what means that the average multiplicity scales as the total mass of the system at the transition point, i.e. $`g=1`$. On the other hand, the size of the largest cluster is a natural order parameter for these kind of phase transitions. In this case we have :
$`<s_{max}>N^{\frac{1}{\tau 1}},`$ (49)
what is a direct consequence of :
$$\underset{s=<s_{max}>}{\overset{N}{}}n_s1,$$
i.e. that there is in the average only one largest cluster. Moreover, the relation (46) derived for the second-order critical phenomenon is correctly recovered. One should emphasize, that the relation (49) is very general and its derivation does not depend on the assumption of thermodynamic equilibrium. In other words, the relation (49) between the anomalous dimension and the exponent $`\tau `$ , is valid also for the off-equilibrium second-order phase transitions. We shall return to this point in Sect. VI.
The cluster multiplicity could be the order parameter whenever $`\tau <2`$, though this cannot happen in the equilibrium phase transitions. Note that this argument to sort among different candidates for the order parameter , requires only the knowledge of $`\tau `$, i.e. the complete information about the critical process is superfluous. We shall use this argument later in the case of percolation model and Smoluchowski model of gelation. Finally, we shall see below in the Mekjian model that we may have power-law size distribution with $`\tau <2`$, in the absence of phase transition governed by the multiplicity as the order parameter.
### D Landau - Ginzburg theory of phase transitions
Let us consider the Landau - Ginzburg (LG) theory as an exactly solvable example of the second-order phase transition. The homogeneous LG free energy density is :
$`f(\eta )=ϵ\eta ^2+b\eta ^4+\mathrm{}`$ (50)
where $`b`$ is a positive constant. The most probable value of the order parameter $`\eta `$ in the disordered phase ($`ϵ<0`$) is implicitly set to 0. It is more convenient to work with the extensive order parameter $`m=N\eta `$ when dealing with the finite systems. The probability of a state $`m`$ for a given $`ϵ`$ is :
$`P_N[m]={\displaystyle \frac{1}{Z_N}}\mathrm{exp}\left[\beta _T(ϵ{\displaystyle \frac{m^2}{N}}+b{\displaystyle \frac{m^4}{N^3}}\mathrm{})\right].`$ (51)
$`Z_N`$ is defined by the normalization of $`P_N[m]`$. To remain consistent with other sections of this paper and without loss of generality, we consider now the case where $`m`$ is positive. We will admit that $`N`$ is so large that the first two terms in the free energy expansion are sufficient to study the phase transition. At the critical point : $`ϵ=0`$, the leading term of the free energy density is proportional to $`m^4`$. Standard integrations yield the values for the partition function $`Z_N`$ and the average value of the order parameter $`<m>`$, both proportional to $`N^{3/4}`$. Introducing them in (51), one finds :
$`<m>P_N[m]`$ $`=`$ $`{\displaystyle \frac{4\sqrt{\pi }}{\mathrm{\Gamma }^2[1/4]}}\times `$ (52)
$`\times `$ $`\mathrm{exp}\left({\displaystyle \frac{\pi ^2}{\mathrm{\Gamma }^4[1/4]}}\left({\displaystyle \frac{m}{<m>}}\right)^4\right),`$ (53)
which has the form of (17). Note that the scaling function : $`\mathrm{\Phi }(z)\mathrm{exp}(z^4)`$, decreases very fast as one moves away from the most probable value. This result is consistent with the analysis done in the previous section.
The pseudo-critical point is the value of $`ϵ`$ for which the finite-size thermal susceptibility reaches its maximum. Writing that the inverse of this susceptibility is the second derivative of the free energy with respect to the order parameter, one finds :
$`ϵ=6{\displaystyle \frac{\mathrm{\Gamma }[3/4]}{\mathrm{\Gamma }[1/4]}}\left({\displaystyle \frac{b}{\beta _TN}}\right)^{1/2}.`$ (54)
This results is correct at the first order in $`N^{1/2}`$. Replacing $`ϵ`$ in (51) by (54), leads to the scaling form of $`P_N[m]`$ :
$`<m>P_N[m]`$ $`=`$ $`A\mathrm{exp}[{\displaystyle \frac{\mathrm{\Gamma }[3/4]^2}{\mathrm{\Gamma }[1/4]^2}}((m/<m>)^4`$ (55)
$``$ $`6(m/<m>)^2)],`$ (56)
where $`A`$ denotes a normalization constant. We recover indeed the first-scaling law with the exponential tail : $`\mathrm{exp}(az^4)`$, for the large arguments.
Outside of the critical point in the disordered phase ($`ϵ>0`$), the leading term of the free energy is proportional to $`m^2`$, and the probability distribution $`P_N[m]`$ is essentially Gaussian. Deriving, as previously, the values of $`Z_N`$ and $`<m>`$ (both behave like $`N^{1/2}`$ in this case), we get the scaling form :
$`<m>P_N[m]={\displaystyle \frac{4}{\pi }}\mathrm{exp}\left({\displaystyle \frac{4}{\pi }}\left({\displaystyle \frac{m}{<m>}}\right)^2\right)`$ (57)
which is still under the form (25) but with a Gaussian scaling function reminiscent of the Gaussian fluctuations.
Finally, in the low temperature regime ($`ϵ<0`$), the most probable value of the order parameter is positive :
$`m^{}=\sqrt{{\displaystyle \frac{ϵ}{2b}}}N.`$ (58)
Developing $`P_N[m]`$ in (51) around this point leads to the expression :
$`m^{1/2}P_N[m]`$ $``$ $`(2{\displaystyle \frac{ϵ^3}{b\pi ^2}})^{1/4}\times `$ (59)
$`\times `$ $`\mathrm{exp}\left(ϵ\sqrt{2{\displaystyle \frac{ϵ}{b}}}{\displaystyle \frac{(mm^{})^2}{m^{}}}\right),`$ (60)
which is no more in the standard form (25). In this case, the average value of the order parameter $`<m>`$ is of the same order of magnitude as its most probable value $`m^{}`$ and one can rewrite (59) in the scaling form :
$`<m>^{1/2}P_N[m]\mathrm{exp}\left(a{\displaystyle \frac{(mm^{})^2}{<m>}}\right),`$ (61)
where $`a`$ is a positive constant. This particular scaling form will be discussed later in details.
### E The $`\mathrm{\Delta }`$-scaling law
One may ask, what happens if the observable quantity is not the order parameter but the $`N`$-dependent function of the order parameter like :
$`m=N^{a_1}N^{a_2}\eta ,`$ (62)
where
$`a_1>g+a_21.`$ (63)
The latter condition ensures that the order parameter does not determine the leading behaviour of $`m`$. For large $`N`$ :
$$<m>N^{a_1}.$$
Writing (27) with $`m`$ instead of $`\eta `$ and taking into account that :
$$P_N[\eta ]d\eta =P_N[m]dm$$
one finds the generalized law :
$`<m>^\mathrm{\Delta }P_N[m]`$ $`=`$ $`\mathrm{\Phi }(z_{(\mathrm{\Delta })})\mathrm{\Phi }\left({\displaystyle \frac{mm^{}}{<m>^\mathrm{\Delta }}}\right),`$ (64)
where :
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{g+a_21}{a_1}}<1.`$ (65)
This generalized law will be called in the following the $`\mathrm{\Delta }`$-scaling law. The scaling function $`\mathrm{\Phi }(z_{(\mathrm{\Delta })})`$ depends only on one scaled variable :
$`z_{(\mathrm{\Delta })}={\displaystyle \frac{mm^{}}{<m>^\mathrm{\Delta }}}.`$ (66)
The normalization of the probability distribution $`P_N[m]`$ and the definition of the average value of $`m`$ provide two constraints :
$`lim{\displaystyle _{<m>^{1\mathrm{\Delta }}}^{\mathrm{}}}\mathrm{\Phi }(z_\mathrm{\Delta })𝑑z_\mathrm{\Delta }=1`$ (67)
(68)
$`lim{\displaystyle _{<m>^{1\mathrm{\Delta }}}^{\mathrm{}}}z_\mathrm{\Delta }\mathrm{\Phi }(z_\mathrm{\Delta })𝑑z_\mathrm{\Delta }=0`$ (69)
which are consistent with : $`\mathrm{\Delta }1`$, because the scaling function $`\mathrm{\Phi }`$ is positive. The scaling function $`\mathrm{\Phi }(z_{(\mathrm{\Delta })})`$ in (64) has an identical form as $`\mathrm{\Phi }(z_{(1)})`$, except for the inversion of the abscissa axis. In particular, its tail for large $`z_{(\mathrm{\Delta })}`$ has the same form :
$`\mathrm{\Phi }(z_{(\mathrm{\Delta })})\mathrm{exp}\left(z_{(\mathrm{\Delta })}^{\widehat{\nu }}\right)=\mathrm{exp}\left(z_{(\mathrm{\Delta })}^{\frac{1}{1g}}\right)`$ (70)
as given in (37). One should mention in passing that if
$$a_1<g+a_21$$
in (62), then
$$<m>N^{a_2}<\eta >$$
and $`\mathrm{\Delta }=1`$, following the remark of Sect. II.A.
According to (26), the logarithm of scaling function $`\mathrm{\Phi }(z_{(\mathrm{\Delta })})`$ :
$`\mathrm{ln}\mathrm{\Phi }(z_{(\mathrm{\Delta })})=\beta _Tf(az_{(\mathrm{\Delta })},c)`$ (71)
is related to the non-critical free energy $`f`$, in either ordered ($`c>0`$) or disordered $`(c<0)`$ phase.
As an important example, we see from (62), (64), that the $`\mathrm{\Delta }`$-scaling of the extensive variable :
$`\widehat{m}=N(1\eta )N\widehat{\eta }`$ (72)
can be used to determine the anomalous dimension since in this case : $`\mathrm{\Delta }=g`$. For this reason , $`\widehat{m}`$ is a very useful variable in all phenomenological studies. At the phase transition :
$$<N\widehat{\eta }>N$$
but the finite-size corrections are algebraic.
### F Off-critical scaling
$`\mathrm{\Delta }=1/2`$ with $`\mathrm{\Phi }(z_{(\mathrm{\Delta })})`$ nearly Gaussian, is a particular case of a $`\mathrm{\Delta }`$-scaling associated with the non-critical systems . This limit :
$`<m>^{1/2}P_N[m]`$ $`=`$ $`\mathrm{\Phi }\left({\displaystyle \frac{mm^{}}{<m>^{1/2}}}\right)`$ (73)
$``$ $`\mathrm{\Phi }(z_{(1/2)}),`$ (74)
which is called the second-scaling law, has been found in the shattering phase of the non-equilibrium FIB process and in the ’gaseous’ phase of equilibrium percolation process . We should recall also that this form of scaling function has been seen for LG model in the low temperature regime (see (61)).
More generally, let us suppose now that the extensive parameter $`m`$ is not critical, i.e. either the system is in a critical state but the parameter $`m`$ is not critical, or the system is outside of the critical region. The value of $`m`$ at the equilibrium is obtained by minimizing the free energy. The free energy $`F`$ is analytical in the variable $`m`$ close to its most probable value $`m^{}`$ :
$`FN^1(mm^{})^2.`$ (75)
Using (75) one obtains
$$<m>\mu ^{}N,$$
where $`\mu ^{}`$ is a positive (finite) number independent of $`N`$ and :
$`Z_NN^{1/2}<m>^{1/2}.`$ (76)
The probability density $`P_N[m]`$ verifies the second-scaling law (64) :
$`<m>^{1/2}P_N[m]`$ $`=`$ $`\mathrm{exp}\left(\beta _T\mu ^{}\left({\displaystyle \frac{m<m>}{<m>^{1/2}}}\right)^2\right)`$ (77)
$``$ $`\mathrm{\Phi }(z_{(1/2)}).`$ (78)
This is a particular case of the $`\mathrm{\Delta }`$-scaling law ($`\mathrm{\Delta }=1/2`$) and the scaling function is now Gaussian . This scaling (73) holds for $`<m>N`$ but now with the exponential finite-size corrections. This is a principle difference from the finite-size corrections and/or $`\mathrm{\Delta }`$-scaling. The above arguments apply to any second order phase transition. In particular, they are not limited to LG theory of phase transitions (see Eq. (61)).
### G Finite-size cross-over effects
The discussion of previous section is valid for systems at the critical (and pseudo-critical) point or far from the critical point in the ordered phase. Let us suppose now that the system is prepared such that :
$`<m>N^g^{^{}},g^{^{}}<1`$ (79)
and $`g^{^{}}`$ is not the anomalous exponent. Here, we would like to study how the finite system evolves when the control parameter $`ϵ`$ tends slowly to 0, namely :
$`ϵN^{2g^{^{}}2}.`$ (80)
We shall address this question in the mean-field approximation using LG theory. Let us first write down the average value :
$`<m>={\displaystyle \frac{_0^{\mathrm{}}m\mathrm{exp}(ϵm^2/Nbm^4/N^3)𝑑m}{_0^{\mathrm{}}\mathrm{exp}(ϵm^2/Nbm^4/N^3)𝑑m}}.`$ (81)
Hence, writing this definition with the new driving parameter $`ϵ^{}=ϵN^{1/2}`$ and using the rescaled variable $`m^{}=m/N^{3/4}`$, the average value of $`m`$ can be put in the form :
$`<m>=N^{3/4}\psi (ϵN^{1/2}),`$ (82)
while its most probable value is :
$$m^{}=\sqrt{\frac{ϵ}{2b}}N.$$
If the exponent $`g^{^{}}`$ is not too small, i.e. if $`ϵ`$ does not vanish too fast, the two quantities : $`<m>`$ and $`m^{}`$, have to coincide. This is because the exponential weight term in (81) diverges as $`\mathrm{exp}(ϵ^2N/4b)`$ when
$`ϵN^{1/2}N^{\frac{4g^{^{}}3}{2}}`$ (83)
becomes large with increasing $`N`$. As a consequence, the common behaviour of $`<m>`$ and $`m^{}`$ is : $`N^g^{^{}}`$. The scaling form (59) in this case is :
$`<m>^\mathrm{\Delta }P_N[m]\mathrm{exp}\left(c{\displaystyle \frac{(mm^{})^2}{<m>^{2\mathrm{\Delta }}}}\right)`$ (84)
with $`c`$ a positive constant, and :
$$\mathrm{\Delta }=\frac{3}{2}g^{^{}}1.$$
We recover here the two cases previously discussed in Sect. II.B. When $`g^{^{}}=1`$, i.e. when $`ϵ=\text{const}`$, then this is the second-scaling law. When $`g^{^{}}=3/4`$, then this is the first-scaling law since the finite system is yet in the critical region ($`g^{^{}}=g`$). In between these two limiting cases, $`\mathrm{\Delta }`$-scaling holds with $`1/2<\mathrm{\Delta }<1`$. Note also that the scaling function in (84) has a Gaussian form, even for $`\mathrm{\Delta }>1/2`$, which is quite different from the case (37) of Sect. II.C.
### H Summary : panorama of the $`\mathrm{\Delta }`$-scaling for thermodynamic systems
Several features of finite systems are important if one wants to study either the criticality of the corresponding infinite system or the distance to the critical point. One should name here : the $`\mathrm{\Delta }`$-scaling (this includes the first-scaling law $`\mathrm{\Delta }=1`$ as well), the form of the tail of the scaling function $`\mathrm{\Phi }`$, and the anomalous exponent. All these features are closely related with the properties of the scaling function which characterizes finite system at the equilibrium. If the infinite system experiences a second-order phase transition, and if $`m`$ is the scalar order parameter or the shifted scalar order parameter (62), then :
At the critical point, the corresponding finite system exhibits first-scaling law if $`m`$ is the order parameter, or $`\mathrm{\Delta }`$-scaling law if $`m`$ is the shifted order parameter. In both cases, the tail of the scaling function $`\mathrm{exp}(z^{\widehat{\nu }})`$ is characterized by a large value of the exponent $`\widehat{\nu }=1/(1g)>2`$, with $`g`$ being the anomalous exponent, i.e. the exponent characterizing decrease of the extensive order parameter with the size $`N`$ of the finite system. The values of $`\mathrm{\Delta }`$ are restricted to $`0<\mathrm{\Delta }1`$, and the anomalous exponent $`g`$ takes values in between 1/2 and 1 for second-order at-equilibrium phase transition.
Far from the critical point, finite system exhibits second-scaling law with the Gaussian tail of the scaling function.
Close to the critical point when $`ϵ0`$ if $`N\mathrm{}`$, the finite system exhibits the cross-over phenomenon from the first-scaling to the second-scaling law by a continuous $`\mathrm{\Delta }`$-scaling law with Gaussian shape of the scaling function. One should remind here that the precise dependence of $`ϵ(N)`$ is irrelevant provided that : $`g<g^{^{}}<1`$, i.e. that the conditional point approaches 0 not faster than the pseudocritical point. This last remark is important in the phenomenological applications of the scaling theory to the situations where the $`N`$ \- dependence of the conditional point is governed by an external control parameter with unknown relation to the system size.
Last but not least, if the parameter $`m`$ is not singular at the transition, then all properties of its probability distribution are the same as in the case of non-critical systems.
In phenomenological applications, it is often difficult to get the probability distribution with sufficient accuracy for the values of scaling variable which are far from the most probable value since this corresponds to very small probabilities. It is then more judicious to work with moments of the distribution instead of the distribution itself. For example, when the system undergoes $`\mathrm{\Delta }`$-scaling, the properly normalized cumulant moments (9) :
$`{\displaystyle \frac{\kappa _q}{(\kappa _1)^{q\mathrm{\Delta }}}}\text{const},`$ (85)
are independent of the size of the system. An important consequence is that the generating function of the $`m`$-distribution : $`\widehat{G}(u)=P[m]\mathrm{exp}(mu)`$, is a function of the reduced variable $`<m>^\mathrm{\Delta }u`$ only, generalizing a remark of Sect. II.A for the generating function in the first-scaling case.
## III A non-critical model : the Mekjian model
The Mekjian fragmentation model is an equilibrium model which describes the decomposition of system into an ensemble of fragments. The statistical weights for every configuration of fragments are given explicitly in this model. If $`n_s`$ denotes the number of fragments of size $`s`$ with the size-conservation : $`N=_ssn_s`$, the weight function for the configuration $`\{n_s\}`$ is given by :
$`W_N(\{n_s\},x)={\displaystyle \underset{s=1}{\overset{N}{}}}{\displaystyle \frac{sx^{n_s}}{n_s!s^{n_s}(x+s1)}},`$ (86)
with $`x`$ being a real control parameter. Many exact results can be obtained in this simple model. Here, we are interested in the multiplicity distribution $`P_N[m]`$, where the fragment multiplicity is : $`m=_sn_s`$. We can show that :
$`P_N[m]=x^m|S_N^{(m)}|{\displaystyle \frac{\mathrm{\Gamma }(x)}{\mathrm{\Gamma }(N+x)}},`$ (87)
where $`|S_N^{(m)}|`$ are the signless Stirling numbers of the first kind. Knowing then the generating function for these Stirling numbers :
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}P_N[m]e^{mu}={\displaystyle \frac{\mathrm{\Gamma }(x)\mathrm{\Gamma }(xe^u+N)}{\mathrm{\Gamma }(xe^u)\mathrm{\Gamma }(x+N)}}`$ (88)
one obtains the average value of $`m`$ :
$`<m>`$ $`=`$ $`x{\displaystyle \underset{s=1}{\overset{N}{}}}{\displaystyle \frac{1}{x+s}}`$ (89)
$`=`$ $`x\mathrm{ln}N+(x1)\gamma \psi (x)+𝒪(1/N).`$ (90)
Moreover, making an asymptotic development of (88) for large $`N`$ and small $`s`$, one obtains :
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}P_N[m]e^{mu}N^{x(e^u1)}.`$ (91)
The latter approximation is known to be correct for finite values of $`u`$ . This means that $`P_N[m]`$ is approximately a Poisson $`m`$-distribution with parameter $`x\mathrm{ln}N`$. In the leading order we have then : $`<m>m^{}`$ . Inverting (88) to get $`P_N[m]`$ as a Fourier transform, and making $`N`$ large yields the scaling formula :
$`<m>^{1/2}P_N[m]`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi }}}\mathrm{exp}({\displaystyle \frac{(m<m>)^2}{2<m>}}`$ (92)
$``$ $`(x1)(\gamma \psi (x)){\displaystyle \frac{m<m>}{<m>}}+`$ (93)
$`+`$ $`𝒪\left({\displaystyle \frac{1}{<m>}}\right)).`$ (94)
This is nothing else but the second-scaling law (73) for the multiplicity distribution when $`N`$ becomes large enough, because $`<m>m^{}`$ . When $`<m>`$ is large enough, the second term in (92) is always very small compared to the first one for a finite $`x`$.
Different fixed values of the control parameter $`x`$ mimic different situations of the fragmentation. For $`x1`$, one has the situation of a fused system. For $`x0.5`$, the fragmentation resembles the evaporation of light fragments. The limit $`x1`$ corresponds to the complete dissociation of the mass into light fragments (monomers). Each of this situation is characterized by a different fragment-size distribution. The case $`x=1`$ is particular in this model since it leads to the power-law size-distribution (42) with the exponent $`\tau =1`$. Following discussion in Sect. II.C, cluster multiplicity could be the order parameter. On the other hand, the second-order equilibrium phase transition is associated with $`2<\tau <3`$, what implies that the equilibrium model of Mekjian is a non-critical model. Indeed, that is what can be seen also in the cluster-multiplicity scaling law (92). Hence, the power-law cluster-size distribution alone does not guarantee that the system exhibits the critical behaviour of any kind .
## IV Example : The Potts model
A generalization of the magnetic Ising spin model has been proposed by Domb and studied in details by Potts . In this model, one considers a system of $`N`$ sites in the $`d`$-dimensional space. The magnetic state of each site $`i`$ is characterized by a quantity called : a spin (say : $`s_i`$). Each spin is of the same constant modulus and points to one of the $`q`$ equally spaced directions, labelled from 0 to $`q1`$. The ferromagnetic short-ranged Potts Hamiltonian is then :
$`H_q=J{\displaystyle \underset{i,j}{}}\delta (s_i,s_j),`$ (95)
where $`\delta `$ is the Kronecker symbol, and $`J`$ is the positive coupling constant. The sum is restricted to nearest-neighbour pairs. The site percolation corresponds to the $`q=1`$ Potts model, and ferromagnetic Ising model to the $`q=2`$ case. This model is one of the simplest non-trivial critical thermodynamic N-body system, and many exact or accurate results are known for standard values of $`(d,q)`$. In particular, there exists a value $`q_c(d)`$ (for example : $`q_c(2)=4`$), for which when $`qq_c(d)`$, then such an interacting system experiences a second-order phase transition at a finite critical temperature (for example : $`\beta _cJ=\mathrm{ln}(1+\sqrt{q})`$ at $`d=2`$), while for $`q>q_c(d)`$ the transition is a first-order one.
Let us consider here the case of the second order phase transition. All the scalings described above should hold. We have first to define the order parameter for the system. If for a given configuration of the system, we call $`N_k`$ the number of sites in the state $`k`$, where $`k`$ varies from 0 to $`q1`$, then the order parameter $`m`$ is given by :
$`m={\displaystyle \frac{q(N_{max}/N)1}{q1}}.`$ (96)
$`N_{max}`$ in (96) is defined as the maximum of all $`N_k`$’s. Fig. 1 shows the $`m`$-distribution at the critical temperature in the $`(d,q)=(2,3)`$ case, in the first-scaling form. The scaling is recovered very precisely, even for such small system sizes as $`64\times 64`$. Note also the complicated shape of the scaling curve.
We can discuss this scaling here in the slightly different context of correlated variables . Consider for simplicity, the Ising model (i.e. $`(d,q)=(d,2)`$ case). The extensive order parameter is just the sum of $`N`$ correlated variables : $`M=s_i`$. When the system is disordered, the spins are correlated on the short distance $`\xi `$ ($`\xi /N0`$ at the thermodynamic limit) and their mean value is zero. The central limit theorem tells then that the distribution of the random variable $`M/\sqrt{N}`$ is Gaussian when $`N`$ becomes large, with zero mean and finite variance. This can also be expressed by the asymptotic law :
$`<M^2>P[M^2]=\sqrt{{\displaystyle \frac{<M^2>}{2\pi M^2}}}\mathrm{exp}\left({\displaystyle \frac{M^2}{2<M^2>}}\right)`$ (97)
which is correctly in the first-scaling form, with a Gaussian shape and the trivial anomalous exponent $`g=1/2`$ (see (57)).
On the other hand, if the system is in the ordered phase, the average value of the individual spins is finite, say : $`<s_i>=m`$, and the same reasoning can be used for the variable : $`(MNm)/\sqrt{N}`$. This variable is of zero mean , finite variance and short-ranged correlated. So, its fluctuations are Gaussian and can be put in the second-scaling form :
$`<M>^{1/2}P[M]={\displaystyle \frac{1}{\sqrt{2\pi }}}\mathrm{exp}\left({\displaystyle \frac{(M<M>)^2}{2<M>}}\right).`$ (98)
Of course, the most interesting case corresponds to the critical temperature. At this point, the spins are correlated throughout the system, and the magnetization cannot be evaluated by the central limit theorem. Instead, we can remark that the spin-correlations are power-law :
$$<s_{\stackrel{}{r}_o}s_{\stackrel{}{r}+\stackrel{}{r}_o}>\frac{1}{r^{d2+\eta }},$$
with $`\eta `$ a critical exponent whose value should be between $`2d`$ and 2. Looking at the total magnetization as the sum of $`N`$ correlated variables, one gets
$`<M^2>`$ $`=`$ $`{\displaystyle \underset{i,j}{}}<s_is_j>={\displaystyle \underset{i}{}}<s_i^2>+`$ (99)
$`+`$ $`N{\displaystyle _1^L}<s_\stackrel{}{0}s_\stackrel{}{r}>r^{d1}𝑑r`$ (100)
with $`LN^{1/d}`$ the typical macroscopic length of the system. This means that :
$$<M^2>N^{1+\frac{2\eta }{d}}$$
as the leading behaviour. This non-trivial anomalous exponent
$$g=\frac{1}{2}+\frac{2\eta }{2d}$$
between $`1/2`$ and 1, is here the sign of the criticality. The first-scaling law should hold in this case, as for the $`(d,q)=(2,3)`$ Potts model discussed above, but the scaling function should be different since it depends on the precise form of the interactions. Only the tail can be linked to another critical exponent, as it has been written in Sect. II.C.
## V Reversible aggregation - the percolation model
Percolation model can be defined as follows. In a box (a part of the regular lattice), each site corresponds to a monomer and a proportion $`p`$ of active bonds is set randomly between sites (the bond percolation model). Such a network results in a distribution of clusters defined as ensemble of occupied sites connected by active bonds. For a definite value of $`p`$, say $`p_{cr}`$, a giant cluster almost surely spans the whole box. The sol-gel transition corresponds to the appearance of ’infinite’ cluster (gel) at a finite time. ’Infinite’ in this context means that gel contains the finite fraction of total mass of the system. The sol - gel transition in finite systems can be suitably studied using moments of the number-size-distribution $`n_s`$, i.e. the number of finite clusters of size $`s`$ :
$`M_k^{}={\displaystyle \underset{s<s_{max}}{}}s^kn_s,`$ (101)
where the summation is performed over all clusters with the exception of the largest cluster $`ss_{max}`$. The superscript recalls this constraint on summation in Eq. (101). The mass of the largest cluster is then : $`NM_1^{}`$ , with :
$`N={\displaystyle \underset{\text{a}lls}{}}sn_s.`$ (102)
In infinite systems, one works with the normalized moments of the concentration-size-distribution $`c_s`$, i.e. the concentration of clusters of size $`s`$ :
$`m_k^{}={\displaystyle \underset{s}{}}s^kc_s,`$ (103)
where the summation in (103) runs over all finite clusters. Generally, concentrations are normalized such that :
$`c_s=\underset{N\mathrm{}}{lim}{\displaystyle \frac{n_s}{N}}.`$ (104)
The probability that a monomer belongs to the infinite cluster (gel) is equal to $`1m_1^{}`$ , with :
$`m_k^{}=\underset{N\mathrm{}}{lim}{\displaystyle \frac{M_k^{}}{N}}.`$ (105)
For example, in the thermodynamic limit when the size of the box becomes infinite , a finite fraction of the total number of vertices belongs to this cluster. Therefore, we get the results : $`m_1^{}=1`$ for $`p<p_{cr}`$ and $`m_1^{}<1`$ for $`p>p_{cr}`$. Moreover, $`m_1^{}`$ is a decreasing function of the occupation probability. This typical behaviour is commonly (and incorrectly) called : ’the failure of mass conservation’, but, as stated before, $`m_1^{}`$ is more simply the probability for a vertex to belong to some finite cluster.
### A Percolation on the Bethe lattice
The bond-percolation on the Bethe lattice with coordination number $`\widehat{z}`$, has been solved by Fisher and Essam . Here the main result we are interested in, is the concentration-size-distribution :
$`c_s=\widehat{z}{\displaystyle \frac{((\widehat{z}1)s)!}{((\widehat{z}2)s+2)!s!}}p^{s1}(1p)^{(\widehat{z}2)s+\widehat{z}}`$ (106)
and the first normalized moment :
$`m_1^{}=\left({\displaystyle \frac{1p}{1p^{}}}\right)^{2\widehat{z}2},`$ (107)
with $`p^{}`$ being the smallest solution of equation :
$`p^{}(1p^{})^{\widehat{z}2}=p(1p)^{\widehat{z}2}.`$ (108)
Let us define :
$$p_{cr}\frac{1}{\widehat{z}1}.$$
For $`p<p_{cr}`$, the only solution of the above equation is : $`p^{}=p`$, but when $`p`$ is larger than $`p_{cr}`$, then there is a smaller non-trivial solution which behaves as $`p_{cr}|pp_{cr}|`$ near $`p_{cr}`$. Above this threshold, the moment $`m_1^{}`$ is smaller than 1 and behaves approximately as :
$$m_1^{^{}}1\frac{2(pp_{cr})}{1p_{cr}}.$$
The marginal case $`\widehat{z}=2`$ corresponds to the linear-chain case.
Coming back to the concentrations, we can see that for large values of the size $`s`$, the following Stirling approximation holds :
$`c_ss^{5/2}\mathrm{exp}(\alpha s)`$ (109)
with $`\alpha `$ given by :
$`\alpha =\mathrm{ln}\left[{\displaystyle \frac{p}{p_{cr}}}\left({\displaystyle \frac{1p}{1p_{cr}}}\right)^{\widehat{z}2}\right].`$ (110)
For this model, a power-law behaviour of the concentrations is seen at the threshold $`p_{cr}`$, namely : $`c_ss^\tau `$ with $`\tau =5/2`$. Outside this threshold, an exponential cut-off is always present . This sort of critical behaviour at equilibrium is analogous to the thermal critical phenomena, and in particular, there exist two independent critical exponents, for example $`\tau `$ and $`\sigma `$. The latter one is the exponent of the mean cluster-size divergence. Together, those two critical exponents : $`\tau =5/2`$ and $`\sigma =1`$, describe completely the critical features.
This singular behaviour is due to the appearance of a giant cluster, the so-called percolation cluster, at the transition point. More precisely, in the infinite system the probability for a given site to belong to this infinite cluster is zero below the critical threshold $`p_{cr}`$ and positive above it. This probability is non-analytical at the critical point. Because of this behaviour, the extensive order parameter defined for finite systems is just the size of the largest cluster $`s_{max}`$.
As discussed in Sect. II.C, the corresponding finite-size order parameter scales as
$$s_{max}N^{2/3}.$$
Even though the system experiences the second-order critical phenomenon, fluctuations of the multiplicity distribution remain small and the KNO-scaling does not hold. Of course, $`m_0^{}`$ is not in this case an order parameter since $`\tau >2`$ even though there is some irregularity in its behaviour passing the threshold. This non-analyticity can be illustrated by the exact result for bond-percolation on the Bethe lattice. In this mean-field case, the normalized $`0^{th}`$-moment is :
$`m_0^{}`$ $`=`$ $`(1{\displaystyle \frac{\widehat{z}}{2}}p^{})\left({\displaystyle \frac{1p}{1p^{}}}\right)^{2\widehat{z}2}`$ (111)
$``$ $`{\displaystyle \frac{\widehat{z}2}{2(\widehat{z}1)}}(\widehat{z}1)ϵ+(1{\displaystyle \frac{\widehat{z}}{2}})|ϵ|`$ (112)
with : $`ϵ=pp_{cr}`$, and $`ϵ1`$. There is a jump of the first derivative of $`m_0^{}`$ with respect to $`p`$ : $`\widehat{z}/2`$ for $`pp_{cr}^{}`$, and $`(43\widehat{z})/2`$ for $`pp_{cr}^+`$.
### B 3-dimensional percolation
As shown by Botet et al , the multiplicity distribution for the 3-dimensional bond percolation model on the cubic lattice at the infinite-network percolation threshold exhibits the $`\mathrm{\Delta }`$-scaling with $`\mathrm{\Delta }=1/2`$ , and hence the fragment multiplicity is not related to the order parameter in this process. This is shown in Fig. 2c as a typical example of non-critical parameter scaling. Note that the multiplicity distributions in Fig. 2c are plotted in semi-logarithmic form to show clearly the Gaussian behaviour (parabolic shape on the figure). The proper order parameter for this model is the normalized mass of the gel-phase, i.e. the mass of the largest cluster divided by the total mass of the system $`s_{max}/N`$ . Different probability distributions $`P_N[s_{max}/N]`$ for different system sizes $`N`$ can be all compressed into a unique characteristic function (see Fig. 2.a) :
$`<{\displaystyle \frac{s_{max}}{N}}>P_N\left[{\displaystyle \frac{s_{max}}{N}}\right]=\mathrm{\Phi }\left({\displaystyle \frac{s_{max}<s_{max}>}{<s_{max}>}}\right)`$ (113)
which is analogous to the KNO-scaling function. As an application of results developed in Sect. II.E, Fig.2b shows the $`\mathrm{\Delta }`$-scaling for the shifted order parameter : $`M_1^{}=Ns_{max}`$. The value of $`\mathrm{\Delta }`$ (= 0.8), is consistent with the value of the anomalous dimension (25) : $`g=0.8435`$, for the accepted values of the critical exponents $`\beta ,\gamma `$ in the 3-dimensional percolation . One should also remember, that $`\mathrm{\Delta }`$ has been extracted from the small size ($`N=14^3,20^3,32^3`$) percolation network calculations at the infinite-network percolation threshold. This explains a small difference between the value for $`\mathrm{\Delta }`$ from the scaling analysis and the expected value : $`\mathrm{\Delta }=g`$, in the infinite network.
According to the results derived above for the second order phase transition, the second-scaling should hold outside of the critical point. This is correctly realized with the three variables $`s_{max}`$, $`M_1^{}`$ and $`M_0`$ for large or small values of the probability $`p`$. Fig. 3 shows such results for the value $`p=0.35`$.
Finally, it is instructive to see how the first-scaling is disappearing when the value of $`p`$ is slightly shifted away from its critical value. Fig. 4 illustrates the deviations from the first-scaling for the values of parameter $`p`$ close to $`p_{cr}`$, on both sides of $`p_{cr}`$. Even very close to the critical point, these deviations are quite significant and can be easily seen in this representation.
## VI Irreversible aggregation - example of Smoluchowski kinetic model
The irreversible sol-gel transition can be modelled using the coupled non-linear differential equations in distributions $`c_s`$ of clusters of mass $`s`$ per unit volume (the Smoluchowski equations ) :
$`{\displaystyle \frac{dc_s}{dt}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i+j=s}{}}K_{i,j}c_ic_j{\displaystyle \underset{j}{}}K_{s,j}c_sc_j.`$ (114)
Coefficients $`K_{i,j}`$ represent the probability of aggregation per unit of time between two clusters of mass $`i`$ and $`j`$ . The Smoluchowski equations are derived from the master equation in the mean-field approximation :
$$<c_kc_l>=<c_k><c_l>.$$
The time $`t`$ includes both diffusion and reaction times. Eqs. (114) suppose irreversibility of the aggregation, i.e. the cluster fragmentation is excluded. One should notice however that the sum over $`j`$ in (114) does not include the infinite cluster (gel) because :
$$c_{j=\mathrm{}}=1/\mathrm{}=0.$$
Experimentally known aggregation kernels $`K_{ij}`$ are homogeneous functions :
$`K_{ai,aj}=a^\lambda K_{i,j}`$ (115)
with $`\lambda `$ being the homogeneity index. Perhaps the simplest physically relevant example of a homogeneous kernel is : $`K_{i,j}=(ij)^\mu `$. It has been shown in this case that if $`\mu `$ is larger than 1/2, then there exists a time $`t_{cr}`$ ($`t_{cr}<\mathrm{}`$) such that $`m_1^{}=1`$ for $`tt_{cr}`$ but $`m_1^{}<1`$ for $`t>t_{cr}`$ .
Let us consider now the case : $`K_{i,j}=(ij)^\mu `$ with $`\mu =1`$ in more details. It was shown in this case that the critical gelation time is : $`t_{cr}=1`$, and the solution for size-distribution of the Smoluchowski equations with the monodisperse initial condition is :
$`c_s`$ $`=`$ $`{\displaystyle \frac{s^{s2}}{s!}}e^{st}t^{s1}\text{ for}t1`$ (116)
$`c_s`$ $`=`$ $`{\displaystyle \frac{s^{s2}}{s!}}{\displaystyle \frac{\mathrm{exp}(s)}{t}}\text{ for}t>1`$ (118)
The asymptotic solution for large $`s`$ is :
$`c_s`$ $``$ $`{\displaystyle \frac{1}{t\sqrt{2\pi }}}s^{5/2}\mathrm{exp}\left[s(t1+\mathrm{ln}t)\right]\text{for}t1`$ (119)
$`c_s`$ $``$ $`{\displaystyle \frac{1}{t\sqrt{2\pi }}}s^{5/2}\text{for}t>1`$ (121)
Note that the power-law behaviour ($`\tau =5/2`$) is present for $`t>1`$ and not only at the threshold. The whole distribution of finite-size clusters evolves self-similarly and the appearance of a power-law behaviour is not here a sign of the critical behaviour but a specific characteristics of the gelation phase.
The solution for the first normalized moment is :
$`m_1^{}`$ $`=`$ $`1\text{ for}t1`$ (122)
$`m_1^{}`$ $`=`$ $`{\displaystyle \frac{1}{t}}\text{ for}t>1`$ (124)
With those asymptotic forms of $`c_s`$ one can calculate the gel fraction in the infinite system before and after the critical point :
$`m_G`$ $`=`$ $`0\text{for}t1`$ (125)
$`m_G`$ $`=`$ $`1{\displaystyle \frac{1}{t}}\text{for}t>1`$ (127)
It has been shown that gelation is analogous to a dynamical critical phenomenon with :
$`m_G=\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N}}<s_{max}>`$ (128)
as the order parameter. For one realization, $`s_{max}`$ corresponds to the mass of the gel above $`t_c=1`$.
For finite sizes, one makes the usual assumption that there exists a characteristic size diverging at the transition, say :
$`N_c|t1|^{1/\sigma _N}`$ (129)
such that for the mass gel in the finite system one has :
$`{\displaystyle \frac{1}{N}}<s_{max}>(t1)f\left({\displaystyle \frac{N}{N_c}}\right)\text{for}t1`$ (130)
In particular, at the gelation time one has :
$`<s_{max}>N^{1\sigma _N}N^g.`$ (131)
Using the formula (49), which is valid both for the equilibrium and non-equilibrium systems, one can calculate the anomalous dimension. Given the value of $`\tau `$ (see (119)), one finds $`g=2/3`$. Hence : $`\sigma _N=1/3`$ can be deduced from (131). The average value of the order parameter $`<s_{max}>`$ increases logarithmically for $`t<1`$, and is a finite proportion of the system size when $`t>1`$.
The illustration of the above discussion is shown in Figs. 5 and 6. Fig. 5 shows the distribution of $`s_{max}`$ in the first-scaling variables for systems of different sizes. The results have been obtained in the Smoluchowski model with the kernel : $`K_{ij}=ij`$, at the critical time : $`t=t_{cr}=1`$. Fluctuation properties of $`s_{max}`$ outside of the critical time : $`t=2t_{cr}`$, are shown in Fig. 6. The remaining parameters of the Smoluchowski calculations are the same as used in the calculations shown in Fig. 5. In this case, the data for different system sizes collapse into the universal curve in the scaling variables with $`\mathrm{\Delta }=1/2`$. One should keep in mind that in both cases, the fragment-size distribution is a power-law with $`\tau =5/2`$ (see (119)).
The relation between the form of tail of the scaling function and the anomalous dimension (41) was derived analytically in Sect. II.C for the equilibrium systems at the second-order phase transition. For non-equilibrium systems, we do not know equally rigorous derivation (see also Sect. VI.A). On the other hand, one may expect that the relation between the $`N`$-dependence of the average value of the order parameter and the asymptotic form of the scaling function in the limit $`N\mathrm{}`$, i.e. between $`\widehat{\nu }`$ and $`g`$, is connected to the asymptotic stability of the limit distributions. Actually, there is a very close connection of the renormalization group ideas and the limit theorems in the probability theory . If true, then the relation (41) could be valid in a more general framework than the one provided by the equilibrium statistical mechanics. To check this assertion, in Fig. 7 we show the plot of the logarithm of the scaling function $`\mathrm{\Phi }(z_{(1)})`$ (see Fig. 5) versus $`z_{(1)}^3`$ for large values of $`z_{(1)}`$. If the relation (41) is valid also for the non-equilibrium sol - gel second-order phase transition, then : $`\mathrm{\Phi }(z_{(1)})\mathrm{exp}(z^3)`$, and the tail of the scaling function should be a straight line in Fig. 7. That is indeed the case.
Figs. 8 and 9, show the $`\mathrm{\Delta }`$-scaling for the shifted order parameter variable : $`M_1^{^{}}=Ns_{max}`$. Results of the Smoluchowski calculations with the kernel : $`K_{ij}=ij`$, are shown at $`t=t_{cr}`$ (see Fig. 8), and at $`t=2t_{cr}`$ (see Fig. 9). One sees that the $`M_1^{^{}}`$ distribution exhibits a qualitative change while going from the critical time $`t=t_{cr}`$, where $`\mathrm{\Delta }=0.67`$, to $`t=2t_{cr}`$ for which $`\mathrm{\Delta }=1/2`$. At $`t=t_{cr}`$, the value of $`\mathrm{\Delta }`$ obtained by superposing different $`M_1^{^{}}`$ distributions in the scaling plot (64), agrees perfectly with the value of the anomalous dimension $`g`$ (=2/3).
By comparing Figs. 5, 6 with Figs. 8, 9, one may see also that the effect of changing the variable : $`s_{max}M_1^{^{}}`$, is seen only at the critical time (compare Figs. 5 and 8) where $`(\mathrm{\Delta }=1)(\mathrm{\Delta }=0.67)`$, and is absent above the critical time (compare Figs. 6 and 9) where $`\mathrm{\Delta }`$ (=1/2) remains unchanged.
Finally, in Fig. 10 we show the size-dependence of the $`M_1^{^{}}`$-distributions at $`t=t_{cr}`$, when the distributions are plotted in the ’wrong’ variables of the second-scaling $`\mathrm{\Delta }=1/2`$. The distributions for two system sizes are clearly displaced , showing the sensitivity of the scaling analysis and failure of the second-scaling.
### A The origin of fluctuations and the argument of Van Kampen
$`\mathrm{\Omega }`$ \- expansion is a systematic expansion of the master equations in powers of $`1/N`$ . Lushnikov was the first to express the generating functions as the contour integrals for quantities like the moments $`M_k^{}`$ . Then Van Dongen and Ernst used $`\mathrm{\Omega }`$ \- expansion to calculate explicitly these integrals for the moments $`M_k^{}`$ in some simple cases like : $`K_{ij}=ij`$. For example, the result for $`M_1^{}`$ can be expressed in terms of the generating function for the $`s_{max}`$-distribution $`P_N[s_{max}]`$, as :
$`{\displaystyle \underset{s_{max}}{}}P_N[s_{max}]e^{s_{max}u}`$ $`=`$ $`{\displaystyle \frac{N!}{2i\pi }}e^{Nu}{\displaystyle }{\displaystyle \frac{dz}{z^{N+1}}}\times `$ (132)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \underset{s=1}{\overset{N}{}}}{\displaystyle \frac{c_s}{N^{s1}}}e^{s(Ns)/(2N)}(ze^u)^s\right]`$ (133)
Using then the identity :
$`{\displaystyle \frac{^n}{u^n}}`$ $`\left[\mathrm{exp}\left({\displaystyle \underset{s}{}}\alpha _se^{su}\right)\right]={\displaystyle }{\displaystyle \frac{n!}{1!^{a_1}a_1!\mathrm{}n!^{a_n}a_n!}}\times `$ (134)
$`\times `$ $`\left[{\displaystyle \underset{s}{}}\alpha _ss^1e^{su}\right]^{a_1}\mathrm{}\left[{\displaystyle \underset{s}{}}\alpha _ss^ne^{su}\right]^{a_n}\mathrm{exp}\left({\displaystyle \underset{s}{}}\alpha _se^{su}\right)`$ (135)
, where the sum runs over different sets $`\{a_1,\mathrm{}a_n\}`$ with the constraint : $`a_1+\mathrm{}+a_n=n`$, and the particular result written down by Van Dongen and Ernst for the $`K_{ij}=ij`$-case :
$`\mathrm{exp}`$ $`\left[{\displaystyle \underset{s=1}{\overset{N}{}}}{\displaystyle \frac{c_s}{N^{s1}}}\mathrm{exp}[{\displaystyle \frac{s(Ns)}{2N}}]z^s\right]=`$ (136)
$`=`$ $`{\displaystyle \underset{s=1}{\overset{N}{}}}{\displaystyle \frac{z^s}{s!}}\mathrm{exp}[{\displaystyle \frac{1}{2}}st{\displaystyle \frac{1s}{N}}]+O(1/N)`$ (137)
we can find :
$`{\displaystyle \underset{s_{max}}{}}P_N[s_{max}]`$ $`\mathrm{exp}`$ $`(s_{max}u)=`$ (138)
$`\mathrm{exp}`$ $`\left[N\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{m_k^{}(u)^k}{k!}}+u\right)\right].`$ (139)
Having the $`\mathrm{\Omega }`$ \- expansion of the generating function of the $`s_{max}`$-distribution, we can conclude about the scaling at the gelation point. The moments $`m_k^{}`$ of the size-distribution for infinite systems are known to diverge near the gelation time as :
$`m_k^{}|t1|^{32k}.`$ (140)
For finite systems, using : $`\sigma _N=1/3`$, one obtains :
$`m_k^{}|t1|^{2k+3}f_k(N|t1|^3)N^{\frac{2k3}{3}}`$ (141)
at the gelation time. We have then found the asymptotic result :
$`Nm_k^{}a_k<s_{max}>^k,`$ (142)
where $`a_k`$’s are some positive constants. Inserting (142) in (138), one can show that the generating function for $`s_{max}`$ is the function of a single variable : $`<s_{max}>u`$, what is a sufficient condition for the validity of the first-scaling law (18).
We can have informations on similar scalings for various moment-distributions. $`\mathrm{\Omega }`$ \- expansion leads to the results :
$`<M_k^2><M_k^{}>^2`$ $`=`$ $`Nm_{2k}^{}(1t)m_{k+1}^2`$ (143)
$`<M_k^{}>`$ $`=`$ $`Nm_k^{}`$ (145)
for the values of $`k`$ when all the quantities are defined. At the transition ($`t=1`$), the relation (141) allows to calculate the reduced moments $`m_k^{}`$. The results can be written in the compact form :
$`{\displaystyle \frac{<M_k^2><M_k^{}>^2}{<M_k^{}>^{2\mathrm{\Delta }}}}\text{const},`$ (146)
with the following values of exponent $`\mathrm{\Delta }`$ :
$`\mathrm{\Delta }`$ $`=`$ $`1/2\text{for}k3/4`$ (147)
$`\mathrm{\Delta }`$ $`=`$ $`2k/3\text{for}3/4k3/2`$ (148)
$`\mathrm{\Delta }`$ $`=`$ $`1\text{for}3/2k`$ (149)
These are indications for $`\mathrm{\Delta }`$-scaling according to the remark (85) in Sect. II.H. More precisely : the moments of order $`k<3/4`$ are not critical (the second-scaling law), the moments of order $`k`$ between 3/4 and 3/2 exhibit the $`\mathrm{\Delta }`$-scaling with $`\mathrm{\Delta }=2k/3`$. In particular, for $`k=1`$, one recovers the correct value : $`\mathrm{\Delta }=g=2/3`$, corresponding to the general argument of the shifted order-parameter (62) with $`a_1=a_2=1`$. Finally, when the value of $`k`$ is larger than 3/2, we obtain the first-scaling law for the distribution of moments $`M_k^{}`$. This is also a consequence of the shifted order parameter argument, since in these cases :
$$<M_k^{}><s_{max}^k>.$$
Far from the critical point, all the reduced moments $`m_k^{}`$ are independent of $`N`$, since the correlation size :
$$\frac{1}{t1+\mathrm{ln}t}$$
(see (119)) is finite. Then, for any value of $`k`$, the second-scaling law holds, as expected from the general theory.
The above results about $`\mathrm{\Delta }`$-scaling for various moments of the size-distribution in the Smoluchowski model with kernel $`K_{ij}=ij`$, are not complete since the arguments involve only the second cumulant moment $`\kappa _2`$. In principle, as shown in Sect. II.H, all cumulants should be investigated. So, even though many exact results are known in this model, the complete analytical solution is not yet available.
The same study as presented above for gelling systems, can be performed also for non-gelling systems. An example of this kind is obtained for : $`K_{ij}=i+j`$. In this case, the size-distribution is power-law with the exponent $`\tau =3/2`$ and, following the discussion in Sect. II.C, the cluster multiplicity can be the order parameter. One can derive analytically, that the multiplicity distribution is binomial :
$`P_N[M_0,t]`$ $`=`$ $`\left(\begin{array}{c}N1\\ M_01\end{array}\right)(1\mathrm{exp}(Nt))^{NM_0}\times `$ (152)
$`\times `$ $`\mathrm{exp}[(M_01)Nt]`$ (153)
and can be approximated for $`N\mathrm{}`$ and for a finite value of $`<M_0>/N`$ by :
$`<M_0>^{1/2}`$ $`P_N[M_0,t]{\displaystyle \frac{1}{\sqrt{2\pi }(1e^{Nt})}}\times `$ (154)
$`\times `$ $`\mathrm{exp}\left({\displaystyle \frac{1}{2(1e^{Nt})}}{\displaystyle \frac{(M<M_0>)^2}{<M_0>}}\right)`$ (155)
what corresponds to the second-scaling. One may notice, that this binomial distribution is exactly equivalent to the bond percolation on a Bethe lattice with the occupation probability :
$$p=1\mathrm{exp}(Nt).$$
In spite of self-similar features in the fragment-size distribution at the infinite time, one does not see any critical behaviour in the cluster multiplicity distribution at any time in the non-gelling aggregation systems. This confirms the observation made in Sect. III.A in the Mekjian equilibrium model, that the power-law size-distribution alone does not guarantee that the system exhibits the critical behaviour.
The insight gained from the numerical simulations of Smoluchowski equations and the evidences from exact results for both gelling and non-gelling aggregation systems, provide strong hints that the discussion of Sect. II.H is valid not only for the equilibrium systems but also for the non-equilibrium ones. We see the same significance of the $`\mathrm{\Delta }`$-scaling in non-equilibrium systems as found in thermodynamic systems, and not only at the critical point but also close to the critical point or even far from it. We believe that this universality, which is common to equilibrium and non-equilibrium systems, has deeper foundation in the relation between renormalization group ideas for self-similar systems and the limit theorems of probability theory for the asymptotic scaling laws of order-parameter distributions. The concept of statistical equilibrium does not intervene at this level. One should also remember that the universality discussed in this work, is associated with only one critical exponent, and certainly does not exhausts all the singularity properties of the thermodynamical potential in second-order thermal phase-transition.
## VII Conclusions
We have presented in this paper the theory of universal scaling laws of the order parameter fluctuations in any system in which the second-order critical behaviour exist. These scaling laws, called $`\mathrm{\Delta }`$-scaling laws, are rigorously derived for the equilibrium systems. Moreover, both analytical and numerical evidence is presented in favour of the general validity of the $`\mathrm{\Delta }`$-scaling laws also for the off-equilibrium processes which exhibit the critical phenomenon of the second-order.
In this work we have discussed different aggregation models, both reversible and irreversible, finding the same connexion between scaling function properties and the anomalous dimension (critical exponents). These results can be important in the phenomenological analysis of ’critical behaviour’ in finite systems, where the critical exponent analysis is dubious and, moreover, the precise mechanism of the process may be unknown. In these cases, the $`\mathrm{\Delta }`$-scaling analysis allows to select both the relevant observable and the interesting initial conditions, which lead to the ’pseudocritical’ behaviour in the studied process. Another interesting aspect of the $`\mathrm{\Delta }`$-scaling analysis is the possibility of compressing data and, hence, the elimination of redundant dependences in the data on parameters like the system size, the total energy, or the total momentum etc.. This provides an obligatory intermediate step in any phenomenological analysis before the laws governing complicated dynamics can be found. Finally, one should stress that the scaling laws discussed in this paper are independent of whether one deals with an equilibrium or an off-equilibrium process. This is a crucial advantage in the studies of short-lived systems. Examples of the multifragmentation processes in the collisions of atomic nuclei or atomic clusters well illustrate this problematic . In the absence of thermal equilibrium, which is a theoretical hypothesis difficult to verify in dynamically formed short-lived systems, we simply do not have at our disposal any other tool to address reliably the question of possible ’criticality’ of the studied process.
As said before, the $`\mathrm{\Delta }`$-scaling analysis developed in this work provides an alternative to the critical exponents analysis in the equilibrium systems and is the only tool for the analysis of the non-equilibrium systems. All essential information can be deduced from the scaling function, the value of $`\mathrm{\Delta }`$ parameter, the form of the tail of the scaling function and the value of the anomalous exponent. With these informations it is possible to find out whether the studied process is at the critical point, in its neighbourhood or far away from it. Reference point in this analysis is the (approximate) self-similarity of the system. Generalization of the above scaling theory to discontinuous phase-order transitions for which the characteristic length can be defined, is in progress.
We thank R. Paredes V. for providing us with the Potts model data which were used in preparing Fig. 1.
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# A Monte Carlo Formulation of the Bogolubov Theory
## 1 Introduction
Since the first experimental demonstrations of Bose-Einstein condensation in atomic gases , we have assisted to a fantastic development of the field both on experimental and theoretical ground. From the theoretical point of view the Gross-Pitaevskii equation (GPE) or non-linear Schrödinger equation for the condensate wave function, in its steady-state and time dependent versions, revealed itself to be an extremely powerful tool to explain most of the experimental results .
Generally speaking the Gross-Pitaevskii equation is obtained by neglecting the interaction between the condensate and non condensed particles . This approximation is justified at temperatures much below the transition temperature where the non condensed cloud is very small and much less dense that the condensate. This is often the case in experimental conditions where almost pure condensates can be obtained.
There are nevertheless situations in which the rich physics coming from the non condensed fraction of the cloud and its coupling to the condensate cannot be neglected. A first example concerns the collective modes of the gas induced by a modulation of the trapping potential . Different theories going beyond the Gross-Pitaevskii equation were put forward in order to explain the temperature dependence of the frequencies and the damping rates measured in experiments (for more references see , pages 489-500).
A second example where the non condensed fraction may play a significant role concerns condensates in thermodynamically unstable states, like vortices and dark solitons recently created experimentally . In for example one or more vortices are created by stirring the condensate with a slightly anisotropic rotating potential: vortices appear provided that the rotation frequency exceeds a critical value. If the rotation is then stopped the vortices are not any more a local minimum of energy and are expected to spiral out of the condensate . It is indeed found experimentally that the vortices leave the condensate . It is predicted theoretically that the dynamics and lifetime of the thermodynamically unstable Bose-Einstein condensate depend crucially on the interaction between the condensate and the non condensed cloud .
The examples we have given above concern time dependent properties of the condensate. Another class of interesting finite temperature issues is given by the thermal equilibrium properties of the gas. To cite some example, one might be interested in the non condensed spatial density of the atoms, or in the thermal fluctuations of the number of particles in the condensate . Condensates with vortices give rise again to interesting questions in this respect such as the thermal fluctuations in the position of the vortex core and the temperature dependence of the critical rotation velocity .
A well established technique to study Bose-Einstein condensates at thermal equilibrium is the Bogolubov method, which is a perturbative method valid for low temperatures when the density of condensate particles is much larger than the density of non condensed particles. This method relies on the quadratization of the Hamiltonian in the field operator representing the non condensed fraction, and on the Bogolubov transform . Another method which could be valid for a larger range of temperatures is the Hartree-Fock-Bogolubov (HFB) method , which is instead a variational method, based on a Gaussian Ansatz for the system density operator. For time dependent problems, one disposes of the former two methods, plus more refined techniques including kinetic equations for the non condensed particles .
A common point to all these methods is the necessity to calculate mean values using a density operator which is an exponential of a quadratic form in the field operator: the exponential of the quadratized Hamiltonian in the case of the Bogolubov approach or the Gaussian Ansatz for the HFB approach. The procedure that is usually followed to perform those calculations is to diagonalize the quadratic form by the Bogolubov transform, turning it into a sum of decoupled harmonic oscillator energy terms.
Unfortunately the corresponding calculations are quite heavy for 3D non-homogeneous systems such as trapped gases, and this constitutes a serious limitation to the use of these methods. E.g. in practice, by discretizing the real space on a grid with a modest number of $`64`$ points per spatial dimension, and in the absence of rotational symmetry (as in a multi-vortex configuration) one has to diagonalize a matrix which is $`[2(64)^n\times 2(64)^n]`$ large, where $`n=1,2,3`$ is the dimension of the space. For $`n=3`$ the matrix to diagonalize is $`\mathrm{524\hspace{0.17em}288}\times \mathrm{524\hspace{0.17em}288}`$.
In this paper we propose an alternative formulation of the Bogolubov theory which can be implemented numerically more efficiently in two or three spatial dimensions since it avoids the diagonalization of big matrices. Our method relies on the use of the Wigner representation of the density operator and on a Monte Carlo simulation to sample the steady-state Wigner distribution of the system at finite temperature. Since our method allows to sample a “Gaussian” density operator, it applies both to the Bogolubov and to the HFB theory, in the thermal equilibrium as well as in the time dependent case. We give in this paper a detailed derivation of the method in the frame of a particle number conserving Bogolubov theory . We consider first the thermal equilibrium case and then we show how to extend the treatment to the time dependent case.
The paper is organized as follows: in section 2 we recall the number-conserving version of the Bogolubov theory put forward in , which we adopt in the following. In section 3 we reformulate the theory by expressing the Bogolubov density operator in terms of the Wigner quasi-distribution function. In section 4 we describe the stochastic simulation that we use to sample the equilibrium Wigner function. In section 5 we compare numerically our stochastic method to the direct diagonalization used in the traditional Bogolubov approach in the case of a 1D trapped Bose gas, where the explicit diagonalization is feasible, and we calculate the spatial distribution of the non condensed fraction. In section 6 we show how a time dependent problem would be treated with our method. Conclusions and perspectives are presented in section 7.
## 2 Bogolubov theory conserving the total number of particles
In this section we recall briefly the main ideas and results of where a number-conserving version of the Bogolubov theory is established. The reader already familiar with this theory can jump directly to section 3.
We consider a cloud of atoms trapped in the potential $`U(\stackrel{}{r})`$ interacting through an effective low energy contact potential $`V=g\delta (\stackrel{}{r_1}\stackrel{}{r_2})`$:
$$H=d^3\stackrel{}{r}\widehat{\psi }^{}(\stackrel{}{r})h_1\widehat{\psi }(\stackrel{}{r})+\frac{g}{2}\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }^{}(\stackrel{}{r})\widehat{\psi }(\stackrel{}{r})\widehat{\psi }(\stackrel{}{r})$$
(1)
where $`\widehat{\psi }`$ is the atomic field operator, $`h_1`$ is the one-body Hamiltonian:
$$h_1=\frac{\mathrm{}^2}{2m}\mathrm{\Delta }+U(\stackrel{}{r}).$$
(2)
The coupling constant $`g`$ is expressed in terms of the $`s`$-wave scattering length $`a`$ of the true interaction potential:
$$g=\frac{4\pi \mathrm{}^2a}{m}.$$
(3)
Let $`\rho _1`$ be the one-body density operator of our system and let $`\varphi _{\mathrm{ex}}`$ be the eigenvector of $`\rho _1`$ with the largest eigenvalue $`N_0`$:
$$\rho _1|\varphi _{\mathrm{ex}}=N_0|\varphi _{\mathrm{ex}}.$$
(4)
We normalize $`\varphi _{\mathrm{ex}}`$ to unity. Physically $`\varphi _{\mathrm{ex}}`$ is the most populated single particle state. We assume here the presence of a Bose-Einstein condensate so that $`\varphi _{\mathrm{ex}}`$ is the condensate wave function and $`N_0N`$. This motivates the splitting of the atomic field operator $`\widehat{\psi }`$ into condensate and non condensed modes:
$$\widehat{\psi }(\stackrel{}{r})=\varphi _{\mathrm{ex}}(\stackrel{}{r})\widehat{a}_{\varphi _{\mathrm{ex}}}+\delta \widehat{\psi }(\stackrel{}{r})$$
(5)
where $`\widehat{a}_{\varphi _{\mathrm{ex}}}`$ annihilates a particle in the condensate wavefunction $`\varphi _{\mathrm{ex}}`$, and $`\delta \widehat{\psi }(\stackrel{}{r})`$ is orthogonal to $`\varphi _{\mathrm{ex}}`$:
$$\varphi _{\mathrm{ex}}^{}\delta \widehat{\psi }=0.$$
(6)
An important property that comes from (4) and (6) it that there is no single particle coherence between the condensate and the non condensed modes:
$$\widehat{a}_{\varphi _{\mathrm{ex}}}^{}\delta \widehat{\psi }=0.$$
(7)
The interest of the splitting (5) is that the typical matrix element of $`\widehat{a}_{\varphi _{\mathrm{ex}}}`$ is of order $`N_0^{1/2}N^{1/2}`$ while the matrix elements of $`\delta \widehat{\psi }`$ scale as $`(\delta N)^{1/2}N^{1/2}`$ where $`\delta N`$ is the number of non condensed particles.
An essential ingredient of the number conserving Bogolubov approach is the operator $`\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}`$ transferring one non condensed particle into the condensate:
$$\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}=\widehat{N}^{1/2}\widehat{a}_{\varphi _{\mathrm{ex}}}^{}\delta \widehat{\psi }.$$
(8)
This operator plays here the role played by $`\delta \widehat{\psi }`$ in the symmetry breaking approach commonly adopted. In particular due to (7), the expectation value of $`\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}`$ vanishes.
In the weakly interacting Bose gas limit
$$N\mathrm{}\text{and}gN=\mathrm{const},$$
(9)
and for a constant trapping potential, the mean interaction energy per particle $`ng`$ (where $`n`$ is the density) tends to a finite value whereas the small gaseous parameter $`(na^3)^{1/2}`$ tends to zero. If the temperature, the trapping potential and the mean interaction energy are fixed, the number of non condensed particles $`\delta N`$ is bounded from above so that $`\delta NN`$ in the large $`N`$ limit . One can then make a systematic expansion of the exact condensate wavefunction $`\varphi _{\mathrm{ex}}`$ and of the fields $`\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}`$ and $`\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}^{}`$ in powers of the small parameter:
$$ϵ=\sqrt{\delta N/N}1/\sqrt{N}.$$
(10)
Formally:
$`\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_{(0)}+{\displaystyle \frac{1}{\sqrt{N}}}\widehat{\mathrm{\Lambda }}_{(1)}+{\displaystyle \frac{1}{N}}\widehat{\mathrm{\Lambda }}_{(2)}+\mathrm{}`$ (11)
$`\varphi _{\mathrm{ex}}`$ $`=`$ $`\varphi _{(0)}+{\displaystyle \frac{1}{\sqrt{N}}}\varphi _{(1)}+{\displaystyle \frac{1}{N}}\varphi _{(2)}+\mathrm{}`$ (12)
One then calculates $`d\widehat{\mathrm{\Lambda }}_{\mathrm{ex}}/dt=0`$ order by order in $`ϵ`$. At the order $`1`$ in $`ϵ`$ (the lowest approximation order), one finds $`\widehat{\mathrm{\Lambda }}_{(0)}=0`$ and $`\varphi _{(0)}=\varphi `$ where $`\varphi `$ is solution of the time-dependent Gross-Pitaevskii equation:
$$i\mathrm{}_t\varphi =[h_1+Ng|\varphi |^2\mu ]\varphi .$$
(13)
The quantity $`\mu `$ in (13) is such that the wavefunction of the condensate at equilibrium is a time independent solution of the above equation. Physically $`\mu `$ is the lowest order approximation to the chemical potential of the gas.
The order $`0`$ in $`ϵ`$ provides the equations of motion for $`\widehat{\mathrm{\Lambda }}\frac{1}{\sqrt{N}}\widehat{\mathrm{\Lambda }}_{(1)}`$ and $`\widehat{\mathrm{\Lambda }}^{}`$:
$$i\mathrm{}_t\left(\begin{array}{c}\widehat{\mathrm{\Lambda }}\\ \widehat{\mathrm{\Lambda }}^{}\end{array}\right)=\left(\begin{array}{c}\widehat{\mathrm{\Lambda }}\\ \widehat{\mathrm{\Lambda }}^{}\end{array}\right).$$
(14)
The first non zero correction to $`\varphi _{(0)}`$ is $`\frac{1}{N}\varphi _{(2)}`$, whose equation of motion is obtained in the next order of the expansion. The operator $``$ in (14) is given by:
$$=\left(\begin{array}{cc}H_{\mathrm{gp}}\mu +QNg|\varphi |^2Q& QNg\varphi ^2Q^{}\\ Q^{}Ng(\varphi ^{})^2Q& [H_{\mathrm{gp}}^{}\mu +Q^{}Ng|\varphi |^2Q^{}]\end{array}\right)$$
(15)
where
$$H_{\mathrm{gp}}=h_1+Ng|\varphi |^2$$
(16)
and $`Q`$ projects orthogonally to $`\varphi `$:
$$Q=1|\varphi \varphi |.$$
(17)
The complex conjugation indicated by the star is taken in the $`r`$-basis so that e.g. $`Q^{}`$ projects orthogonally to the state with wavefunction $`\varphi ^{}`$. The projection operator appears as well in the commutation relations of $`\widehat{\mathrm{\Lambda }}`$ and $`\widehat{\mathrm{\Lambda }}^{}`$:
$$[\widehat{\mathrm{\Lambda }}(\stackrel{}{r}),\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{s})]=\stackrel{}{r}|Q|\stackrel{}{s}.$$
(18)
The linear equations of motion (14) for $`\widehat{\mathrm{\Lambda }}`$ correspond to an approximation of the Hamiltonian which is quadratic in $`\widehat{\mathrm{\Lambda }}`$. This quadratic approximation can be obtained by inserting the decomposition (5) in the Hamiltonian, by neglecting terms cubic or quartic in $`\delta \widehat{\psi }`$, and by using:
$`a_{\varphi _{\mathrm{ex}}}^{}a_{\varphi _{\mathrm{ex}}}`$ $`=`$ $`\widehat{N}\delta \widehat{N}`$ (19)
$`\delta \widehat{N}`$ $``$ $`{\displaystyle \delta \widehat{\psi }^{}\delta \widehat{\psi }}{\displaystyle \widehat{\mathrm{\Lambda }}^{}\widehat{\mathrm{\Lambda }}}.`$ (20)
As equation (13) is satisfied, no term linear in $`\widehat{\mathrm{\Lambda }}`$ is left in the Hamiltonian. One is left with the quadratic Hamiltonian :
$$\widehat{H}_{\mathrm{quad}}=f(\widehat{N})+\frac{1}{2}(\widehat{\mathrm{\Lambda }}^{},\widehat{\mathrm{\Lambda }})\left(\begin{array}{c}\widehat{\mathrm{\Lambda }}\\ \widehat{\mathrm{\Lambda }}^{}\end{array}\right)$$
(21)
At thermal equilibrium in the canonical ensemble, the non condensed atoms are then described by the density operator:
$$\widehat{\sigma }_{\mathrm{nc}}=\frac{1}{Z}e^{\beta \widehat{H}_{\mathrm{quad}}}.$$
(22)
The usual way of proceeding at this point is to diagonalize the operator $``$. Let us introduce the eigenvalues and eigenvectors of $``$:
$$\left(\begin{array}{c}u_k\\ v_k\end{array}\right)=ϵ_k\left(\begin{array}{c}u_k\\ v_k\end{array}\right).$$
(23)
We suppose that all the eigenvalues $`ϵ_k`$ are real so that the equilibrium condensate wavefunction is a dynamically stable solution of the Gross-Pitaevskii equation. Moreover we restrict the notation $`(u_k,v_k)`$ to the modes satisfying the normalization and orthogonality conditions:
$$u_k^{}u_k^{}v_k^{}v_k^{}=\delta _{k,k^{}}.$$
(24)
To each mode $`(u_k,v_k)`$ is associated a “dual mode” $`(v_k^{},u_k^{})`$, of energy $`ϵ_k`$ and satisfying obviously $`v_k^{}v_k^{}u_k^{}u_k^{}=\delta _{k,k^{}}`$. In general the $`(u_k,v_k)`$’s plus the dual modes form a complete family in the subspace orthogonal to $`(\varphi ,0)`$ and $`(0,\varphi ^{})`$ so that the field $`\mathrm{\Lambda }`$ can then be expressed as:
$$\left(\begin{array}{c}\widehat{\mathrm{\Lambda }}(\stackrel{}{r})\\ \widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r})\end{array}\right)=\underset{k}{}\widehat{b}_k\left(\begin{array}{c}u_k(\stackrel{}{r})\\ v_k(\stackrel{}{r})\end{array}\right)+\widehat{b}_k^{}\left(\begin{array}{c}v^{}_k(\stackrel{}{r})\\ u^{}_k(\stackrel{}{r})\end{array}\right).$$
(25)
The normalization condition of the $`(u_k,v_k)`$’s ensures that the $`b_k,b_k^{}`$ satisfy bosonic commutation relations. If (25) is injected in (21), the quadratic Hamiltonian takes the form:
$$H_{\mathrm{quad}}=E_0+\underset{k}{}ϵ_k\widehat{b}_k^{}\widehat{b}_k$$
(26)
corresponding to a set of decoupled harmonic oscillators. For thermodynamical stability of the condensate we must require that all the $`ϵ_k`$ are positive. Physical quantities can then be readily extracted from $`\widehat{\sigma }_{\mathrm{nc}}`$. The inconvenient of this method is that the diagonalization of $``$ in two and three dimensions is a very heavy numerical task in the absence of particular symmetries of the problem.
## 3 Formulation of the Bogolubov theory in terms of the Wigner function
We wish to deal with complex numbers and complex functions instead of operators and fields
$$\widehat{b}_k^{},\widehat{b}_kb_k^{},b_k$$
$$\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r}),\widehat{\mathrm{\Lambda }}(\stackrel{}{r})\mathrm{\Lambda }^{}(\stackrel{}{r}),\mathrm{\Lambda }(\stackrel{}{r}).$$
For this reason we introduce the Wigner quasi distribution function $`W(\{b_k\},\{b_k^{}\})`$ defined as the Fourier transform of the characteristic function $`\chi (\{\alpha _k\})`$:
$$\chi (\{\alpha _k\},\{\alpha _k^{}\})=\text{Tr}\left[\widehat{\sigma }_{\mathrm{nc}}\underset{k}{}\mathrm{exp}(\alpha _k\widehat{b}_k^{}\alpha _k^{}\widehat{b_k})\right]$$
(27)
and
$$W(\{b_k\},\{b_k^{}\})=\chi (\{\alpha _k\},\{\alpha _k^{}\})\underset{k}{}\mathrm{exp}\left[(\alpha _k^{}b_k\alpha _kb_k^{})\right]\frac{(d\mathrm{}e\alpha _k)(d\mathrm{}m\alpha _k)}{\pi }.$$
(28)
With the diagonalized form of the quadratic Hamiltonian (26), the density operator (22) corresponds to decoupled harmonic oscillators at thermal equilibrium, so that:
$$W(\{b_k\},\{b_k^{}\})=\underset{k}{}2\mathrm{tanh}\left(\frac{\beta ϵ_k}{2}\right)\mathrm{exp}\left[2|b_k|^2\text{tanh}(\beta ϵ_k/2)\right].$$
(29)
In the high temperature limit $`k_BTϵ_k`$ the Wigner distribution for mode $`k`$ is simply the Boltzmann distribution for a classical field. In the low temperature limit $`k_BTϵ_k`$ it is worth noting that the Wigner distribution has a finite width tending to 1/2: even at zero temperature, the c-number quantities $`b_k`$ fluctuate mimicking the quantum fluctuations of the operators $`\widehat{b}_k`$.
Since we want to avoid the diagonalization of the operator $``$, we wish now to express the Wigner function (29) in terms of the functions $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ corresponding to the field operators $`\widehat{\mathrm{\Lambda }}`$ and $`\widehat{\mathrm{\Lambda }}^{}`$:
$$\left(\begin{array}{c}\mathrm{\Lambda }(\stackrel{}{r})\\ \mathrm{\Lambda }^{}(\stackrel{}{r})\end{array}\right)\underset{k}{}b_k\left(\begin{array}{c}u_k(\stackrel{}{r})\\ v_k(\stackrel{}{r})\end{array}\right)+b_k^{}\left(\begin{array}{c}v^{}_k(\stackrel{}{r})\\ u^{}_k(\stackrel{}{r})\end{array}\right).$$
(30)
By using (30) and the properties (24) of the $`u_k(\stackrel{}{r})`$ and $`v_k(\stackrel{}{r})`$, one can show that the Wigner function takes the form:
$$W(\mathrm{\Lambda },\mathrm{\Lambda }^{})=\text{det}\left[4\eta \mathrm{tanh}\frac{\beta }{2}\right]^{1/2}\mathrm{exp}\left\{(\mathrm{\Lambda }^{},\mathrm{\Lambda })\eta \mathrm{tanh}\frac{\beta }{2}\left(\begin{array}{c}\mathrm{\Lambda }\\ \mathrm{\Lambda }^{}\end{array}\right)\right\}$$
(31)
where $`\eta `$ is the matrix:
$$\eta =\left(\begin{array}{cc}\text{Id}& \text{0}\\ \text{0}& \mathrm{Id}\end{array}\right),$$
(32)
and the determinant is defined in the space orthogonal to $`(\varphi ,0)`$ and $`(0,\varphi ^{})`$.
The quantum mechanical mean values of totally symmetrized operators are then simply given as averages over the Wigner distribution ; e.g.:
$$\mathrm{\Lambda }(\stackrel{}{r_1})\mathrm{\Lambda }^{}(\stackrel{}{r_2})_W=\frac{1}{2}\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r_1})\widehat{\mathrm{\Lambda }}(\stackrel{}{r_2})+\widehat{\mathrm{\Lambda }}(\stackrel{}{r_2})\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r_1}).$$
(33)
Note that according to (18) the previous expression gives infinity for $`\stackrel{}{r_1}=\stackrel{}{r_2}`$! In practice we will not encounter this problem because we will work on a grid with a finite number of modes.
## 4 Brownian motion simulation
Since we quadratized the Hamiltonian with respect to $`\widehat{\mathrm{\Lambda }}`$ and $`\widehat{\mathrm{\Lambda }}^{}`$, the Wigner distribution (31) is Gaussian and can be recovered as the steady-state probability distribution of a conveniently chosen Brownian motion of the functions $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$. In this section we show that one can find such Brownian motion in a convenient way for numerical simulations. First of all we discretize our problem on a grid. To this aim we introduce the vector $`\stackrel{}{\mathrm{\Lambda }}`$ representing the set of values $`\{\mathrm{\Lambda }(\stackrel{}{r}_l)\}`$ where $`\stackrel{}{r_l}=(x_i,y_j,z_k)`$ are points of the grid:
$$(\stackrel{}{\mathrm{\Lambda }})_l\mathrm{\Lambda }(\stackrel{}{r_l}).$$
(34)
If $`dV`$ is the volume element of the grid, the Wigner function (31) takes the form:
$$W(\mathrm{\Lambda },\mathrm{\Lambda }^{})=𝒜\mathrm{exp}\left[dV(\stackrel{}{\mathrm{\Lambda }}^{},\stackrel{}{\mathrm{\Lambda }})\eta \mathrm{tanh}\frac{\beta }{2}\left(\begin{array}{c}\stackrel{}{\mathrm{\Lambda }}\\ \stackrel{}{\mathrm{\Lambda }}^{}\end{array}\right)\right]$$
(35)
where now $``$ is represented by a matrix, and where we did not write explicitly the normalization factor. The vector $`\stackrel{}{\mathrm{\Lambda }}`$ is then submitted to the following stochastic evolution during the time step $`dt`$:
$$d\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{F}dt+Y_{11}d\stackrel{}{\xi }+Y_{12}d\stackrel{}{\xi }^{}$$
(36)
where $`\stackrel{}{F}`$ is a linear friction force
$$\stackrel{}{F}=\alpha _{11}\stackrel{}{\mathrm{\Lambda }}\alpha _{12}\stackrel{}{\mathrm{\Lambda }}^{}$$
(37)
and $`d\stackrel{}{\xi }`$ is a noise term satisfying:
$`\overline{d\xi (\stackrel{}{r})}`$ $`=`$ $`0`$ (38)
$`\overline{d\xi (\stackrel{}{r}_1)d\xi ^{}(\stackrel{}{r}_2)}`$ $`=`$ $`2{\displaystyle \frac{dt}{dV}}\left[\delta _{\stackrel{}{r}_1,\stackrel{}{r}_2}dV\varphi (\stackrel{}{r}_1)\varphi ^{}(\stackrel{}{r}_2)\right]`$ (39)
$`\overline{d\xi (\stackrel{}{r}_1)d\xi (\stackrel{}{r}_2)}`$ $`=`$ $`0,`$ (40)
where the overline denotes the average over the time interval $`dt`$. In (36) and (37), $`\stackrel{}{\mathrm{\Lambda }}`$ and $`\stackrel{}{\xi }`$ are vectors whose length $`𝒩`$ is the size of the grid, while $`\alpha _{ij}`$ and $`Y_{ij}`$ $`(i,j=1,2)`$ are $`𝒩\times 𝒩`$ matrices. The prescription (39) is equivalent to the usual prescription for spatially $`\delta `$-correlated Gaussian white noise with the additional requirement that the noise should be orthogonal to the condensate wavefunction $`\varphi `$, so that $`\mathrm{\Lambda }`$ remains within the subspace orthogonal to $`\varphi `$. It is important to note that the choice of the ‘time’ variable $`t`$ in the Brownian motion simulation is totally arbitrary; in what follows $`dt`$ will have the dimension of an energy.
Let $`P(\stackrel{}{\mathrm{\Lambda }},\stackrel{}{\mathrm{\Lambda }}^{},t)`$ be the probability distribution of $`\stackrel{}{\mathrm{\Lambda }}`$ and $`\stackrel{}{\mathrm{\Lambda }}^{}`$ at time $`t`$. The stochastic equations (36) can then be turned into a Fokker-Plank equation for $`P(\stackrel{}{\mathrm{\Lambda }},\stackrel{}{\mathrm{\Lambda }^{}},t)`$:
$$_tP+\underset{\stackrel{}{r}}{}\left[_{\mathrm{\Lambda }(\stackrel{}{r})}J(\stackrel{}{r})+\text{c.c.}\right]=0$$
(41)
where $`J(\stackrel{}{r})`$ is a density current:
$$J(\stackrel{}{r})=F(\stackrel{}{r})P\frac{1}{dV}\underset{\stackrel{}{s}}{}_{\mathrm{\Lambda }(\stackrel{}{s})}(D_{12}(\stackrel{}{r},\stackrel{}{s})P)+_{\mathrm{\Lambda }^{}(\stackrel{}{s})}(D_{11}(\stackrel{}{r},\stackrel{}{s})P)$$
(42)
and $`D_{ij}`$ are constant diffusion matrices:
$`(D_{11})(\stackrel{}{r},\stackrel{}{s})`$ $`=`$ $`{\displaystyle \frac{dV}{2dt}}\overline{d\mathrm{\Lambda }(\stackrel{}{r})d\mathrm{\Lambda }^{}(\stackrel{}{s})}`$ (43)
$`(D_{12})(\stackrel{}{r},\stackrel{}{s})`$ $`=`$ $`{\displaystyle \frac{dV}{2dt}}\overline{d\mathrm{\Lambda }(\stackrel{}{r})d\mathrm{\Lambda }(\stackrel{}{s})}.`$ (44)
Neglecting terms on the order of $`dt`$, we can express the identities (43-44) more synthetically in terms of a diffusion matrix:
$$D\left(\begin{array}{cc}D_{11}& D_{12}\\ D_{12}^{}& D_{11}^{}\end{array}\right)=YY^{}$$
(45)
where we have introduced the noise matrix
$$Y=\left(\begin{array}{cc}Y_{11}& Y_{12}\\ Y_{12}^{}& Y_{11}^{}\end{array}\right).$$
(46)
The steady-state solution of (41) with $`J(\stackrel{}{r})=0`$ is found to be:
$$P\mathrm{exp}\left[\frac{dV}{2}(\stackrel{}{\mathrm{\Lambda }}^{},\stackrel{}{\mathrm{\Lambda }})D^1\alpha \left(\begin{array}{c}\stackrel{}{\mathrm{\Lambda }}\\ \stackrel{}{\mathrm{\Lambda }}^{}\end{array}\right)\right]$$
(47)
with the friction matrix $`\alpha `$ defined as:
$$\alpha =\left(\begin{array}{cc}\alpha _{11}& \alpha _{12}\\ \alpha _{12}^{}& \alpha _{11}^{}\end{array}\right).$$
(48)
If we now require that the steady-state distribution (47) coincides with the Wigner distribution (31) describing thermal equilibrium, we obtain:
$$\alpha =2D\eta \mathrm{tanh}\left(\frac{\beta }{2}\right).$$
(49)
Equation (49) is simply the generalization to a field of Einstein’s relation for Brownian motion of a particle: in the high temperature limit the diffusion term divided by the friction term is proportional to temperature.
Provided that the Einstein’s relation (49) is satisfied, the choice of $`\alpha `$ and $`Y`$ (which determines $`D`$) is not unique. It turns out that a convenient choice of the friction and noise terms from the point of view of a numerical simulation is given by:
$$\alpha =\left[\mathrm{cosh}\left(\frac{\beta }{2}\right)\right]\eta \left[\mathrm{sinh}\left(\frac{\beta }{2}\right)\right]/\left(\frac{\beta }{2}\right)$$
(50)
and:
$$Y=\left[\mathrm{cosh}\left(\frac{\beta }{2}\right)\right]\frac{1}{\beta ^{1/2}}.$$
(51)
One can check that the choice (50) and (51) satisfies (49) by using the following properties of the operators $``$ and $`\eta `$:
$$\eta \eta =^{}\text{and}\eta ^1=\eta .$$
(52)
One also has to check that the choice (50) and (51) reproduces the structure of the matrices (46) and (48), that is the second line of the matrix is obtained from the first line by complex conjugation and exchange of indices 1 and 2, a property that can be formalized as
$$S\alpha S=\alpha ^{}\text{and}SYS=Y$$
(53)
where $`S`$ is the matrix
$$S=S^1=\left(\begin{array}{cc}\text{0}& \text{Id}\\ \mathrm{Id}& \text{0}\end{array}\right).$$
(54)
From the properties $`SS=^{}`$ and $`S\eta S=\eta `$ one finds that (50) and (51) have the right structure indeed.
The convenience of the choice (50) and (51) relies in the fact that one then is brought to calculate the action of the operators $`\mathrm{exp}\left(\frac{\beta }{2}\right)`$ on the functions $`\mathrm{\Lambda }`$, which is easily implemented by using a splitting plus Fourier transform technique (more details about the numerical simulation will be given in the end of the next section).
## 5 Test of the method for a 1D trapped Bose gas
As a test of the method we have made a numerical comparison between the Monte Carlo simulation and the direct diagonalization of the operator $``$ for the case of a Bose-Einstein condensate in a harmonic potential with frequency $`\omega `$ in one spatial dimension. Physically this situation would correspond to a very elongated cigar shaped trap with a longitudinal trap frequency $`\omega `$ and a strong transversal confinement $`\mathrm{}\omega _{}>k_BT`$, so that most of the atoms are in the transverse ground state $`\varphi _{}`$ of the trap. In these conditions the field operator can be approximated by:
$$\widehat{\psi }(\stackrel{}{r})=\varphi _{}(x,y)\widehat{\psi }(z)$$
(55)
where the operators $`\widehat{\psi }^{}(z)`$ and $`\widehat{\psi }(z)`$ have the usual bosonic commutation relations for field operators in one dimension. We have then reduced our problem to one spatial dimension with an effective coupling constant $`g_{1D}`$ :
$$g_{1D}=g𝑑x𝑑y|\varphi _{}(x,y)|^4.$$
(56)
We wish to calculate the spatial density of non condensed atoms. In order to calculate the average $`\widehat{\mathrm{\Lambda }}^{}(z)\widehat{\mathrm{\Lambda }}(z)`$ we use the formula (33) and equation (18) in its discretized version, to obtain:
$$\widehat{\mathrm{\Lambda }}^{}(z)\widehat{\mathrm{\Lambda }}(z)=\mathrm{\Lambda }^{}(z)\mathrm{\Lambda }(z)\frac{1}{2}\left[\frac{1}{dz}|\varphi (z)|^2\right].$$
(57)
The result for the non condensed cloud spatial density as a function of $`z`$ is reported in figure 1. We have chosen a total number of atoms $`N=10^4`$, a temperature $`k_BT=30\mathrm{}\omega `$, and a coupling constant $`g_{1D}=0.01\mathrm{}(\mathrm{}\omega /m)^{1/2}`$ leading to a chemical potential $`\mu =14.1\mathrm{}\omega `$. The full line is the direct diagonalization result and the points with the error bars are obtained with our Monte Carlo simulation with 2000 realizations. The non condensed cloud is repelled by the condensate which is sitting in the center of the harmonic trap, which explains the double peak structure of figure 1.
We can also get information on the number of condensate particles by calculating the number of non-condensed particles, since the total number of particles is fixed. The mean number of non-condensed particles is given by the spatial integral of Eq.(57):
$$\delta \widehat{N}=\delta N_W\frac{𝒩1}{2}$$
(58)
where the operator $`\delta \widehat{N}`$, defined in (20), is replaced at the present order by its approximation in terms of $`\widehat{\mathrm{\Lambda }}`$, $`𝒩`$ is the number of modes in the simulation (that is the number of points on the spatial grid) and
$$\delta Ndz\underset{z}{}\mathrm{\Lambda }^{}(z)\mathrm{\Lambda }(z).$$
(59)
After some algebra, involving the total symmetrization of a product of 4 operators $`\widehat{\mathrm{\Lambda }}`$ or $`\widehat{\mathrm{\Lambda }}^{}`$, we also obtain
$$\delta \widehat{N}^2\delta \widehat{N}^2=\delta N^2_W\delta N_W^2\frac{𝒩1}{4}.$$
(60)
For the parameters of the figure we get the mean values and the standard deviations
$`(\delta \widehat{N})_{\mathrm{MC}}=385\pm 6`$ and $`(\delta \widehat{N})_{\mathrm{diag}}=391,`$ (61)
$`\sigma (\delta \widehat{N})_{\mathrm{MC}}=269\pm 10`$ and $`\sigma (\delta \widehat{N})_{\mathrm{diag}}=279,`$ (62)
where the subscript ‘MC’ represents the results of the Monte Carlo simulation and ‘diag’ indicates the predictions of a direct diagonalization of $``$.
We give in the following some details and tricks about the numerical simulation. The first thing do to is to solve the steady-state Gross-Pitaevskii equation to find the function $`\varphi `$. This is done by imaginary time evolution. In practice we iterate the action of the evolution operator $`\mathrm{exp}\left(d\tau H_{\mathrm{gp}}\right)`$ (where $`H_{\mathrm{gp}}`$ is given in (16)) on an initial wavefunction, renormalizing the wavefunction and calculating $`\mu `$ at each time step $`d\tau `$, until $`\mu `$ reaches a stationary state within the required accuracy. If $`d\tau `$ is small enough, we can approximate the evolution of the wavefunction in each interval $`d\tau `$ by splitting the exponential $`\mathrm{exp}\left(d\tau H_{\mathrm{gp}}\right)`$ into the exponential of kinetic energy and the exponential of potential energy:
$$\varphi (\tau +d\tau )=\mathrm{exp}\left(d\tau H_{\mathrm{gp}}\right)\varphi (\tau )\mathrm{exp}\left(d\tau E_{\mathrm{cin}}\right)\mathrm{exp}\left(d\tau E_{\mathrm{pot}}\right)\varphi (\tau )$$
(64)
where:
$`E_{\mathrm{cin}}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}^2}{2m}}`$ (65)
$`E_{\mathrm{pot}}`$ $`=`$ $`U(z)+Ng|\varphi (z,\tau )|^2.`$ (66)
The advantage of the splitting in (64) is that the action of $`E_{\mathrm{cin}}`$ and $`E_{\mathrm{pot}}`$ on a vector can be easily calculated by going back and forth in the Fourier space and in the real space. This procedure is called the splitting method and is it a common technique to solve the steady-state Gross-Pitaevskii equation. Let us now consider the Brownian motion simulation: because of our choice (50) and (51) we have to calculate the action of the exponential of $`\beta /2`$ over a vector. It is convenient to write:
$$\mathrm{exp}\left(\frac{\beta }{2}\right)=\left[\mathrm{exp}\left(d\tau \right)\right]^n\text{with}d\tau =\frac{\beta /2}{n}.$$
(67)
The idea underlying equation (67) is to choose $`n`$ large enough so that the splitting approximation sketched above can be applied. Before implementing the splitting however, we still have to perform some transformation. We notice that if acting on functions orthogonal to the condensate wavefunctions, $`\mathrm{exp}\left(d\tau \right)`$ can be expressed as:
$$\mathrm{exp}\left(d\tau \right)=\left(\begin{array}{cc}Q& \text{0}\\ \text{0}& Q^{}\end{array}\right)\mathrm{exp}\left(d\tau _{gp}\right)$$
(68)
where $`_{gp}`$ is the operator obtained by linearization of the Gross-Pitaevskii equation, which is the same as $``$ given by (15) without the projectors. We are then led to the action of $`\mathrm{exp}\left(d\tau _{gp}\right)`$ on a vector and will do it approximately by using the splitting technique. In order to be consistent however we must take care that in our simulation the approximate action of $`_{gp}`$ corresponds exactly to the linearization of the approximate Gross-Pitaevskii equation in imaginary time (64) that we have used to find the condensate wavefunction. We then obtain:
$$\mathrm{exp}\left(d\tau _{gp}\right)\mathrm{exp}\left(d\tau E_{\mathrm{cin}}\right)\mathrm{exp}\left[d\tau (E_{\mathrm{pot}}\mu )\right]\mathrm{exp}\left[Ngd\tau \left(\begin{array}{cc}|\varphi |^2& \varphi ^2\\ \varphi ^{}^2& |\varphi |^2\end{array}\right)\right],$$
(69)
where the matrix in the last factor has the nice property of giving zero if squared.
A last point concerns the choice of the time step $`dt`$ of the Brownian motion, and of the total simulation time $`t_{\mathrm{max}}`$ guaranteeing that the steady-state is reached.
The time step $`dt`$ should satisfy $`\alpha _{\mathrm{max}}dt1`$ where $`\alpha _{\mathrm{max}}`$ is the largest eigenvalue of the friction matrix Eq.(50). To estimate $`\alpha _{\mathrm{max}}`$ we consider the eigenmode $`(u_{\mathrm{max}},v_{\mathrm{max}})`$ of $``$ with the largest eigenvalue $`e_{\mathrm{max}}`$. This value is well above the chemical potential $`\mu `$, so that $`v_{\mathrm{max}}`$ can be neglected as compared to $`u_{\mathrm{max}}`$, in the so-called Hatree-Fock approximation. In this case this largest energy eigenmode is an approximate eigenvector of the friction matrix, and we get the condition
$$\alpha _{\mathrm{max}}dtdt\mathrm{sinh}(\beta e_{\mathrm{max}})/\beta 1.$$
(70)
In a one-dimensional harmonic trap one expects $`e_{\mathrm{max}}e_{\mathrm{max}}^{\mathrm{ho}}\mu `$ in the Hartree-Fock approximation, where $`e_{\mathrm{max}}^{\mathrm{ho}}`$ is the maximal energy level of the bare trapping potential.
The condition to be satisfied by the duration $`t_{\mathrm{max}}`$ of the simulation is $`\alpha _{\mathrm{min}}t_{\mathrm{max}}1`$ where $`\alpha _{\mathrm{min}}`$ is the smallest eigenvalue of the friction matrix. We assume here that the corresponding eigenmode of $`\alpha `$ has components on energy modes of $``$ with energy much smaller than $`k_BT`$ only. In this case $`\alpha _{\mathrm{min}}`$ corresponds approximately to the minimum eigenvalue of $`\eta `$. In the case of real condensate wavefunction, the determination of the minimum eigenvalue of $`\eta `$ orthogonal to the condensate wavefunction boils down to the calculation of the first excited state energy of $`H_{gp}\mu `$ ; following we then find the condition:
$$t_{\mathrm{max}}\frac{\mathrm{}^2}{2m}\left(|\varphi |^2z^2\right)^11.$$
(71)
In the weakly interacting regime one finds $`t_{\mathrm{max}}\mathrm{}\omega `$. In the opposite Thomas-Fermi regime one finds that $`t_{\mathrm{max}}`$ scales as $`\mu \mathrm{}\omega `$. In practice, in the case of figure 1 we had: $`dt=1.15\times 10^4/\mathrm{}\omega `$ and $`t_{\mathrm{max}}=56.9/\mathrm{}\omega `$.
## 6 The time-Dependent Case
Let us consider the following situation: initially we have a Bose-Einstein condensate in a thermodynamically stable state at thermal equilibrium. Suddenly a parameter of the system is changed (e.g. the trapping potential or the scattering length) so that the system is no more in an equilibrium state and it undergoes a dynamical evolution during a given time, after which we perform a measurement.
In order to study this situation we need to recall some more results of the time-dependent number conserving Bogolubov theory of . It is shown in this paper that the expectation value of a one-particle observable $`\widehat{X}=_{i=1}^NX(i)`$ is given by:
$`X`$ $``$ $`(N\delta \widehat{N})\varphi |X(1)|\varphi +\varphi |X(1)|\varphi _{(2)}+\varphi _{(2)}|X(1)|\varphi `$ (72)
$`+{\displaystyle 𝑑\stackrel{}{r}𝑑\stackrel{}{r^{}}\stackrel{}{r}|X(1)|\stackrel{}{r^{}}\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r})\widehat{\mathrm{\Lambda }}(\stackrel{}{r^{}})}+O\left({\displaystyle \frac{1}{\sqrt{N}}}\right).`$
In this equation the first term contains the leading contribution to $`X`$ scaling as $`N`$ in the limit (9), where $`\varphi `$ is solution of the Gross-Pitaevskii equation (13), and
$$\delta \widehat{N}=𝑑\stackrel{}{r}\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r},t)\widehat{\mathrm{\Lambda }}(\stackrel{}{r},t),$$
(73)
so that $`N\delta \widehat{N}`$ is the mean number of particles in the condensate. The other terms involve the first correction to condensate wavefunction $`\varphi _{(2)}`$ and the contribution of the non condensed atoms to the expectation value of the observable. The equation of motion of $`\varphi _{(2)}`$ given in is:
$$\left(i\mathrm{}\frac{d}{dt}(t)\right)\left(\begin{array}{c}\hfill \varphi _{(2)}(t)\\ \hfill \varphi _{(2)}^{}(t)\end{array}\right)=\left(\begin{array}{c}\hfill Q(t)R(t)\\ \hfill Q^{}(t)R^{}(t)\end{array}\right)$$
(74)
where $``$ is given by Eq.(15) and
$`R(\stackrel{}{r})`$ $`=`$ $`gN|\varphi (\stackrel{}{r})|^2\varphi (\stackrel{}{r})1+{\displaystyle 𝑑\stackrel{}{s}\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{s})\widehat{\mathrm{\Lambda }}(\stackrel{}{s})}`$ (75)
$`+`$ $`2gN\varphi (\stackrel{}{r})\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{r})\widehat{\mathrm{\Lambda }}(\stackrel{}{r})+gN\varphi ^{}(\stackrel{}{r})\widehat{\mathrm{\Lambda }}(\stackrel{}{r})\widehat{\mathrm{\Lambda }}(\stackrel{}{r})`$
$``$ $`gN{\displaystyle 𝑑\stackrel{}{s}|\varphi (\stackrel{}{s})|^2\left[\widehat{\mathrm{\Lambda }}^{}(\stackrel{}{s})\varphi (\stackrel{}{s})+\widehat{\mathrm{\Lambda }}(\stackrel{}{s})\varphi ^{}(\stackrel{}{s})\right]\widehat{\mathrm{\Lambda }}(\stackrel{}{r})}.`$
The first term in Eq.(75) corrects the overestimation of the number of condensed particles in calculating their mutual interaction ($`NN(1+\delta \widehat{N}`$) in the Gross-Pitaevskii equation. The terms in the second line describe the interaction of the condensed particles and the non condensed ones. We refer to and for more details.
The implementation of our method for the time-dependent case then consists of three parts: (i) we generate many stochastic realizations of the initial state sampling the correct equilibrium distribution. (ii) We perform a deterministic evolution in real time of the stochastic initial states. (iii) We calculate the expectation value by averaging over the stochastic realizations.
Let us first concentrate on part (i): $`\varphi `$ which is the condensate wavefunction to the lowest approximation, is found by solving the Gross-Pitaevskii equation (13) in steady-state. Once we have $`\varphi `$, we can use the Monte Carlo simulation described in section 4 to generate many realizations of the stochastic functions $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ representing the non condensed fraction. To calculate $`\varphi _{(2)}`$, we calculate the source term (75) and solve (74) at steady-state with an imaginary time evolution:
$$\frac{d}{d\tau }\left(\begin{array}{c}\hfill \varphi _{(2)}(t)\\ \hfill \varphi _{(2)}^{}(t)\end{array}\right)=\eta (t=0)\left(\begin{array}{c}\hfill \varphi _{(2)}(t)\\ \hfill \varphi _{(2)}^{}(t)\end{array}\right)\left(\begin{array}{c}\hfill Q(t=0)R(t=0)\\ \hfill Q^{}(t=0)R^{}(t=0)\end{array}\right).$$
(76)
We note in this respect that the last term in (75) for a given realization of $`\mathrm{\Lambda }`$ factorizes as $`\mathrm{\Lambda }(\stackrel{}{r})`$ times an $`\stackrel{}{r}`$-independent integral over $`\stackrel{}{s}`$.
The real time evolution part (ii) is more straightforward since it is deterministic. The wavefunction $`\varphi `$ evolves according to (13). Each stochastic realization of $`\mathrm{\Lambda }`$, $`\mathrm{\Lambda }^{}`$ evolves according to (14), and the corrections $`\varphi _{(2)}`$, $`\varphi _{(2)}^{}`$ evolve according to (74). The averages in the source term $`R`$ involving $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ must be calculated at each time step.
## 7 Conclusions and Perspectives
We have put forward a Monte Carlo method to sample numerically the Gaussian density operator for the non condensed modes obtained in the Bogolubov approach, once it is expressed in term of the Wigner quasi-distribution function. This method allows to implement Bogolubov theory both in steady-state and in the time-dependent case without having to diagonalize big matrices, which makes it readily scalable to two or three spatial dimensions. We plan to use this method to study thermodynamics and finite temperature dynamics of Bose-Einstein condensates with vortices.
In this paper we have sticked to the Bogolubov theory as presented in . In this frame, for the time-dependent case, the Gross-Pitaevskii condensate wavefunction, the non condensed fraction and the first correction to the condensate wavefunction must be first obtained in the initial equilibrium state and then they must be evolved separately.
A natural development of the method beyond the Bogolubov approximation consists in a Monte Carlo simulation which samples the Wigner distribution for the whole field $`\psi `$ in equilibrium at temperature $`T`$ followed by the real time evolution of the field according to the classical field equation:
$$i\mathrm{}_t\psi =h_0\psi +g|\psi |^2\psi .$$
(77)
The classical field equation (77), which is formally identical to the Gross-Pitaevskii equation, corresponds in fact to the approximate evolution in real time of the Wigner quasi distribution function (truncated Wigner) .
The classical equations of motion for the field (77) has the disadvantage of considering only stimulated processes. Consider for example two colliding condensates in three dimensions, corresponding to an initial field of the form
$$\psi (\stackrel{}{r})=A(\mathrm{exp}(ikz)+\mathrm{exp}(ikz)).$$
(78)
By evolving this field according to (77) we will populate only modes that are plane waves of momentum $`\mathrm{}nk`$, $`n`$ integer (e.g. $`n=2`$ corresponds to the stimulated process $`k+kk+2k`$). In reality we expect to have binary collisions producing emerging atoms along any spatial direction. At the level of $`\psi `$ this would require a broken translational symmetry in the $`x`$-$`y`$ transverse plane. Fortunately, the advantage of being in the Wigner point of view is that all the modes are initially “filled” by quantum noise. In this case all processes can be stimulated and all possible symmetries of $`\psi `$ are broken.
We acknowledge useful discussions with Gora Shlyapnikov, Yuri Kagan and Anthony Leggett in the early stage of this work. Laboratoire Kastler Brossel is a unité de recherche de l’École normale supérieure et de l’Université Pierre et Marie Curie, associée au CNRS.
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# Perturbations of self-gravitating, ellipsoidal superfluid-normal fluid mixtures
## I Introduction
The problem of the equilibrium and stability of rotating neutron stars is encountered in various astrophysical contexts, ranging from the limiting frequencies of rapidly rotating isolated millisecond pulsars, emission of gravitational waves in neutron star-neutron star or neutron star-black hole binaries, to the generation of $`\gamma `$-rays in the neutron star mergers and $`X`$-rays in accreting systems. Considerable current interest is attached to the problem of neutron star binary inspiral, which would be the primary source of gravitational wave radiation for detection by future laser interferometers. Such a detection, apart from testing the general theory of relativity, potentially could provide useful information on the equation of state of superdense matter. Also the stability criteria for rapidly rotating neutron stars are essential for placing firm upper limits on the frequencies to which millisecond pulsars can be spun up thereby constraining the range of admissible equations of states.
High precision, fully relativistic treatments of rapidly rotating isolated neutron stars and binaries comprising two neutron stars have become available in recent years. Nevertheless, the development of simpler models that provide a fast and transparent insight into the underlying physics is needed when the basic set of equations is modified to include new effects.
A systematic framework for the treatment of the equilibrium and stability of rotating liquid masses bound by self-gravitation in the Newtonian theory is contained in Chandrasekhar’s Ellipsoidal Figures of Equilibrium (hereafter EFE). The tensor virial method, developed most extensively by Chandrasekhar and co-workers, transforms the local hydrodynamical equations into global virial equations that contain the full information on the structure and stability of the Newtonian self-gravitating system as a whole. The method describes, in a coherent manner, the properties of solitary ellipsoids with and without intrinsic spin, and ellipsoids in binaries subject to Newtonian tidal fields. It is especially useful for studying divergence-free displacements of uniform ellipsoids from equilibrium, in which case the each perturbed virial equation yields (in the absence of viscous dissipation) a different set of normal modes.
Recent alternative formulations of the theory of ellipsoids are based mainly on either the energy variation method or the affine star model or the (Eulerian) two potential formalism. A large class of incompressible and compressible ellipsoidal models has been studied using these methods. The energy variation method employs the observation that an equilibrium configuration is possible if the energy of the ellipsoid is an extremum for variations of the ellipsoidal semiaxis at a constant volume; the ellipsoidal figure is stable only if the energy is a true local minimum. In the affine star method the figures are described by a time-dependent Lagrangian as a function of a deformation matrix and its derivative. The structure of the star at any particular time is related conformally to the initial unperturbed sphere via a quadratic form constructed from the deformation matrix.
The purpose of this paper is to extend previous studies to a treatment of the oscillation modes of ellipsoidal figures of equilibrium that contain a mixture of normal fluid and superfluid. Many-body studies of the pair correlations in neutron star matter show that the baryon fluids in their ground state form superfluid condensates in the bulk of the star. The superfluid phases, in the hydrodynamic limit, can be treated as a mixture of superfluid condensate and normal matter. The superfluid rotation is supported by the Feynman-Onsager vortex lattice state, and on the average leads to quasi-rigid body rotation of the superfluid component. The corresponding time-dependent two-fluid hydrodynamic equations are completely specified by two independent velocities for the superfluid and the normal component and corresponding densities of the constituents. The two fluids are coupled to one another by mutual friction forces, which we model phenomenologically according the prescription given by Khalatnikov . Because of the additional degrees of freedom in this system, there are twice as many modes as for a single fluid. A natural question is whether the new modes affect the stability criteria previously deduced from studying perturbations of a single self-gravitating fluid.
In our treatment of perturbations of a mixture of normal fluid and superfluid, we shall follow most closely Chandrasekhar’s formulation. However, since the basic equations of motion for the two fluids will include mutual friction between them, the system we study is inherently dissipative. Nevertheless, since the frictional forces only depend linearly on the relative velocity between the two fluids and vanish in the background, where the two fluids move together, one can still derive relations that resemble Chandrasekhar’s tensor virial equations. Because these equations include dissipation, we shall prefer to regard them as moments of the fluid equations, rather than tensor virial equations. In fact, we shall relegate the derivation of perturbation equations from these moment equations to the Appendix of the paper, and instead derive the necessary equations for the fluid displacements directly by taking moments of linearized equations of motion for the two fluid components. In this paper, we concentrate on two-fluid variants of the classical Maclaurin, Jacobi and Roche ellipsoids.
Although our main aim is to access the oscillation modes and instabilities of neutron stars within the ellipsoidal approximation, the results obtained here may be of significance in other contexts (Ref. , the epilogue). One example is the understanding of rapidly rotating nuclei in the spirit of the Bohr-Wheeler model of a charged incompressible liquid droplet. In this case, the stability is determined by the competition between the attractive surface tension, the repulsive Coulomb potential, and the centrifugal stretching due to the rotation. Another example is the stability of rotating superfluid liquid drops of Bose condensed atomic gases, where the stability is determined through an interplay among the pressure of the condensate, the confining potential of the magnetic trap and the centrifugal potential .
Previous work on the oscillations of superfluid neutron stars concentrated mainly on perturbations of non-rotating or slowly rotating isolated neutron stars , and used methods entirely different from the one adopted here. The propagation of acoustic waves in neutron star interiors, including those related to the relative motion of neutron-proton superfluids, was studied by Epstein, who found the compressional and shear modes related to short-wavelength oscillations of neutron star matter. The small-amplitude pulsation modes of superfluid neutron stars were derived by Lindblom and Mendell , who found that the lowest frequency modes were almost indistinguishable from the normal modes of a single fluid star. Their analytical solutions also reveal the existence of a spectrum of modes which are absent in a single fluid star. Subsequent work concentrated on numerical solutions for the radial and non-radial pulsations of the two-fluid stars and identified distinct superfluid modes in the absence of rotation. The r modes of slowly rotating two-fluid neutron stars have been derived by Lindblom and Mendell , who find that they are identical to their ordinary-fluid counterparts to the lowest order in their small-angular-velocity expansion. The linear oscillations of general relativistic stars composed of two non-interacting fluids in a non-rotating static background have been studied by Comer et al. .
Our calculations allow arbitrary fast rotation, in the context of (incompressible) Newtonian fluid models. One may anticipate that the effects of superfluidity on oscillation modes, if any, should be affected by the underlying vortex structure of the rotating superfluid. In our treatment, dissipation arises because of the drag forces experienced by the vortex lines as they move through the normal fluid (and there is no dissipation if the drag force is zero). We ignore the motions related to the isospin degrees of freedom in the core of a neutron star and, hence, the mutual entrainment of the neutron and proton condensates, as well as forces arising from deviations from $`\beta `$ equilibrium. The two-fluid equations used in the remainder of this work can adequately describe the mutual friction of a two-condensate fluid in the core of a neutron star, since the entrainment effect renormalizes the effective superfluid densities and the frictional coefficients, i.e. the phenomenological input in the two-fluid equations. While the ellipsoidal approximation to a superfluid neutron star is restrictive, it allows us to study the effects of vorticity on the oscillation modes of a self-gravitating star in a transparent manner, avoiding complications due to the star’s inhomogeneity (multi-layer composition).
## II Perturbation Equations
The equations of motion for a mixture of two fluids may be summarized simply as
$$\rho _\alpha D_\alpha u_{\alpha ,i}=\frac{p_\alpha }{x_i}\rho _\alpha \frac{\varphi }{x_i}+\frac{1}{2}\rho _\alpha \frac{|𝛀\mathbf{\times }𝐱|^2}{x_i}+2\rho _\alpha ϵ_{ilm}u_{\alpha ,l}\mathrm{\Omega }_m+F_{\alpha \beta ,i},$$
(1)
where the subscript $`\alpha \{S,N\}`$ identifies the fluid component, and Latin subscripts denote coordinate directions; $`\rho _\alpha `$, $`p_\alpha `$, and $`𝐮_\alpha `$ are the density, pressure, and velocity of fluid $`\alpha `$, $`\varphi `$ is the gravitational potential, and $`𝐅_{\alpha \beta }`$ is the mutual friction force on fluid $`\alpha `$ due to fluid $`\beta `$. These equations have been written in a frame rotating with angular velocity $`𝛀`$ relative to some inertial coordinate reference system. The total time derivative operator
$$D_\alpha \frac{}{t}+u_{\alpha ,j}\frac{}{x_j}.$$
(2)
The gravitational potential $`\varphi `$ is derived from
$$^2\varphi =^2(\varphi _S+\varphi _N)=4\pi G[\rho _S(𝐱)+\rho _N(𝐱)];$$
(3)
the individual fluid potentials $`\varphi _\alpha `$ obey $`^2\varphi _\alpha =4\pi G\rho _\alpha `$. The two fluids are coupled to one another via the frictional force $`𝐅_{\alpha \beta }`$ which is antisymmetric on interchange of $`\alpha `$ and $`\beta `$. For a normal-superfluid mixture
$$𝐅_{SN}=𝐅_{NS}\rho _S\omega _S\{\beta ^{}\nu \mathbf{\times }(𝐮_S𝐮_N)+\beta 𝝂\mathbf{\times }[𝝂\mathbf{\times }(𝐮_S𝐮_N)]\beta ^{\prime \prime }𝝂\mathbf{}(𝐮_S𝐮_N)\},$$
(4)
where $`\beta `$, $`\beta ^{}`$ and $`\beta ^{\prime \prime }`$ are coupling coefficients, and $`𝝎_S=𝝂\omega _S\mathbf{}\mathbf{\times }𝐮_S`$; in components we have
$$F_{SN,i}=\rho _S\omega _S\beta _{ij}(u_{S,j}u_{N,j}),$$
(5)
where, from Eq. (4),
$$\beta _{ij}=\beta \delta _{ij}+\beta ^{}ϵ_{ijm}\nu _m+(\beta ^{\prime \prime }\beta )\nu _i\nu _j.$$
(6)
The net rate at which this force does work is
$$𝐮_S\mathbf{}𝐅_{SN}+𝐮_N\mathbf{}𝐅_{NS}=\rho _S\omega _S\{\beta |𝝂\mathbf{\times }(𝐮_S𝐮_N)|^2\beta ^{\prime \prime }[𝝂\mathbf{}(𝐮_S𝐮_N)]^2\};$$
(7)
there is no dissipation associated with the term proportional to $`\beta ^{}`$ in $`𝐅_{\alpha \beta }`$. Throughout this paper, we assume that $`\beta `$, $`\beta ^{}`$ and $`\beta ^{\prime \prime }`$ are independent of position in the fluid mixture. In Eq. (4), we have neglected the effects of the vortex tension, and expressed the mutual friction force in terms of the phenomenological coefficients $`\beta `$, $`\beta ^{}`$ and $`\beta ^{\prime \prime }`$. While these parameters determine the macroscopic behavior of the fluid system, they are not the optimal ones for connecting microscopic parameters of the mixture to its macroscopic motion. Instead, the macroscopic results can be parametrized in terms of frictional coefficients $`\eta `$ and $`\eta ^{}`$, which connect $`\beta `$ and $`\beta ^{}`$ to the drag on individual superfluid vortices via the relations
$$\beta =\frac{\eta \rho _S\omega _S}{\eta ^2+(\rho _S\omega _S\eta ^{})^2}\beta ^{}=1\frac{\rho _S\omega _S(\rho _S\omega _S\eta ^{})}{\eta ^2+(\rho _S\omega _S\eta ^{})^2}.$$
(8)
The physical meaning of $`\eta `$’s is apparent from the equation of motion of a single vortex line
$$\rho _S\omega _S\left[\left(𝐮_S𝐮_L\right)\times 𝝂\right]\eta \left(𝐮_L𝐮_N\right)+\eta ^{}\left[\left(𝐮_L𝐮_N\right)\times 𝝂\right]=0,$$
(9)
where $`𝐮_L`$ is the vortex velocity. Equation (9) states that the Magnus force, which represents a lifting force due to the superflow imposed on the vortex circulation, is balanced by the viscous friction forces along the vortex motion (the term $`\eta `$) and perpendicular to the vortex motion (the term $`\eta ^{}`$); these latter forces arise from the scattering of the normal quasiparticles off the vortex line<sup>*</sup><sup>*</sup>*The inertial mass of the vortex is neglected in the standard formulation of the two-fluid superfluid hydrodynamics .. The characteristic dynamical relaxation time scale related to the vortex motion can be defined as
$$\tau _\mathrm{D}=\frac{1}{\omega _S}\left(\frac{\eta }{\rho _S\omega _S}+\frac{\rho _S\omega _S}{\eta }\right),$$
(10)
where $`\omega _S`$ is the superfluid circulation averaged over macroscopic scales; e.g. for uniformly rotating superfluid $`\omega _S=2\mathrm{\Omega }_S`$, where $`\mathrm{\Omega }_S`$ is the superfluid rotation frequency. For fixed density $`\rho _S`$, $`\tau _D\mathrm{}`$ asymptotically, when $`\eta \rho _S\omega _S`$ (strong coupling limit) and $`\rho _S\omega _S\eta `$ (weak coupling limit). Its minimal value is attained when $`\eta /\rho _S\omega _S=1`$.
The relations (8) do not determine $`\beta ^{\prime \prime }`$, which, if nonzero, implies friction along vortex lines, which would be possible if vortices oscillate or are deformed in the plane perpendicular to the rotation axis. Generally, we assume in this paper that $`\beta ^{\prime \prime }\beta `$ and $`\beta ^{}`$, but occasionally we retain nonzero (and not necessarily negligible) $`\beta ^{\prime \prime }`$ to examine its effects on the modes. Moreover, if we assume that the drag force on a vortex line is principally opposite to its velocity relative to the normal fluid, then $`\eta ^{}\eta `$, and the relationships between $`\beta `$ and $`\beta ^{}`$ and the drag force on vortex lines simplifies accordingly.
Ultimately, we want to compute perturbations of Eq. (1) around some presumed background state. In general, the two fluids need not occupy the same volume, and we shall suppose that fluid $`\alpha `$ occupies a volume $`V_\alpha `$. However, we shall restrict ourselves to background states in which the two fluids occupy identical volumes and have densities $`\rho _\alpha (𝐱)=f_\alpha \rho (𝐱)`$, where $`\rho (𝐱)`$ is the total density and $`f_\alpha `$ does not depend on $`𝐱`$. The perturbation equations derived below can be applied to nonuniform background states satisfying these conditions, but we shall only consider uniformly dense backgrounds, as was done by Chandrasekhar. Note, though, that we shall not assume that the two fluids must occupy the same volumes in the perturbed state. In fact, at least for low order perturbations, there are no nontrivial perturbations that leave the volumes of the two fluids identical.
Following the example set by Chandrasekhar, we could take moments of the fluid equations (1) to obtain tensor virial theorems of various orders, and then perturb them to find linear modes for uniform ellipsoids. We have derived the necessary moment equations in this way (see the Appendix), by analogy to Chandrasekhar’s treatment for a single fluid, but here we present a somewhat different (and possibly more transparent approach) to their derivation. We begin with the equation of motion for the displacement of fluid $`\alpha `$ from equilibrium, a direct generalization of Eq. (107) in Chap. 2, Sec. 14 in EFE:
$`\rho _\alpha {\displaystyle \frac{d^2\xi _{\alpha ,i}}{dt^2}}`$ $`=`$ $`{\displaystyle \frac{\xi _{\alpha ,l}}{x_l}}{\displaystyle \frac{p_\alpha }{x_i}}{\displaystyle \frac{\mathrm{\Delta }_\alpha p_\alpha }{x_i}}\rho _\alpha {\displaystyle \frac{\mathrm{\Delta }_\alpha \varphi }{x_i}}+\rho _\alpha {\displaystyle \frac{}{x_i}}\left[{\displaystyle \frac{\xi _{\alpha ,l}}{2}}\left({\displaystyle \frac{|𝛀\mathbf{\times }𝐱|^2}{x_l}}\right)\right]+2\rho _\alpha ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{d\xi _{\alpha ,l}}{dt}}`$ (12)
$`+{\displaystyle \frac{\xi _{\alpha ,l}}{x_i}}\left({\displaystyle \frac{p_\alpha }{x_l}}+\rho _\alpha {\displaystyle \frac{\varphi }{x_l}}{\displaystyle \frac{\rho _\alpha }{2}}{\displaystyle \frac{|𝛀\mathbf{\times }𝐱|^2}{x_l}}\right)+F_{\alpha \beta ,i},`$
where the Lagrangian variation $`\mathrm{\Delta }_\alpha Q`$ denotes the change in $`Q`$ seen by a moving element of fluid $`\alpha `$; in particular, it is easy to show that $`\mathrm{\Delta }_\alpha \varphi (𝐱)`$ separates into an Eulerian part, $`\delta \varphi `$, plus $`\xi _{\alpha ,l}\varphi /x_l`$, where
$$\delta \varphi G\underset{\gamma =\alpha ,\beta }{}_Vd^3x^{}\rho _\gamma (𝐱^{})\xi _{\gamma ,l}(𝐱^{})\frac{}{x_l^{}}\left(\frac{1}{|𝐱𝐱^{}|}\right)$$
(13)
is the Eulerian potential perturbation. We assume that the mutual friction force is $`F_{\alpha \beta ,i}=0`$ in the background solution, i.e., either the two fluids are stationary in the rotating frame for the background or have identical fluid velocities in this frame. Our strategy will be to take moments of Eq. (12) by multiplying by appropriate factors of $`x_i`$ and integrating over the (common) volume of the unperturbed background configuration. Note that this does not restrict the volumes of the perturbed fluids in any way.
Although the method we shall use to derive the perturbations resembles Chandrasekhar’s, there is an important distinction due to the mutual friction force. Chandrasekhar’s method of solution yields exact modes only in the dissipationless limit, where the modes of uniform ellipsoids are, respectively, linear, quadratic, cubic, etc. functions of the coordinates. Viscous terms would prevent exact solution in this manner, and one resorts to an approximation in which they are evaluated by substitution of the inviscid eigenfunctions (EFE, Chap. 5, §37(b)). However, it is possible to employ the moment method to find exact modes for the normal fluid-superfluid mixture coupled by mutual friction because $`𝐅_{\alpha \beta }`$ is a linear function of the velocity difference between the two fluids, not of their spatial derivatives (as is the case for viscous dissipation). By taking moments, we can derive the analogue of tensor virial theorems of various orders, but because the equations involve manifestly dissipative mutual friction forces, we prefer to think of these merely as moments of the original equations of motion. Of course, if we also choose to include the effects of viscous dissipation in the normal fluid, which we shall not do in this paper, we must resort to a low Reynolds number approximation, as was done in EFE. We discuss this briefly in Sec. IV.
Using the above result for the perturbed gravitational potential, we can rewrite Eq. (12) as
$`\rho _\alpha {\displaystyle \frac{d^2\xi _{\alpha ,i}}{dt^2}}`$ $`=`$ $`{\displaystyle \frac{\xi _{\alpha ,l}}{x_l}}{\displaystyle \frac{p_\alpha }{x_i}}\xi _{\alpha ,l}{\displaystyle \frac{^2p_\alpha }{x_ix_l}}{\displaystyle \frac{\delta p_\alpha }{x_i}}\rho _\alpha {\displaystyle \frac{\delta \varphi }{x_i}}`$ (15)
$`\rho _\alpha \xi _{\alpha ,l}{\displaystyle \frac{^2\varphi }{x_ix_l}}+\rho _\alpha \xi _{\alpha ,l}(\mathrm{\Omega }^2\delta _{il}\mathrm{\Omega }_i\mathrm{\Omega }_l)+2\rho _\alpha ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{d\xi _{\alpha ,l}}{dt}}+F_{\alpha \beta ,i}.`$
Then if we define
$$𝝃_+=f_S𝝃_S+f_N𝝃_N𝝃_{}=𝝃_S𝝃_N,$$
(16)
we find
$`\rho {\displaystyle \frac{d^2\xi _{+,i}}{dt^2}}`$ $`=`$ $`{\displaystyle \frac{\xi _{+,l}}{x_l}}{\displaystyle \frac{p}{x_i}}\xi _{+,l}{\displaystyle \frac{^2p}{x_ix_l}}{\displaystyle \frac{\delta p}{x_i}}\rho {\displaystyle \frac{\delta \varphi }{x_i}}`$ (18)
$`\rho \xi _{+,l}{\displaystyle \frac{^2\varphi }{x_ix_l}}+\rho \xi _{+,l}(\mathrm{\Omega }^2\delta _{il}\mathrm{\Omega }_i\mathrm{\Omega }_l)+2\rho ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{d\xi _{+,l}}{dt}}`$
where $`\delta p=\delta p_S+\delta p_N`$, and
$`\rho {\displaystyle \frac{d^2\xi _{,i}}{dt^2}}`$ $`=`$ $`{\displaystyle \frac{\xi _{,l}}{x_l}}{\displaystyle \frac{p}{x_i}}\xi _{,l}{\displaystyle \frac{^2p}{x_ix_l}}{\displaystyle \frac{1}{f_S}}{\displaystyle \frac{\delta p_S}{x_i}}+{\displaystyle \frac{1}{f_N}}{\displaystyle \frac{\delta p_N}{x_i}}`$ (21)
$`+\rho \xi _{,l}(\mathrm{\Omega }^2\delta _{il}\mathrm{\Omega }_i\mathrm{\Omega }_l)\rho \xi _{,l}{\displaystyle \frac{^2\varphi }{x_ix_l}}+2\rho ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{d\xi _{,l}}{dt}}`$
$`\rho \omega _S\left(1+{\displaystyle \frac{f_S}{f_N}}\right)\beta _{ik}{\displaystyle \frac{d\xi _{,k}}{dt}}.`$
Equation (18) is identical to what is found for a single fluid, and therefore contains the well-known modes documented by Chandrasekhar: if we define
$$V_{i;j}=_Vd^3x\rho \xi _{+,i}x_j,$$
(22)
then we find
$`{\displaystyle \frac{d^2V_{i;j}}{dt^2}}`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{dV_{l;j}}{dt}}+\mathrm{\Omega }^2V_{ij}\mathrm{\Omega }_i\mathrm{\Omega }_kV_{kj}+\delta _{ij}\delta \mathrm{\Pi }`$ (23)
$``$ $`\pi G\rho \left(2B_{ij}V_{ij}a_i^2\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}V_{ll}\right),`$ (24)
where $`\delta \mathrm{\Pi }\delta \mathrm{\Pi }_S+\delta \mathrm{\Pi }_N`$ and all other quantities are defined exactly as in EFE. All of the new modes of a mixture of normal fluid and superfluid are contained in Eq. (21) for their relative displacements. One noteworthy feature of Eq. (21) is that the Eulerian gravitational potential does not appear. Consequently, the new normal modes of the system only depend on the unperturbed gravitational potential; for perturbations of homogeneous ellipsoids, only the coefficients $`A_i`$ defined by EFE, Chap. 3, Eqs. (18) and (40), will appear.
For the most part, we shall be interested in displacements that are linear functions of $`x_i`$ in this paper. For the homogeneous ellipsoids, we can find the new modes that result from the differential displacements of normal fluid and superfluid by taking the first moment of Eq. (21), e.g. by multiplying by $`x_j`$ and integrating over the unperturbed volume.In the nomenclature of EFE, this is the second order virial equation corresponding to Eq. (21). If we define
$$U_{i;j}=_Vd^3x\rho \xi _{,i}x_j,$$
(25)
then we find
$`{\displaystyle \frac{d^2U_{i;j}}{dt^2}}`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{dU_{l;j}}{dt}}+\mathrm{\Omega }^2U_{ij}\mathrm{\Omega }_i\mathrm{\Omega }_kU_{kj}+\delta _{ij}\left({\displaystyle \frac{\delta \mathrm{\Pi }_S}{f_S}}{\displaystyle \frac{\delta \mathrm{\Pi }_N}{f_N}}\right)`$ (26)
$``$ $`2\pi G\rho A_iU_{ij}\omega _S\left(1+{\displaystyle \frac{f_S}{f_N}}\right)\beta _{ik}{\displaystyle \frac{dU_{k;j}}{dt}},`$ (27)
where
$$\delta \mathrm{\Pi }_\alpha _Vd^3x\delta p_\alpha .$$
(28)
To obtain Eq. (27), various surface terms can be eliminated using the conditions that $`p_\alpha `$ and $`\mathrm{\Delta }_\alpha p_\alpha =\delta p_\alpha +\xi _{\alpha ,l}p_\alpha /x_l`$ vanish on the boundary; also, the equation of hydrostatic equilibrium for the unperturbed configuration,
$$0=\frac{p}{x_i}+\rho \frac{\varphi }{x_i}\rho \frac{}{x_i}\left(\frac{|𝛀\mathbf{\times }𝐱|^2}{2}\right),$$
(29)
must be used. It is straightforward to compute higher moments of Eq. (21). For example, taking its second momentThis is the third order virial equation in the nomenclature of EFE. by multiplying by $`x_jx_k`$ and integrating over the unperturbed volume gives
$`{\displaystyle \frac{d^2U_{i;jk}}{dt^2}}`$ $`=`$ $`\delta _{ij}\left({\displaystyle \frac{\delta \mathrm{\Pi }_{S,k}}{f_S}}{\displaystyle \frac{\delta \mathrm{\Pi }_{N,k}}{f_N}}\right)+\delta _{ik}\left({\displaystyle \frac{\delta \mathrm{\Pi }_{S,j}}{f_S}}{\displaystyle \frac{\delta \mathrm{\Pi }_{N,j}}{f_N}}\right)`$ (32)
$`2\pi G\rho A_iU_{ijk}+(\mathrm{\Omega }^2\delta _{il}\mathrm{\Omega }_i\mathrm{\Omega }_l)U_{ljk}`$
$`+\left[2ϵ_{ilm}\mathrm{\Omega }_m\omega _S\left(1+{\displaystyle \frac{f_S}{f_N}}\right)\beta _{il}\right]{\displaystyle \frac{dU_{l;jk}}{dt}},`$
where, by analogy to definitions in EFE for a single fluid,
$`U_{i;jk}`$ $`=`$ $`{\displaystyle _V}d^3x\rho \xi _{,i}x_jx_k`$ (33)
$`U_{ijk}`$ $`=`$ $`U_{i;jk}+U_{j;ki}+U_{k;ij}`$ (34)
$`\delta \mathrm{\Pi }_{\alpha ,k}`$ $`=`$ $`{\displaystyle _V}d^3xx_k\delta p_\alpha .`$ (35)
Equations (27) and (32) only apply to unperturbed states that are static in the rotating frame.
## III New Modes of a Two-Fluid Mixture
In this section, we derive the characteristic equations for the normal modes for relative fluid displacements implied by Eq. (27) assuming $`U_{i;j}\mathrm{exp}(\lambda t)`$. We list the results separately for perturbations of Maclaurin, Jacobi and Roche ellipsoids. Note that the modes described by Eq. (24) are identical to those treated in EFE.
The superfluid state of the matter does not affect the equilibrium figure.<sup>§</sup><sup>§</sup>§Here, we have taken $`p_\alpha =f_\alpha p`$ in the background state, which is a mathematically convenient idealization. More realistically, contributions from the pair condensation energy and the energy density of the superfluid vortex lattice, which distinguish the superfluid from the normal fluid, could play a role in both the equilibria and the perturbations. We need not restrict our attention to displacements for which the pressure perturbations of the two fluids are still proportional to one another, although, as argued in §IV, for the solenoidal displacements considered here, this turns out to be the case. In equilibrium, Eq. (27) is satisfied trivially, for all virials $`U_{ij}=0`$ in the absence of relative motion between the superfluid and normal components. The equilibrium figure follows from Eq. (24) by dropping the temporal variations and is identical to its non-superfluid counterpart.
For irrotational ellipsoids we can set $`\omega _S=2\mathrm{\Omega }`$ in Eq. (27) and using the compact notations
$$\delta \stackrel{~}{\mathrm{\Pi }}=\left(\frac{\delta \mathrm{\Pi }_S}{f_S}\frac{\delta \mathrm{\Pi }_N}{f_N}\right),\stackrel{~}{\beta }_{ik}=\left(1+\frac{f_S}{f_N}\right)\beta _{ik},$$
(36)
we rewrite the Eq. (37) as
$`{\displaystyle \frac{d^2U_{i;j}}{dt^2}}`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{dU_{l;j}}{dt}}+\mathrm{\Omega }^2U_{ij}\mathrm{\Omega }_i\mathrm{\Omega }_kU_{kj}+\delta _{ij}\delta \stackrel{~}{\mathrm{\Pi }}2\pi G\rho A_iU_{ij}2\mathrm{\Omega }\stackrel{~}{\beta }_{ik}{\displaystyle \frac{dU_{k;j}}{dt}}.`$ (37)
For time-dependent Lagrangian displacements of the form
$$𝝃_\alpha (x_i,t)=𝝃_\alpha (x_i)e^{\lambda t},$$
(38)
Eq. (37) becomes
$`\lambda ^2U_{i;j}2ϵ_{ilm}\mathrm{\Omega }_m\lambda U_{l;j}`$ $`=`$ $`+\mathrm{\Omega }^2U_{ij}\mathrm{\Omega }_i\mathrm{\Omega }_kU_{kj}+(\pi \rho G)^1\delta _{ij}\delta \stackrel{~}{\mathrm{\Pi }}2A_iU_{ij}2\mathrm{\Omega }\lambda \stackrel{~}{\beta }_{ik}U_{k;j};`$ (39)
here all frequencies are measured in the units $`(\pi \rho G)^{1/2}`$. Equation (39) contains all the second harmonic modes of the relative oscillation of Maclaurin and Jacobi ellipsoids, and only requires a slight modification for the application to Roche ellipsoids.
### A Superfluid Maclaurin spheroid
Next, we specialize Eq. (37) to Maclaurin spheroids, the equilibrium figures of a self-gravitating fluid with two equal semi-major axis, say $`a_1`$ and $`a_2`$, uniformly rotating about the third semi-major axis $`a_3`$ (i.e. the $`x_3`$ axis). The sequence of quasi-equilibrium figures can be parametrized by the eccentricity $`ϵ^2=1a_3^2/a_1^2`$, with (squared) angular velocity $`\mathrm{\Omega }^2=2ϵ^2B_{13}`$, in units of $`(\pi \rho G)^{1/2}`$.
Surface deformations related to various modes can be classified by corresponding terms of the expansion in surface harmonics labeled by the indexes $`l,m`$. Second order harmonic deformations correspond to $`l=2`$ with five distinct values of $`m`$, $`2m2`$. The 18 equations represented by Eq. (37) separate into two independent subsets which are odd and even with respect to index 3. The corresponding oscillation modes can be treated separately.
#### 1 Relative transverse shear modes
These modes correspond to surface deformations with $`|m|=1`$ and represent relative shearing of the northern and southern hemispheres of the ellipsoid. The components of Eq. (39), which are odd in index 3, are
$`\lambda ^2U_{3;1}`$ $`=`$ $`2A_3U_{31}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;1},`$ (40)
$`\lambda ^2U_{3;2}`$ $`=`$ $`2A_3U_{32}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;2},`$ (41)
$`\lambda ^2U_{1;3}2\mathrm{\Omega }\lambda U_{2;3}`$ $`=`$ $`\left(2A_1+\mathrm{\Omega }^2\right)U_{13}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;3}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;3},`$ (42)
$`\lambda ^2U_{2;3}+2\mathrm{\Omega }\lambda U_{1;3}`$ $`=`$ $`\left(2A_1+\mathrm{\Omega }^2\right)U_{23}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;3}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;3}.`$ (43)
Note that because of the degeneracy in indexes 1 and 2 for the Maclaurin spheroid $`A_1=A_2`$. We sum Eqs. (40), (42) and (41), (43), respectively, and use the symmetry properties of $`U_{ij}`$ combined with Eqs. (41), (43). We find
$`\left[\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)+2\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)A_3+\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)(2A_1\mathrm{\Omega }^2)\right]U_{13}`$ (44)
$`2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3\right)U_{23}=0,`$ (45)
$`\left[\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)+2\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)A_3+\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)(2A_1\mathrm{\Omega }^2)\right]U_{23}`$ (46)
$`+2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3\right)U_{13}=0.`$ (47)
It is instructive to consider first the limit of zero mutual friction, in which case Eqs. (45)-(47) reduce to
$`\lambda \left[\lambda ^2+(2A_1+2A_3\mathrm{\Omega }^2)\right]U_{13}2\mathrm{\Omega }\left(\lambda ^2+2A_3\right)U_{23}=0,`$ (48)
$`\lambda \left[\lambda ^2+(2A_1+2A_3\mathrm{\Omega }^2)\right]U_{23}+2\mathrm{\Omega }\left(\lambda ^2+2A_3\right)U_{13}=0,`$ (49)
excluding the trivial mode $`\lambda =0`$. The characteristic equation can be factorized by substituting $`\lambda =i\sigma `$
$$\sigma \left[\sigma ^22(A_1+A_3)+\mathrm{\Omega }^2\right]\pm 2\mathrm{\Omega }(\sigma ^22A_3)=0.$$
(50)
The purely rotational mode $`\sigma =\mathrm{\Omega }`$ decouples only in the spherical symmetric limit where $`A_1=A_3`$. If only axial symmetry is imposed then the characteristic equation is third order:
$`\sigma ^3\pm 2\mathrm{\Omega }\sigma ^2+\left[2\left(A_1+A_3\right)+\mathrm{\Omega }^2\right]\sigma 4A_3\mathrm{\Omega }=0.`$ (51)
Along the entire sequence parametrized in terms of the eccentricity the three modes are real This is easy to prove directly from Eq. (49). Write the dispersion relation as $`f(\lambda ^2)=0`$. Then show that (i) $`f(\lambda ^2)\pm \mathrm{}`$ as $`\lambda ^2\pm \mathrm{}`$, (ii) $`f(0)>0`$, and (iii) the two extrema of $`f(\lambda ^2)`$ are both at $`\lambda ^2<0`$. Thus, the zeros of $`f(\lambda ^2)`$ are all at $`\lambda ^2<0`$, so $`\sigma =i\lambda `$ must be real. and are given by
$`\sigma _1={\displaystyle \frac{2\mathrm{\Omega }}{3}}+(s_++s_{}),\sigma _{2,3}={\displaystyle \frac{2\mathrm{\Omega }}{3}}{\displaystyle \frac{1}{2}}(s_++s_{})\pm {\displaystyle \frac{i\sqrt{3}}{2}}(s_+s_{}),`$ (52)
where
$`s_\pm ^3`$ $`=`$ $`{\displaystyle \frac{8\mathrm{\Omega }^3}{27}}+{\displaystyle \frac{2\mathrm{\Omega }(A_12A_3)\mathrm{\Omega }^3}{3}}`$ (53)
$`\pm `$ $`\left[\left({\displaystyle \frac{4\mathrm{\Omega }^2}{9}}{\displaystyle \frac{2\left(A_1+A_3\right)\mathrm{\Omega }^2}{3}}\right)^3+\left({\displaystyle \frac{8\mathrm{\Omega }^3}{27}}+{\displaystyle \frac{2\mathrm{\Omega }(A_12A_3)\mathrm{\Omega }^3}{3}}\right)^2\right]^{1/2}.`$ (54)
Three complementary modes follow from Eqs.(52)-(53) via the replacement $`\mathrm{\Omega }\mathrm{\Omega }`$. The frictionless modes are real and are shown in Fig. 1. The two high frequency, frictionless modes are roughly twice as large as the transverse shear modes of ordinary Maclaurin spheroids. The third low frequency mode corresponds to the nearly rotational mode and indeed coincides with $`\mathrm{\Omega }`$ in the limits $`ϵ0`$ and $`A_1A_3`$, but not generally. In the dissipative case the characteristic equation is of sixth order:
$`\lambda ^6+4\mathrm{\Omega }(\stackrel{~}{\beta }+\stackrel{~}{\beta }^{\prime \prime })\lambda ^5+\left[4A_3+2\mathrm{\Omega }(1\stackrel{~}{\beta }^{})^2+4\mathrm{\Omega }^2(\stackrel{~}{\beta }+\stackrel{~}{\beta }^{\prime \prime })^22(2A_1\mathrm{\Omega }^2)\right]\lambda ^4`$ (55)
$`+[16A_3\stackrel{~}{\beta }\mathrm{\Omega }+8A_3\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }+8\stackrel{~}{\beta }^{\prime \prime }(1\stackrel{~}{\beta }^{})^2\mathrm{\Omega }^2+16\stackrel{~}{\beta }^2\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^3`$ (56)
$`+16\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^34\stackrel{~}{\beta }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)8\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)]\lambda ^3`$ (57)
$`+[4A_3^2+8A_3\mathrm{\Omega }(1\stackrel{~}{\beta }^{})^2+16A_3\stackrel{~}{\beta }^2\mathrm{\Omega }^2+32A_3\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2+8\stackrel{~}{\beta }^{\prime \prime 2}(1\stackrel{~}{\beta }^{})^2\mathrm{\Omega }^3+16\stackrel{~}{\beta }^2\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^4`$ (58)
$`4A_3(2A_1\mathrm{\Omega }^2)16\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)8\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)+(2A_1\mathrm{\Omega }^2)^2]\lambda ^2`$ (59)
$`+[16A_3^2\stackrel{~}{\beta }\mathrm{\Omega }+16A_3\stackrel{~}{\beta }^{\prime \prime }(1\stackrel{~}{\beta }^{})^2\mathrm{\Omega }^2+32A_3\stackrel{~}{\beta }^2\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^38A_3\stackrel{~}{\beta }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)`$ (60)
$`8A_3\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)16\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^3(2A_1\mathrm{\Omega }^2)+4\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)^2]\lambda `$ (61)
$`+\left[8A_3^2\mathrm{\Omega }(1\stackrel{~}{\beta }^{})^2+16A_3^2\stackrel{~}{\beta }^2\mathrm{\Omega }^216A_3\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)+4\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)^2\right]=0.`$ (62)
The real and imaginary parts of the relative transverse shear modes are shown in the Fig. 1 for several values of $`\eta `$ and $`\eta ^{}=0=\beta ^{\prime \prime }`$ (here and below we scale $`\eta `$ in units of $`\rho _S\omega _S`$). The real part of the low frequency rotational modes is diminished as $`\eta `$ is increased; the high frequency modes are unaffect except in the stong coupling limit $`\eta 50`$ \[the strong and weak coupling limits are discussed after Eq. (10)\]. The damping of the modes is maximal for $`\eta =1`$ and decreases to zero for $`\eta 0`$ and $`\eta \mathrm{}`$. Note that in the limiting cases the vortices are locked either in the superfluid ($`\eta 0`$) or the normal component ($`\eta \mathrm{}`$) and hence the damping is ineffective. The transverse shear modes are stable for arbitrary values of the eccentricity of the spheroid.
#### 2 Relative toroidal modes
These modes correspond to $`|m|=2`$ and the motions in this case are confined to the planes parallel to the equatorial plane. The components of Eq. (39), which are even in index 3, are:
$`\lambda ^2U_{3;3}`$ $`=`$ $`(\pi \rho G)^1\delta \stackrel{~}{\mathrm{\Pi }}2A_3U_{33}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;3},`$ (63)
$`\lambda ^2U_{1;1}2\mathrm{\Omega }\lambda U_{2;1}`$ $`=`$ $`(\pi G\rho )^1\delta \stackrel{~}{\mathrm{\Pi }}+(\mathrm{\Omega }^22A_1)U_{11}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;1}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;1},`$ (64)
$`\lambda ^2U_{2;2}+2\mathrm{\Omega }\lambda U_{1;2}`$ $`=`$ $`(\pi G\rho )^1\delta \stackrel{~}{\mathrm{\Pi }}+(\mathrm{\Omega }^22A_1)U_{22}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;2}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;2},`$ (65)
$`\lambda ^2U_{1;2}2\mathrm{\Omega }\lambda U_{2;2}`$ $`=`$ $`(2A_1+\mathrm{\Omega }^2)U_{12}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;2}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;2},`$ (66)
$`\lambda ^2U_{2;1}+2\mathrm{\Omega }\lambda U_{1;1}`$ $`=`$ $`(2A_1+\mathrm{\Omega }^2)U_{21}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;1}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;1}.`$ (67)
We add Eqs. (66) and (67) and subtract Eqs. (64) and (65) to find the following coupled equations for the toroidal modes (note that $`A_1=A_2`$ for Maclaurin spheroids)
$`\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda +4A_12\mathrm{\Omega }^2\right)(U_{11}U_{22})4\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})U_{12}`$ $`=`$ $`0,`$ (68)
$`\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda +4A_12\mathrm{\Omega }^2\right)U_{12}+\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})(U_{11}U_{22})=0.`$ (69)
The characteristic equation for the toroidal modes is
$`\lambda ^4+4\stackrel{~}{\beta }\lambda ^3\mathrm{\Omega }+\lambda ^2(8A_1+4\stackrel{~}{\beta }^2\mathrm{\Omega }^28\stackrel{~}{\beta }^{}\mathrm{\Omega }^2+4\stackrel{~}{\beta }^2\mathrm{\Omega }^2)`$ (70)
$`+\lambda (16A_1\stackrel{~}{\beta }\mathrm{\Omega }8\stackrel{~}{\beta }\mathrm{\Omega }^3)+16A_1^216A_1\mathrm{\Omega }^2+4\mathrm{\Omega }^4=0.`$ (71)
In the frictionless limit this can be written as
$$(\lambda ^2+4A_12\mathrm{\Omega }^2)^2+4\mathrm{\Omega }^2\lambda ^2=0,$$
(72)
which is factorized by writing $`\lambda =i\sigma `$. The two solutions are then
$$\sigma _{1,2}=\mathrm{\Omega }\pm \sqrt{4A_1\mathrm{\Omega }^2},$$
(73)
and there are two complementary modes which are found by substituting $`\mathrm{\Omega }`$ for $`\mathrm{\Omega }`$. The modes are always real because $`4A_1>\mathrm{\Omega }^2`$ for incompressible Maclaurin spheroids. This result is in contrast to the modes of ordinary Maclaurin spheroids which become dynamically unstable at $`4B_{12}=\mathrm{\Omega }^2`$, $`ϵ=0.953`$. Our model of the superfluid Maclaurin spheroid also becomes dynamically unstable at the same point, but only via the toroidal modes derived from the perturbation equations for $`V_{ij}`$, just as in EFE.
The real and imaginary parts of the dissipative toroidal modes are shown in the Fig. 2, for the same values of $`\eta `$ as in Fig. 1. The real parts of the modes tend towards each other and merge in the large friction limit. Note that there are no neutral points associated with these modes and the necessary condition for a point of bifurcation is not satisfied. The damping of the modes is maximal, as in the case of the transverse shear modes for $`\eta =1`$, and decreases in both limits of $`\eta 0`$ and $`\eta \mathrm{}`$. In contrast to ordinary Maclaurin spheroids, which become secularly unstable at the bifurcation point where $`2B_{12}=\mathrm{\Omega }^2`$ and $`ϵ=0.813`$, the new toroidal modes are stable at all values of the eccentricity. Our main conclusion is that the toroidal modes associated with the relative motions of the superfluid and the normal components always remain stable for incompressible Maclaurin spheroids. In the case of the compressible Maclaurin spheroids the point of the onset of secular instability may vary as a function of the adiabatic index (in the case of a polytropic type of equation of state) and hence the conclusions reached above should be verified for these models separately.
#### 3 The relative pulsation mode
To find the pulsation modes, which correspond to $`m=0`$, we first add Eqs. (64)-(65) and subtract from the result the Eq. (63). In this manner we find that
$`\left(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+2A_1\right)(U_{11}+U_{22})`$ (74)
$`+2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})(U_{1;2}U_{2;1})(\lambda ^2+4A_3+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda )U_{33}=0.`$ (75)
Subtracting Eqs. (67) and (66) one finds
$`\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)(U_{1;2}U_{2;1})\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})(U_{11}+U_{22})=0.`$ (76)
Equations (74) and (76) can be further combined to a single equation:
$`\left[\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda 2\mathrm{\Omega }^2+4A_1\right)(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda )+4\mathrm{\Omega }^2\lambda ^2(1\stackrel{~}{\beta }^{})^2\right](U_{11}+U_{22})`$ (77)
$`2\left[\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)\left(\lambda ^2+2\mathrm{\Omega }\lambda \stackrel{~}{\beta }^{\prime \prime }+4A_3\right)\right]U_{33}=0.`$ (78)
The solution is found by supplementing these equations by the divergence free condition
$$\frac{U_{11}}{a_1^2}+\frac{U_{22}}{a_2^2}+\frac{U_{33}}{a_3^2}=0$$
(79)
or, in terms of the eccentricity $`ϵ=1a_3^2/a_1^2,`$
$$(U_{11}+U_{22})(1ϵ^2)+U_{33}=0.$$
(80)
The third order characteristic equation is
$`(32ϵ^2)\lambda ^3+(8\stackrel{~}{\beta }\mathrm{\Omega }+4\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }4\stackrel{~}{\beta }ϵ^2\mathrm{\Omega }4\stackrel{~}{\beta }^{\prime \prime }ϵ^2\mathrm{\Omega })\lambda ^2+(4A_1+8A_38A_3ϵ^2`$ (81)
$`+2\mathrm{\Omega }^2+4\stackrel{~}{\beta }^2\mathrm{\Omega }^28\stackrel{~}{\beta }^{}\mathrm{\Omega }^2+4\stackrel{~}{\beta }^2\mathrm{\Omega }^2+8\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^28\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }ϵ^2\mathrm{\Omega }^2)\lambda `$ (82)
$`+8A_1\stackrel{~}{\beta }\mathrm{\Omega }+16A_3\stackrel{~}{\beta }\mathrm{\Omega }16A_3\stackrel{~}{\beta }ϵ^2\mathrm{\Omega }4\stackrel{~}{\beta }\mathrm{\Omega }^3=0,`$ (83)
where the trivial mode $`\lambda =0`$ is neglected. In the frictionless limit we find ($`\lambda =i\sigma `$ as before)
$$\sigma =\pm \left[\frac{2\mathrm{\Omega }^2+4A_1+8A_3(1ϵ^2)}{(32ϵ^2)}\right]^{1/2}.$$
(84)
The pulsation modes for a sphere follow in the limit $`ϵ,\mathrm{\Omega }0`$: for a sphere $`A_i/(\pi \rho G)=2/3`$, and Eq. (84) reduces to $`\sigma ^2=8/3`$ \[$`\sigma `$ is given in units of $`(\pi \rho G)^{1/2}]`$. This result could be compared with the pulsation modes of an ordinary sphere: $`\sigma ^2=16/15`$. Thus a superfluid sphere, apart form the ordinary pulsations, shows pulsations at frequencies roughly twice as large as the ordinary ones.
The real and imaginary parts of the dissipative pulsation modes of a superfluid Maclaurin spheroid are shown in the Fig. 3. The real parts of the modes are weakly affected by mutual friction and closely resemble those of an ordinary Maclaurin spheroid in the frictionless limit. These are located, however, at higher frequencies. The symmetry of the damping rate as a function of $`\eta `$ observed for the transverse shear and toroidal modes is again observed in Fig. 3. Note that the results in Fig. 6 were obtained in the case $`\beta ^{\prime \prime }=0.`$ The pulsation modes of the superfluid Maclaurin spheroid are stable, as is the case for the ordinary Maclaurin spheroids.
### B Modes of superfluid Jacobi ellipsoid
The sequence of the Jacobi ellipsoids emerges from the Maclaurin sequence at the bifurcation point $`ϵ=0.813`$ via a spontaneous breaking of symmetry in the plane perpendicular to the rotation ($`a_1a_2`$). The superfluid equilibrium figures are again identical to their ordinary fluid counterparts and the defining relations $`a_1^2a_2^2A_{12}=a_3^2A_3`$ and $`\mathrm{\Omega }^2=2B_{12}`$ are unchanged. Ordinary Jacobi ellipsoids are known to be stable against second order harmonic perturbations while they become dynamically unstable against transformation into Poincarè’s pear shaped figures through a mode belonging to third order harmonic perturbations. If the sequence of Jacobi ellipsoids is parametrized in terms of the variable $`\mathrm{cos}^1(a_3/a_1)`$, it is stable between the point of bifurcation from the Maclaurin sequence, $`\mathrm{cos}^1(a_3/a_1)=54.36`$, and the point where Poincarè’s figures bifurcate, $`\mathrm{cos}^1(a_3/a_1)=69.82`$. Here, by an explicit calculation, we verify that superfluid Jacobi ellipsoids do not develop new instabilities via second order harmonic modes of the relative displacements.
#### 1 Relative odd modes
The treatment of the oscillations of the Maclaurin spheroid of the previous sections can be readily extended to the Jacobi ellipsoids by lifting the degeneracy in indexes 1 and 2 and imposing $`A_1A_2`$. The equations odd in index 3 are
$`\lambda ^2U_{3;1}`$ $`=`$ $`2A_3U_{31}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;1},`$ (85)
$`\lambda ^2U_{3;2}`$ $`=`$ $`2A_3U_{32}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;2},`$ (86)
$`\lambda ^2U_{1;3}2\mathrm{\Omega }\lambda U_{2;3}`$ $`=`$ $`\left(2A_1+\mathrm{\Omega }^2\right)U_{13}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;3}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;3},`$ (87)
$`\lambda ^2U_{2;3}+2\mathrm{\Omega }\lambda U_{1;3}`$ $`=`$ $`\left(2A_2+\mathrm{\Omega }^2\right)U_{23}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;3}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;3}.`$ (88)
Combining Eqs. (85) and (87) and, similarly, Eqs. (86) and (88), after some manipulation we find
$`\left[\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)+2\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)A_3+\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)(2A_1\mathrm{\Omega }^2)\right]U_{13}`$ (89)
$`2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3\right)U_{23}=0,`$ (90)
$`\left[\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)+2\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda \right)A_3+\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda \right)(2A_2\mathrm{\Omega }^2)\right]U_{23}`$ (91)
$`+2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})\left(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3\right)U_{13}=0.`$ (92)
Equations (90) and (92) are sufficient to determine the symmetric parts of the virials, and any two of Eqs. (85)-(88) can be used to find the antisymmetric parts. The sixth order characteristic equation is
$`\lambda ^6+4\mathrm{\Omega }(\stackrel{~}{\beta }+\stackrel{~}{\beta }^{\prime \prime })\lambda ^5+\left[4A_3+2\mathrm{\Omega }(1\stackrel{~}{\beta }^{})^2+4\mathrm{\Omega }^2(\stackrel{~}{\beta }+\stackrel{~}{\beta }^{\prime \prime })^2(2A_1\mathrm{\Omega }^2)(2A_2\mathrm{\Omega }^2)\right]\lambda ^4`$ (93)
$`+[16A_3\stackrel{~}{\beta }\mathrm{\Omega }+8A_3\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }+8\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(1\stackrel{~}{\beta }^{})^2+16\stackrel{~}{\beta }^2\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^3+16\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^3`$ (94)
$`2\stackrel{~}{\beta }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)2\stackrel{~}{\beta }\mathrm{\Omega }(2A_2\mathrm{\Omega }^2)4\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)4\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_2\mathrm{\Omega }^2)]\lambda ^3`$ (95)
$`+[4A_3^2+8A_3\mathrm{\Omega }(1\stackrel{~}{\beta }^{})^2+16A_3\stackrel{~}{\beta }^2\mathrm{\Omega }^2+32A_3\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2+8\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^3(1\stackrel{~}{\beta }^{})^2+16\stackrel{~}{\beta }^2\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^4`$ (96)
$`2A_3(2A_1\mathrm{\Omega }^2)8\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)4\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)2A_3(2A_2\mathrm{\Omega }^2)`$ (97)
$`8\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(2A_2\mathrm{\Omega }^2)4\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^2(2A_2\mathrm{\Omega }^2)+(2A_1\mathrm{\Omega }^2)(2A_2\mathrm{\Omega }^2)]\lambda ^2`$ (98)
$`+[16A_3^2\stackrel{~}{\beta }\mathrm{\Omega }+16A_3\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(1\stackrel{~}{\beta }^{})^2+32A_3\stackrel{~}{\beta }^2\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^34A_3\stackrel{~}{\beta }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)4A_3\stackrel{~}{\beta }\mathrm{\Omega }(2A_2\mathrm{\Omega }^2)`$ (99)
$`4A_3\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)8\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^3(2A_1\mathrm{\Omega }^2)4A_3\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_2\mathrm{\Omega }^2)`$ (100)
$`8\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^3(2A_2\mathrm{\Omega }^2)+4\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }(2A_1\mathrm{\Omega }^2)(2A_2\mathrm{\Omega }^2)]\lambda `$ (101)
$`+[8A_3^2\mathrm{\Omega }(1\stackrel{~}{\beta }^{})^2+16A_3^2\stackrel{~}{\beta }^2\mathrm{\Omega }^28A_3\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)`$ (102)
$`8A_3\stackrel{~}{\beta }\stackrel{~}{\beta }^{\prime \prime }\mathrm{\Omega }^2(2A_2\mathrm{\Omega }^2)+4\stackrel{~}{\beta }^{\prime \prime 2}\mathrm{\Omega }^2(2A_1\mathrm{\Omega }^2)(2A_2\mathrm{\Omega }^2)]=0.`$ (103)
The real and imaginary parts of the dissipative odd parity modes are shown in the Fig. 4. The Jacobi sequence is parametrized in terms of $`\mathrm{cos}^1(a_3/a_1)`$ starting off from the point of bifurcation of the Jacobi ellipsoid from the Maclaurin sequence. The low frequency mode resembles the rotational mode of the ellipsoid; its frequency decreases with increasing friction. One of the remaining two distinct high frequency modes is almost unaffected by the dissipation, while the other is suppressed close to the bifurcation point in a monotonic manner. The damping rates of the odd modes are maximal at $`\eta =1`$ and tend to zero for both large and small friction. The modes are damped along the entire sequence; hence, we conclude that superfluid Jacobi ellipsoids are stable against the odd second harmonic modes of oscillations.
#### 2 Relative even modes
The explicit form of the even parity modes for the Jacobian sequence is
$`\lambda ^2U_{3;3}`$ $`=`$ $`(\pi \rho G)^1\delta \stackrel{~}{\mathrm{\Pi }}2A_3U_{33}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;3},`$ (104)
$`\lambda ^2U_{1;1}2\mathrm{\Omega }\lambda U_{2;1}`$ $`=`$ $`(\pi G\rho )^1\delta \stackrel{~}{\mathrm{\Pi }}2A_1U_{11}+\mathrm{\Omega }^2U_{11}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;1}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;1},`$ (105)
$`\lambda ^2U_{2;2}+2\mathrm{\Omega }\lambda U_{1;2}`$ $`=`$ $`(\pi G\rho )^1\delta \stackrel{~}{\mathrm{\Pi }}2A_2U_{22}+\mathrm{\Omega }^2U_{22}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;2}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;2},`$ (106)
$`\lambda ^2U_{1;2}2\mathrm{\Omega }\lambda U_{2;2}`$ $`=`$ $`(\mathrm{\Omega }^22A_1)U_{12}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;2}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;2},`$ (107)
$`\lambda ^2U_{2;1}+2\mathrm{\Omega }\lambda U_{1;1}`$ $`=`$ $`(\mathrm{\Omega }^22A_2)U_{12}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;1}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;1}.`$ (108)
These equations can be reduced to a simpler set of equations through manipulations which eliminate the variations of the pressure tensor. Explicitly, in the first step we subtract the Eqs. (105) and (106); in the second we sum Eqs. (105) and (106) and subtract twice Eq. (104). The result is
$`(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+2A_1)U_{11}(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+2A_2)U_{22}`$ (109)
$`2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})U_{12}=0,`$ (110)
$`(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+2A_1)U_{11}+(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+2A_2)U_{22}`$ (111)
$`(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +4A_3)U_{33}+2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})(U_{1;2}U_{2;1})=0.`$ (112)
Further we add and subtract Eqs. (107) and (108) to find
$`\left[\lambda ^2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda 4B_{12}+2(A_1+A_2)\right]U_{12}+\mathrm{\Omega }(1\stackrel{~}{\beta }^{})\lambda (U_{11}U_{22})=0,`$ (113)
$`(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda )(U_{1;2}U_{2;1})\mathrm{\Omega }(1\stackrel{~}{\beta }^{})\lambda (U_{11}+U_{22})+2(A_1A_2)U_{12}=0.`$ (114)
Equations (109)-(114), supplemented by the divergence free condition, Eq. (79), are sufficient to determine the modes. The characteristic equation is of seventh order, excluding the trivial root $`\lambda =0`$; in the frictionless limit the characteristic equation is of third order. The explicit form of these equations is cumbersome and will not be given here.
The real and imaginary parts of the dissipative even parity modes are shown in Fig. 5. For each member of the sequence, the eigenvalues of the two high frequency modes are suppressed and that of the low-frequency mode is amplified as the dissipation increases. As in the case of the odd modes the damping rates of the even parity modes are maximal at $`\eta =1`$ and tend to zero for both large and small friction. The damping rates are again positive along the entire sequence and we conclude that superfluid Jacobi ellipsoids are stable against the even parity second harmonic modes of oscillations.
### C Modes of superfluid Roche ellipsoid
In this section we extend the previous discussion of isolated ellipsoids to binary star systems, and consider the simplest case – the Roche problem. The classical Roche binary consists of a finite size ellipsoid (primary of mass $`M`$) and a point mass (secondary of mass $`M^{}`$) rotating about their common center of mass with an angular velocity $`\mathrm{\Omega }`$. The new ingredient in the problem of the equilibrium and stability of the primary is the tidal Newtonian gravitational field of the secondary. Place the center of the coordinate system at the center of mass of the primary with the $`x_1`$-axis pointing to the center of mass of the secondary and $`x_3`$-axis along the vector $`𝛀`$. The equation of motion for a fluid element of the primary in the frame rotating with angular velocity $`\mathrm{\Omega }`$ is, then (EFE, Chap. 8, Sec. 55)
$$\rho _\alpha D_\alpha u_{\alpha ,i}=\frac{p_\alpha }{x_i}\rho _\alpha \frac{(\varphi +\varphi ^{})}{x_i}+\frac{1}{2}\rho _\alpha \frac{|𝛀\mathbf{\times }𝐱|^2}{x_i}+2\rho _\alpha ϵ_{ilm}u_{\alpha ,l}\mathrm{\Omega }_m+F_{\alpha \beta ,i},$$
(115)
where the tidal potential of the secondary, up to quadratic terms in $`x_i/R`$, is
$$\varphi ^{}=\frac{GM^{}}{R}\left(1+\frac{x_1}{R}+\frac{2x_1^2x_2^2x_3^2}{2R^2}\right).$$
(116)
The modified Keplerian rotation frequency for circular orbits, which is consistent with the first order virial equations, is
$$\mathrm{\Omega }^2=(1+P)\varphi _0(1+\delta ),$$
(117)
where $`P=M/M^{}`$ is the mass ratio, $`\varphi _0=GM^{}/R^3`$ is the tidal potential at the origin of the primary, and $`\delta `$ is the quadrupole part of the tidal field. The latter correction to the Keplerian frequency is maximal at the Roche limit where $`\delta 0.13`$ . For the sake of simplicity this correction is dropped in the following, as it does not enter into the analysis of the stability of the Roche ellipsoid for displacements that are linear functions of the coordinates. However, the relation between the frequency $`\mathrm{\Omega }`$ and the orbital separation is now determined within an accuracy $`\delta 1`$. As in the case of the solitary ellipsoids, we find that the equilibrium figure of the superfluid Roche ellipsoid is identical to its ordinary fluid counterpart.
To treat Roche ellipsoids we modify Eq. (39) to
$`\lambda ^2U_{i;j}2ϵ_{il3}\mathrm{\Omega }_3\lambda U_{l;j}`$ $`=`$ $`\delta _{ij}(\pi \rho G)^1\delta \stackrel{~}{\mathrm{\Pi }}A_iU_{ij}2\mathrm{\Omega }\lambda \stackrel{~}{\beta }_{ik}U_{k;j}`$ (118)
$`+`$ $`(\mathrm{\Omega }^2\varphi _0)U_{ij}\mathrm{\Omega }^2\delta _{i3}U_{3j}+3\varphi _0\delta _{i1}U_{1j},`$ (119)
where all frequencies are measured in units of $`(\pi \rho G)^{1/2}`$. Equation (119) is appropriate for finding the second harmonic modes of oscillations of Roche ellipsoids.
#### 1 Relative odd modes
The equations determining the modes even and odd in index 3 form separate sets. We start with the modes belonging to $`l=2`$ and $`m=1,1`$ displacements, which are odd in index 3; for these,
$`\lambda ^2U_{3;1}`$ $`=`$ $`(2A_3+\varphi _0)U_{31}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;1},`$ (120)
$`\lambda ^2U_{3;2}`$ $`=`$ $`2(A_3+\varphi _0)U_{32}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;2},`$ (121)
$`\lambda ^2U_{1;3}2\mathrm{\Omega }\lambda U_{2;3}`$ $`=`$ $`\left(2A_1+\mathrm{\Omega }^2+2\varphi _0\right)U_{13}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;3}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;3},`$ (122)
$`\lambda ^2U_{2;3}+2\mathrm{\Omega }\lambda U_{1;3}`$ $`=`$ $`\left(2A_2+\mathrm{\Omega }^2\varphi _0\right)U_{23}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;3}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;3}.`$ (123)
On combining Eqs. (120) and (122) and, similarly, Eqs. (121) and (123) we obtain
$`\left[(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3+\varphi _0)(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda )+(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda )(2A_1\mathrm{\Omega }^22\varphi _0)\right]U_{13}`$ (124)
$`2\mathrm{\Omega }(1\stackrel{~}{\beta }^{})\lambda \left[\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3+\varphi _0\right]U_{32}=0,`$ (125)
$`\left[(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3+\varphi _0)(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda )+(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda )(2A_2\mathrm{\Omega }^2+\varphi _0)\right]U_{23}`$ (126)
$`+2\mathrm{\Omega }(1\stackrel{~}{\beta }^{})\lambda \left[\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2A_3+\varphi _0\right]U_{13}=0.`$ (127)
The $`U_{ij}`$ are symmetric under interchange of their indexes, and Eqs. (124) and (126) completely determine the modes \[any two of Eqs. - may be used to find the antisymmetric parts of $`U_{i;j}`$\]. The real and imaginary parts of the dissipative odd parity modes are shown in Fig. 6, for the case of an equal mass binary ($`P=1`$). The relative modes of Roche ellipsoids for other values of the mass ratio display behavior similar to the $`P=1`$ case. The Roche sequence is parametrized in terms of $`\mathrm{cos}^1(a_3/a_1)`$. The low frequency mode resembles the rotational mode of the ellipsoid; its frequency decreases with increasing friction. The high frequency modes tend towards each other and merge in the limit of slow rotation; in the opposite limit the modes remain unaffected by the dissipation. There are three distinct rates for the damping of oscillations. These are maximal at $`\eta =1`$ and tend to zero in both limits of large and small friction, as was the case for the Maclaurin and Jacobi ellipsoids. The damping rates are positive along the entire sequence. We conclude that superfluid Roche ellipsoids do not develop instabilities via the second order harmonic odd modes of relative oscillation.
#### 2 Relative Even Modes
As is well known, Roche ellipsoids develop a dynamical instability via the second order even parity modes beyond the Roche limit, which is the point of closest approach of the primary to the secondary. If viscous dissipation is allowed for, Roche ellipsoids become secularly unstable at the Roche limit via an even parity mode and before dynamical instability sets in. We have seen that Maclaurin spheroids do not develop any instabilities (i.e. neither dynamical nor secular) via the modes associated with the $`U_{ij}`$ in the presence of superfluid dissipation. The extension of the theory above to superfluid Roche ellipsoids, as we show now, does not reveal any new instabilities in the presence of superfluid dissipation, again in contrast to the analysis based on the ordinary viscous dissipation.
The explicit form of the even parity modes for the Roche sequence is
$`\lambda ^2U_{3;3}`$ $`=`$ $`(\pi \rho G)^1\delta \stackrel{~}{\mathrm{\Pi }}(2A_3+\varphi _0)U_{33}2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda U_{3;3},`$ (128)
$`\lambda ^2U_{1;1}2\mathrm{\Omega }\lambda U_{2;1}`$ $`=`$ $`(\pi G\rho )^1\delta \stackrel{~}{\mathrm{\Pi }}(2A_1\mathrm{\Omega }^22\varphi _0)U_{11}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;1}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;1},`$ (129)
$`\lambda ^2U_{2;2}+2\mathrm{\Omega }\lambda U_{1;2}`$ $`=`$ $`(\pi G\rho )^1\delta \stackrel{~}{\mathrm{\Pi }}(2A_2\mathrm{\Omega }^2+\varphi _0)U_{22}+\mathrm{\Omega }^2U_{22}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;2}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;2},`$ (130)
$`\lambda ^2U_{1;2}2\mathrm{\Omega }\lambda U_{2;2}`$ $`=`$ $`(\mathrm{\Omega }^2+2\mu 2A_1)U_{12}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{1;2}2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{2;2},`$ (131)
$`\lambda ^2U_{2;1}+2\mathrm{\Omega }\lambda U_{1;1}`$ $`=`$ $`(\mathrm{\Omega }^2\varphi _02A_2)U_{12}2\mathrm{\Omega }\stackrel{~}{\beta }\lambda U_{2;1}+2\mathrm{\Omega }\stackrel{~}{\beta }^{}\lambda U_{1;1}.`$ (132)
These equations can be reduced to a simpler set of equations through manipulations which eliminate variations of the pressure tensor. Using the symmetry properties of the virials we first subtract Eqs. (129) and (130), then sum Eqs. (129) and (130) and subtract twice Eq. (128) to obtain
$`(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^22\varphi _0+2A_1)U_{11}(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+\varphi _0+2A_2)U_{22}`$ (133)
$`2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})U_{12}=0,`$ (134)
$`(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^22\varphi _0+2A_1)U_{11}+(\lambda ^2/2+\mathrm{\Omega }\stackrel{~}{\beta }\lambda \mathrm{\Omega }^2+\varphi _0+2A_2)U_{22}`$ (135)
$`(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }^{\prime \prime }\lambda +2\varphi _0+4A_3)U_{33}+2\mathrm{\Omega }\lambda (1\stackrel{~}{\beta }^{})(U_{1;2}U_{2;1})=0.`$ (136)
Further we add and subtract Eqs. (131) and (132) to find
$`\left[\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda +2(A_1+A_2)2\mathrm{\Omega }^2\varphi _0\right]U_{12}+\mathrm{\Omega }(1\stackrel{~}{\beta }^{})\lambda (U_{11}U_{22})=0,`$ (137)
$`(\lambda ^2+2\mathrm{\Omega }\stackrel{~}{\beta }\lambda )(U_{1;2}U_{2;1})\mathrm{\Omega }(1\stackrel{~}{\beta }^{})\lambda (U_{11}+U_{22})[3\varphi _02(A_1A_2)]U_{12}=0.`$ (138)
Equations (137)-(138), supplemented by the divergence free condition, Eq. (79), are sufficient to determine the unknown virials.
The real and imaginary parts of the dissipative even parity modes are shown in Fig. 7. For each member of the sequence the eigenvalues of the two high frequency modes are suppressed and that of the low frequency one is amplified with increasing dissipation. In effect these modes merge in the slow rotation limit. The modes do not become neutral at any point along the frictionless sequence and, hence, the necessary condition for the onset of dynamical instability is not achieved. Note that our model for the Roche ellipsoids is dynamically unstable as is its classical counterpart, but via the modes governed by the Eq. (24) modified appropriately to include the external tidal potential. We do not repeat the mode analysis for the virials $`V_{ij}`$ as it is a complete analogue of the analysis in EFE. As in the case of odd modes the damping rates of even parity modes are maximal at $`\eta =1`$ and tend to zero in both limits of large and small friction. The damping rates are positive along the entire sequence and we conclude that superfluid Roche ellipsoids are secularly stable against the even parity second harmonic modes of oscillations associated with the relative motions between the superfluid and normal components.
## IV Viscosity and Gravitational Radiation
Above, we neglected viscous dissipation in computing normal modes of a normal fluid-superfluid mixture. However, for a single fluid, viscous dissipation is important for understanding stability, for it is responsible for secular instability. Although viscous terms spoil the calculation of the modes of uniform ellipsoids from moment equations formally, when the dissipative time-scale is long, one can include them perturbatively (e.g. EFE, Chap. 5, §37b).
The inclusion of viscosity is slightly more complicated for a mixture of superfluid and normal fluid because viscous dissipation only operates on the normal fluid. The separation of Eq. (12) for the fluid displacements into Eqs. (18) and (21) for the common and differential fluid displacements, $`𝝃_\pm `$, is possible because the only form of dissipation, the mutual friction force, included in Eq. (12) only depends on $`d𝝃_{}/dt`$. Viscous dissipation depends on $`𝝃_N`$ only, and, in a formal sense, the dynamics no longer separate into the independent dynamics of $`𝝃_\pm `$.
If we assume that the time-scale associated with viscous dissipation is relatively long, then we can include it perturbatively. The calculation is a bit more subtle than for a single fluid, because we have to deduce $`𝝃_N`$ for the modes. This brings up an issue that we glossed over earlier, in setting up the calculation of modes in §II: even when computing the modes that arise from the equation for $`𝝃_{}`$ alone, the equation for $`𝝃_+`$ must be satisfied, and viceversa. The simplest way for this to work is for $`𝝃_+`$ to vanish when $`𝝃_{}`$ is nonzero and vice versa. In fact it is easy to show that this is a reasonable solution provided that the Eulerian pressure perturbations are $`\delta p_\alpha =\xi _{\alpha ,l}p_\alpha /x_l`$, a situation that arises naturally for adiabatic perturbations, where the Lagrangian pressure perturbations are $`\mathrm{\Delta }_\alpha p_\alpha =\mathrm{\Gamma }_\alpha \xi _{\alpha ,l}/x_l`$, and the perturbations are solenoidal, so that $`\xi _{\alpha ,l}/x_l=0`$, as is true for all modes considered in this paper. More realistically, one would also need to consider non-adiabatic effects, such as might arise from perturbations from $`\beta `$equilibrium; see e.g. Lindblom and Mendell . These would tend to couple $`𝝃_\pm `$, but if small, could be computed perturbatively, as we do here for viscosity, which also couples $`𝝃_\pm `$. To see how this works, consider a mode of Eq. (21), and examine under what conditions Eq. (18) will be satisfied. Then, using
$$\delta p_\alpha =\xi _{\alpha ,l}\frac{p_\alpha }{x_l}=f_\alpha \xi _{\alpha ,l}\frac{p}{x_l}$$
(139)
and \[substituting the definition of $`𝝃_+`$, and $`\rho _\alpha =f_\alpha \rho `$ into Eq. (13)\]
$$\delta \varphi G_Vd^3x^{}\rho (𝐱^{})\xi _{+,l}(𝐱^{})\frac{}{x_l^{}}\left(\frac{1}{|𝐱𝐱^{}|}\right)$$
(140)
it is easy to see that Eq. (18) only depends on $`𝝃_+`$. But since the normal modes of Eq. (21) have different frequencies than normal modes of Eq. (18), we must have $`𝝃_+=0`$ when $`𝝃_{}0`$. Since, in general,
$$𝝃_S=𝝃_++f_N𝝃_{}𝝃_N=𝝃_+f_S𝝃_{},$$
(141)
we conclude that, for modes of Eq. (21), $`𝝃_S=f_N𝝃_{}`$ and $`𝝃_N=f_S𝝃_{}`$, when Eulerian pressure perturbations are given by Eq. (139). Similarly, since Eq. (21) only depends on $`𝝃_{}`$, for modes with $`𝝃_+0`$, we must have $`𝝃_{}=0`$ and, therefore, $`𝝃_S=𝝃_N=𝝃_+`$, assuming Eq. (139). In particular, if the kinematic viscosity $`\nu `$ is held constant \[see EFE, Chap. 5, Sec. 36, Eq. (111)\], for modes with $`𝝃_+0`$, the viscous dissipation rate is smaller by a factor of $`f_N`$ than it is for a single fluid with the same background and displacement $`𝝃_+`$, as might have been expected qualitatively (i.e. for small normal fluid density, the viscous dissipation must be diminished). For perturbations with $`𝝃_{}0`$ and displacements that are linear functions of the coordinates, we must add
$$5f_S\nu \left(\frac{1}{a_j^2}\frac{dU_{i;j}}{dt}+\frac{1}{a_i^2}\frac{dU_{j;i}}{dt}\right)$$
(142)
to the right hand side of Eq. (27) to include the effects of viscous dissipation.
Perturbations with $`𝝃_{}0`$ emit no gravitational radiation because $`𝝃_+=0`$ for them, and therefore there are no perturbations of the quadrupole moment or any other net mass currents associated with them. Gravitational radiation is emitted by the modes in which the two fluids move together at the same rate as for a single fluid (e.g. ref. ). Thus, none of the new modes of a superfluid-normal fluid mixture found here is affected by gravitational radiation at all.
## V Conclusions
Despite a number of simplifying assumptions, the study of the oscillation modes of uniform ellipsoids is useful for understanding the equilibrium and stability of real neutron stars, at least qualitatively. Moreover, the theory is simple enough that it can be extended readily to include modifications to the underlying physics; here, we have considered new features that arise because a neutron star contains a mixture of normal fluid and superfluid coupled by mutual friction. The theory of uniform ellipsoids is interesting from the viewpoint of mathematical physics because it is solvable exactly, and it may also have applications to other physical systems, such as the physics of trapped, rotating Bose-Einstein condensates.
In this paper we have extended previous treatments of the oscillation modes of ellipsoidal figures of equilibrium to the case of a mixture of normal fluid and superfluid. The basic equations of motion for the two fluid hydrodynamics include mutual friction exactly, since the frictional forces depend linearly on the relative velocity between the two fluids, and vanish in the background where the two fluids move together. In addition the fluids are coupled via mutual gravitational attraction, which is also treated without further approximations. As a result our relations closely resemble Chandrasekhar’s tensor virial equations, even though they are intrinsically dissipative due to the mutual friction. While we have developed these moment equations for general underlying equilibria, they are most useful for perturbations around uniform backgrounds, in which case the moment equations of various orders decouple and yield exact solutions for the normal modes.
Quite generally, there are two classes of modes for small perturbations, one class in which the two fluids move together and the other in which there is relative motion between them. The former are identical to the modes found for a single fluid. As a result our models of superfluid Maclaurin and Roche ellipsoids undergo dynamical instabilities with respect to these modes which are indistinguishable from what is found for their classical counterparts. When ordinary viscous dissipation is included they are also subject to secular instabilities related to the modes in which two fluids move together. If the kinematic viscosity is held constant, then the rate of viscous dissipation, when computed in the “low Reynolds number” approximation, is diminished by a factor $`f_N`$, the fraction of the total mass in the normal fluid (see however below).
The modes involving the relative motion between the fluids are completely new and are shown to be stable along the entire sequences of the incompressible Maclaurin, Jacobi and Roche ellipsoids independent of the magnitude of the phenomenological mutual friction. These modes also do not become neutral at selected points along any sequence and the necessary condition for the point of bifurcation is not achieved. Our main conclusion is that mutual friction does not drive secular instabilities in incompressible and irrotational ellipsoids. In addition we find that even though the new modes are subject to viscous dissipation (a consequence of viscosity of the normal matter), they do not emit gravitational radiation, and are therefore immune to any instabilities associated with gravitational radiation, irrespective of their modal frequencies.
The results summarized above hold within a combined framework of two-fluid superfluid hydrodynamics, Newtonian gravity, and the ellipsoidal approximation, as formulated in EFE. Each of these elements of our approach contains a number of simplifying approximations which need to be relaxed in realistic applications to neutron stars. For example, to treat the mixture of neutron and proton superfluids in the neutron star cores, the one-constituent two-fluid superfluid hydrodynamics must be replaced by the hydrodynamics of the multi-constituent superfluid mixtures, in which case the mutual entrainment of the superfluids and deviations from $`\beta `$-equilibrium must be accounted for (see Ref. for a discussion of these effects and their impact on neutron star oscillations within the real energy functional method). We anticipate that these effects can be incorporated in the tensor virial approach in a perturbative manner and the results of previous sections will hold in leading order of the perturbation expansion. On the other hand, relaxing the incompressible approximation and, hence, including the partial pressures of the superfluid and the normal fluid, will lead to non-perturbative effects as the pressure terms significantly alter the balance between gravitational attraction and centrifugal stretching. As is well known, the points of the onset of the dynamical and secular instabilities of compressible ellipsoids depend on the adiabatic index of the underlying polytropic equation of state. As noted in Sec. 3, the conclusions reached with respect to the stability of the incompressible ellipsoids should be verified for the compressible models anew. The differences in the pressures (or equations of states) of the normal and superfluid phases, however, are typically small in neutron stars, since the condensation energy is negligible compared to the degeneracy pressures of interacting Fermi-liquids. The coupling between the partial pressures of the noraml fluid and superfluid, therefore, can be treated perturbatively.
Another effect that needs to be included is the density dependence of the mutual friction coefficients and the kinematic viscosity. For example, in neutron stars the kinematic viscosity is density dependent in general, explicitly as a result of to the density dependence of the phase space of normal quasiparticles undergoing collisions, implicitly because of the density dependence of the in-medium scattering amplitudes (see for further details Ref. ). The ramification for comparison of the secular instabilities of the normal fluid and superfluid ellipsoids is that the modifications of the rate of the viscous dissipation will depend in a non-trivial manner on the fraction of the normal fluid in the system.
The entrainment, $`\beta `$-nonequilibrium, compressibility, e. t. c. will couple the relative and center-of-mass modes in general. Such a “mixing”, as discussed in Sec. 5, implies $`𝝃_+0`$ for the relative modes and, similarly, $`𝝃_{}0`$ for the center-of-mass modes. Therefore, the mutual friction might tend to drive the center-of-mass modes secularly unstable; if they emit gravitational radiation, the mutual friction will suppress the gravitational radiation induced instabilities. One of the important issues to be addressed by the future work is the magnitude of the “mixing” of the modes for general equilibria and the corresponding times scales.
## ACKNOWLEDGMENTS
A.S. acknowledges the Nederlandse Organisatie voor Wetenschappelijk Onderzoek for its support at KVI Groningen via the Stichting voor Fundamenteel Onderzoek der Materie, the Institute for Nuclear Theory at the University of Washington for its hospitality and the Department of Energy for partial support. I.W. acknowledges partial support for this project from NASA.
## A “Virial” Equations and Perturbations
The two fluids need not occupy the same volume, and we shall suppose that fluid $`\alpha `$ occupies a volume $`V_\alpha `$. Taking the zeroth moment of Eq. (1) amounts to integrating over $`V_\alpha `$; doing so, we obtain<sup>\**</sup><sup>\**</sup>\**We assume that $`p_\alpha =0`$ on the bounding surface of $`V_\alpha `$. the “first order ‘virial’ equation”<sup>††</sup><sup>††</sup>††We put the word “virial” in quotes because the equations are dissipative.
$`{\displaystyle \frac{d}{dt}}\left({\displaystyle _{V_\alpha }}d^3x\rho _\alpha u_{\alpha ,i}\right)`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle _{V_\alpha }}d^3x\rho _\alpha u_{\alpha ,l}+(\mathrm{\Omega }^2\delta _{ij}\mathrm{\Omega }_i\mathrm{\Omega }_j){\displaystyle _{V_\alpha }}d^3x\rho _\alpha x_j`$ (A2)
$`(1\delta _{\alpha \beta }){\displaystyle _{V_\alpha }}d^3x\rho _\alpha {\displaystyle \frac{\varphi _\beta }{x_i}}+{\displaystyle _{V_\alpha }}d^3xF_{\alpha \beta ,i}.`$
Apart from inertial forces, which do not couple the two fluids, there are two forces that do couple them: gravity and friction. The net, mutual gravitational force between the fluids only vanishes if they (i) occupy the same volume and (ii) have densities that are proportional to one another (i.e. $`\rho _S\rho _N`$). The mutual friction force is nonzero as long as the fluids move relative to one another. Thus we see that, for a two fluid mixture, the zeroth moment of Eq. (1) is not trivial, as it would be for a single fluid (as in EFE). Note, though, that the center of mass motion of the combined system is trivial; i.e., if the center of mass of the combined system starts out at $`𝐱=0`$ with zero velocity, it does not move.<sup>‡‡</sup><sup>‡‡</sup>‡‡There is a slight subtlety that has not been stated explicitly in deriving Eq. (A2). The mutual friction force is nonzero only in the overlap volume of the two fluids. This restriction is necessary to derive conservation of total momentum for the combined fluids. We might be interested in situations involving long range coupling, such as between the core superfluid and crustal normal fluid. For such a coupling to occur, we need to introduce long range fields (e.g. magnetic fields) capable of transmitting forces between fluid elements at different points in space.
Taking the first moment of Eq. (1) results in the second order “virial” equation
$`{\displaystyle \frac{d}{dt}}\left({\displaystyle _{V_\alpha }}d^3x\rho _\alpha x_ju_{\alpha ,i}\right)`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_m\left({\displaystyle _{V_\alpha }}d^3x\rho _\alpha x_ju_{\alpha ,l}\right)+\mathrm{\Omega }^2I_{\alpha ,ij}\mathrm{\Omega }_i\mathrm{\Omega }_kI_{\alpha ,kj}`$ (A4)
$`+2𝒯_{\alpha ,ij}+\delta _{ij}\mathrm{\Pi }_\alpha +_{\alpha ,ij}+(1\delta _{\alpha \beta })_{\alpha \beta ,ij}+_{\alpha \beta ,ij},`$
where
$`I_{\alpha ,ij}`$ $``$ $`{\displaystyle _{V_\alpha }}d^3x\rho _\alpha x_ix_j`$ (A5)
$`\mathrm{\Pi }_\alpha `$ $``$ $`{\displaystyle _{V_\alpha }}d^3xp_\alpha `$ (A6)
$`𝒯_{\alpha ,ij}`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{V_\alpha }}d^3x\rho _\alpha u_{\alpha ,i}u_{\alpha ,j}`$ (A7)
$`_{\alpha ,ij}`$ $``$ $`{\displaystyle \frac{G}{2}}{\displaystyle _{V_\alpha }}{\displaystyle \frac{d^3xd^3x^{}\rho _\alpha (𝐱)\rho _\alpha (𝐱^{})(x_ix_i^{})(x_jx_j^{})}{|𝐱𝐱^{}|^3}}`$ (A8)
$`_{\alpha \beta ,ij}`$ $``$ $`G{\displaystyle _{V_\alpha }}d^3x{\displaystyle _{V_\beta }}{\displaystyle \frac{d^3x^{}\rho _\alpha (𝐱)\rho _\beta (𝐱^{})x_j(x_ix_i^{})}{|𝐱𝐱^{}|^3}}`$ (A9)
$`_{\alpha \beta ,ij}`$ $``$ $`{\displaystyle _{V_\alpha }}d^3xx_jF_{\alpha \beta ,i}.`$ (A10)
When there is only one fluid present, this equation reduces to the results found in Chap. 2 of EFE. There are two new terms here: there is a term that arises from the mutual gravitational forces of the two fluids ($`_{\alpha \beta ,ij}`$), and also a term from the mutual friction ($`_{\alpha \beta ,ij}`$).
Consider first the variation of the first order virial equation under the influence of perturbations. Most terms are simple to compute, but we must take special care in computing
$$\delta _{V_\alpha }d^3x\rho _\alpha \frac{\varphi _\beta }{x_i}=\delta G_{V_\alpha }d^3x_{V_\beta }\frac{d^3x^{}\rho _\alpha (𝐱)\rho _\beta (𝐱^{})(x_ix_i^{})}{|𝐱𝐱^{}|^3}.$$
(A11)
In computing the necessary variation, think of $`𝐱`$ as having a label $`\alpha `$ and $`𝐱^{}`$ as having a label $`\beta `$. It is then easy to find that
$`\delta {\displaystyle _{V_\alpha }}d^3x\rho _\alpha {\displaystyle \frac{\varphi _\beta }{x_i}}=G{\displaystyle _{V_\alpha }}d^3x\rho _\alpha (𝐱)\xi _{\alpha ,l}(𝐱){\displaystyle \frac{}{x_l}}{\displaystyle _{V_\beta }}{\displaystyle \frac{d^3x^{}\rho _\beta (𝐱^{})(x_ix_i^{})}{|𝐱𝐱^{}|^3}}`$ (A12)
$`+G{\displaystyle _{V_\beta }}d^3x\rho _\beta (𝐱)\xi _{\beta ,l}(𝐱){\displaystyle \frac{}{x_l}}{\displaystyle _{V_\alpha }}{\displaystyle \frac{d^3x^{}\rho _\alpha (𝐱^{})(x_ix_i^{})}{|𝐱𝐱^{}|^3}}`$ (A13)
which is manifestly antisymmetric on $`\alpha \beta `$. Assuming $`V_\alpha =V_\beta =V`$ and $`\rho _\alpha =f_\alpha \rho (𝐱)`$ in the background equilibrium, we can simplify this to
$$\delta _{V_\alpha }d^3x\rho _\alpha \frac{\varphi _\beta }{x_i}=Gf_\alpha f_\beta _Vd^3x\rho (𝐱)[\xi _{\beta ,l}(𝐱)\xi _{\alpha ,l}(𝐱)]\frac{}{x_l}_V\frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})}{|𝐱𝐱^{}|^3}.$$
(A14)
Gathering terms, we find that the perturbed first order virial equation is
$`{\displaystyle \frac{d^2}{dt^2}}\left(f_\alpha {\displaystyle _V}d^3x\rho \xi _{\alpha ,i}\right)`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_m{\displaystyle \frac{d}{dt}}\left(f_\alpha {\displaystyle _V}d^3x\rho \xi _{\alpha ,l}\right)+(\mathrm{\Omega }^2\delta _{ij}\mathrm{\Omega }_i\mathrm{\Omega }_j)f_\alpha {\displaystyle _V}d^3x\rho \xi _{\alpha ,j}`$ (A17)
$`+Gf_\alpha f_\beta {\displaystyle _V}d^3x\rho (𝐱)[\xi _{\beta ,l}(𝐱)\xi _{\alpha ,l}(𝐱)]{\displaystyle \frac{}{x_l}}{\displaystyle _V}{\displaystyle \frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})}{|𝐱𝐱^{}|^3}}`$
$`+\delta {\displaystyle _V}d^3xF_{\alpha \beta ,i}.`$
Although we simplified the final answer by assuming that the fluids occupy identical volumes and have proportional densities in the background state, we could not have derived the correct perturbation of the first order virial theorem if we had not allowed the volumes to differ.
For uniform ellipsoids, we can simplify the mutual gravitational term further. First, we recognize that
$$G_V\frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})}{|𝐱𝐱^{}|^3}=\frac{\varphi (𝐱)}{x_i},$$
(A18)
so we can write the gravitational term generally as
$$f_\alpha f_\beta _Vd^3x\rho (𝐱)[\xi _{\beta ,l}(𝐱)\xi _{\alpha ,l}(𝐱)]\frac{^2\varphi (𝐱)}{x_lx_i}.$$
(A19)
Second, recall that the potential at any interior point of a homogeneous ellipsoid is (Theorem 3 in Chap. 3 of EFE)
$$\varphi (𝐱)=\pi G\rho \left(I\underset{k=1}{\overset{3}{}}A_kx_k^2\right),$$
(A20)
where $`I`$ is a constant; consequently
$$\frac{^2\varphi }{x_lx_i}=2\pi G\rho A_i\delta _{il}.$$
(A21)
Thus, the mutual gravitational contribution to the equation of motion for the perturbed center of mass is
$$2\pi G\rho ^2A_if_\alpha f_\beta _Vd^3x(\xi _{\beta ,i}\xi _{\alpha ,i}).$$
(A22)
A sufficient condition for this to vanish is
$$_Vd^3x𝝃_\alpha =_Vd^3x𝝃_\beta ,$$
(A23)
which is just the statement that the perturbations have the same the center of mass, specialized to the case of uniform density.
In Eq. (A17), we did not compute the variation in the last term. To do this, let us write
$$F_{SN,i}=\rho _S\omega _S\beta _{ij}(u_{S,j}u_{N,j});$$
(A24)
in the background state, the two fluids move together (and may even be stationary in the rotating frame) so we have
$$\delta _Vd^3xF_{\alpha \beta ,i}=𝒮_{\alpha \beta }\frac{d}{dt}\left(f_S_Vd^3x\rho (𝐱)\omega _S\beta _{ij}(\xi _{S,j}\xi _{N,j})\right),$$
(A25)
where $`𝒮_{\alpha \beta }=0`$ if $`\alpha =\beta `$, 1 if $`\alpha =S`$ and $`\beta =N`$, and $`1`$ if $`\alpha =N`$ and $`\beta =S`$, and \[see Eq. (4)\] $`\beta _{ij}=\beta \delta _{ij}+\beta ^{}ϵ_{ijm}\nu _m+(\beta ^{\prime \prime }\beta )\nu _i\nu _j`$. It is clear that the centers of mass of the two fluids remain stationary if the fluid displacements are identical. However, there may be other conditions under which they remain stationary. For example, the integrated mutual friction force will be zero as long as
$$f_S_Vd^3x\rho (𝐱)\omega _S\beta _{ij}(\xi _{S,j}\xi _{N,j})=0;$$
(A26)
for a background with uniform density, vorticity and frictional coupling coefficients, this is guaranteed if
$$_Vd^3x𝝃_S=_Vd^3x𝝃_N,$$
(A27)
which is less restrictive than the requirement of identical displacements. We found the same condition for the vanishing of the integrated, mutual gravitational force for perturbations of uniformly dense ellipsoids. Equation (A27) may be true, in fact, for all of the perturbations considered in our paper, since both sides may vanish identically.
Next, consider variations of the second order virial equation. Most of the terms are varied exactly as for single fluids; one exception is
$`\delta _{\alpha \beta ,ij}`$ $`=`$ $`Gf_\alpha f_\beta \{{\displaystyle _V}d^3x\rho (𝐱)\xi _\alpha (𝐱){\displaystyle \frac{}{x_l}}{\displaystyle _V}{\displaystyle \frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})(x_jx_j^{})}{|𝐱𝐱^{}|^3}}`$ (A29)
$`+{\displaystyle _V}d^3x\rho (𝐱)[\xi _{\alpha ,l}(𝐱)\xi _{\beta ,l}(𝐱)]{\displaystyle \frac{}{x_l}}{\displaystyle }{\displaystyle \frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})x_j^{}}{|𝐱𝐱^{}|^3}}\},`$
where we have specialized to backgrounds with proportional densities and identical bounding volumes. The first term in the brackets can be combined with $`\delta _{\alpha ,ij}`$ and we find
$`\delta _{\alpha ,ij}+(1\delta _{\alpha \beta })\delta _{\alpha \beta ,ij}=Gf_\alpha {\displaystyle _V}d^3x\rho (𝐱)\xi _\alpha (𝐱){\displaystyle \frac{}{x_l}}{\displaystyle _V}{\displaystyle \frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})(x_jx_j^{})}{|𝐱𝐱^{}|^3}}`$ (A30)
$`Gf_\alpha f_\beta {\displaystyle _V}d^3x\rho (𝐱)[\xi _{\alpha ,l}(𝐱)\xi _{\beta ,l}(𝐱)]{\displaystyle \frac{}{x_l}}{\displaystyle _V}{\displaystyle \frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})x_j^{}}{|𝐱𝐱^{}|^3}}.`$ (A31)
The last equation can be written more compactly in terms of the functions
$$_{ij}G_V\frac{d^3x^{}\rho (𝐱^{})(x_ix_i^{})(x_jx_j^{})}{|𝐱𝐱^{}|^3}\frac{𝒟_j}{x_i}G_V\frac{d^3x^{}\rho (𝐱^{})x_j^{}(x_ix_i^{})}{|𝐱𝐱^{}|^3}$$
(A32)
defined in EFE, Chap. 2, Eqs. (14) and (27); using them, we find
$$\delta _{\alpha ,ij}+(1\delta _{\alpha \beta })\delta _{\alpha \beta ,ij}=f_\alpha _Vd^3x\rho \xi _{\alpha ,l}\frac{_{ij}}{x_l}+f_\alpha f_\beta _Vd^3x\rho (\xi _{\alpha ,l}\xi _{\beta ,l})\frac{^2𝒟_j}{x_lx_i}.$$
(A33)
We can rewrite this last result using EFE, Chap. 2 Eq. (28) i.e.
$$\frac{𝒟_j}{x_i}=_{ij}x_j\frac{\varphi }{x_i};$$
(A34)
with this substitution we get
$`\delta _{\alpha ,ij}+(1\delta _{\alpha \beta })\delta _{\alpha \beta ,ij}`$ $`=`$ $`f_\alpha {\displaystyle _V}d^3x\rho \xi _{\alpha ,l}{\displaystyle \frac{_{ij}}{x_l}}`$ (A35)
$`+`$ $`f_\alpha f_\beta {\displaystyle _V}d^3x\rho (\xi _{\alpha ,l}\xi _{\beta ,l})\left({\displaystyle \frac{_{ij}}{x_l}}\delta _{lj}{\displaystyle \frac{\varphi }{x_i}}x_j{\displaystyle \frac{^2\varphi }{x_lx_i}}\right).`$ (A36)
For the uniform ellipsoids \[EFE, Chap. 3, Eqs. (125) and (126)\],
$$\frac{𝒟_j}{\pi G\rho }=a_j^2x_j\left(A_j\underset{k=1}{\overset{3}{}}A_{jk}x_k^2\right)\frac{_{ij}}{\pi G\rho }=2B_{ij}x_ix_j+a_i^2\delta _{ij}\left(A_i\underset{i=1}{\overset{3}{}}A_{ik}x_k^2\right),$$
(A38)
$`(\pi G\rho )^1{\displaystyle \frac{^2𝒟_j}{x_lx_i}}`$ $`=`$ $`2a_j^2(\delta _{ij}x_lA_{jl}+\delta _{jl}x_iA_{ji}+\delta _{il}x_jA_{ji})`$ (A39)
$`=`$ $`2a_j^2[A_{il}(\delta _{ij}x_l+\delta _{jl}x_i)+\delta _{il}x_jA_{ji}],`$ (A40)
$`(\pi G\rho )^1{\displaystyle \frac{_{ij}}{x_l}}`$ $`=`$ $`2B_{ij}(\delta _{il}x_j+\delta _{jl}x_i)2a_i^2\delta _{ij}A_{il}x_l.`$ (A41)
Using these results and the definitions \[EFE, Chap. 2, Eqs. (122), (124), and (125)\]
$$V_{\alpha ,i;j}_Vd^3x\rho \xi _{\alpha ,i}x_jV_{\alpha ,ij}=V_{\alpha ,i;j}+V_{\alpha ,j;i},$$
(A42)
we find
$`{\displaystyle \frac{\delta _{\alpha ,ij}+(1\delta _{\alpha \beta })\delta _{\alpha \beta ,ij}}{\pi G\rho }}=f_\alpha \left(2B_{ij}V_{\alpha ,ij}a_i^2\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}V_{\alpha ,ll}\right)`$ (A43)
$`a_j^2f_\alpha f_\beta \left[2A_{ij}(V_{\alpha ,ij}V_{\beta ,ij})+\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}(V_{\alpha ,ll}V_{\beta ,ll})\right].`$ (A44)
It is possible to write this more compactly, but the form of Eq. (A44) makes clear which terms depend on the differences between the displacements of the two fluids and which do not.
The other new moment we need is
$$\delta _{\alpha \beta ,ij}=\delta _{V_\alpha }d^3xx_jF_{\alpha \beta ,i}.$$
(A45)
Since the two fluids move together in the unperturbed state, this moment is first order in the perturbations at the largest. We then find
$$\delta _{V_\alpha }d^3xx_jF_{\alpha \beta ,i}=𝒮_{\alpha \beta }f_S_Vd^3x\rho \omega _Sx_j\beta _{ik}\left(\frac{d\xi _{S,k}}{dt}\frac{d\xi _{N,k}}{dt}\right).$$
(A46)
For perturbations of uniform ellipsoids, $`\omega _S`$ and $`\rho `$ are independent of position in the unperturbed background, and we may also assume that $`\beta _{ij}`$ is constant; then,
$$\delta _{V_\alpha }d^3xx_jF_{\alpha \beta ,i}=𝒮_{\alpha \beta }f_S\omega _S\beta _{ik}\left[\frac{d}{dt}\left(_Vd^3x\rho x_j(\xi _{S,k}\xi _{N,k})\right)_Vd^3x\rho u_j(\xi _{S,k}\xi _{N,k})\right].$$
(A47)
For backgrounds in which there are no fluid motions, the last term is absent and
$$\delta _{V_\alpha }d^3xx_jF_{\alpha \beta ,i}=𝒮_{\alpha \beta }f_S\rho \omega _S\beta _{ik}\left(\frac{dV_{\alpha ,k;j}}{dt}\frac{dV_{\beta ,k;j}}{dt}\right),$$
(A48)
using the definition in Eq. (A42).
When there are no fluid motions of the unperturbed star in the rotating frame, the second order virial equations are
$`f_S{\displaystyle \frac{d^2V_{S,i;j}}{dt^2}}`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_mf_S{\displaystyle \frac{dV_{S,l;j}}{dt}}+\mathrm{\Omega }^2f_SV_{S,ij}\mathrm{\Omega }_i\mathrm{\Omega }_kf_SV_{S,kj}+\delta _{ij}\delta \mathrm{\Pi }_S`$ (A49)
$``$ $`f_S\pi G\rho \left(2B_{ij}V_{S,ij}a_i^2\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}V_{S,ll}\right)`$ (A50)
$``$ $`a_j^2f_Sf_N\pi G\rho \left[2A_{ij}(V_{S,ij}V_{N,ij})+\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}(V_{S,ll}V_{N,ll})\right]`$ (A51)
$``$ $`f_S\omega _S\beta _{ik}\left({\displaystyle \frac{dV_{S,k;j}}{dt}}{\displaystyle \frac{dV_{N,k;j}}{dt}}\right),`$ (A52)
$`f_N{\displaystyle \frac{d^2V_{N,i;j}}{dt^2}}`$ $`=`$ $`2ϵ_{ilm}\mathrm{\Omega }_mf_N{\displaystyle \frac{dV_{N,l;j}}{dt}}+\mathrm{\Omega }^2f_NV_{N,ij}\mathrm{\Omega }_i\mathrm{\Omega }_kf_NV_{N,kj}+\delta _{ij}\delta \mathrm{\Pi }_N`$ (A53)
$``$ $`f_N\pi G\rho \left(2B_{ij}V_{N,ij}a_i^2\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}V_{N,ll}\right)`$ (A54)
$``$ $`a_j^2f_Nf_S\pi G\rho \left[2A_{ij}(V_{N,ij}V_{Sij})+\delta _{ij}{\displaystyle \underset{l=1}{\overset{3}{}}}A_{il}(V_{N,ll}V_{S,ll})\right]`$ (A55)
$``$ $`f_S\omega _S\beta _{ik}\left({\displaystyle \frac{dV_{N,k;j}}{dt}}{\displaystyle \frac{dV_{S,k;j}}{dt}}\right).`$ (A56)
We can replace these equations with a different set by defining
$$V_{i;j}f_SV_{S,i;j}+f_NV_{N,i;j},U_{i;j}V_{S,i;j}V_{N,i;j}.$$
(A57)
In terms of these new quantities we find Eqs. (24) and (27).
Fig 1: The real (upper panel) and imaginary (lower panel) parts of the relative transverse-shear modes of a superfluid Maclaurin spheroid as a function of eccentricity for three values of $`\eta =0.5,1,50`$. The $`\eta `$ parameter is scaled in units of $`\omega _S\rho _S`$. The imaginary parts of the modes for $`\eta =50`$ are magnified by a factor of 10. The grey lines show the frictionless solutions. To relate the $`\beta `$-coefficients to the rescaled $`\stackrel{~}{\beta }`$-coefficients we have set $`f_S/f_N=0.2`$. The results are insensitive to the choice of this ratio.
Fig 2: The relative toroidal modes of superfluid Maclaurin spheroids. Conventions are the same as in Fig. 1.
Fig 3: The relative pulsation modes of superfluid Maclaurin spheroids. Conventions are the same as in Fig. 1.
Fig 4: The relative odd-parity modes of superfluid Jacobi ellipsoids as a function of $`\mathrm{cos}^1(a_3/a_1)`$. Conventions are the same as in Fig. 1.
Fig 5: The relative even-parity modes of superfluid Jacobi ellipsoids as a function of $`\mathrm{cos}^1(a_3/a_1)`$. Conventions are the same as in Fig. 1.
Fig 6: The relative odd-parity modes of superfluid Roche ellipsoids as a function of $`\mathrm{cos}^1(a_3/a_1)`$ for $`P=1`$. Conventions are the same as in Fig. 1.
Fig 7: The relative even-parity modes of superfluid Roche ellipsoids as a function of $`\mathrm{cos}^1(a_3/a_1)`$ for $`P=1`$. Conventions are the same as in Fig. 1.
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# Magnetic Screening in Hot Non-Abelian Gauge Theory
## Abstract
We analyze the large distance and low-momentum behavior of the magnetic gluon propagator of the SU(2) gauge theory at finite temperature. Lattice calculations within the $`4`$-dimensional as well as the effective, dimensionally reduced $`3`$-dimensional gauge theories in generalized Landau gauges and MAG show that the magnetic propagator is strongly infrared suppressed in Landau gauges but stays large and finite in MAG. Despite these differences in the low-momentum behavior of the propagator calculated in different gauges the magnetic fields are exponentially screened in all gauges considered. From the propagator calculated in maximally Abelian gauge we find for the screening mass, $`m_M=(1.48\pm 0.17)T`$ at $`T=2T_c`$.
PACS numbers: 11.10.Wx, 11.15.Ha, 12.38.Mh
Magnetic fields in hot non-abelian gauge theories have to be screened at high temperature for any perturbative description of the plasma phase to make sense . In fact, there is plenty of evidence from analytic as well as numerical calculations for the screening of electric and magnetic fields in non-abelian gauge theories at high temperature. The way this is realized on the partonic level, however, is poorly understood. Our intuitive understanding of screening at high temperature is largely influenced by the perturbative concept of electric and magnetic screening masses . However, while the former is calculable at least to leading order in perturbation theory the latter is intrinsically non-perturbative. Its existence has been postulated to render the perturbative expansion finite .
In the vacuum the gluon propagator in Landau gauge is expected to be infrared suppressed. The existing rigorous bounds on the unrenormalized gluon propagator in momentum space imply that it is less singular than $`p^2`$ in $`4`$ dimension and $`p^1`$ in $`3`$ dimension and it has further been argued that it is likely to vanish at zero momentum in the thermodynamic limit . Numerical evidence for this has recently been found in lattice simulations of the SU(N) gauge theory in $`3`$ and in $`4`$ dimensions . As it is well established that the $`4d`$, $`SU(N)`$ gauge theory at finite temperature is well described by an effective, $`3d`$ gauge theory with an adjoint Higgs field and that, moreover, the magnetic sector is little influenced by the scalar Higgs fields , the findings for the gluon propagator in $`3`$ dimensions are of direct relevance for the behavior of the magnetic gluon propagator in the $`4d`$ gauge theory at finite temperature. In fact, it recently has been suggested that also at finite temperature the gluon propagator might be infrared suppressed .
It is the purpose of this letter to clarify the long-distance, low-momentum structure of the magnetic gluon propagator of the finite temperature $`SU(2)`$ gauge theory in $`4d`$, establish the connection to the $`3d`$ gauge theory and analyze to what extent these findings are gauge dependent. To this end we will analyze the magnetic gluon propagator in Landau as well as maximally Abelian gauges (MAG). Both gauges are complementary in so far as the bounds on the zero momentum limit of the gluon propagator derived in Landau gauge are not valid in MAG because in this case the Fadeev-Popov matrix turns out to be quadratic in the gauge fields. Although details of the infrared behavior will turn out to be different in different gauges we show that in all cases considered magnetic fields are exponentially screened at large distances. In particular, we will make contact to earlier numerical calculations, which gave evidence for exponential damping of the magnetic propagator at large distances and led to the determination of a magnetic screening mass .
We have performed numerical calculations for the $`4d`$, $`SU(2)`$ gauge theory at finite temperature as well as for the effective, $`3d`$ gauge theory with an adjoint Higgs field and the pure $`SU(2)`$ gauge theory in $`3d`$. The magnetic gluon propagator, $`D_M(z)`$, and its Fourier transform, $`\stackrel{~}{D}_M(p)`$, have been calculated in Landau gauge and MAG on lattices of size $`\mathrm{\Omega }=N_\sigma ^2N_z`$, in $`3d`$ and $`\mathrm{\Omega }=N_\sigma ^2N_zN_\tau `$ in $`4d`$, respectively. In our $`3d`$ simulations we also have considered an anisotropic generalization of the Landau gauge defined by the gauge condition,
$$_1A_1+_2A_2+\lambda _3_3A_3=0,$$
(1)
which for $`\lambda _3=1`$ reduces to the ordinary Landau gauge condition. For the gauge fields we use the straightforward lattice definition, $`A_\mu (x)=\left[U_\mu (x)U_\mu ^{}(x)\right]/2i`$, with $`U_\mu (x)SU(2)`$, which reproduces the continuum gauge fields up to $`𝒪(a^2)`$ discretization errors. We have checked previously that, up to an overall normalization, other lattice definitions with formally smaller discretization errors lead to identical results .
The correlation function for the $`\mu `$-component of the gauge fields $`A_\mu `$ is defined in terms of the sum over gauge fields in a hyperplane orthogonal to $`x_3`$,
$`D_{\mu \mu }(z)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Omega }}}{\displaystyle \underset{b=1}{\overset{3}{}}}{\displaystyle \underset{x_3}{}}Q_\mu ^b(x_3+z)Q_\mu ^b(x_3)\text{,}`$ (2)
with $`Q_\mu ^b(x_3)=_x_{}\mathrm{Tr}[A_\mu (x_{}\text{,}x_3)\sigma ^b]/2`$. Here $`\sigma ^b`$ are the usual Pauli matrices; $`x_{}=(x_1,x_2)`$ in $`3d`$ and $`x_{}=(x_0,x_1,x_2)`$ in $`4d`$, respectively. In terms of this we define the magnetic gluon propagator in coordinate space,
$`D_M(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(D_{11}(z)+D_{22}(z)),`$ (3)
and the momentum space propagator
$$\stackrel{~}{D}_M(p)=\frac{\beta a^2}{12}\underset{\mu =1}{\overset{3}{}}\underset{z=0}{\overset{N_z1}{}}\mathrm{e}^{2\pi ikz/N_z}D_{\mu \mu }(z),$$
(4)
with $`p=2|\mathrm{sin}(\pi k/N_z)|`$ and $`k=0,1,\mathrm{},N_z1`$. Here $`\beta `$ is the coupling appearing in the Euclidean lattice action and $`a`$ is the lattice spacing.
The $`4d`$ calculations at finite temperature have been performed in the high temperature phase at twice the critical temperature for the deconfinement transition. All calculations have been performed with the standard Wilson gauge action with gauge coupling $`\beta =2.512`$ for $`N_\tau =4`$ and $`\beta =2.740`$ for $`N_\tau =8`$. These values correspond to $`T=2T_c`$ . Corresponding calculations in $`3d`$ have been performed at $`\beta =8`$ on lattices of size $`16^2\times 32`$, $`24^2\times 48`$, $`28^2\times 56`$, $`32^2\times 64`$, $`48^2\times 64`$ and $`64^3`$. Moreover, we have performed $`3d`$ calculations at a smaller coupling, $`\beta =5`$, on lattices of size $`32^3`$ up to $`96^3`$. This choice of parameters for calculations within the reduced theory allows to analyze the behavior of correlation functions in large physical volumes which are 40 times larger than what is feasible within our $`4d`$ simulations. The gauge coupling of the pure gauge theory was chosen according to $`g_3^2(T)=g^2(\mu )T`$, where $`g^2(\mu )`$ is 1-loop running coupling constant of the $`4d`$ $`SU(2)`$ gauge theory in the $`\overline{MS}`$ scheme, $`\mu =18.86T`$ and for $`\mathrm{\Lambda }_{\overline{MS}}`$ the numerically determined relation to the deconfinement temperature has been used, i.e. $`T_c=1.06\mathrm{\Lambda }_{\overline{MS}}`$ . This choice of $`g_3^2`$ ensures a perfect description of the temperature dependence of the spatial string tension of the $`4d`$ finite temperature theory in terms of the effective $`3d`$ theory down to temperatures $`T2T_c`$. Furthermore, we have checked on our smaller lattices that the magnetic propagators calculated in the pure $`3d`$ $`SU(2)`$ gauge theory and in the complete $`3d`$ effective theory, i.e. in the presence of an adjoint Higgs field, agree within statistical errors. We therefore will restrict our discussion here to results coming from the pure gauge sector alone. Our $`3d`$ calculations have been performed at $`\beta =5`$ and 8 which correspond to lattice spacings $`a=0.2752/T`$ and $`0.172/T`$, respectively. Varying $`\beta `$ in our $`3d`$ calculations and converting to physical units with the above relation for $`g_3^2`$ we have checked that cut-off effects are small. To verify the cut-off independence of our results in $`4d`$ the calculations have been performed at two different values of the lattice cut-off, i.e. we used lattices of size $`N_\sigma ^2N_zN_\tau `$ with $`N_\tau =4`$ and 8, $`N_z=2N_\sigma `$ and $`8N_\sigma 32`$. Keeping the physical volume $`VT^3=2(N_\sigma /N_\tau )^3`$ constant we have verified that $`D_M(z)`$ is cut-off independent within the statistical errors of our calculation.
The gluon propagator has been calculated on 1000 to 8000 gauge fixed gauge field configurations in the case of $`3d`$ and about 400 configurations in $`4d`$ calculations. In Fig. 1 we show the magnetic propagator, $`D_M(z)`$, calculated in Landau and maximally Abelian gauges on various size lattices. We note that the $`4d`$ and $`3d`$ calculations performed on lattices of similar physical size are in good agreement. While the magnetic gluon propagator calculated in MAG shows no significant volume dependence the Landau gauge propagators are strongly volume dependent for $`zT>\mathrm{\hspace{0.33em}1}`$.
A similarly strong volume dependence we also observe in other generalized Landau gauges, i.e. gauges based on the gauge condition given in Eq. 1 with $`\lambda _31`$. In fact, at large distances the propagators in generalized Landau gauges, including the Coulomb gauge limit, become negative.
The virtue of the effective $`3d`$ theory is that we can reach here a much larger physical volume and thus can control the volume dependence still apparent in the comparison of $`3d`$ and $`4d`$ data shown in Fig. 1. We have analyzed the large distance behavior of $`D_M(z)`$ in these gauges in our $`3d`$ calculations on large lattices of size $`L^3`$, $`L=32`$, 40, 48, 56, 64, 72 and 96 at $`\beta =5`$. This corresponds to physical volumes of size $`VT^3=683`$ up to $`VT^3=5464`$. In all cases we clearly observe negative correlation functions at large distances. For some of these lattices we show the long-distance part of the $`D_M(z)`$ in Fig. 2 for $`\lambda _3=1`$. Similar results we find for $`\lambda _31`$. We note that the volume dependence is small on these large lattices. Moreover, the correlation functions become smaller with increasing lattice size. This is in accordance with our expectation that finite size effects are due to zero mode contributions which disappear in the infinite volume limit. On finite lattices, however, they add a volume-dependent, positive constant to the correlation function in coordinate space.
Our analysis thus gives clear evidence that the magnetic gluon correlation function in coordinate space calculated in generalized Landau gauges becomes negative at large distances, i.e. for $`zT>1.5`$. At $`T=0`$ such a behavior is, in fact, expected to occur for the gluon propagator as a consequence of the conjecture made by Zwanziger that the Landau gauge gluon propagator vanishes in the infrared limit. In coordinate space it has been suggested that the propagator takes on the form
$$D_M(z)=A\mathrm{e}^{m_Mz}\mathrm{cos}(bm_Mz+c),$$
(5)
Although we can at present not quantify the functional form of the correlation function in generalized Landau gauges, our results are consistent with this ansatz up to distances $`zT3`$. To judge whether the correlation function will indeed continue to oscillate as suggested by Eq. 5 will, however, require a more detailed analysis of the large distance behavior of $`D_M(z)`$. In general, for $`\stackrel{~}{D}_M(p)`$ to vanish at $`p=0`$ in momentum space one should find $`\mathrm{tan}(c)=1/b`$ with the above ansatz. This conjecture, $`lim_{p0}\stackrel{~}{D}_M(p)=0`$, may also be valid at non-zero temperature .
In contrast to the complicated structure of the Landau gauge propagator the propagator calculated in MAG does show a simple exponential decay at large distances and does not show any significant volume dependence. We could verify that it stays positive at least up to distance $`zT=3.5`$. For $`zT>1.5`$ simple exponential fits gave $`\chi ^2/(d.o.f)=1.7`$ and we also found that local masses became independent of $`z`$ within statistical errors for $`zT>1`$. Best fits with the ansatz given in Eq. 5 were thus obtained for $`b=0`$. From this we find for the magnetic screening mass of the magnetic propagator calculated in MAG,
$$\frac{m_M}{T}=1.48\pm 0.17\mathrm{at}T=2T_c.$$
(6)
When taking the absolute value of the Landau gauge propagator this is found to be bounded by the exponentially decaying MAG-propagator. This confirms that also in Landau gauge correlations of magnetic fields are exponentially screened.
The fact that the magnetic gluon propagator calculated in Landau gauge becomes negative at large distances translates into a suppression of the momentum space propagator at small momentum. This is shown in Fig. 3 for the Landau gauge propagators, similar results hold for $`\lambda _31`$. For $`p<T`$ the momentum space propagator is sensitive to the volume, while for large momenta ($`p>T`$) it is essentially independent of the lattice size. Moreover, we note that for small momenta the finite volume effects lead to a decrease of $`\stackrel{~}{D}(p)`$ with increasing volume while the volume dependence is already negligible for $`p/T1`$. The occurrence of a maximum in $`\stackrel{~}{D}(p)`$ for $`p/T1`$ and the infrared suppression of the momentum space propagator at $`p=0`$ thus is firmly established by our data.
As expected the volume dependence of the magnetic propagator in momentum space is strongest at $`p=0`$. On large lattices $`\stackrel{~}{D}_M(p)`$ reaches a maximum at a non-zero but small value of $`p/T`$ and keeps decreasing at $`p=0`$ with increasing lattice size. This might indicate that $`\stackrel{~}{D}_M(0)`$ could actually vanish in the infinite volume limit. The propagator at vanishing momentum reflects the contribution of zero momentum fluctuations of the gauge fields,
$$\stackrel{~}{D}_M(0)=\frac{\beta a^2}{12}\mathrm{\Omega }\underset{\mu ,b}{}(\varphi _\mu ^b)^2\mathrm{with}\varphi _\mu ^b=\frac{1}{\mathrm{\Omega }}\underset{x}{}A_\mu ^b(x).$$
(7)
In order for it to vanish in the infinite volume limit the zero mode fluctuations $`(\varphi _\mu ^b)^2`$ thus have to drop faster than $`\mathrm{\Omega }^1`$. To quantify this behavior we have fitted $`\stackrel{~}{D}_M(0)T^2`$ to a simple ansatz that assumes a power-like dependence on the volume, $`\stackrel{~}{D}_M(0)T^2=b(VT^3)^z+c`$. At present such fits to our $`4d`$ and $`3d`$ data can, however, not rule out a non-vanishing value for the constant $`c`$. We stress, however, that the value of $`\stackrel{~}{D}_M(0)`$ itself is of little importance for the observed complex structure of the gluon propagator in Landau gauge in coordinate as well as in momentum space and in any case is gauge dependent. Nonetheless, the behavior found here for the finite-T magnetic gluon propagator is consistent with findings of Ref. for the $`T=0`$ gluon propagator of the $`3d`$, $`SU(2)`$ gauge theory. There it was shown that the magnetic propagator calculated in Landau gauge is less singular then $`p^1`$ in the infrared and is likely to vanish in $`p0`$, $`V\mathrm{}`$ limit.
The momentum dependence of $`\stackrel{~}{D}_M(p)`$ shown in Fig. 3 also makes clear why earlier calculations on medium size lattices let to the determination of a non-zero magnetic mass. The infrared suppression becomes significant only for momenta $`p/T<0.5`$. The existence of a simple pole mass becomes unlikely on the basis of our analysis. However, we stress that magnetic fields in coordinate space are, of course, strongly screened. The magnetic gluon propagator calculated in Landau gauge as well as MAG is exponentially damped. The value for the screening mass extracted from the propagator in MAG gauge is in good agreement with earlier in $`4d`$ and $`3d`$ lattice calculations in Landau gauge, although the data analysis in these cases was based on a simple pole ansatz for the magnetic mass.
In the course of the calculations reported here we also have reanalyzed the behavior of the electric gluon propagator in different gauges. In that case no suppression in the infrared has been observed and the electric screening mass has been found to be gauge independent within statistical errors. Details of this calculation as well as a more detailed discussion of our results for the magnetic gluon propagator will be presented elsewhere.
Acknowledgements:
The work has been supported by the TMR network ERBFMRX-CT-970122 and by the DFG under grant Ka 1198/4-1. Our calculations have partly been performed at the HLRS in Stuttgart and the $`(PC)^2`$ in Paderborn.
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# 1. Introduction
## 1. Introduction
The observation of the large tunneling magnetoresistance effect at room temperature in tunnel junctions of the form M/O/M’ (where M and M’ are magnetic metals and O is an oxide tunnel barrier) has stimulated a renewed interest for these systems . Besides the fundamental interest for spin-polarized transport, these structures are also foreseen as potential candidates for sensitive magnetic sensors and memory cells in random access memory devices. The first model of spin-dependent tunneling in the framework of classical quantum mechanics was proposed by Słonczewski . However, in this approach no scattering of electrons in the magnetic metallic electrodes was taken into account. This model has been subsequently developed in Refs. by using the Kubo formalism of linear response. The effects of elastic impurity scattering inside the metallic layers and at interfaces between the dielectric and conductive layers could then be incorporated in the model. On the other hand, it is well known that the presence of impurities inside the potential barrier can lead to the mechanism of resonant tunneling when the localized electronic states within the gap of the insulator formed by embedded atoms lie close to the chemical potential of the system. This situation was qualitatively studied in mesoscopic semiconductor system in case of one- and two-impurity resonant channels by means of classical quantum mechanical treatment. The same approach afterwards has been used in with the application to the impurity-assisted tunneling magnetoresistance (TMR). The numerical analysis of this problem which was carried out in Refs. should also be mentioned. In paper only the case of spinless impurities was considered, and the author came to the conclusion that the TMR amplitude decreased due to the impurity assisted tunneling. The problem of the paramagnetic impurity assisted tunneling in tunnel magnetic junctions was investigated recently in Ref. . The author claim, that he investigated resonant tunneling through impurity resonance level, however he did not introduce the line-width of this resonance, which, as it will be shown below, does depend on the position of the impurity atom inside the barrier as well on the magnetic configuration of the magnetic layers. As it will be shown below, namely these line-widths define the value of the tunneling conductance and the amplitude of the TMR for spin-conserving and spin-flip resonant tunneling. An attempt of the analysis the same problem has also been undertaken in Ref. , but nevertheless the microscopic mechanism of electron scattering on the paramagnetic impurity has not been taken into account.
In this paper, we propose a renewed study of the problem of impurity-assisted tunneling in spin-valve junctions of the form F<sub>1</sub>/O/F<sub>3</sub>, where F’s are ferromagnetic electrodes and O is an insulating barrier with embedded paramagnetic impurities, that incorporates the effect of both elastic and non-elastic spin-flip scattering due to the exchange interaction between the itinerant electrons forming the tunneling current and the localized spins of impurities. It will be shown that non-elastic scattering has an essential impact not only on the temperature variation of the TMR (which is a well established result ) but also on the I–V characteristics of considered structures. The latter effect was predicted in Ref. , where the TMR dependence on the electron scattering on interfacial magnons was investigated.
## 2. Model
### 2.1. Kubo formula and general expression for the conductivity of the system
The following simplified model is adopted throughout the paper. First of all, the thickness of an oxide layer is supposed to be much smaller than its in-plane dimension, so that the system may be considered as homogeneous in the $`xy`$–plane (parallel to the interfaces) and inhomogeneous only in the $`z`$–direction (growth direction). Within each layer, the electrons are described as a free-electron gas and they undergo scattering on the 3-D $`\delta `$-function impurity potential within the insulating barrier. Within these approximations, the Hamiltonian of the system has the form
$$\widehat{H}=\widehat{H_0}+\widehat{H}_{\mathrm{int}},$$
where
$`\widehat{H}_0={\displaystyle \frac{\mathrm{}^2}{2m(z)}}\mathrm{\Delta }+U(z)2\mu _BH_z^{\mathrm{eff}}(z)(\widehat{s_z}+\widehat{S_z})`$ (1)
$`\widehat{H}_{\mathrm{int}}={\displaystyle \underset{i}{}}a_0^3\delta (𝐫𝐜_i)\left\{\epsilon _0J(\mathrm{𝐬𝐒})\right\}.`$
Here the summation is performed over the location of impurities $`𝐜_i`$ inside the barrier, $`a_0`$ is the lattice constant, $`\epsilon _0`$ denotes the scattering potential amplitude on the impurity, $`J`$ is responsible for the s–d type exchange interaction between a conduction electron spin $`𝐬`$ and the impurity spin $`𝐒`$, $`U(z)`$ is a model step-like potential seen by the conduction electron as it is represented in Fig. 1. We take into account the exchange splitting of the $`d`$–band by introducing different values $`V_\mu ^{1,3}`$ for the position of the bottom of the conduction band in F<sub>1</sub> and F<sub>3</sub>, depending on the mutual orientation of magnetization in the layers and the spin $`\mu =,`$ of the conduction electron. $`H_z^{\mathrm{eff}}(z)`$ represents the effective field acting on impurity and electron spins inside the barrier. The origin of this field is the super-exchange between the spins in the bulk of ferromagnetic layer and in insulating layer. We suppose that $`H_z^{\mathrm{eff}}(z)`$ decreases exponentially with the distance from the interface in the depth of the oxide layer. $`m(z)`$ corresponds to the effective electron mass that we suppose is equal to $`m`$ in the ferromagnetic layers and to $`m_0`$ in the insulator. Throughout the paper, it is expressed in units of bare electron mass $`m_e`$. We also assume that the mass of free-like electrons in ferromagnet is slightly differs from $`m_e`$, i.e. $`m1`$ and we will eliminate it from all subsequent expressions.
We start from the Keldysh technique for Green functions together with Kubo exact formula of linear response theory for the static conductivity which relates its real part with the current-current correlation function and may be written in the form
$`\sigma _{\mu \rho }(𝐫,𝐫^{})={\displaystyle \frac{1}{2k_BT}}{\displaystyle _{\mathrm{}}^+\mathrm{}}j_\rho (𝐫^{},t^{})j_\mu (𝐫,t)d(tt^{})`$
$`={\displaystyle \frac{1}{2k_BT}}\left({\displaystyle \frac{e\mathrm{}}{2m}}\right)^2{\displaystyle _{\mathrm{}}^+\mathrm{}}G_{\mu \rho }^<(𝐫,t,𝐫^{},t^{})\underset{𝐫}{\overset{}{}}\underset{𝐫^{}}{\overset{}{}}G_{\rho \mu }^>(𝐫^{},t^{},𝐫,t)d(tt^{}),`$ (2)
where $`\mu `$, $`\rho `$ denote the projections of the spin of the electrons, $`\underset{𝐫}{\overset{}{}}=\frac{1}{2}(\underset{𝐫}{\overset{}{}}\underset{𝐫}{\overset{}{}})`$ is the asymmetric gradient operator and $`G_{\mu \rho }^<`$ and $`G_{\rho \mu }^>`$ are corresponding Green functions in Keldysh formalism . $`\mathrm{}`$ represents the quantum statistical averaging over the distribution of impurities and degrees of freedom of impurity spin. This expression is most general and holds both for the elastic impurity and defect scattering or inelastic, including magnon and phonon, scattering. To evaluate expression (2) one needs to introduce the retarded Green function $`G_{\mu \rho }^R(z,z^{},\kappa ,\epsilon )`$ that in our case is defined by the following differential equation:
$$\left\{\epsilon +\frac{\mathrm{}^2}{2m}\frac{^2}{z^2}\frac{\kappa ^2}{2m}U(z)\right\}G_{\mu \rho }^R(z,z^{},\kappa ,\epsilon )=\delta _{\mu \rho }\delta (zz^{})$$
in the mixed real-space momentum representation , where $`\kappa =(\kappa _x,\kappa _y)`$ is the component of the electron momentum in $`xy`$-plane of the layers and $`z`$ is the coordinate perpendicular to $`xy`$-plane. We should note, that by definition the conductivity (2) is defined as a linear response on the externally applied electric field and does not depend on $`z`$ and $`z^{}`$ because of the obvious condition $`{\displaystyle \frac{j(z)}{z}}=0`$.
Let now denote $`k_1^\mu =\sqrt{2(\epsilon V_1^\mu )}`$, $`k_3^\mu =\sqrt{2(\epsilon V_3^\mu )}`$ momenta of electrons with energy $`\epsilon `$ and spin $`\mu `$ in the ferromagnetic layers and $`q_0=\sqrt{2m_0(U\epsilon )}`$ is an imaginary momentum inside the barrier. By introducing the following functions on $`x=\kappa /q_0`$:
$$c_1^\mu (x)=\sqrt{k_1^{\mu \mathrm{\hspace{0.17em}2}}q_0^2x^2},c_3^\mu (x)=\sqrt{k_3^{\mu \mathrm{\hspace{0.17em}2}}q_0^2x^2},q_2(x)=q_0\sqrt{1+x^2},$$
the final expression for the conductance of the system, comprising only one impurity, located at point $`𝐜`$, at given temperature $`T`$ is written as follows:
$$\sigma (T,𝐜)=\sigma _0(T)+\sigma ^{\mathrm{imp}}(T,𝐜).$$
The first term is given by
$$\sigma _0(T)=\frac{q_0^2e^2}{2\pi \mathrm{}}\underset{\mu }{}_{\mathrm{}}^+\mathrm{}𝑑\epsilon \left(\frac{f(\epsilon )}{\epsilon }\right)_0^{x_0^\mu }\frac{xdx}{2\pi }\frac{16c_1^\mu c_3^\mu m_0^2q_2^2e^{2q_2w}}{(m_0^2c_1^{\mu \mathrm{\hspace{0.17em}2}}+q_2^2)(m_0^2c_3^{\mu \mathrm{\hspace{0.17em}2}}+q_2^2)},$$
(3)
where $`x_0^\mu =\mathrm{min}\{\frac{k_1^\mu }{q_0},\frac{k_3^\mu }{q_0}\}`$, $`f(\epsilon )=\left[1+e^{\beta (\epsilon \epsilon _F)}\right]^1`$ is Fermi function, and $`w=ba`$ is the width of the insulating spacer. It represents the well known result for the pure tunneling conductance . The second term $`\sigma ^{\mathrm{imp}}(T,𝐜)`$ is directly related to the impurity assisted tunneling. It is convenient to write it down as a sum of two contributions:
$$\sigma ^{\mathrm{imp}}(T,𝐜)=\sigma _{\mathrm{el}}^{\mathrm{imp}}(T,𝐜)+\sigma _{\mathrm{sf}}^{\mathrm{imp}}(T,𝐜),$$
where the first term corresponds to the conductivity due to elastic spin-conserving processes of scattering electron on the impurity and the second one summarizes all other events when the conduction electron changes its spin after tunneling through the barrier. We have derived the analytical expressions for these two terms, which are valid under two assumptions. Namely, under domination of single electron scattering on impurities over multiple scattering of two and more electrons on the same center and under absence of polarization of impurity spin induced by the ejection of spin-polarized electrons. Then the final result for these terms is written as (the details of its derivation are outlined further)
$`\sigma _{\mathrm{el}}^{\mathrm{imp}}(T,𝐜)={\displaystyle \frac{1}{S}}\left({\displaystyle \frac{2e^2}{\pi \mathrm{}}}\right){\displaystyle _{\mathrm{}}^+\mathrm{}}d\epsilon \{{\displaystyle \frac{f_{}(\epsilon \mu _BH_z^{\mathrm{eff}})}{\epsilon }}(\widehat{t}_z^{}(\epsilon ))^{}\widehat{t}_z^{}(\epsilon )\mathrm{\Phi }_{}^L(𝐜)\mathrm{\Phi }_{}^R(𝐜)`$
$`{\displaystyle \frac{f_{}(\epsilon +\mu _BH_z^{\mathrm{eff}})}{\epsilon }}(\widehat{t}_z^{}(\epsilon ))^{}\widehat{t}_z^{}(\epsilon )\mathrm{\Phi }_{}^L(𝐜)\mathrm{\Phi }_{}^R(𝐜)\},`$ (4)
$`\sigma _{\mathrm{sf}}^{\mathrm{imp}}(T,𝐜)={\displaystyle \frac{1}{S}}\left({\displaystyle \frac{2e^2}{\pi \mathrm{}}}\right){\displaystyle \frac{1}{k_BT}}{\displaystyle _{\mathrm{}}^+\mathrm{}}d\epsilon \{f_{}(\epsilon \mu _BH_z^{\mathrm{eff}})[1f_{}(\epsilon +\mu _BH_z^{\mathrm{eff}})]\widehat{t}_{}(\epsilon )\widehat{t}_+(\epsilon )+`$
$`f_{}(\epsilon +\mu _BH_z^{\mathrm{eff}})[1f_{}(\epsilon \mu _BH_z^{\mathrm{eff}})]\widehat{t}_+(\epsilon )\widehat{t}_{}(\epsilon )\}\times {\displaystyle \frac{1}{2}}\{\mathrm{\Phi }_{}^L(𝐜)\mathrm{\Phi }_{}^R(𝐜)+\mathrm{\Phi }_{}^L(𝐜)\mathrm{\Phi }_{}^R(𝐜)\}.`$
Here $`S`$ is the junction area, $`\mathrm{\Phi }_{()}^L(𝐜)`$ and $`\mathrm{\Phi }_{()}^R(𝐜)`$ are the probabilities of tunneling of the electron from the left or from the right electrode to impurity, located at point $`𝐜`$. Omitting the exponentially small terms, the expression for these probabilities can be written as
$`\mathrm{\Phi }_\mu ^L(𝐜)`$ $`=`$ $`{\displaystyle _0^{x_{\mathrm{max}}^\mu }}{\displaystyle \frac{xdx}{2\pi }}{\displaystyle \frac{2c_1^\mu m_0^2q_0^2}{(m_0^2c_1^{\mu \mathrm{\hspace{0.17em}2}}+q_2^2)}}e^{2q_2(ca)},`$
$`\mathrm{\Phi }_\rho ^R(𝐜)`$ $`=`$ $`{\displaystyle _0^{x_{\mathrm{max}}^\rho }}{\displaystyle \frac{xdx}{2\pi }}{\displaystyle \frac{2c_3^\rho m_0^2q_0^2}{(m_0^2c_3^{\rho \mathrm{\hspace{0.17em}2}}+q_2^2)}}e^{2q_2(bc)}.`$
The quantities $`(\widehat{t}_z^{()})^{}\widehat{t}_z^{()}(\epsilon )`$ and $`\widehat{t}_{}\widehat{t}_+(\epsilon )`$, $`\widehat{t}_+\widehat{t}_{}(\epsilon )`$ in (4) represent the scattering amplitudes of electron on the impurity center for the case of spin conserving $`(|in,|out,)`$ or $`|in,|out,)`$ and spin-flip $`(|in,|out,`$ or $`|in,|out,)`$ transitions averaged over the distribution of paramagnetic impurity spin. Here $`|in`$ and $`|out`$ denote the initial and final states of impurity. Operators $`\widehat{t}_z^{()}`$ and $`\widehat{t}_\pm `$ form a one-center matrix $`\widehat{t}=\left(\genfrac{}{}{0pt}{}{\widehat{t}_z^{}}{\widehat{t}_+}\genfrac{}{}{0pt}{}{\widehat{t}_{}}{\widehat{t}_z^{}}\right)`$ in the direct product of the linear subspaces of electron’s and impurity’s spins and are expressed as
$`\widehat{t}_z^{()}(\epsilon )={\displaystyle \frac{1}{1\widehat{V}_z^{()}(\epsilon )G_{()}(\epsilon )}}\widehat{V}_z^{()}(\epsilon ),`$
$`\widehat{t}_\pm (\epsilon )={\displaystyle \frac{1}{1\widehat{V}_z^{()}(\epsilon )G_{()}(\epsilon )}}\widehat{S}_\pm {\displaystyle \frac{a_0^3J/2}{1a_0^3(\epsilon _0+\frac{1}{2}J\widehat{S}_z)G_{()}(\epsilon )}},`$ (5)
where effective potentials $`\widehat{V}_z^{()}`$ are given by
$$\widehat{V}_z^{()}(\epsilon )=a_0^3\left\{\epsilon _0\frac{1}{2}J\widehat{S}_z+\frac{1}{4}\widehat{S}_{}\frac{a_0^3G_{()}(\epsilon )J^2}{1a_0^3(\epsilon _0\pm \frac{1}{2}J\widehat{S}_z)G_{()}(\epsilon )}\widehat{S}_\pm \right\}.$$
Here $`G_{()}(\epsilon )`$ is the electron Green function at point $`𝐜`$:
$$G_\mu (\epsilon ,𝐜)=_0^{\kappa _{\mathrm{max}}}G_{\mu \kappa }(\epsilon ,𝐜)\frac{\kappa d\kappa }{2\pi },$$
where $`\kappa _{\mathrm{max}}=\frac{2\sqrt{\pi }}{a_0}`$ is a cut-off of in plane momentum that stems from the finite size of Brillouin zone (we substitute the Brillouin zone’s projection onto $`(\kappa _x,\kappa _y)`$ plane by the circle of radius $`\kappa _{\mathrm{max}}`$ of the same square in $`\kappa _{}`$–plane). The real and imaginary part of $`G_\kappa ^\mu (\epsilon ,𝐜)`$ ($`\mu `$ is the spin index) in the leading order of magnitude are given by
$`\mathrm{Re}G_\kappa ^\mu (\epsilon ,𝐜)={\displaystyle \frac{m_0}{q_2}},`$ (6)
$`\mathrm{Im}G_\kappa ^\mu (\epsilon ,𝐜)=\left(\mathrm{\Phi }_\mu ^L(𝐜)+\mathrm{\Phi }_\mu ^R(𝐜)\right).`$
Let us now explain the derivation of expression (4) and clarify the two assumptions under which this formula is valid. To derive (4) from the starting point (2) one can first of all examine two diagrams (a) and (b) (see Fig. 2) that contribute to spin-conserving and spin-flip part of $`\sigma ^{\mathrm{imp}}(T,𝐜)`$ at second order of $`J`$, respectively. One may easy verify that a general structure of these diagrams is just the same as the final result in form (4) with the mere difference that the one-center $`t`$–matrix is reduced at first order of $`J`$ to the initial potential $`\epsilon _0\frac{J}{2}\left(\genfrac{}{}{0pt}{}{\widehat{S}_z}{\widehat{S}_+}\genfrac{}{}{0pt}{}{\widehat{S}_{}}{\widehat{S}_z}\right)`$.
Moreover, the diagram (b) contains both direct and indirect processes in equal proportion with common factor 1/2 for any of the possible channels $`|in,|out,`$ or $`|in,|out,`$. The thermodynamic averaging $`\mathrm{}`$ in the second-order expansion is simply reduced to the averaging over the Boltzman distribution of the impurity spin in the ”external” effective magnetic field $`H_z^{\mathrm{eff}}`$ which was introduced in (1), i.e. with the density matrix $`\widehat{\rho }_0=Z^1\mathrm{exp}\left\{\frac{2\mu _BH_z^{\mathrm{eff}}\widehat{S}_z}{k_BT}\right\}`$, where $`Z=2\mathrm{cosh}\left(\frac{2\mu _BH_z^{\mathrm{eff}}}{k_BT}\right)+1`$. After that, it is easy to check that the total probabilities (with account of Fermi factors of electron states) of direct and inverse processes are equal which means that the principle of detailed equilibrium holds. In particular it leads to the vanishing of spin-flip processes in a system at zero temperature and vanishing voltage bias.
After that preliminary discussion two assumptions should be made to justify the result (4):
i) We assume that the occupation of the given impurity center simultaneously by two electrons with different spin (due to Pauli principle) is a rather rare event or, in other words, we do not take into account many-electrons effects. It may be justified: a) by Coulomb interaction between electrons that make unprofitable their arrangement at the same site of the lattice; b) by the large number of impurity centers that provides a sufficiently large number of one-step channels so that electrons may be considered as independent.
ii) We also neglect the influence of electron current on the statistical distribution of paramagnetic spins inside the oxide barrier. This assumption is valid for practical intensity of tunneling current which is low enough not to produce a spin polarization of impurities by injection of spin-polarized charge carriers.
Under these assumptions the expression (4) can be obtained by simple substitution of scattering potential $`\widehat{H}_{\mathrm{int}}`$ at site $`𝐜_i`$ on Fig. 2 by the corresponding one-center $`t`$–matrix in accordance with (LABEL:t-matrix) and assuming that the averaging over the degrees of freedoms of the impurity is carried out by means of unperturbed density matrix $`\widehat{\rho }_0=Z^1\mathrm{exp}\left\{\frac{2\mu _BH_z^{\mathrm{eff}}\widehat{S}_z}{k_BT}\right\}`$. In this form the structure of the result (4) is similar to the one obtained in Ref. , where the spin-flip scattering of electrons at interfaces of tunnel junctions was investigated in the framework of tunneling Hamiltonian and the second order perturbation theory.
### 2.2. Resonant tunneling in the case of nonmagnetic impurities
To extract the physical nature of resonant tunneling through the impurity states contained in expression (4) we proceed as follows. For the sake of clarity and simplicity we consider first the case of zero-spin impurity. Then only one element of $`t`$–matrix at cite $`𝐜`$ survives
$$t_0^{()}(\epsilon )=\frac{a_0^3\epsilon _0}{1a_0^3\epsilon _0G_{()}(\epsilon )}.$$
It defines the position of a resonant level inside the gap of the dielectric band structure by finding the root of the equation $`a_0^3\epsilon _0\mathrm{Re}G_{()}(\epsilon _i)=1`$. From expression (LABEL:G\_c), it follows that the real part of the Green function $`\mathrm{Re}G_\mu (\epsilon ,𝐜)`$ is independent on c and spin $`\mu `$ up to exponentially small terms. Therefore, the position of level $`\epsilon _i`$ is weakly depend both on the position of impurity inside the barrier and on the direction of the spin of tunneling electron. Evidently, only those impurities for which $`\epsilon _i`$ is close to the chemical potential $`\epsilon _F`$ contribute to a significant extent of the total current at low bias voltage. Therefore, it is possible to expand the denominator in $`t_0^{()}(\epsilon )`$ in powers of $`(\epsilon _i\epsilon )`$. If we now introduce the position dependent line-widths
$$\mathrm{\Gamma }_\mu ^L(𝐜)=\mathrm{\Phi }_\mu ^L(𝐜)/\mathrm{Re}G^{}(\epsilon _F),\mathrm{\Gamma }_\mu ^R(𝐜)=\mathrm{\Phi }_\mu ^R(𝐜)/\mathrm{Re}G^{}(\epsilon _F),$$
(7)
where $`\mathrm{Re}G^{}(\epsilon _F)=\frac{}{\epsilon }\mathrm{Re}G(\epsilon )|_{\epsilon =\epsilon _F}`$ is the energy derivative of the electron Green function at Fermi level, then we obtain the general formula for the resonant case of impurity assisted tunneling
$$\sigma ^{\mathrm{imp}}(𝐜)\frac{1}{S}\left(\frac{2e^2}{\pi \mathrm{}}\right)\underset{\mu }{}_{\mathrm{}}^+\mathrm{}\underset{i}{}\frac{\mathrm{\Gamma }_\mu ^L(𝐜)\mathrm{\Gamma }_\mu ^R(𝐜)}{(\epsilon _F\epsilon _i)^2+\mathrm{\Gamma }_\mu ^2(𝐜)}\left(\frac{f(\epsilon )}{\epsilon }\right)d\epsilon ,$$
(8)
where $`\mathrm{\Gamma }_\mu (𝐜)=\mathrm{\Gamma }_\mu ^L(𝐜)+\mathrm{\Gamma }_\mu ^R(𝐜)`$ and the summation by $`i`$ is performed over all resonant levels. For the qualitative analysis, one may evaluate the expressions (7) for $`\mathrm{\Gamma }_\mu ^{R(L)}(𝐜)`$ approximately by considering the case $`\kappa =0`$ which is valid if $`e^{2q_0w}1`$. In this approximation
$`\mathrm{\Gamma }_\mu ^L(𝐜)`$ $`=`$ $`{\displaystyle \frac{2k_{1\mu }^Fm_0}{m_0^2k_{1\mu }^{F\mathrm{\hspace{0.17em}2}}+q_0^2}}\left({\displaystyle \frac{q_0^2}{2m_0}}\right){\displaystyle \frac{e^{2q_0(ca)}}{ca}},`$
$`\mathrm{\Gamma }_\mu ^R(𝐜)`$ $`=`$ $`{\displaystyle \frac{2k_{3\mu }^Fm_0}{m_0^2k_{3\mu }^{F\mathrm{\hspace{0.17em}2}}+q_0^2}}\left({\displaystyle \frac{q_0^2}{2m_0}}\right){\displaystyle \frac{e^{2q_0(bc)}}{bc}},`$ (9)
and expression (8) reproduces the result of Ref. . To proceed further, we discuss some assumptions concerning the parameters of the model. We consider the case of Co electrode and Al<sub>2</sub>O<sub>3</sub> as an oxide layer and take typical values of $`k_{}^F=1.09`$ Å<sup>-1</sup>, $`k_{}^F=0.42`$ Å<sup>-1</sup>, $`m1`$ for itinerant electrons in Co and a typical barrier height for Al<sub>2</sub>O<sub>3</sub> (measured from the Fermi level $`\epsilon _F`$) $`U_0\epsilon _F=3\mathrm{eV}`$ with an effective mass $`m_0=0.4`$ (Ref. ), that gives $`q_00.56`$ Å<sup>-1</sup>. Assuming the thickness of the barrier $`w20`$ Å<sup>-1</sup>, one may estimate the conductance $`\sigma _0`$ of the system without impurity from (3) by means of approximate formula
$$\sigma _0\frac{2e^2}{\pi ^2\mathrm{}}\left(\frac{q_0}{w}\right)\underset{\mu }{}\frac{k_{1\mu }^Fk_{3\mu }^Fq_0^2m_0^2e^{2q_0w}}{(m_0^2k_{1\mu }^{F2}+q_0^2)(m_0^2k_{3\mu }^{F2}+q_0^2)},$$
(10)
that leads to GMR $``$ 16 %. To estimate the value of the line-width (9) we consider impurities located close to left interface at a distance, say, of two atomic layers which corresponds to $`(ca)4`$ Å. For spin up electrons it gives $`\mathrm{\Gamma }_{}(𝐜)0.02`$ eV. Further in this paper we restrict ourselves to the case of temperature interval from $`4.2300`$ K (0.025 eV). We assume that the impurity levels $`\epsilon _i`$ in the band gap of the insulator form a narrow impurity band of width $`\mathrm{\Delta }\epsilon `$ which spreads symmetrically with respect to Fermi level $`\epsilon _F`$ and, following Ref. , we introduce its density of states $`\nu (\epsilon )`$ per unit volume and unit energy interval. We assume that $`\mathrm{\Delta }\epsilon `$ is of the order $`0.1`$ to $`0.2`$ eV, i.e. an order of magnitude greater that the above estimated line-width. In this context, with a good accuracy, the impurity conductance (8) rewrites as follows:
$$\sigma ^{\mathrm{imp}}(𝐜)\frac{2e^2}{\mathrm{}}\nu (\epsilon _F)\underset{\mu }{}_{\mathrm{}}^+\mathrm{}\left(\frac{f}{\epsilon }\right)\frac{\mathrm{\Gamma }_\mu ^L(𝐜)\mathrm{\Gamma }_\mu ^R(𝐜)}{\mathrm{\Gamma }_\mu (𝐜)}\rho (\epsilon ,𝐜)𝑑\epsilon ,$$
(11)
where factor
$$\rho (\epsilon ,𝐜)=\frac{2}{\pi }\mathrm{arctan}\left(\frac{\mathrm{\Delta }\epsilon }{2\mathrm{\Gamma }_\mu (𝐜)}\right)$$
arises from the integration of exp. (8) over impurity levels $`\epsilon _i`$ in the range of impurity band. Due to above mentioned estimations equality $`\rho (\epsilon ,𝐜)1`$ holds with a good degree of accuracy and in this case exp. (11) becomes in agreement with Refs. .
### 2.3. Resonant tunneling in the case of paramagnetic impurities
To investigate the general case of paramagnetic impurity we follow the same procedure as in the previous section. Let $`𝐉=𝐬+𝐒`$ be the total magnetic moment of the system. We may state that $`[H,J_z]=0`$ and, therefore, $`J_z`$ is a good quantum number. We regard the $`\widehat{t}`$–matrix (LABEL:t-matrix) as an operator acting on the spinor subspace $`|\sigma ,m`$, where $`\sigma =\pm \frac{1}{2}`$ and $`m=\pm 1,0`$ corresponds to the projection of the $`z`$–component of electron and impurity spin, respectively (we consider the case $`𝐒=1`$). As far as its total magnetic moment along the $`z`$–axis $`J_z=s_z+S_z`$ is conserved, the matrix elements $`\sigma _1m_1|\widehat{t}|\sigma _2m_2`$ are non-zero only if $`m_1+\sigma _1=m_2+\sigma _2`$ and, therefore, it is convenient to introduce the notation $`t_{m_j}^{\sigma _1\sigma _2}=\sigma _1m_1|\widehat{t}|\sigma _2m_2`$, where $`m_j=m_1+\sigma _1=m_2+\sigma _2`$. These elements are simply calculated from (LABEL:t-matrix). The non-zero ones are written as follows
$`t_{3/2}^{}`$ $`=`$ $`{\displaystyle \frac{a_0^3(\epsilon _0J/2)}{1a_0^3(\epsilon _0J/2)G_{}(\epsilon )}},`$
$`t_{3/2}^{}`$ $`=`$ $`{\displaystyle \frac{a_0^3(\epsilon _0J/2)}{1a_0^3(\epsilon _0J/2)G_{}(\epsilon )}},`$ (12)
and $`\widehat{t}_{\pm 1/2}=\left(\genfrac{}{}{0pt}{}{t_{\pm 1/2}^{}}{t_{\pm 1/2}^{}}\genfrac{}{}{0pt}{}{t_{\pm 1/2}^{}}{t_{\pm 1/2}^{}}\right)`$ corresponding to the subspace $`m_j=\pm \frac{1}{2}`$ with
$`t_{1/2}^{}(t_{1/2}^{})`$ $`=`$ $`{\displaystyle \frac{a_0^3}{\mathrm{\Delta }_{\pm 1/2}(\epsilon )}}\left\{\epsilon _0+J/2a_0^3G_{()}(\epsilon )(\epsilon _0J/2)(\epsilon _0+J)\right\},`$
$`t_{1/2}^{}(t_{1/2}^{})`$ $`=`$ $`{\displaystyle \frac{a_0^3}{\mathrm{\Delta }_{\pm 1/2}(\epsilon )}}\left\{\epsilon _0a_0^3G_{()}(\epsilon )(\epsilon _0J/2)(\epsilon _0+J)\right\},`$ (13)
$`t_{\pm 1/2}^{}`$ $`=`$ $`t_{\pm 1/2}^{}={\displaystyle \frac{a_0^3}{\sqrt{2}\mathrm{\Delta }_{\pm 1/2}}}J,`$
where denominators are
$$\mathrm{\Delta }_{\pm 1/2}(\epsilon )=\left(1a_0^3G_{}(\epsilon )(\epsilon _0J/2)\right)\left(1a_0^3G_{}(\epsilon )(\epsilon _0+J)\right)\pm a_0^3J(G_{}(\epsilon )G_{}(\epsilon )).$$
As can be seen, two poles of the $`\widehat{t}`$–matrix defined from equations $`a_0^3\mathrm{Re}G(\epsilon _{3/2})(\epsilon _0J/2)=1`$ and $`a_0^3\mathrm{Re}G(\epsilon _{1/2})(\epsilon _0+J)=1`$ correspond to two multiplets $`\epsilon _{3/2}`$ and $`\epsilon _{1/2}`$ with a total angular momentum $`j=3/2`$ and $`j=1/2`$, respectively. If $`J>0`$, then $`\epsilon _{3/2}<\epsilon _{1/2}`$, i.e. the multiplet with $`j=3/2`$ has a lower energy than one with $`j=1/2`$. As for non-magnetic impurity, we restrict ourselves by considering the regime of only one-channel resonant tunneling. We assume that $`J>0`$ and the lowest impurity levels $`\epsilon _i=\epsilon _{3/2}`$ corresponding to the multiplet with $`j=3/2`$ lie close to $`\epsilon _F`$. We note that the typical value of exchange coupling $`J`$ is of order 1 eV and due to this fact we may eliminate the resonant level $`\epsilon _{1/2}`$ from further consideration. Then, as in the previous analysis for a non- magnetic impurity, only the resonant part of the $`\widehat{t}`$– matrix (12,13) at energies close to chosen $`\epsilon _{3/2}`$ is essential for the subsequent calculations. Expressions (12,13) can be easily written as follows:
$`t_{\pm 3/2}^{()}(\epsilon )={\displaystyle \frac{1}{G^{}(\epsilon )}}{\displaystyle \frac{1}{\epsilon \epsilon _i+i\mathrm{\Gamma }_{()}(𝐜)}},`$
$`\widehat{t}_{1/2}(\epsilon )={\displaystyle \frac{1}{G^{}(\epsilon )}}{\displaystyle \frac{1}{\epsilon \epsilon _i+i\gamma _{}(𝐜)}}\left(\begin{array}{cc}\frac{1}{3}& \frac{\sqrt{2}}{3}\\ \frac{\sqrt{2}}{3}& \frac{2}{3}\end{array}\right),`$ (16)
$`\widehat{t}_{1/2}(\epsilon )={\displaystyle \frac{1}{G^{}(\epsilon )}}{\displaystyle \frac{1}{\epsilon \epsilon _i+i\gamma _{}(𝐜)}}\left(\begin{array}{cc}\frac{2}{3}& \frac{\sqrt{2}}{3}\\ \frac{\sqrt{2}}{3}& \frac{1}{3}\end{array}\right),`$ (19)
where $`\gamma _{}(𝐜)=\frac{2}{3}\mathrm{\Gamma }_{}(𝐜)+\frac{1}{3}\mathrm{\Gamma }_{}(𝐜)`$, $`\gamma _{}(𝐜)=\frac{1}{3}\mathrm{\Gamma }_{}(𝐜)+\frac{2}{3}\mathrm{\Gamma }_{}(𝐜)`$ are the inverse lifetimes of the resonant states with $`m_j=\pm 1/2`$. This result allows simple qualitative interpretation. Let us look, for a example, at quantum states with $`m_j=1/2`$. From elementary quantum mechanical theory one may conclude that
$`\varphi _{1/2}^{}=|,m_s=0=\sqrt{{\displaystyle \frac{2}{3}}}|j={\displaystyle \frac{3}{2}},m_j={\displaystyle \frac{1}{2}}+\sqrt{{\displaystyle \frac{1}{3}}}|j={\displaystyle \frac{1}{2}},m_j={\displaystyle \frac{1}{2}},`$
$`\varphi _{1/2}^{}=|,m_s=1=\sqrt{{\displaystyle \frac{1}{3}}}|j={\displaystyle \frac{3}{2}},m_j={\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{2}{3}}}|j={\displaystyle \frac{1}{2}},m_j={\displaystyle \frac{1}{2}}.`$ (20)
As we have assumed, only $`|j=\frac{3}{2},m_j=\frac{1}{2}`$ gets into resonance and, therefore, e.g. $`t_{1/2}^{}\varphi _{1/2}^{}|\widehat{t}|\varphi _{1/2}^{}\sqrt{2}/3`$ in agreement with (LABEL:t\_res). On the other hand,
$`|j=\frac{3}{2},m_j=\frac{1}{2}=\sqrt{\frac{2}{3}}|,m_s=0+\sqrt{\frac{1}{3}}|,m_s=1`$
and, hence, its inverse life time is given by $`\tau _{1/2}^1=\frac{2}{3}\tau _{}^1+\frac{1}{3}\tau _{}^1`$.
We substitute all $`\widehat{t}`$–matrix elements in (4) by its resonance expansion (LABEL:t\_res). To proceed further, one has to perform in (4) the configuration averaging over all impurity centers and thermodynamic one over all possible channels. Suppose, that the impurities are distributed uniformly in the space in the interval $`[z_0\mathrm{\Delta },z_0+\mathrm{\Delta }]`$ along $`z`$–direction with the width of $`2\mathrm{\Delta }`$ and center $`z_0`$ which we have chosen to be close to left (L) ferromagnetic contact. After that, it is possible, first of all, to average the Lorentzian peaks over the distribution of impurity levels by averaging them over $`\epsilon _i`$ in the range of impurity band with factor $`\nu (\epsilon _F)`$ and to perform the thermodynamic averaging by integrating over $`\epsilon `$ and neglecting the dependencies of $`\mathrm{\Gamma }^{L(R)}(𝐜)`$ on energy. On the third step, the averaging over the space distribution of impurities along $`z`$–direction should be made. Following the outlined procedure, the total conductance (4) as a function of temperature for the parallel and antiparallel alignment of magnetization of the ferromagnetic layers is written as a sum of factorized terms over all possible scattering channels:
$`\sigma ^P(T)={\displaystyle \frac{2e^2}{\mathrm{}}}{\displaystyle \underset{\mu \rho m_j}{}}P_{m_j}^{\mu \rho }\left({\displaystyle \frac{\mu _BH_z^{\mathrm{eff}}}{k_BT}}\right){}_{}{}^{P}\sigma _{m_j}^{\mu \rho }(z_0,\mathrm{\Delta })\nu (\epsilon _F)+\sigma _0^P,`$
$`\sigma ^{AP}(T)={\displaystyle \frac{2e^2}{\mathrm{}}}{\displaystyle \underset{\mu \rho m_j}{}}P_{m_j}^{\mu \rho }\left({\displaystyle \frac{\mu _BH_z^{\mathrm{eff}}}{k_BT}}\right){}_{}{}^{AP}\sigma _{m_j}^{\mu \rho }(z_0,\mathrm{\Delta })\nu (\epsilon _F)+\sigma _0^{AP},`$ (21)
where
$$\sigma _{m_j}^{\mu \rho }(z_0,\mathrm{\Delta })=\frac{1}{2\mathrm{\Delta }}_{z_0\mathrm{\Delta }}^{z_0+\mathrm{\Delta }}\sigma _{m_j}^{\mu \rho }(c)\rho (\epsilon _F,c)𝑑c.$$
The origine of $`\rho (\epsilon _F,c)`$ is the same as in exp. (11) and explicit form of functions $`P_{m_j}^{\mu \rho }(h)`$ and $`\sigma _{m_j}^{\mu \rho }(c)`$ is given in Appendix A. We have also used the same notation of matrix indexes as it was previously introduced for $`\widehat{t}`$–matrix elements. $`\sigma _0^P`$ and $`\sigma _0^{AP}`$ are the tunnel conductances of the pure system in accordance with (3). Factors $`P_{m_j}^{\mu \rho }`$ and $`\sigma _{m_j}^{\mu \rho }`$ represent the thermodynamic and quantum mechanical probabilities of the given process, respectively. Expressions (LABEL:s\_T) are the final results of this section and their analysis is presented below (see Sec. III).
### 2.4. Dependence of conductivity on bias-voltage
We are also interested in $`I(V)`$ characteristic of the considered system. To derive the general formula for the current, one may simply extend the expressions (4) to the case of finite applied bias voltage. Consider, for example, the contribution to the total current $`I`$, coming from all possible channels of the form $`|in,|out,`$ for tunnel electrons moving from the left electrode to the right electrode and of the form $`|in,|out,`$ for electrons moving from the right to the left, respectively, i.e. in both cases an electron has an ”up” projection of spin in the left contact and ”down” projection of spin in the right one after or before scattering. From the general concept, one may conclude that this contribution to the current can be written as
$`I^{}(V)={\displaystyle \frac{1}{S}}\left({\displaystyle \frac{2e^2}{\pi \mathrm{}}}\right){\displaystyle _{\mathrm{}}^+\mathrm{}}d\epsilon \{f_{}(\epsilon \mu _BH_z^{\mathrm{eff}}eV)[1f_{}(\epsilon +\mu _BH_z^{\mathrm{eff}})]\widehat{t}_{}(\epsilon )\widehat{t}_+(\epsilon )`$
$`f_{}(\epsilon +\mu _BH_z^{\mathrm{eff}})[1f_{}(\epsilon \mu _BH_z^{\mathrm{eff}}eV)]\widehat{t}_+(\epsilon )\widehat{t}_{}(\epsilon )\}\mathrm{\Phi }_{}^L(𝐜)\mathrm{\Phi }_{}^R(𝐜),`$ (22)
where it is assumed that the voltage bias is applied from the left to the right direction. It is important to notice that inelastic spin-flip processes of the electron scattering on the impurity were taken into account in derivation of the exp. (22) but they were omitted in Ref. . Analogous expressions can be written for all other channels. In the case under consideration, expression (22) contains two regimes of non-linear behavior of $`I(V)`$ characteristic. The first one reproduces a zero bias anomaly due to excitation of spin-flip processes at low bias voltages of order of magnitude $`\mu _BH_z^{\mathrm{eff}}`$ (we believe that it is of order 5 mV). In this range, as before, one may assume that the resonance amplitudes $`(\widehat{t}_z^{()})^{}(\widehat{t}_z^{()})`$ and $`\widehat{t}_{}\widehat{t}_+`$ are nearly independent of the energy after averaging over all possible configurations of impurities. As a result, the voltage dependence of total currents for parallel and antiparallel configurations are given by formulae similar to (LABEL:s\_T):
$`I^P(V,T)`$ $`=`$ $`{\displaystyle \frac{2e}{\mathrm{}}}{\displaystyle \underset{\mu \rho m_j}{}}I_{m_j}^{\mu \rho }(V,H_z^{\mathrm{eff}}){}_{}{}^{P}\sigma _{m_j}^{\mu \rho }(z_0,\mathrm{\Delta })\nu (\epsilon _F)+\sigma _0^PV,`$
$`I^{AP}(V,T)`$ $`=`$ $`{\displaystyle \frac{2e}{\mathrm{}}}{\displaystyle \underset{\mu \rho m_j}{}}I_{m_j}^{\mu \rho }(V,H_z^{\mathrm{eff}}){}_{}{}^{AP}\sigma _{m_j}^{\mu \rho }(z_0,\mathrm{\Delta })\nu (\epsilon _F)+\sigma _0^{AP}V.`$ (23)
The expressions for $`I_{m_j}^{\mu \rho }(V,H_z^{\mathrm{eff}})`$ are given in Appendix B. The voltage dependent conductances $`\sigma ^P(V,T)`$ and $`\sigma ^{AP}(V,T)`$ can be obtained from (22) by derivation with respect to $`V`$. The detailed analysis of this physical situation is presented in the next section.
The second source of possible non-linear character of $`I(V)`$ dependence is the variation of potential profile $`U(z)`$ (see Fig. 1) under applied bias voltage. It follows, then, that the latter introduces the correction to exp. (9) and (10) and they can be calculated with the use of Wentzel–Kramers–Brillouin (WKB) approximation assuming that the applied voltage produces the uniform electrical field inside the insulating layer. In the case of pure tunnel conductance it is known that both conductances for the parallel and antiparallel configurations increase with the increase of applied voltage so that the TMR as a function of $`V`$, defined as $`{\displaystyle \frac{I^P(V)I^{AP}(V)}{I^{AP}(V)}}`$, drops significantly at the voltages of order 1 eV. The contribution of impurity assisted tunneling may change considerably this situation in the case of non-uniform spatial distribution of impurities, e.g. when they are distributed in the vicinity of only one electrode. In this particular situation, as we will show, the essential variation of TMR amplitude in the case of magnetic impurities (in contrast to non-magnetic ones) will take place at the bias voltages compared with impurity band width $`\mathrm{\Delta }\epsilon `$.
For the sake of simplicity we consider, first, the case of non-magnetic impurities. In the WKB approximation the contribution from all impurities, located at given point $`𝐜`$, to the total current $`I(V)`$ will have the form similar to (9) and (11):
$$j^{\mathrm{imp}}(𝐜)=\frac{2e^2}{\mathrm{}}\nu (\epsilon _F)\underset{\mu }{}_{\mathrm{}}^+\mathrm{}\left\{f(\epsilon eV)f(\epsilon )\right\}\frac{\mathrm{\Gamma }_\mu ^L(𝐜)\mathrm{\Gamma }_\mu ^R(𝐜)}{\mathrm{\Gamma }_\mu ^L(𝐜)+\mathrm{\Gamma }_\mu ^R(𝐜)}\rho (\epsilon ,V)𝑑\epsilon ,$$
(24)
where
$$\mathrm{\Gamma }_\mu ^L(𝐜)=\frac{k_{1\mu }^Fq_am_0\tau _a^1}{(q_a^{})^2+k_{1\mu }^{F2}m_0^2}e^{S_a/\mathrm{}};$$
$$\mathrm{\Gamma }_\mu ^R(𝐜)=\frac{k_{3\mu }^Fq_bm_0\tau _b^1}{(q_b^+)^2+k_{3\mu }^{F2}m_0^2}e^{S_b/\mathrm{}}.$$
(25)
Here $`q_b^2=q_0^2=2m_0(U\epsilon )`$, $`q_a^2=q_0^2+2m_0eV`$ are imaginary momenta of electron with the energy $`\epsilon `$ in the vicinity of the right and the left electrode, $`q_{a(b)}^\pm =q_0\pm {\displaystyle \frac{1}{2}}{\displaystyle \frac{eEm_0}{q_{a(b)}^2}}`$, $`E`$ is the electric field in the device. We also introduce $`q_c=q_0+2m_0eV(bc)/w`$ imaginary momenta of electron on the impurity center. Then $`S_a={\displaystyle \frac{q_a^3q_c^3}{3m_0eE}}`$, $`S_b={\displaystyle \frac{q_c^3q_b^3}{3m_0eE}}`$ represent the classical actions along the path from the left contact to the point $`𝐜`$ in the barrier and, afterwards, from this point to the right contact, respectively; $`\tau _a={\displaystyle \frac{q_aq_c}{eE}}`$ and $`\tau _b={\displaystyle \frac{q_cq_b}{eE}}`$ denote the passage times associated with these paths. Factor
$$\rho (\epsilon ,V)=\frac{1}{\pi }\left\{\mathrm{arctan}\left[\frac{\epsilon \epsilon _FeV\left(\frac{bc}{w}\right)+\frac{\mathrm{\Delta }\epsilon }{2}}{\mathrm{\Gamma }_\mu ^L(𝐜)+\mathrm{\Gamma }_\mu ^R(𝐜)}\right]\mathrm{arctan}\left[\frac{\epsilon \epsilon _FeV\left(\frac{bc}{w}\right)\frac{\mathrm{\Delta }\epsilon }{2}}{\mathrm{\Gamma }_\mu ^L(𝐜)+\mathrm{\Gamma }_\mu ^R(𝐜)}\right]\right\},$$
as before, arises from the summation over all impurity levels $`\epsilon _i`$ and gives the relative weight of all resonant channels with energy $`\epsilon `$. To clarify the situation, it is sufficiently to consider the most resonant channel with energy $`\epsilon _r=\epsilon _F+eV(bc)/w`$ at which $`\rho (\epsilon ,V)`$ reaches its maximum. One may note that $`\epsilon _r`$ corresponds to the resonant impurity level that exactly coincides with Fermi energy at vanishing voltage and it shifts linearly with the increase of applied bias depending on the position $`𝐜`$ of impurity inside the barrier. As it was stated earlier, the most interesting case takes place when the point $`𝐜`$ is situated close to the left contact. Then one can see that $`\mathrm{\Gamma }^L(𝐜)\mathrm{\Gamma }^R(𝐜)`$ and, thus, $`j^{\mathrm{imp}}(𝐜)\mathrm{\Gamma }^R(𝐜)\rho (\epsilon _r,V)`$. At bias voltages much more lower than the height of the barrier $`\phi =(U\epsilon _F)`$, $`S_b`$ can be expanded in powers of $`V`$:
$$S_b=q_0(bc)\left\{1\frac{m_0eV}{2q_0^2}\left(\frac{bc}{w}\right)+\mathrm{}\right\}$$
which shows that $`\mathrm{\Gamma }^R(𝐜)\mathrm{exp}(S_b/\mathrm{})`$ is an increasing function of $`V`$ in the vicinity of $`V=0`$. Hence, it leads to increase of differential conductivity $`\sigma (V)=I/V`$ under direct bias voltage, applied to the barrier from the left to the right direction, and to decrease of $`\sigma (V)`$ under inverse bias voltage. The physical meaning of such behavior is rather obvious. From expression for $`S_b`$ it follows that electrons tunneling under forward bias due to resonant levels lying close to $`\epsilon _r`$ will propagate through the potential barrier with the height which is effectively less compared with that in case of inverse bias.
The expression for the paramagnetic impurity assisted current at finite voltages has the structure similar to exp. (24) with Fermi distribution factors written in accordance with the general formula (22) and integrand expression has the form given in Appendix A, where line-widths $`\mathrm{\Gamma }_{()}`$ have to be substituted by WKB approximation (25). In the case of magnetic impurities the above outlined mechanism of asymmetry in $`I(V)`$ characteristic due to the resonance levels will essentially contribute to the voltage dependence of TMR in question, and it will be discussed in the next section.
## 3. Results and discussion
In this section we consider the temperature and bias voltage dependencies of the conductances and TMR effect of the considered structures. We investigate the case of Co/Al<sub>2</sub>O<sub>3</sub>/Co junction with the typical parameters that have been already mentioned in section 2: $`k_{}^F=1.09`$ Å<sup>-1</sup>, $`k_{}^F=0.42`$ Å<sup>-1</sup> are the Fermi momenta of itinerant electrons in Co, $`q_0=0.56`$ Å<sup>-1</sup> is the imaginary momentum in the barrier, $`m_0=0.4`$ is the effective mass in the insulator and $`w=20`$ Å is its thickness. We focus on the most interesting situation when impurities are introduced at the vicinity of the left electrode at a depth $`w_1`$ inside the insulator layer. We chose width $`w_1=4.06`$ Å that corresponds to two atomic monolayers. The essential parameter of the model that has to be defined is the effective molecular field $`H_z^{\mathrm{eff}}`$ acting on impurity spins. One may suppose that it should exponentially decay in the depth of the barrier. We have, therefore, set it to $`\mu _BH_z^{\mathrm{eff}}=5`$ meV (58 K) that is of two order less than the critical temperature in the bulk Co.
Consider, first, the case $`T=0`$. It is possible to estimate the concentration of impurity atoms so that its contribution to the resonance conductivity is comparable with the ordinary tunnel conductance of the system. One can write that the impurity density of states (see exp.(11)) $`\nu (\epsilon _F)=\frac{N_i}{Sw_1\left(\mathrm{\Delta }\epsilon \right)}`$, where $`N_i`$ is a total number of impurities and $`\mathrm{\Delta }\epsilon `$ is the width of its energy distribution. On the other hand $`N_i=xN`$ and $`N=Sw_1/a_0^3`$ where $`N`$ is the total number of atoms in the layer which contains the impurities, $`x`$ is the local concentration of impurities in this volume and $`a_0`$ is the lattice constant. This yields $`\nu (\epsilon _F)=x\left(\frac{w_1}{a_0}\right)\frac{1}{a_0^2\mathrm{\Delta }\epsilon }`$. We introduce the characteristic concentration $`x_0`$ defined so that in the case of nonmagnetic impurity, the impurity conductance (11) is equal to the tunnel conductance (10) of spin $``$ channel in the parallel magnetic configuration of the ferromagnetic layers. Such definition leads to
$$x_0\frac{\mathrm{\Delta }\epsilon }{\pi ^2}\left(\frac{2m_0a_0^2}{\mathrm{}^2}\right)\left(\frac{w}{ww_1}\right)\left(\frac{a_0}{w_1}\right)\frac{k_{}^Fq_0m_0}{k_{}^{F}{}_{}{}^{2}m_0^2+q_0^2}\frac{(q_0w_1)\mathrm{exp}(2q_0w_1)}{1\mathrm{exp}(2q_0w_1)}$$
(26)
If we choose $`\mathrm{\Delta }\epsilon =0.2`$ eV, then $`x_0=6.5\times 10^5`$. The conductance of the system at $`T=0`$ can be extract from the general expression (14). We suppose that for both parallel and antiparallel configurations the left electrode has ”up” magnetization and, hence, $`H_z^{\mathrm{eff}}`$ is positive in both cases. At zero temperature, all spin-flip processes are frozen and due to the above assumption, only the configuration of impurity spin with $`m_s=1`$ is possible. As a result, only two resonance channels from many possible ones have nonzero contribution to conductivity, namely $`|,m_s=1|,m_s=1`$ and $`|,m_s=1|,m_s=1`$ with $`m_j=3/2`$ and $`1/2`$ respectively. From exp. (A1) (see Appendix A) it follows that the channel with $`m_j=3/2`$ gives the main contribution into the conductivity at low temperatures and $`{}_{}{}^{P}\sigma _{3/2}^{}\mathrm{\Gamma }_{}`$ and $`{}_{}{}^{AP}\sigma _{3/2}^{}\mathrm{\Gamma }_{}`$. So these contributions depend on the mutual orientation of the magnetization of the ferromagnetic layers, therefore they increase the total amplitude of the TMR. The total expression for $`TMR=\frac{\sigma ^P\sigma ^{AP}}{\sigma ^{AP}}`$ including all possible channels may be written as
$$TMR=\frac{(\mathrm{\Gamma }_{}\mathrm{\Gamma }_{})(\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}+\frac{x}{x_0}\mathrm{\Gamma }_{}(1\frac{1}{9}\mathrm{\Gamma }_{}/\gamma _{}))}{\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}(2+\frac{x}{x_0}(1+\frac{1}{9}\mathrm{\Gamma }_{}/\gamma _{}))}$$
(27)
where $`x`$ is concentration, $`x_0`$ is defined by (26), $`\gamma _{}=\frac{2}{3}\mathrm{\Gamma }_{}+\frac{1}{3}\mathrm{\Gamma }_{}`$ and $`\mathrm{\Gamma }_{()}=k_{()}^Fq_0^2/(k_{()}^{F2}m_0^2+q_0^2)`$ are the tunneling density of states for $`()`$ spin electrons. The dependence of TMR effect versus the polarization $`P=(k_{}^Fk_{}^F)/(k_{}^F+k_{}^F)`$ is shown in Fig. 3 in comparison with the non-resonant tunnel conductance at $`x=0`$. The total TMR amplitude is larger than the TMR due to direct tunneling, in accordance with considerations written above. For a given polarization $`P=0.44`$ in case of chosen parameters, the contribution of the impurity assisted tunneling leads to strong enhancement of TMR amplitude (typically by a factor 2, see Fig. 3).
We note, that in the case of nonmagnetic impurities, distributed in the vicinity of only one contact, the resonant impurity conductance $`\sigma ^{\mathrm{imp}}\mathrm{\Gamma }_{}+\mathrm{\Gamma }_{}`$ is equal for both parallel and antiparallel configurations and, therefore, in this case the mechanism of impurity assisted tunneling is not able to enhance the TMR effect. The enhancement of TMR amplitude in the case of paramagnetic impurities is essentially due to the presence of ferromagnetic exchange coupling between the magnetization in ferromagnetic electrode and the impurity spins which tends to induce a ferromagnetic order in the plane of impurities and, as the result, leads to the preference of impurity spin to be found in the quantum state with $`m_s=1`$.
The temperature dependences of resonant conductances for parallel and antiparallel configurations in the interval from 4.2 K – 300 K are presented in Fig. 4. In the case of parallel alignment, $`\sigma _{\mathrm{imp}}^P(T)`$ is nearly independent on the temperature, but in antiparallel situation there is a 50% increase of impurity conductance $`\sigma _{\mathrm{imp}}^{AP}(T)`$. This originates from the thermal excitation of both spin-flip and spin-conserving processes which are frozen at zero temperature. For AP configuration the process $`|,m_s=0|,m_s=1`$ was forbidden at $`0^{}`$K but now it is allowed and gives a large contribution into the current as it is proportional to the product of the largest density of states $`\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}`$. As a consequence, the TMR effect decreases with the increase of the temperature.
We have also calculated the dependence of the differential conductances on bias voltage according to (23) and (24). These dependences at $`T=4.2`$ K and $`T=77`$ K are presented in Fig. 5. A new effect is predicted: the voltage dependence of the conductance in the antiparallel alignment of magnetization in ferromagnetic layers is asymmetric under forward and inverse bias voltage when the paramagnetic impurities inside the insulator layer are distributed close to only one of the interfaces and are bound by exchange interaction with magnetization of the nearest ferromagnetic layer.
One can distinguish two different mechanisms that give rise to the presented asymmetrical behavior with respect to inversion of bias voltage. The first one, which we refer as zero bias anomaly, manifests itself at low voltages of the order of 10 mV (for particular chosen parameters in our model) and is strongly pronounced only at low temperatures (see Fig. 5, the case of $`T=4.2`$ K). It originates from the excitations of spin-flip processes on the impurity centers. One may look at the general expression (23) and consider the case of low temperature. An electron undergoing spin-flip scattering, may transfer an amount of energy $`\omega _0=2\mu _BH_z^{\mathrm{eff}}`$ to the impurity spin thus exiting it at an higher energy level or on the contrary may acquire this quantum of energy from it. The latter process is impossible at low temperature. The former one is possible only if an electron moving, say, from the left contact possesses an excess energy of at least $`\omega _0`$ with respect to Fermi level in the right contact. The only one process that contributes to this anomaly at low temperature is $`\varphi _{1/2}^{}\varphi _{1/2}^{}`$ (see exp. LABEL:phi). For antiparallel alignment of the magnetization its quantum mechanical probability is proportional to $`{}_{}{}^{AP}\sigma _{1/2}^{}\frac{2}{9}\mathrm{\Gamma }_{}^2/\gamma _{}`$ for electrons moving from the left ferromagnetic layer into the right one and is proportional to $`{}_{}{}^{AP}\sigma _{1/2}^{}\frac{2}{9}\mathrm{\Gamma }_{}^2/\gamma _{}`$ in the case of electrons moving from the right to the left. For parallel configuration of magnetizations these probabilities are equal in both directions and are proportional to $`{}_{}{}^{P}\sigma _{1/2}^{}={}_{}{}^{P}\sigma _{1/2}^{}\frac{2}{9}\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}/\gamma _{}`$. As a result, zero bias anomaly at $`T=4.2`$ K looks asymmetrically in the case of antiparralel configuration and is symmetrical in the case of parallel alignment of magnetizations.
The differential conductances as a functions of the bias voltage at $`T=77`$ K have been calculated using two different approximations. Thick dashed and solid lines correspond to the conductances in the parallel and antiparallel configurations, respectively, that have been calculated by means of WKB approximation in accordance with expressions (22) and (24). For the sake of comparison, the same dependences, indicated by thin lines, have been calculated with the use approximate formulae (23), where the dependence of $`\widehat{t}`$–matrix elements on the applied voltage has been neglected. In course of this, the latter curves demonstrate the only zero bias anomaly discussed above, which is substantially smoothed, compared with the case of $`T=4.2`$ K. On the contrary, the WKB scheme of calculation takes into account the variation of the potential profile inside the insulating barrier under applied bias voltage. In view of this, the differential conductances calculated by this scheme exhibit the tendency to increase at the direct bias voltage and to decrease at the reverse one. As it was shown above (see section 2.4), this behavior originates from the shift of the resonant levels inside the insulator due to externally applied electric field. This second mechanism in the origin of non-linear voltage dependence of impurity assisted conductance does not relate with the excitation of spin-flip processes. It becomes apparent at the voltages of the order of 50 mV and leads to the asymmetric voltage bias dependences in both cases of parallel and antiparallel configurations.
Finally, the TMR amplitude as a function of the bias voltage is shown in Fig. 6 for the broad range of applied voltage. Its non-linear and asymmetric behavior in the range of 0.2 V originates from the asymmetry of the shifts of resonant impurity levels with respect to forward and inverse bias. The low bias voltage anomaly at 10 mV is also strongly pronounced at the curve corresponding to 4.2 K. The relative contribution of the impurity assisted conductance to the total current of electrons is diminished after the value of applied bias voltage exceeds the half width of impurity band $`\mathrm{\Delta }\epsilon /2=0.1`$ eV. Therefore, the TMR amplitude drops to value $`20\%`$ at $`0.5`$ V corresponding primary to the pure tunnel conductance.
## Acknowledgments
A. Vedyayev and D. Bagrets are grateful to CEA/Grenoble/DRFMC/SP2M/NM for hospitality. This work was partially supported by Russian Foundation for Basic Research (grant No. 98–02–16806).
## Appendix A
Let $`w_1=ca`$ and $`w_2=bc`$ be the position of impurity with respect to the left and right interfaces, respectively, and $`w=ba`$ be the width of the insulator layer. We introduce tunneling densities of states for spin $`()`$ electrons $`\mathrm{\Gamma }_{()}=k_{()}^Fq_0^2/(k_{()}^Fm_0^2+q_0^2)`$ and denote $`\gamma _{}=\frac{2}{3}\mathrm{\Gamma }_{}+\frac{1}{3}\mathrm{\Gamma }_{}`$, $`\gamma _{}=\frac{1}{3}\mathrm{\Gamma }_{}+\frac{2}{3}\mathrm{\Gamma }_{}`$. Then, the position-dependent quantum mechanical probabilities $`\sigma _{m_j}^{\mu \rho }(c)`$ can be found as follows.
(a) In case of parallel configuration:
$$\sigma _{\frac{3}{2}(\frac{3}{2})}^{()}(c)=\frac{\mathrm{\Gamma }_{()}^2e^{2q_0w}/\left(w_1w_2\right)}{\mathrm{\Gamma }_{()}\frac{e^{2q_0w_1}}{w_1}+\mathrm{\Gamma }_{()}\frac{e^{2q_0w_2}}{w_2}},$$
$$\sigma _{\frac{1}{2}(\frac{1}{2})}^{()}(c)=\frac{4}{9}\frac{\mathrm{\Gamma }_{()}^2e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{()}\frac{e^{2q_0w_1}}{w_1}+\gamma _{()}\frac{e^{2q_0w_2}}{w_2}},$$
$$\sigma _{\frac{1}{2}(\frac{1}{2})}^{()}(c)=\frac{1}{9}\frac{\mathrm{\Gamma }_{()}^2e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{()}\frac{e^{2q_0w_1}}{w_1}+\gamma _{()}\frac{e^{2q_0w_2}}{w_2}},$$
$`(A1)`$
$$\sigma _{\frac{1}{2}}^{}(c)=\sigma _{\frac{1}{2}}^{}(c)=\frac{2}{9}\frac{\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{}\frac{e^{2q_0w_1}}{w_1}+\gamma _{}\frac{e^{2q_0w_2}}{w_2}},$$
$$\sigma _{\frac{1}{2}}^{}(c)=\sigma _{\frac{1}{2}}^{}(c)=\frac{2}{9}\frac{\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{}\frac{e^{2q_0w_1}}{w_1}+\gamma _{}\frac{e^{2q_0w_2}}{w_2}};$$
(b) In case of antiparallel configuration:
$$\sigma _{\frac{3}{2}(\frac{3}{2})}^{()}(c)=\frac{\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}e^{2q_0w}/\left(w_1w_2\right)}{\mathrm{\Gamma }_{()}\frac{e^{2q_0w_1}}{w_1}+\mathrm{\Gamma }_{()}\frac{e^{2q_0w_2}}{w_2}},$$
$$\sigma _{\frac{1}{2}(\frac{1}{2})}^{()}(z)=\frac{4}{9}\frac{\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{()}\frac{e^{2q_0w_1}}{w_1}+\gamma _{()}\frac{e^{2q_0w_2}}{w_2}},$$
$$\sigma _{\frac{1}{2}(\frac{1}{2})}^{()}(z)=\frac{1}{9}\frac{\mathrm{\Gamma }_{}\mathrm{\Gamma }_{}e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{()}\frac{e^{2q_0w_1}}{w_1}+\gamma _{()}\frac{e^{2q_0w_2}}{w_2}},$$
$`(A2)`$
$$\sigma _{\frac{1}{2}}^{()}(z)=\frac{2}{9}\frac{\mathrm{\Gamma }_{()}^2e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{}\frac{e^{2q_0w_1}}{w_1}+\gamma _{}\frac{e^{2q_0w_2}}{w_2}},$$
$$\sigma _{\frac{1}{2}}^{()}(z)=\frac{2}{9}\frac{\mathrm{\Gamma }_{()}^2e^{2q_0w}/\left(w_1w_2\right)}{\gamma _{}\frac{e^{2q_0w_1}}{w_1}+\gamma _{}\frac{e^{2q_0w_2}}{w_2}}.$$
The statistical probabilities $`P_j^{\mu \rho }(h)`$ are independent of the configuration of the system. We denote $`h=\mu _BH_z^{\mathrm{eff}}/kT`$ and $`Z=2\mathrm{cosh}(2h)+1`$, then
$$P_{\frac{3}{2}}^{}(h)=P_{\frac{1}{2}}^{}(h)=Z^1e^{2h},P_{\frac{1}{2}}^{}(h)=P_{\frac{1}{2}}^{}(h)=Z^1,$$
$$P_{\frac{1}{2}}^{}(h)=P_{\frac{3}{2}}^{}(h)=Z^1e^{2h},$$
$`(A3)`$
$$P_{\frac{1}{2}}^{}(h)=P_{\frac{1}{2}}^{}(h)=Z^1\frac{he^h}{\mathrm{sinh}(h)},P_{\frac{1}{2}}^{}(h)=P_{\frac{1}{2}}^{}(h)=Z^1\frac{he^h}{\mathrm{sinh}(h)}$$
## Appendix B
The non-trivial functions $`I_{m_j}^{\mu \rho }`$ are written as follows
$`I_{1/2}^{}(V,H_z^{\mathrm{eff}})`$ $`=`$ $`{\displaystyle \frac{(eV2\mu _BH_z^{\mathrm{eff}})(e^{eV/kT}1)}{(e^{(eV2\mu _BH_z^{\mathrm{eff}})/kT}1)(2\mathrm{cosh}(2\mu _BH_z^{\mathrm{eff}}/kT)+1)}},`$
$`I_{1/2}^{}(V,H_z^{\mathrm{eff}})`$ $`=`$ $`{\displaystyle \frac{(eV+2\mu _BH_z^{\mathrm{eff}})(e^{eV/kT}1)e^{2\mu _BH/kT}}{(e^{(eV+2\mu _BH_z^{\mathrm{eff}})/kT}1)(2\mathrm{cosh}(2\mu _BH_z^{\mathrm{eff}}/kT)+1)}},`$
$`I_{1/2}^{}(V,H_z^{\mathrm{eff}})`$ $`=`$ $`{\displaystyle \frac{(eV+2\mu _BH_z^{\mathrm{eff}})(e^{eV/kT}1)}{(e^{(eV+2\mu _BH_z^{\mathrm{eff}})/kT}1)(2\mathrm{cosh}(2\mu _BH_z^{\mathrm{eff}}/kT)+1)}},`$
$`I_{1/2}^{}(V,H_z^{\mathrm{eff}})`$ $`=`$ $`{\displaystyle \frac{(eV2\mu _BH_z^{\mathrm{eff}})(e^{eV/kT}1)e^{2\mu _BH/kT}}{(e^{(eV2\mu _BH_z^{\mathrm{eff}})/kT}1)(2\mathrm{cosh}(2\mu _BH_z^{\mathrm{eff}}/kT)+1)}}.`$
All the other ones, not written above, are equal to $`eV`$.
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# R-Parity Violation and the Decay 𝑏→𝑠𝛾
## I Introduction
The decay $`bs\gamma `$, forbidden at tree-level in the Standard Model (SM), is an excellent candidate for exploring the influence of new physics beyond the SM. However, the experimentally measured branching ratio $`\text{Br}(BX_s\gamma )=(3.15\pm 0.93)10^4`$ is in perfect agreement with the SM prediction computed at the next-to-leading order:$`\text{Br}(bs\gamma )_{SM}=(3.28\pm 0.30)10^4`$ . This leads to the conclusion that the influence of new physics on this decay is either very limited or the new contributions cancel among each other to a large extent. The most serious and attractive extention of the Standard Model is Supersymmetry (SUSY). It has a variety of very appealing features. For instance it provides a natural solution to the hierarchy problem. SUSY doubles the particle spectrum of the SM providing every fermion with a bosonic partner and vice versa. If SUSY were an exact symmetry of nature the new particles would be of the same mass as their partners. This is definitely excluded from experiment. Therefore SUSY must be broken. To retain the solution of the hierarchy problem one allows only a breaking which does not introduce quadratic divergences in loop diagrams. Although this reduces the set of possible breaking terms substantially to the so-called *soft-breaking* terms the number of free parameters of softly broken SUSY still exceeds one hundred. This theory, i.e., SUSY with soft breaking terms *and no others*, emerges naturally as low energy limit of local supersymmetry, *Supergravity (SUGRA)* by breaking it at some high scale $`10^{11}`$ GeV and taking the so-called *flat limit* . The connection to supergravity eliminates much of the freedom in choosing the parameters for the soft breaking terms, enhancing the predictive power of the model.
Viewing naturalness as a first principle in model building one runs immediately into a problem: The Yukawa interactions are fixed by the so-called *Superpotential $`W`$*, which is the most general third degree polynomial that can be built of gauge invariant combinations of the (left-handed) superfields of the theory. In the case of the minimal extension of the standard model with gauge group $`SU(3)\times SU(2)\times U(1)`$ it shows the form
$`W`$ $`=`$ $`\lambda _{ij}^uQ_iH_2U_j^c+\lambda _{ij}^dH_1Q_iD_j^c+\lambda _{ij}^eH_1L_iE_j^c+\mu H_1H_2`$ (2)
$`+\mu _{2i}L_iH_2+{\displaystyle \frac{1}{2}}\lambda _{ijk}L_iL_jE_k^c+\lambda _{ijk}^{}L_iQ_jD_k^c+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c,`$
where $`Q,U^c,D^c,L,E^c,H_1`$ and $`H_2`$ label the left-handed superfields that describe left- and right-handed (s)quarks and (s)leptons and the Higgs-bosons (-fermions) respectively. The terms on the second line of (2) lead to unwanted baryon- and lepton-number violating vertices. The product $`\lambda _{ijk}^{}\lambda _{ijk}^{\prime \prime }`$, for instance, is restricted to be smaller than $`10^{10}`$ . It is not obvious why the coefficients of these terms should not be of order unity if the corresponding interaction is not protected by any symmetry. Hence one way to avoid these unwanted couplings is the invention of a new discrete symmetry, called *R-Parity* . The multiplicative quantum number $`R`$ is then defined as
$$R=(1)^{3B+L+2s},$$
(3)
where $`s`$ is the spin of the particle. The particles of the standard model are then R-parity even fields while their supersymmetric partners are R-parity odd fields. The superfields adopt the value of R from their scalar components. The terms on the second line of (2) are then forbidden by this symmetry. One ends up with the *Minimal Supersymmetric Standard Model (MSSM)* . One could even go one step further and promote this new symmetry to a $`U(1)`$ gauge symmetry (*R-symmetry*). The R-charges must then be chosen such that the “nice” terms remain in the Lagrangian and the unwanted terms will not be allowed anymore. Further requirement is the vanishing of possible additional anomalies . This scenario is preferred from a string theoretical view because in string theory one has many additional $`U(1)`$’s “floating around”, what makes the introduction of this symmetry more natural.
It is interesting to explore what the constraints on the R-parity violating couplings, especially $`\lambda _{ijk}`$, $`\lambda _{ijk}^{}`$, and $`\lambda _{ijk}^{\prime \prime }`$, are from the experimental point of view. Bounds on these couplings have been found by many authors . Sometimes the bounds on products of couplings are more restrictive than the products of the individual bounds. Depending on the reactions that the couplings are involved in the constraints from experiment are very strong or rather poor.
In this paper we want to explore the (theoretical) influence of $`\lambda `$, $`\lambda ^{}`$, and $`\lambda ^{\prime \prime }`$ on the decay $`bs\gamma `$. This has been done before by other authors . However, here we will include the full operator basis at leading-log of the effective low-energy theory.
The comparison with the experiment must result in bounds that are model dependent. The relevant couplings are proportional to the inverse mass-squared of particles (Higgs, SUSY-partners) that have not yet been detected. However, it is clear from next-to-leading-log calculations that the prediction of the MSSM with a realistic mass spectrum lies within the current experimental bounds.
Because of the strong model dependence we do not perform our calculations at highest precision. Nevertheless we try to include all the possibly relevant contributions at leading-log. We do not claim our results to be very accurate. It is the aim of this paper to explore if $`bs\gamma `$ has the potential power to reduce the bounds of (products of) some of the R-parity breaking couplings substantially, i.e., by some order of magnitude.
This article is divided as follows: Section 2 introduces the model we are working with. We try to describe as precise as possible what our assumptions are. The following section deals with the effective Hamiltonian approach. The enhanced operator basis is presented and the $`\gamma `$-matrix as well as the matching conditions at $`M_W`$ calculated. The comparison with the experiment is performed in section 4. Section 5 contains our conclusions. In the appendix we present some technical details of our computations, namely the mixing matrices,the relevant part of the interaction Lagrangian and the RGE’s that are needed.
## II Framework
In supersymmetry the matter fields are described by left-handed chiral superfields $`\mathrm{\Sigma }^i`$. They contain a scalar boson $`z^i`$ and a two-component fermion $`\psi ^i`$. Real vector superfields $`V^a`$ are needed to form the gauge bosons $`A_\mu ^a`$ and the gauginos $`\lambda ^a`$. The minimal supersymmetric standard model is the model with the smallest particle content that is able to mimic the features of the standard model, i.e., including all the observed particles, gauge group $`SU(3)_{\mathrm{colour}}\times SU(2)_{\mathrm{weak}}\times U(1)_Y`$, spontaneous symmetry breaking and the Higgs mechanism. Its superfields together with their components are collected in table I. A few comments are in order:
* Because the theory only deals with left-handed chiral superfields the $`SU(2)`$-singlet matter fields must be defined via their charged conjugated (anti-) fields.
* No conjugated superfields are allowed in the superpotential $`W`$. Therefore we need to introduce a second Higgs field to give the up and the down quarks a mass when the neutral components gain a vacuum expectation value (vev).
* The $`L`$’s and $`R`$’s in the names of squarks and sleptons only identify the fermionic partners. These fields are just normal complex scalar bosons.
It must be mentioned that the fields in table I are *not* the physical fields.
* In the Higgs sector three degrees of freedom are eaten by the gauge bosons in analogy to the SM. We end up with one charged and three neutral Higgs bosons.
* Higgsinos and gauginos of $`SU(2)\times U(1)`$ mix to form *charginos* and *neutralinos*.
* Photon, $`W`$\- and $`Z`$-boson form when the electroweak symmetry breaks down.
* The three generations of quarks and leptons mix via the Cabibbo-Kobayashi-Maskawa-matrix $`K`$ to give the mass eigenstates in complete analogy to the standard model.
* The family mixing also takes place in the squark and slepton sector. However, there is an additional mixing between the partners of left- and right-handed fermions due to the soft-breaking terms.
Appendix A gives a more detailed description of the different mixings including a complete listing of all relevant mixing matrices.
The component expression of the Lagrangian we base our model on can be written as
$$=_{\mathrm{kin}}+_{\mathrm{int}}+_WV_{\mathrm{soft}},$$
(4)
where
* $`_{\mathrm{kin}}`$ and $`_{\mathrm{int}}`$ stand for the kinetic energies, the interactions between chiral and gauge fields and part of the scalar potential.
* $`_W`$ contains the rest of the scalar potential and the Yukawa interactions:
$$_W=\underset{i}{}\left|\frac{W}{z_i}\right|^2\frac{1}{2}\underset{ij}{}\left[\psi _i\frac{^2W}{z_iz_j}\psi _j+\text{h.c.}\right],$$
(5)
where $`W`$ should be viewed as a function of the scalar fields.
* $`W`$ is the superpotential which contains all possible gauge invariant combinations of the left-handed superfields (but not their conjugated right-handed partners). In our case, this results in equation (2). It should be noted that every term is a gauge invariant combination of the corresponding superfields, for instance, $`Q_iH_2U_j^c`$ is abbreviated for $`ϵ_{\alpha \beta }Q_i^{\alpha A}H_2^\beta U_j^{cA}`$, where $`\alpha `$ and $`\beta `$ are $`SU(2)`$-indices, $`A`$ is an $`SU(3)`$-index and $`ϵ_{\alpha \beta }`$ is the completely antisymmetric tensor with $`ϵ_{12}=1`$ <sup>*</sup><sup>*</sup>*The antisymmetry of the $`SU(2)`$ product is the reason why a term $`H_1H_1E_i^c`$ is not introduced. The second line of (2) represents the R-parity breaking sector. As a consequence of the antisymmetry in the fields $`\lambda `$ is antisymmetric in its first two indices and $`\lambda ^{\prime \prime }`$ is antisymmetric in its last two indices. Therefore the trilinear R-parity breaking couplings of $`W`$ contain $`9+27+9=45`$ new parameters.
* $`V_{\mathrm{soft}}`$ includes the soft-breaking trilinear terms of the scalar potential and mass terms for the scalar fields and the gauginos. It has the form
$`V_{\mathrm{soft}}`$ $`=`$ $`[h_{ij}^u\stackrel{~}{q}_i\stackrel{~}{h}_2\stackrel{~}{u}_j^{}+h_{ij}^d\stackrel{~}{h}_1\stackrel{~}{q}_i\stackrel{~}{d}_j^{}+h_{ij}^e\stackrel{~}{h}_1\stackrel{~}{\mathrm{}}_i\stackrel{~}{e}_j^{}`$ (9)
$`+{\displaystyle \frac{1}{2}}C_{ijk}\stackrel{~}{\mathrm{}}_i\stackrel{~}{\mathrm{}}_j\stackrel{~}{e}_k^{}+C_{ijk}^{}\stackrel{~}{\mathrm{}}_i\stackrel{~}{q}_j\stackrel{~}{d}_k^{}+{\displaystyle \frac{1}{2}}C_{ijk}^{\prime \prime }\stackrel{~}{u}_i^{}\stackrel{~}{d}_j^{}\stackrel{~}{d}_k^{}`$
$`+\stackrel{~}{\mu }\stackrel{~}{h}_1\stackrel{~}{h}_2+\stackrel{~}{\mu }_{2i}\stackrel{~}{\mathrm{}}_i\stackrel{~}{h}_2]+\text{h.c.}+m_{ab}^2z^az^b`$
$`+{\displaystyle \frac{1}{2}}M_a[\lambda _a\lambda _a+\text{h.c.}].`$
The MSSM in its full generality with an additional R-parity breaking sector involves over 150 free parameters. These are by far too many for the model to be predictive. In the following we will reduce the parameter space substantially by making some assumptions that are, hopefully, well motivated.
As a first step it is important to mention that we see our softly broken global SUSY at a low energy scale $`M_Z`$ emerging from a spontaneously broken local supersymmetry at a high scale $`M_X10^{16}`$ GeV taking the flat limit $`M_{\mathrm{Planck}}\mathrm{}`$, $`m_0:=m_{\mathrm{gravitino}}=`$constant. This fixes most of the parameters of $`V_{\mathrm{soft}}`$ at $`M_X`$:
* All the coefficients of the trilinear terms in $`V_{\mathrm{soft}}`$ are related to the corresponding terms of $`W`$ by a multiplication with a universal factor $`Am_0`$:
$$\begin{array}{ccc}h_{ij}^u=Am_0\lambda _{ij}^u\hfill & C_{ijk}=Am_0\lambda _{ijk}& \\ & & \\ h_{ij}^d=Am_0\lambda _{ij}^d\hfill & C_{ijk}^{}=Am_0\lambda _{ijk}^{}& \\ & & \\ h_{ij}^e=Am_0\lambda _{ij}^e\hfill & C_{ijk}^{\prime \prime }=Am_0\lambda _{ijk}^{\prime \prime }& \end{array}$$
(10)
* An analogous statement holds for the bilinear terms:
$`\stackrel{~}{\mu }=Bm_0\mu `$ $`\stackrel{~}{\mu }_{2i}=Bm_0\mu _{2i},`$ (11)
where usually
$$B=A1.$$
(12)
* The mass term of the scalars are diagonal and universally equal to the gravitino mass $`m_0`$:
$$m_{ab}^2=m_0^2\delta _{ab}$$
(13)
We assume unification of the gauge group at $`M_X`$. As a consequence all the gaugino masses are equal at that scale:
$$M_i(M_X)=Mi$$
(14)
Not all entries of the Yukawa-matrices $`\lambda _{ij}^u`$, $`\lambda _{ij}^d`$, and $`\lambda _{ij}^e`$ are observable in the SM. One usually chooses two of them (in most cases $`\lambda ^d`$ and $`\lambda ^e`$) to be diagonal. Although this is in principle not possible in our model we will adopt this choice here for convenience. All the entries at *$`M_W`$* are then fixed by the quark/lepton masses, the vevs $`v_1`$ and $`v_2`$ of the neutral Higgs bosons $`H_1`$ and $`H_2`$ respectively and the Cabibbo-Kobayashi-Maskawa-matrix $`K`$.
$`\mu H_1H_2`$ is the so-called $`\mu `$-term. The mass parameter $`\mu `$ must be of order of the weak scale whereas the natural scale would be the Planck mass $`M_P10^{19}`$ GeV. The question why this parameter is so small is referred to as the *$`\mu `$-problem*.
$`\mu _{2i}L_iH_2`$ and $`\stackrel{~}{\mu }_{2i}\stackrel{~}{\mathrm{}}_i\stackrel{~}{h}_2`$ mix Higgs and leptonic sector. We choose to set $`\mu _{2i}(M_X)=\stackrel{~}{\mu }_{2i}(M_X)=0`$. At $`M_Z`$, $`\mu _{2i}L_iH_2`$ can be rotated away with the help of a field redefinition of the Higgs field whereas $`\stackrel{~}{\mu }_{2i}`$ is, at least in the case of a physically realistic spectrum, small enough to be neglected.
One ends up with the following free parameters:
$$A,m_0,M,\mu $$
(15)
Usually, one replaces one of these parameters by $`\mathrm{tan}\beta =v_2/v_1`$. A second parameter will be fixed by the requirement of a correct electroweak symmetry breaking. This means, the minimum of the scalar Higgs potential must occur at values $`(v_1,v_2)`$ which reproduce the correct mass of the $`Z`$-boson:
$$M_Z^2=\frac{g_1^2+g_2^2}{2}(v_1^2+v_2^2)$$
(16)
It has been realized by many authors that the tree-level potential
$`V_0`$ $`=`$ $`(\mu ^2+m_{H_1H_1}^2)(\stackrel{~}{h}_1^0)^2+(\mu ^2+m_{H_2H_2}^2)(\stackrel{~}{h}_2^0)^2+2\stackrel{~}{\mu }\stackrel{~}{h}_1^0\stackrel{~}{h}_2^0`$ (18)
$`+{\displaystyle \frac{g_1^2+g_2^2}{8}}\left[(\stackrel{~}{h}_1^0)^2(\stackrel{~}{h}_2^0)^2\right]^2`$
is not enough to gain sensible values for $`v_1`$ and $`v_2`$. Thus we have to include the first correction $`\mathrm{\Delta }V`$ to the effective potential. At a mass scale $`Q`$, it has the form
$`\mathrm{\Delta }V`$ $`=`$ $`{\displaystyle \frac{1}{64\pi ^2}}\text{Str}\left[^4\left(\mathrm{ln}{\displaystyle \frac{^2}{Q^2}}{\displaystyle \frac{3}{2}}\right)\right]`$ (19)
$`=`$ $`{\displaystyle \frac{1}{64\pi ^2}}{\displaystyle \underset{p}{}}(1)^{2s_p}(2s_p+1)n_pM_p^4\left(\mathrm{ln}{\displaystyle \frac{|M_p^2|}{Q^2}}{\displaystyle \frac{3}{2}}\right).`$ (20)
Here, Str denotes the supertrace and $`^2`$ is the tree-level mass matrix squared. $`p`$ runs over all particles of the theory with spin $`s_p`$ whereas $`M_p`$ is the corresponding eigenvalue (mass) of the particle. $`n_p`$ counts for the degrees of freedom according to colour and helicity. The eigenvalues $`M_p`$ depend on the neutral components of $`\stackrel{~}{h}_1`$ and $`\stackrel{~}{h}_2`$ and therefore change the shape of the potential and its minimum. For most of the particles, they can only be computed numerically. One comment is in order: A phase rotation of the Higgs fields can turn a negative vev into a positive one. This freedom is reflected in the fact that the *sign* of $`\mu `$ can be chosen freely giving then a different phenomenology.
To compute the mixing matrices and the one-loop effective potential one has to know the mass matrices at $`M_Z`$. Unfortunately, for most of the parameters we know the boundary conditions at the high scale $`M_X`$. Hence one has to set up the complete set of RGE’s to run the parameters from $`M_X`$ to $`M_Z`$. The details of how to get a consistent parameter space are described in section IV, the complete set of RGE’s can be found in appendix B.
## III The effective Hamiltonian
### A The case of the MSSM
The decay $`bs\gamma `$ occurs at energies of a few GeV$`m_b`$. This is much below the weak scale. It makes sense to work with an effective Hamiltonian $`_{\mathrm{eff}}`$ where all the heavy fields (compared to $`m_b`$) are integrated out . In the SM these are the $`Z`$\- and the $`W`$-boson and the top quark whereas for our purposes the $`Z`$-boson does not play any role. A result of integrating out these fields is the appearance of new local operators of dimension higher than four. This can be illustrated by the shrinking of the Feynman diagrams in Fig. 1.
Once one lets the strong interaction come into the game, QCD corrections of these new operators give rise to additional operators. For a consistent treatment of all the corrections at a certain level of QCD we have to include a set of operators $`O_i`$ which closes under these corrections. Our low-energy theory is then described by an effective Hamiltonian
$$_{\mathrm{eff}}=\underset{i}{}C_iO_i.$$
(21)
The QCD renormalization of the operators must be performed at a scale where no large logarithms appear, i.e., at $`\mu M_W`$. However, a calculation of $`\text{Br}(bs\gamma )`$ at energy scales $`m_b`$ requires the knowledge of the Wilson coefficients $`C_i`$ at *that* scale. The $`C_i`$’s depend on a renormalization scale $`\mu `$. They obey the renormalization group equations
$`{\displaystyle \frac{dC_i}{d\mathrm{ln}\mu }}=\gamma _{ji}C_j,`$ (22)
where $`\gamma _{ij}`$ is the gamma-matrix that emerges from the QCD-mixing of the operators $`O_i`$. As initial conditions we find $`C_i(M_W)`$ by matching the effective theory with the full theory at that scale. The RGE’s are then solved to give $`C_i(m_b)`$
In the standard model the relevant set of operators for $`bs\gamma `$ is given by
$`O_1`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\alpha })(\overline{c_L}_\beta \gamma _\mu c_{L\beta })`$ (23)
$`O_2`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\beta })(\overline{c_L}_\beta \gamma _\mu c_{L\alpha })`$ (25)
$`O_3`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\alpha }){\displaystyle \underset{i}{}}(\overline{q_L}_{i\beta }\gamma _\mu q_{Li\beta })`$ (27)
$`O_4`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\beta }){\displaystyle \underset{i}{}}(\overline{q_L}_{i\beta }\gamma _\mu q_{Li\alpha })`$ (29)
$`O_5`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\alpha }){\displaystyle \underset{i}{}}(\overline{q_R}_{i\beta }\gamma _\mu q_{Ri\beta })`$ (31)
$`O_6`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\beta }){\displaystyle \underset{i}{}}(\overline{q_R}_{i\beta }\gamma _\mu q_{Ri\alpha })`$ (33)
$`O_7`$ $`=`$ $`{\displaystyle \frac{e}{16\pi ^2}}m_b\overline{s_L}_\alpha \sigma _{\mu \nu }b_{R\alpha }F^{\mu \nu }`$ (35)
$`O_8`$ $`=`$ $`{\displaystyle \frac{g_s}{16\pi ^2}}m_b\overline{s_L}_\alpha \sigma _{\mu \nu }b_{R\beta }t^{a\alpha \beta }G^{a\mu \nu }.`$ (37)
Here, $`\overline{q_{L/R}}\gamma ^\mu q_{L/R}=\overline{q}\gamma ^\mu (1\gamma ^5)q`$, $`\overline{s}_L\sigma _{\mu \nu }b_R=\overline{s}\sigma _{\mu \nu }(1+\gamma ^5)b`$, the sum runs over the five quarks of the effective theory at $`m_b`$, and $`\alpha `$, $`\beta `$ are colour indices. The operators $`O_1`$ and $`O_3O_6`$ are introduced by QCD corrections through diagrams like those depicted in Fig. 2.
It is interesting to note that in the MSSM without R-parity violation the relevant operator basis does not change although we definitely have new decay channels. The $`\gamma `$-matrix at leading-log can be found in appendix C by picking up the relevant entries. Because the contribution of the diagram in Fig. 3 is not divergent, the mixing of $`O_1O_6`$ with $`O_{7/8}`$ involves two-loop diagrams like those of figures 4 and 5 .
Although we only need the divergent part of these diagrams one has to be very careful in computing the counterterms because one must include certain additional operators , so-called *evanescent operators* that vanish in four dimensions but must be kept in $`D`$ dimensions in all intermediate steps of the calculation. Furthermore, the two-loop results are regularization scheme dependent . This regularization scheme dependence of the $`\gamma `$-matrix cancels with possible finite one-loop \[but $`O(\alpha _s^0)`$\] contributions from $`O_5`$ and $`O_6`$, inserted in the diagram of Fig. 3, to the matrix element of $`bs\gamma `$. As a result of these complications, the calculation of the $`\gamma `$-matrix at leading-log has been finished only few years ago . Meanwhile, also the next-to leading result is known . We will concentrate on the leading-log calculation.
The matching of the effective theory with the full theory at $`M_W`$ must be performed only at order $`\alpha _s^0`$ because the leading QCD corrections are already included in the operator mixing. In the standard model this involves diagrams with a $`W`$-exchange for $`O_2`$ and diagrams with a $`W`$-$`t`$-loop (in the unitary gauge) for $`O_{7/8}`$ corresponding to Fig. 6 (a). The result is then
$`C_{2\mathrm{S}\mathrm{M}}(M_W)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}`$ (38)
$`C_{7\mathrm{S}\mathrm{M}}(M_W)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}3x_{tW}[Q_uF_1(x_{tW})+F_2(x_{tW})]`$ (39)
$`C_{8\mathrm{S}\mathrm{M}}(M_W)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}3x_{tW}F_1(x_{tW}),`$ (40)
where $`G_F=g_2^2/(4\sqrt{2}M_W^2)1.16610^5`$ GeV<sup>-2</sup> is the Fermi constant, $`x_{ab}=m_a^2/m_b^2`$ and the functions $`F_i`$ are given in appendix D.
The matching conditions for $`C_{7/8}`$ become far more complicated in the case of the MSSM, even without R-parity violation. In addition to the $`W`$-$`t`$-loop there are four more combinations of particles in the loop:
* charged Higgs $`H^\pm `$—top $`t`$ ,Fig. 6 (b)
* up-squark $`\stackrel{~}{u}`$—chargino $`\chi ^{\mathrm{ch}}`$ ,Fig. 6 (c)
* down-squark $`\stackrel{~}{d}`$—neutralino $`\chi ^0`$ ,Fig. 6 (d)
* down-squark $`\stackrel{~}{d}`$—gluino $`g`$ ,Fig. 6 (e)
In principle, these contributions give rise to two new operators
$`\stackrel{~}{O}_7`$ $`=`$ $`{\displaystyle \frac{e}{16\pi ^2}}m_b\overline{s_R}_\alpha \sigma _{\mu \nu }b_{L\alpha }F^{\mu \nu }`$ (41)
$`\stackrel{~}{O}_8`$ $`=`$ $`{\displaystyle \frac{g_s}{16\pi ^2}}m_b\overline{s_R}_\alpha \sigma _{\mu \nu }b_{L\beta }t^{a\alpha \beta }G^{a\mu \nu }.`$ (43)
Note that $`\stackrel{~}{O}_{7/8}`$ differ from $`O_{7/8}`$ only by their handedness. The Wilson coefficients of these new operators are usually so small that they can be safely neglected . However, we will include them in our calculations because we need them later on anyway. We only neglect contributions from the Higgs sector which are proportional to some light quark masses.
The matching conditions become rather involved now. They include many mixing matrices, whose definitions we give in appendix A 1. We have
$`C_{7\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}`$ $`=`$ $`C_{7\mathrm{S}\mathrm{M}}`$ (50)
$`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}\{\mathrm{cot}^2\beta x_{tH}[Q_uF_1(x_{tH})+F_2(x_{tH})]`$
$`+x_{tH}[Q_uF_3(x_{tH})+F_4(x_{tH})]\}`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_j}^2}}B_{2j\mathrm{}}^dB_{3j\mathrm{}}^d\left[F_1(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})+Q_uF_2(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})\right]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_j}^2}}{\displaystyle \frac{m_{\chi _{\mathrm{}}^{\mathrm{ch}}}}{m_b}}B_{2j\mathrm{}}^dA_{3j\mathrm{}}^d\left[F_3(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})+Q_uF_4(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})\right]`$
$`+{\displaystyle \frac{Q_d}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[D_{2j\mathrm{}}^dD_{3j\mathrm{}}^dF_2(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})+{\displaystyle \frac{m_{\chi _{\mathrm{}}^0}}{m_b}}D_{2j\mathrm{}}^dC_{3j\mathrm{}}^dF_4(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})\right]`$
$`+{\displaystyle \frac{2}{3}}Q_dg_s^2{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[\mathrm{\Gamma }_{L2j}^d\mathrm{\Gamma }_{Lj3}^dF_2(x_{g\stackrel{~}{d}_j}){\displaystyle \frac{m_g}{m_b}}\mathrm{\Gamma }_{L2j}^d\mathrm{\Gamma }_{Rj3}^dF_4(x_{g\stackrel{~}{d}_j})\right]`$
$`C_{8\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}`$ $`=`$ $`C_{8\mathrm{S}\mathrm{M}}`$ (57)
$`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}[\mathrm{cot}^2\beta x_{tH}F_1(x_{tH})+x_{tH}F_3(x_{tH})]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_j}^2}}\left[B_{2j\mathrm{}}^dB_{3j\mathrm{}}^dF_2(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})+{\displaystyle \frac{m_{\chi _{\mathrm{}}^{\mathrm{ch}}}}{m_b}}B_{2j\mathrm{}}^dA_{3j\mathrm{}}^dF_4(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})\right]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[D_{2j\mathrm{}}^dD_{3j\mathrm{}}^dF_2(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})+{\displaystyle \frac{m_{\chi ^0}}{m_b}}D_{2j\mathrm{}}^dC_{3j\mathrm{}}^dF_4(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})\right]`$
$`+{\displaystyle \frac{2}{3}}g_s^2{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[\mathrm{\Gamma }_{L2j}^d\mathrm{\Gamma }_{Lj3}^dF_2(x_{g\stackrel{~}{d}_j}){\displaystyle \frac{m_g}{m_b}}\mathrm{\Gamma }_{L2j}^d\mathrm{\Gamma }_{Rj3}^dF_4(x_{g\stackrel{~}{d}_j})\right]`$
$`{\displaystyle \frac{3}{4}}g_s^2{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[\mathrm{\Gamma }_{L2j}^d\mathrm{\Gamma }_{Lj3}^dF_1(x_{g\stackrel{~}{d}_j}){\displaystyle \frac{m_g}{m_b}}\mathrm{\Gamma }_{L2j}^d\mathrm{\Gamma }_{Rj3}^dF_3(x_{g\stackrel{~}{d}_j})\right]`$
$`\stackrel{~}{C}_{7\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}{\displaystyle \frac{m_sm_b}{m_t^2}}\mathrm{tan}^2\beta x_{tH}[Q_uF_1(x_{tH})+F_2(x_{tH})]`$ (63)
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_j}^2}}A_{2j\mathrm{}}^dA_{3j\mathrm{}}^d\left[F_1(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})+Q_uF_2(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})\right]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_j}^2}}{\displaystyle \frac{m_{\chi _{\mathrm{}}^{\mathrm{ch}}}}{m_b}}A_{2j\mathrm{}}^dB_{3j\mathrm{}}^d\left[F_3(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})+Q_uF_4(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})\right]`$
$`+{\displaystyle \frac{Q_d}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[C_{2j\mathrm{}}^dC_{3j\mathrm{}}^dF_2(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})+{\displaystyle \frac{m_{\chi _{\mathrm{}}^0}}{m_b}}C_{2j\mathrm{}}^dD_{3j\mathrm{}}^dF_4(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})\right]`$
$`+{\displaystyle \frac{2}{3}}Q_dg_s^2{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[\mathrm{\Gamma }_{R2j}^d\mathrm{\Gamma }_{Rj3}^dF_2(x_{g\stackrel{~}{d}_j}){\displaystyle \frac{m_g}{m_b}}\mathrm{\Gamma }_{R2j}^d\mathrm{\Gamma }_{Lj3}^dF_4(x_{g\stackrel{~}{d}_j})\right]`$
$`\stackrel{~}{C}_{8\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}{\displaystyle \frac{m_sm_b}{m_t^2}}\mathrm{tan}^2\beta x_{tH}F_1(x_{tH})`$ (69)
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_j}^2}}\left[A_{2j\mathrm{}}^dA_{3j\mathrm{}}^dF_2(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})+{\displaystyle \frac{m_{\chi _{\mathrm{}}^{\mathrm{ch}}}}{m_b}}A_{2j\mathrm{}}^dB_{3j\mathrm{}}^dF_4(x_{\chi _{\mathrm{}}^{\mathrm{ch}}\stackrel{~}{u}_j})\right]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[C_{2j\mathrm{}}^dC_{3j\mathrm{}}^dF_2(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})+{\displaystyle \frac{m_{\chi _{\mathrm{}}^0}}{m_b}}C_{2j\mathrm{}}^dD_{3j\mathrm{}}^dF_4(x_{\chi _{\mathrm{}}^0\stackrel{~}{d}_j})\right]`$
$`+{\displaystyle \frac{2}{3}}g_s^2{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[\mathrm{\Gamma }_{R2j}^d\mathrm{\Gamma }_{Rj3}^dF_2(x_{g\stackrel{~}{d}_j}){\displaystyle \frac{m_g}{m_b}}\mathrm{\Gamma }_{R2j}^d\mathrm{\Gamma }_{Lj3}^dF_4(x_{g\stackrel{~}{d}_j})\right]`$
$`{\displaystyle \frac{3}{4}}g_s^2{\displaystyle \frac{1}{m_{\stackrel{~}{d}_j}^2}}\left[\mathrm{\Gamma }_{R2j}^d\mathrm{\Gamma }_{Rj3}^dF_1(x_{g\stackrel{~}{d}_j}){\displaystyle \frac{m_g}{m_b}}\mathrm{\Gamma }_{R2j}^d\mathrm{\Gamma }_{Lj3}^dF_3(x_{g\stackrel{~}{d}_j})\right].`$
### B Including the R-parity breaking terms
#### 1 The new operator basis
As mentioned before, the difference between the standard model and the MSSM does not lie in a change of the operator basis but rather in different matching conditions at $`M_W`$. This situation changes drastically if one includes the R-parity breaking sector. Now our basis has to be enlarged. To find out which are the relevant new operators we first write down the R-parity breaking Yukawa couplings:
$$_{\mathrm{Yukawa}\mathit{}}=_\lambda +_\lambda ^{}+_{\lambda ^{\prime \prime }},$$
(70)
where
$`_\lambda `$ $`=`$ $`{\displaystyle \frac{1}{2}}\lambda _{ijk}\left[\mathrm{\Gamma }_{Li\mathrm{}}^e\stackrel{~}{e}_{\mathrm{}}\overline{e_R}_k\nu _{Lj}+\mathrm{\Gamma }_j\mathrm{}^\nu \stackrel{~}{\nu }_{\mathrm{}}\overline{e_R}_ke_{Li}+\mathrm{\Gamma }_{R\mathrm{}k}^e\stackrel{~}{e}_{\mathrm{}}^{}\overline{e_R^c}_i\nu _{Lj}\right]`$ (72)
$`+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{}\left[\mathrm{\Gamma }_{L\mathrm{}i}^e\stackrel{~}{e}_{\mathrm{}}^{}\overline{\nu _L}_je_{Rk}+\mathrm{\Gamma }_\mathrm{}j^\nu \stackrel{~}{\nu }_{\mathrm{}}^{}\overline{e_L}_ie_{Rk}+\mathrm{\Gamma }_{Rk\mathrm{}}^e\stackrel{~}{e}_{\mathrm{}}\overline{\nu _L}_je_{Ri}^c\right]`$
$`_\lambda ^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\lambda _{ijk}^{}[\mathrm{\Gamma }_{Li\mathrm{}}^eK_{jm}^{}\stackrel{~}{e}_{\mathrm{}}\overline{d_R}_ku_{Lm}+\mathrm{\Gamma }_{Lj\mathrm{}}^u\stackrel{~}{u}_{\mathrm{}}\overline{d_R}_ke_{Li}+\mathrm{\Gamma }_{R\mathrm{}k}^dK_{jm}^{}\stackrel{~}{d}_{\mathrm{}}^{}\overline{e_R^c}_iu_{Lm}`$ (77)
$`\mathrm{\Gamma }_i\mathrm{}^\nu \stackrel{~}{\nu }_{\mathrm{}}\overline{d_R}_kd_{Lj}\mathrm{\Gamma }_{Lj\mathrm{}}^d\stackrel{~}{d}_{\mathrm{}}\overline{d_R}_k\nu _{Li}\mathrm{\Gamma }_{R\mathrm{}k}^d\stackrel{~}{d}_{\mathrm{}}^{}\overline{\nu _R^c}_id_{Lj}]`$
$`+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{}[\mathrm{\Gamma }_{L\mathrm{}i}^eK_{mj}\stackrel{~}{e}_{\mathrm{}}^{}\overline{u_L}_md_{Rk}+\mathrm{\Gamma }_{L\mathrm{}j}^u\stackrel{~}{u}_{\mathrm{}}^{}\overline{e_L}_id_{Rk}+\mathrm{\Gamma }_{Rk\mathrm{}}^dK_{mj}\stackrel{~}{d}_{\mathrm{}}\overline{u_L}_me_{Ri}^c`$
$`\mathrm{\Gamma }_\mathrm{}i^\nu \stackrel{~}{\nu }_{\mathrm{}}^{}\overline{d_L}_jd_{Rk}\mathrm{\Gamma }_{L\mathrm{}j}^d\stackrel{~}{d}_{\mathrm{}}^{}\overline{\nu _L}_id_{Rk}\mathrm{\Gamma }_{Rk\mathrm{}}^d\stackrel{~}{d}_{\mathrm{}}\overline{d_L}_j\nu _{Ri}^c]`$
$`_{\lambda ^{\prime \prime }}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\lambda _{ijk}^{\prime \prime }\left[\mathrm{\Gamma }_{R\mathrm{}i}^u\stackrel{~}{u}_{\mathrm{}}^{}\overline{d_R}_jd_{Lk}^c+\mathrm{\Gamma }_{R\mathrm{}j}^d\stackrel{~}{d}_{\mathrm{}}^{}\overline{u_R}_id_{Lk}^c+\mathrm{\Gamma }_{R\mathrm{}k}^d\stackrel{~}{d}_{\mathrm{}}^{}\overline{u_R}_id_{Lj}^c\right]`$ (80)
$`{\displaystyle \frac{1}{4}}\lambda _{ijk}^{\prime \prime }\left[\mathrm{\Gamma }_{Ri\mathrm{}}^u\stackrel{~}{u}_{\mathrm{}}\overline{d_L^c}_kd_{Rj}+\mathrm{\Gamma }_{Rj\mathrm{}}^d\stackrel{~}{d}_{\mathrm{}}\overline{d_L^c}_ku_{Ri}+\mathrm{\Gamma }_{Rk\mathrm{}}^d\stackrel{~}{d}_{\mathrm{}}\overline{d_L^c}_ju_{Ri}\right].`$
Here, all the fields belong to the mass basis. The colour indices have been omitted.
The next task is to build four-quark operators out of two Yukawa couplings that contribute at $`O(\alpha _s)`$ to $`bs\gamma `$. The boson serves as a bridge between the fermions in analogy to the $`W`$ boson in the standard model. The following points must be taken care of:
* The top quark is not at our disposal in the five flavour effective theory.
* We not only need an ingoing $`b`$ and an outgoing $`s`$. The two remaining quarks must be of the same type because one has to be able to close the loop with these fermions.
It is clear that $`_\lambda `$ cannot participate because it contains no squarks and semi-leptonic operators can be neglected. Also, $`_\lambda ^{}`$ does not mix with $`_{\lambda ^{\prime \prime }}`$. As an example, we take the first term of $`_\lambda ^{}`$ together with his hermitian conjugate. The situation is depicted in Fig. 7:
$`{\displaystyle \frac{i}{2}}\lambda _{ij2}^{}\mathrm{\Gamma }_{Li\mathrm{}}^eK_{jm}^{}\stackrel{~}{e}_{\mathrm{}}\overline{s_R}u_{Lm}{\displaystyle \frac{i}{k^2m_{\stackrel{~}{e}_{\mathrm{}}}^2}}{\displaystyle \frac{i}{2}}\lambda _{ab3}^{}\mathrm{\Gamma }_{L\mathrm{}a}^eK_{mb}\stackrel{~}{e}_{\mathrm{}}^{}\overline{u_L}_mb_R`$ (81)
$`\stackrel{k^2m_{\stackrel{~}{e}_{\mathrm{}}}^2}{}`$ $`{\displaystyle \frac{i}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}\lambda _{ij2}^{}\lambda _{aj3}^{}\mathrm{\Gamma }_{Li\mathrm{}}^e\mathrm{\Gamma }_{L\mathrm{}a}^e(\overline{s_R}u_{Lj})(\overline{u_L}_jb_R)`$ (83)
$`=`$ $`{\displaystyle \frac{i}{8}}{\displaystyle \frac{1}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}\lambda _{ij2}^{}\lambda _{aj3}^{}\mathrm{\Gamma }_{Li\mathrm{}}^e\mathrm{\Gamma }_{L\mathrm{}a}^e(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{u_L}_{j\alpha }\gamma ^\mu u_{Lj\beta })`$ (85)
In the first step we used the fact that the two squarks have to be of the same type and the unitarity of the CKM-matrix. In the second step we performed a Fierz rearrangement. This is done to get the same structure (i.e., two vectors) for the four-fermion operator as in the standard model case. The advantages of this rearrangement will become clear when calculating the $`\gamma `$-matrix. In the last line we put the colour indices $`\alpha `$, $`\beta `$ for clarity.
As one can see clearly, the effect of the (unitary) squark mixing matrix $`\mathrm{\Gamma }_L^e`$ becomes enhanced if the masses of the selectrons are very different for the three generations. If there were a mass degeneracy they would simply give a factor $`\delta _{ia}`$. This is a general feature in our calculations.
The operator that appears in Eq. (85) is of a new type. It consists of a right-handed $`b`$\- and a right-handed $`s`$-quark. (Actually, these are two operators, one with a pair of $`u`$-quarks and the other with two $`c`$-quarks.)
A careful investigation results in the following set of new operators:
From $`_\lambda ^{}`$ one gets
$`P_1`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{u_L}_\beta \gamma _\mu u_{L\alpha })`$ (86)
$`P_2`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{c_L}_\beta \gamma _\mu c_{L\alpha })`$ (88)
$`P_3`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{d_L}_\beta \gamma _\mu d_{L\alpha })`$ (90)
$`P_4`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{s_L}_\beta \gamma _\mu s_{L\alpha })`$ (92)
$`P_5`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{b_L}_\beta \gamma _\mu b_{L\alpha })`$ (94)
$`P_6`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\beta })(\overline{d_R}_\beta \gamma _\mu d_{R\alpha })`$ (96)
$`P_7`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\beta })(\overline{s_R}_\beta \gamma _\mu s_{R\alpha })`$ (98)
$`P_8`$ $`=`$ $`(\overline{s_L}_\alpha \gamma ^\mu b_{L\beta })(\overline{b_R}_\beta \gamma _\mu b_{R\alpha })`$ (100)
$`P_9`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\alpha }){\displaystyle \underset{i}{}}(\overline{q_R}_{i\beta }\gamma ^\mu q_{Ri\beta })`$ (102)
$`P_{10}`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta }){\displaystyle \underset{i}{}}(\overline{q_R}_{i\beta }\gamma ^\mu q_{Ri\alpha })`$ (104)
$`P_{11}`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\alpha }){\displaystyle \underset{i}{}}(\overline{q_L}_{i\beta }\gamma ^\mu q_{Li\beta })`$ (106)
$`P_{12}`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta }){\displaystyle \underset{i}{}}(\overline{q_L}_{i\beta }\gamma ^\mu q_{Li\alpha }).`$ (108)
It is worth noting that the operators $`P_1`$ \- $`P_8`$ emerge directly from the Lagrangian and $`P_9`$ \- $`P_{12}`$ are induced through QCD corrections. One would expect partners of $`P_1`$ \- $`P_8`$ with colour structure $`(\alpha \alpha )(\beta \beta )`$ to be introduced by QCD. This does not happen at leading-log (accidently).
$`_{\lambda ^{\prime \prime }}`$ leads to the following additional operators:
$`R_1`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\alpha })(\overline{u_R}_\beta \gamma _\mu u_{R\beta })`$ (109)
$`R_2`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{u_R}_\beta \gamma _\mu u_{R\alpha })`$ (111)
$`R_3`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\alpha })(\overline{c_R}_\beta \gamma _\mu c_{R\beta })`$ (113)
$`R_4`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{c_R}_\beta \gamma _\mu c_{R\alpha })`$ (115)
$`R_5`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\alpha })(\overline{d_R}_\beta \gamma _\mu d_{R\beta })`$ (117)
$`R_6`$ $`=`$ $`(\overline{s_R}_\alpha \gamma ^\mu b_{R\beta })(\overline{d_R}_\beta \gamma _\mu d_{R\alpha }).`$ (119)
Here, all the operators $`R_1`$ \- $`R_6`$ appear already in the tree-level effective Lagrangian.
In principle, the effective Hamiltonian also contains semi-leptonic Operators. However, these can be neglected because, as a consequence of their Dirac structure, they do not contribute to the decay $`bs\gamma `$.
#### 2 The $`\gamma `$-matrix
The whole basis now consists of 28 operators. Their QCD-mixing is described by a $`28\times 28`$-$`\gamma `$-matrix. It is depicted in appendix C. There are three different blocks in this matrix that have to be treated in separate ways.
* *mixing of four-fermion operators among themselves*
This block involves the one-loop diagrams of Fig. 2. There are no further complications. The mixing of $`O_1`$ \- $`O_8`$ is known already since a long time .
* *mixing of $`O_7`$, $`O_8`$ and $`\stackrel{~}{O}_7`$, $`\stackrel{~}{O}_8`$ among themselves*
These entries need not to be computed. The mixing of $`O_7`$ and $`O_8`$ is known and the new operators mix in exactly the same way such that the corresponding numbers can be copied.
* *mixing of the four-fermion operators with $`O_7`$, $`O_8`$, $`\stackrel{~}{O}_7`$, $`\stackrel{~}{O}_8`$*
This task is more difficult. In principle, one has to compute the divergent part of all the diagrams of Fig. 4 and 5 together with their counterterms and the contributions of the evanescent operators . Moreover, there are four different types of chiralities to be inserted: $`(LL)(LL)`$, $`(LL)(RR)`$, $`(RR)(LL)`$ and $`(RR)(RR)`$. For the first two this calculation had to be performed for the case of the standard model. The detailed results are listed in . The second two types of insertions are new. Fortunately, one can deduce the divergent parts of these types of insertions by making the following two observations:
+ The diagrams of Fig. 4 contain only one fermion line. Here, it is crucial if the two quark pairs of the inserted operator have the same or opposite chirality. Thus, an insertion of a $`(RR)(RR)`$ leads to the same divergence as an insertion of a $`(LL)(LL)`$ as well as the divergence of an insertion of a $`(RR)(LL)`$ is the same as for the case of $`(LL)(RR)`$.
+ One must be careful with the diagrams that contain a closed fermion loop (Fig. 5). At a first sight, one may deduce that the divergence does depend on the chirality of the quarks running in the loop but *not* on the chirality of the $`s`$\- and the $`b`$-quark. However, this is wrong.We thank the authors of for clarifying this point to us. The difference between a right- and a left-handed quark in the loop results in a term $`\mathrm{Tr}(\pm \gamma ^5\mathrm{\Gamma })`$ where $`\mathrm{\Gamma }`$ stands for a collection of at least four Dirac $`\gamma `$-matrices and the sign corresponds to a right- or left-handed quark in the loop, respectively. The trace leads to an $`\epsilon `$-tensor that is contracted with a $`\gamma `$-matrix between the external quarks. This produces an additional $`\gamma ^5`$. Now note that $`\gamma ^5P_R=+P_R`$, whereas $`\gamma ^5P_L=P_L`$. This means, we need to change the chirality in *both* fermion pairs to end up with the same mixing. Hence, if an insertion of an operator of type $`(LL)(LL)`$ \[$`(LL)(RR)`$\] gives a certain contribution to $`O_{7/8}`$ the $`(RR)(RR)`$ \[$`(RR)(LL)`$\] operator will give exactly the same contribution to $`\stackrel{~}{O}_{7/8}`$.
To summarize, the mixing of the four-fermion operators with $`\stackrel{~}{O}_{7/8}`$ can be deduced completely from the results of by interchanging left-handed and right-handed projectors.
All the previous calculations involve $`\gamma ^5`$. It is therefore clear that the results depend on the regularization scheme. This dependence will be cancelled by finite one-loop (but $`O(\alpha _s^0)`$) contributions of some four-fermion operators $`Q_i`$ to the Amplitude $`A`$ of $`bs\gamma `$ through the diagram of Fig. 3. Schematically, the result is then
$$A=C_7s\gamma |O_7|b_{\mathrm{tree}}+\stackrel{~}{C}_7s\gamma |\stackrel{~}{O}_7|b_{\mathrm{tree}}+\underset{i}{}C_{Q_i}s\gamma |Q_i|b_{1\mathrm{loop}}.$$
(120)
An alternative is to define effective coefficients in a way that $`A`$ becomes
$$A=C_7^{\mathrm{eff}}s\gamma |O_7|b_{\mathrm{tree}}+\stackrel{~}{C}_7^{\mathrm{eff}}s\gamma |\stackrel{~}{O}_7|b_{\mathrm{tree}}.$$
(121)
This can be achieved by defining four vectors $`\{y_i\}`$, $`\{z_i\}`$, $`\{\stackrel{~}{y}_i\}`$, and $`\{\stackrel{~}{z}_i\}`$ through
$`s\gamma |Q_i|b_{1\mathrm{loop}}`$ $`=:`$ $`y_is\gamma |O_7|b_{\mathrm{tree}}`$ (122)
$`s\mathrm{gluon}|Q_i|b_{1\mathrm{loop}}`$ $`=:`$ $`z_is\mathrm{gluon}|O_8|b_{\mathrm{tree}}`$ (124)
$`s\gamma |Q_i|b_{1\mathrm{loop}}`$ $`=:`$ $`\stackrel{~}{y}_is\gamma |\stackrel{~}{O}_7|b_{\mathrm{tree}}`$ (126)
$`s\mathrm{gluon}|Q_i|b_{1\mathrm{loop}}`$ $`=:`$ $`\stackrel{~}{z}_is\mathrm{gluon}|\stackrel{~}{O}_8|b_{\mathrm{tree}}.`$ (128)
The effective Wilson-coefficients must be defined as
$`C_7^{\mathrm{eff}}(\mu )`$ $`:=`$ $`C_7(\mu )+{\displaystyle \underset{i}{}}y_iC_i(\mu )`$ (129)
$`C_8^{\mathrm{eff}}(\mu )`$ $`:=`$ $`C_8(\mu )+{\displaystyle \underset{i}{}}z_iC_i(\mu )`$ (131)
$`\stackrel{~}{C}_7^{\mathrm{eff}}(\mu )`$ $`:=`$ $`\stackrel{~}{C}_7(\mu )+{\displaystyle \underset{i}{}}\stackrel{~}{y}_iC_i(\mu )`$ (133)
$`\stackrel{~}{C}_8^{\mathrm{eff}}(\mu )`$ $`:=`$ $`\stackrel{~}{C}_8(\mu )+{\displaystyle \underset{i}{}}\stackrel{~}{z}_iC_i(\mu )`$ (135)
$`C_{Q_i}^{\mathrm{eff}}(\mu )`$ $`:=`$ $`C_{Q_i}(\mu ).`$ (137)
The vector
$$\stackrel{}{C}^{\mathrm{eff}}(\mu ):=\{C_{Q_i}^{\mathrm{eff}}(\mu ),C_7^{\mathrm{eff}}(\mu ),C_8^{\mathrm{eff}}(\mu ),\stackrel{~}{C}_7^{\mathrm{eff}}(\mu ),\stackrel{~}{C}_8^{\mathrm{eff}}(\mu )\}$$
(138)
is then regularization scheme independent. Remember that the index $`i`$ runs over all four fermion operators.
The effective coefficients obey RGE’s which can be derived from the RGE’s for $`C_k(\mu )`$, where $`k`$ labels the whole set of operators. They are
$$\frac{d}{\mathrm{ln}\mu }C_k^{\mathrm{eff}}(\mu )=\frac{\alpha _s}{4\pi }\gamma _{jk}^{\mathrm{eff}}C_j^{\mathrm{eff}}(\mu ),$$
(139)
where
$$\gamma _{jk}^{\mathrm{eff}}=\{\begin{array}{cc}\gamma _{jk}+_{i=1}^{24}y_i\gamma _{ji}y_j\gamma _{O_7O_7}z_j\gamma _{O_8O_7}\hfill & k=O_7,j=1,\mathrm{},24\hfill \\ \gamma _{jk}+_{i=1}^{24}z_i\gamma _{ji}z_j\gamma _{O_8O_8}\hfill & k=O_8,j=1,\mathrm{},24\hfill \\ \gamma _{jk}+_{i=1}^{24}\stackrel{~}{y}_i\gamma _{ji}\stackrel{~}{y}_j\gamma _{\stackrel{~}{O}_7\stackrel{~}{O}_7}\stackrel{~}{z}_j\gamma _{\stackrel{~}{O}_8\stackrel{~}{O}_7}\hfill & k=\stackrel{~}{O}_7,j=1,\mathrm{},24\hfill \\ \gamma _{jk}+_{i=1}^{24}\stackrel{~}{z}_i\gamma _{ji}\stackrel{~}{z}_j\gamma _{\stackrel{~}{O}_8\stackrel{~}{O}_8}\hfill & k=\stackrel{~}{O}_8,j=1,\mathrm{},24\hfill \\ \gamma _{jk}\hfill & \text{otherwise}\hfill \end{array}$$
(140)
For a finite contribution of an operator inserted in the diagram of Fig. 3 we need
* two pairs of fermions with different chirality, i.e., operators of the form $`(LL)(RR)`$ or $`(RR)(LL)`$
* a $`b`$-quark running in the loop.
This reduces the possibilities to $`O_5`$, $`O_6`$, $`P_5`$, $`P_8`$, $`P_{11}`$ and $`P_{12}`$. The results for $`\{y_i\}`$, $`\{z_i\}`$, $`\{\stackrel{~}{y}_i\}`$ and $`\{\stackrel{~}{z}_i\}`$ are then
$$\begin{array}{cccccc}y_i\hfill & =& \{\begin{array}{cc}\frac{2}{3}& i=O_5\hfill \\ 2& i=O_6,P_8\hfill \\ 0& \text{otherwise}\hfill \end{array}\hfill & z_i\hfill & =& \{\begin{array}{cc}2& i=O_5\hfill \\ 0& \text{otherwise}\hfill \end{array}\hfill \\ & & & & & \\ \stackrel{~}{y}_i\hfill & =& \{\begin{array}{cc}\frac{2}{3}& i=P_{11}\hfill \\ 2& i=P_5,P_{12}\hfill \\ 0& \text{otherwise}\hfill \end{array}\hfill & \stackrel{~}{z}_i\hfill & =& \{\begin{array}{cc}2& i=P_{11}\hfill \\ 0& \text{otherwise}\hfill \end{array}\hfill \end{array}$$
(141)
There is one more subtlety concerning the matching at $`M_W`$: In the standard model we have $`s\gamma |H^{\mathrm{eff}}|b=C_7(M_W)s\gamma |O_7|b`$. Now, more operators contribute:
$`s\gamma |H^{\mathrm{eff}}|b`$ $`=`$ $`C_7(M_W)s\gamma |O_7|b+\stackrel{~}{C}_7(M_W)s\gamma |\stackrel{~}{O}_7|b`$ (143)
$`2C_{P_8}(M_W)s\gamma |P_8|b2C_{P_5}(M_W)s\gamma |P_5|b`$
$`=`$ $`C_7^{\mathrm{eff}}(M_W)s\gamma |O_7|b+\stackrel{~}{C}_7^{\mathrm{eff}}(M_W)s\gamma |\stackrel{~}{O}_7|b.`$ (144)
Hence, the matching procedure does not lead to $`C_7(M_W)`$ and $`\stackrel{~}{C}_7(M_W)`$ but directly to $`C_7^{\mathrm{eff}}(M_W)`$ and $`\stackrel{~}{C}_7^{\mathrm{eff}}(M_W)`$.
#### 3 The matching conditions
The matching for the additional four-fermion operators must only be performed at tree-level. The matching conditions are therefore easily derived. An example is given in (85). The complete set is
$`C_{P_1}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i12}^{}\lambda _{j13}^{}\mathrm{\Gamma }_{Li\mathrm{}}^e\mathrm{\Gamma }_{L\mathrm{}j}^e{\displaystyle \frac{1}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}`$ (145)
$`C_{P_2}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i22}^{}\lambda _{j23}^{}\mathrm{\Gamma }_{Li\mathrm{}}^e\mathrm{\Gamma }_{L\mathrm{}j}^e{\displaystyle \frac{1}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}+{\displaystyle \frac{G_F}{\sqrt{2}}}K_{ts}^{}K_{tb}{\displaystyle \frac{m_sm_b}{2m_{H^+}^2}}\mathrm{tan}^2\beta `$ (147)
$`C_{P_3}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i12}^{}\lambda _{j13}^{}\mathrm{\Gamma }_{Li\mathrm{}}^\nu \mathrm{\Gamma }_{L\mathrm{}j}^\nu {\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}`$ (149)
$`C_{P_4}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i22}^{}\lambda _{j23}^{}\mathrm{\Gamma }_{Li\mathrm{}}^\nu \mathrm{\Gamma }_{L\mathrm{}j}^\nu {\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}`$ (151)
$`C_{P_5}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i32}^{}\lambda _{j33}^{}\mathrm{\Gamma }_{Li\mathrm{}}^\nu \mathrm{\Gamma }_{L\mathrm{}j}^\nu {\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}`$ (153)
$`C_{P_6}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i31}^{}\lambda _{j21}^{}\mathrm{\Gamma }_{Li\mathrm{}}^\nu \mathrm{\Gamma }_{L\mathrm{}j}^\nu {\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}`$ (155)
$`C_{P_7}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i32}^{}\lambda _{j22}^{}\mathrm{\Gamma }_{Li\mathrm{}}^\nu \mathrm{\Gamma }_{L\mathrm{}j}^\nu {\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}`$ (157)
$`C_{P_8}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i33}^{}\lambda _{j23}^{}\mathrm{\Gamma }_{Li\mathrm{}}^\nu \mathrm{\Gamma }_{L\mathrm{}j}^\nu {\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}`$ (159)
$`C_{P_9}^{\mathrm{eff}}(M_W)`$ $`=`$ $`C_{P_{10}}^{\mathrm{eff}}(M_W)=C_{P_{11}}^{\mathrm{eff}}(M_W)=C_{P_{12}}^{\mathrm{eff}}(M_W)=0`$ (161)
$`C_{R_1}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{1i2}^{\prime \prime }\lambda _{1j3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}`$ (164)
$`C_{R_2}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{1i2}^{\prime \prime }\lambda _{1j3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}=C_{R_1}^{\mathrm{eff}}(M_W)`$ (166)
$`C_{R_3}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{2i2}^{\prime \prime }\lambda _{2j3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}`$ (168)
$`C_{R_4}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{2i2}^{\prime \prime }\lambda _{2j3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}=C_{R_3}^{\mathrm{eff}}(M_W)`$ (170)
$`C_{R_5}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i12}^{\prime \prime }\lambda _{j13}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^u\mathrm{\Gamma }_{R\mathrm{}i}^u{\displaystyle \frac{1}{m_{\stackrel{~}{u}_{\mathrm{}}}^2}}`$ (172)
$`C_{R_6}^{\mathrm{eff}}(M_W)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\lambda _{i12}^{\prime \prime }\lambda _{j13}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^u\mathrm{\Gamma }_{R\mathrm{}i}^u{\displaystyle \frac{1}{m_{\stackrel{~}{u}_{\mathrm{}}}^2}}=C_{R_5}^{\mathrm{eff}}(M_W).`$ (174)
In $`C_{P_2}^{\mathrm{eff}}(M_W)`$ we included a term coming from Higgs exchange because it can possibly be large in the case of a large $`\mathrm{tan}\beta `$.
There are some more terms to add to $`C_7`$, $`C_8`$, $`\stackrel{~}{C}_7`$, and $`\stackrel{~}{C}_8`$, too
$`C_7^{\mathrm{eff}}(M_W)`$ $`=`$ $`C_{7\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}(M_W)`$ (177)
$`{\displaystyle \frac{1}{4}}{\displaystyle \frac{Q_d}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}\lambda _{i3k}^{}\lambda _{j2k}^{}\mathrm{\Gamma }_i\mathrm{}^\nu \mathrm{\Gamma }_\mathrm{}j^\nu F_1(x_{d_k\stackrel{~}{\nu }_{\mathrm{}}})`$
$`+{\displaystyle \frac{1}{48}}{\displaystyle \frac{Q_d}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\lambda _{k3i}^{}\lambda _{k2j}^{}\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d`$
$`C_8^{\mathrm{eff}}(M_W)`$ $`=`$ $`C_{8\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}(M_W)`$ (180)
$`{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}\lambda _{i3k}^{}\lambda _{j2k}^{}\mathrm{\Gamma }_i\mathrm{}^\nu \mathrm{\Gamma }_\mathrm{}j^\nu F_1(x_{d_k\stackrel{~}{\nu }_{\mathrm{}}})`$
$`+{\displaystyle \frac{1}{48}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\lambda _{k3i}^{}\lambda _{k2j}^{}\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d`$
$`\stackrel{~}{C}_7^{\mathrm{eff}}(M_W)`$ $`=`$ $`\stackrel{~}{C}_{7\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}(M_W)`$ (187)
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}\lambda _{ia2}^{}\lambda _{jb3}^{}K_{an}^{}K_{nb}\mathrm{\Gamma }_{Li\mathrm{}}^e\mathrm{\Gamma }_{L\mathrm{}j}^e\left[Q_uF_1(x_{u_n\stackrel{~}{e}_{\mathrm{}}})+Q_eF_2(x_{u_n\stackrel{~}{e}_{\mathrm{}}})\right]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_{\mathrm{}}}^2}}\lambda _{ki2}^{}\lambda _{kj3}^{}\mathrm{\Gamma }_{Li\mathrm{}}^u\mathrm{\Gamma }_{L\mathrm{}j}^u\left[Q_eF_1(x_{e_k\stackrel{~}{u}_{\mathrm{}}})+Q_uF_2(x_{e_k\stackrel{~}{u}_{\mathrm{}}})\right]`$
$`{\displaystyle \frac{Q_d}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}\lambda _{ik2}^{}\lambda _{jk3}^{}\mathrm{\Gamma }_i\mathrm{}^\nu \mathrm{\Gamma }_\mathrm{}j^\nu F_1(x_{d_k\stackrel{~}{\nu }_{\mathrm{}}})`$
$`+{\displaystyle \frac{Q_d}{48}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\lambda _{ki2}^{}\lambda _{kj3}^{}\mathrm{\Gamma }_{Li\mathrm{}}^d\mathrm{\Gamma }_{L\mathrm{}j}^d`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_{\mathrm{}}}^2}}\lambda _{ik2}^{\prime \prime }\lambda _{jk3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^u\mathrm{\Gamma }_{R\mathrm{}i}^u\left[Q_dF_1(x_{d_k\stackrel{~}{u}_{\mathrm{}}})Q_uF_2(x_{d_k\stackrel{~}{u}_{\mathrm{}}})\right]`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\lambda _{ki2}^{\prime \prime }\lambda _{kj3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d\left[Q_uF_1(x_{u_k\stackrel{~}{d}_{\mathrm{}}})Q_dF_2(x_{u_k\stackrel{~}{d}_{\mathrm{}}})\right]`$
$`\stackrel{~}{C}_8^{\mathrm{eff}}(M_W)`$ $`=`$ $`\stackrel{~}{C}_{8\mathrm{M}\mathrm{S}\mathrm{S}\mathrm{M}}(M_W)`$ (195)
$`{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}\lambda _{ia2}^{}\lambda _{jb3}^{}K_{an}^{}K_{nb}\mathrm{\Gamma }_{Li\mathrm{}}^e\mathrm{\Gamma }_{L\mathrm{}j}^eF_1(x_{u_n\stackrel{~}{e}_{\mathrm{}}})`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_{\mathrm{}}}^2}}\lambda _{ki2}^{}\lambda _{kj3}^{}\mathrm{\Gamma }_{Li\mathrm{}}^u\mathrm{\Gamma }_{L\mathrm{}j}^uF_2(x_{e_k\stackrel{~}{u}_{\mathrm{}}})`$
$`{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{\nu }_{\mathrm{}}}^2}}\lambda _{ik2}^{}\lambda _{jk3}^{}\mathrm{\Gamma }_i\mathrm{}^\nu \mathrm{\Gamma }_\mathrm{}j^\nu F_1(x_{d_k\stackrel{~}{\nu }_{\mathrm{}}})`$
$`+{\displaystyle \frac{1}{48}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\lambda _{ki2}^{}\lambda _{kj3}^{}\mathrm{\Gamma }_{Li\mathrm{}}^d\mathrm{\Gamma }_{L\mathrm{}j}^d`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{u}_{\mathrm{}}}^2}}\lambda _{ik2}^{\prime \prime }\lambda _{jk3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^u\mathrm{\Gamma }_{R\mathrm{}i}^u\left[F_1(x_{d_k\stackrel{~}{u}_{\mathrm{}}})+F_2(x_{d_k\stackrel{~}{u}_{\mathrm{}}})\right]`$
$`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\lambda _{ki2}^{\prime \prime }\lambda _{kj3}^{\prime \prime }\mathrm{\Gamma }_{Rj\mathrm{}}^d\mathrm{\Gamma }_{R\mathrm{}i}^d\left[F_1(x_{u_k\stackrel{~}{d}_{\mathrm{}}})+F_2(x_{u_k\stackrel{~}{d}_{\mathrm{}}})\right]`$
$`2C_{P_5}(M_W).`$
#### 4 $`𝐛𝐮|𝐜𝐞\overline{\nu }`$
It is convenient to express the branching ratio $`\text{Br}(bs\gamma )`$ through the semi-leptonic decay $`bu|ce\overline{\nu }`$ :
$$\text{Br}(bs\gamma )=\frac{\mathrm{\Gamma }(bs\gamma )}{\mathrm{\Gamma }(bu|ce\overline{\nu })}\text{Br}_{\mathrm{exp}}(bu|ce\overline{\nu }),$$
(196)
where we take $`\text{Br}_{\mathrm{exp}}(bu|ce\overline{\nu })=10.5\%`$ . This has the advantage that the large bottom mass dependence ($`m_b^5`$) cancels out. In the SM the semi-leptonic decay is mediated by a $`W`$-boson \[Fig. 8 (a)\] whereas in the case of the MSSM the charged Higgs can take the role of the $`W`$. However, the coupling to the leptons is proportional to the electron mass and hence it can be safely neglected. Introducing the R-parity breaking terms (72) -(80) offers new decay channels depicted in Figs. 8 (b) and (c) which have to be included in the decay width. Please note the following few things:
* In the MSSM the decay $`bue\overline{\nu }`$ is suppressed by the small CKM- matrix element $`K_{ub}`$ and can therefore be neglected. In our case we have to include this decay mode.
* The absence of lepton generation mixing in the SM forces the anti-neutrino to be $`\overline{\nu }_e`$. This restriction is no longer valid in our case. Therefore, we have to sum over the three generations before squaring the amplitude because the generation of the neutrino is not detected.
* Setting the mass of the lepton to zero (which is certainly a valid approximation) the computation in the MSSM does not distinguish between the electron and the muon. Here, the two particles involve different couplings. We give the results for an outgoing electron. For the muon just change the “1” in the relevant coupling to a “2”.
The results are then, to leading order and with $`m_u^2/m_b^2=0`$
$`\mathrm{\Gamma }(bs\gamma )`$ $`=`$ $`{\displaystyle \frac{e^2m_b^5}{192\pi ^5}}{\displaystyle \frac{3}{4}}(|C_7|^2+|\stackrel{~}{C}_7|^2)`$ (197)
$`\mathrm{\Gamma }(bu|ce\overline{\nu })`$ $`=`$ $`{\displaystyle \frac{m_b^5}{192\pi ^3}}{\displaystyle \frac{1}{32}}\{(18ϵ^2+8ϵ^6ϵ^824ϵ^4\mathrm{log}ϵ)\times `$ (199)
$`[|2A+C_2|^2+|B_2|^2]+|B_1|^2+|C_1|^2\},`$
where
$`A`$ $`=`$ $`2\sqrt{2}G_FK_{23}`$ (200)
$`B_r`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{ij1}\lambda _{mn3}^{}}{m_{\stackrel{~}{e}_{\mathrm{}}}^2}}\mathrm{\Gamma }_{L\mathrm{}m}^e\mathrm{\Gamma }_{Li\mathrm{}}^eK_{rn}r=1,2`$ (201)
$`C_r`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{i3k}^{}\lambda _{1mn}^{}}{m_{\stackrel{~}{d}_{\mathrm{}}}^2}}\mathrm{\Gamma }_{R\mathrm{}k}^d\mathrm{\Gamma }_{Rn\mathrm{}}^dK_{rm}`$ (202)
$`ϵ`$ $`=`$ $`{\displaystyle \frac{m_c}{m_b}}.`$ (203)
## IV Results
The formulas of the previous section are far too complicated to be treated “by hand” but it is no problem to feed them to a computer. The $`\gamma `$-matrix is independent of the parameters of supersymmetry. With the help of *Mathematica* it is possible to diagonalize it and find the influence of the QCD effects.
### A General results
Because of the enlarged operator basis the expression for $`C_7`$ changes. In general, the solution of the RGE for the Wilson-coefficients is given by
$$\stackrel{}{C}^{\mathrm{eff}}(\mu )=V\left[\left(\frac{\alpha _s(M_W)}{\alpha _s(\mu )}\right)^{\frac{\stackrel{}{\gamma }_D}{2\beta _0}}\right]_DV^1\stackrel{}{C}^{\mathrm{eff}}(M_W),$$
(204)
where $`V`$ diagonalizes $`\gamma ^T`$
$$\gamma _D=V^1\gamma ^TV.$$
(205)
$`\beta _0=23/3`$ is the one-loop beta-function and $`\stackrel{}{\gamma }_D`$ is the vector containing the eigenvalues of $`\gamma `$. In our case
$$\stackrel{}{\gamma }_D=\begin{array}{ccccccc}(16& 16& 16& 16& 16& 16& 16\\ 16& 8& 8& 8& 8& 4& 4\\ 4& 4& \frac{28}{3}& \frac{28}{3}& \frac{32}{3}& \frac{32}{3}& 2.233\\ 2.233& 6.266& 6.266& 13.791& 13.791& 6.486& 6.486).\end{array}$$
(206)
To have an idea which coefficients are relevant we perform a numerical ana-lysis with $`\alpha _s(M_Z)=0.121`$ and $`\mu =m_b=4.2`$ GeV. The coefficients $`C_7^{\mathrm{eff}}(m_b)`$ and $`\stackrel{~}{C}_7^{\mathrm{eff}}(m_b)`$ are then
$`C_7^{\mathrm{eff}}(m_b)`$ $`=`$ $`0.351C_{O_2}^{\mathrm{eff}}(M_W)0.198C_{P_6}^{\mathrm{eff}}(M_W)0.198C_{P_7}^{\mathrm{eff}}(M_W)`$ (208)
$`0.178C_{P_8}^{\mathrm{eff}}(M_W)+0.665C_7^{\mathrm{eff}}(M_W)+0.093C_8^{\mathrm{eff}}(M_W)`$
$`\stackrel{~}{C}_7^{\mathrm{eff}}(m_b)`$ $`=`$ $`0.510C_{P_1}^{\mathrm{eff}}(M_W)+0.510C_{P_2}^{\mathrm{eff}}(M_W)0.198C_{P_3}^{\mathrm{eff}}(M_W)`$ (212)
$`0.198C_{P_4}^{\mathrm{eff}}(M_W)0.178C_{P_5}^{\mathrm{eff}}(M_W)+0.381C_{R_1}^{\mathrm{eff}}(M_W)`$
$`+0.381C_{R_3}^{\mathrm{eff}}(M_W)0.213C_{R_5}^{\mathrm{eff}}(M_W)+0.665\stackrel{~}{C}_7^{\mathrm{eff}}(M_W)`$
$`+0.093\stackrel{~}{C}_8^{\mathrm{eff}}(M_W),`$
It is clear that the four-fermion operators including a left-handed $`s`$-quark contribute to $`C_{7/8}`$ whereas the ones with a right-handed $`s`$-quark contribute to $`\stackrel{~}{C}_{7/8}`$. The numbers multiplying the different Wilson coefficients are all of the same size, hence there is a priory no term which can be neglected.
### B Specific results for R-parity violation
It is obvious that the Wilson coefficients depend in a very complicated way on the parameters of our supersymmetric model. Changes of $`\lambda `$, $`\lambda ^{}`$, or $`\lambda ^{\prime \prime }`$ not only affect the result in a direct way but also in an indirect fashion through an altered mass spectrum and different mixing matrices. Therefore it is very hard to make general statements on the behaviour of the branching ratio.
As mentioned, it is not our aim to perform a high precision analysis of the parameter space but to explore the influence of the R-parity violating couplings on $`bs\gamma `$ with a reasonable accuracy. We know that our results for the branching ratio are only valid at the 25% level because we do a leading-log approximation with large scale uncertainty . However, we expect the impact of the R-parity breaking terms not to change much when including the next-to-leading corrections. This means that the shape of the curves remain more or less the same whereas the offset where our curves start (i.e., no R-parity breaking) may change significantly when calculating the next-to-leading-log approximation. Solving the RGE’s was performed in the following way:
* Solve the equations for the gauge couplings. The boundary conditions are the physical values at $`M_Z`$. The gauge couplings will meet at $`M_X210^{16}`$ GeV. Choose a common value $`M(M_X)`$ for the gaugino masses at that high scale and solve the RGE’s for $`M_1`$-$`M_3`$.
* Set $`\mathrm{tan}\beta `$ and $`\lambda (M_Z)`$, $`\lambda ^{}(M_Z)`$ and $`\lambda ^{\prime \prime }(M_Z)`$ to the desired value and use them together with the quark and lepton masses as inputs at $`M_Z`$ for the Yukawa couplings. Let these couplings run to $`M_X`$.
* Choose $`A`$, $`\mu _{2i}(M_X)`$ and trial values for $`m_0`$ and $`\mu (M_X)`$ to complete the boundary conditions at $`M_X`$. The whole set of RGE’s is then run down to $`M_Z`$. We choose $`\mu _{2i}(M_X)=0`$. This is an approximation because $`\mu _{2i}`$ and $`\stackrel{~}{\mu }_{2i}`$ will not vanish at $`M_Z`$. However, in our examples their values are so small compared to $`\mu `$ and $`\stackrel{~}{\mu }`$ that we can neglect them avoiding a mixing between $`H_1`$ and $`L_i`$.
* The minimum of the one-loop effective Higgs potential $`V_0+\mathrm{\Delta }V`$ will in general not be at $`(v_1,v_2)`$ which makes it necessary to adjust $`\mu (M_X)`$ and $`m_0`$ in a clever way. We chose the Newton method to converge to the desired position of the minimum. Following we evaluate the minimum of the potential at some average mass scale to avoid large logarithms and therefore get a more reliable result.
The next step is the numerical diagonalization of the mass matrices to find the masses of the physical particles and the relevant mixing matrices.
Changing the value of the R-parity violating couplings makes it necessary to continuously adjust the values of $`m_0`$ and $`\mu (M_X)`$ which alters the mass spectrum of the particles at $`M_Z`$. All the following examples correspond to mass spectra within the current bounds for the masses of the supersymmetric particles. We encountered two critical situations, namely too small masses for the lightest selectron and/or the lightest Higgs boson. An idea would then be to constrain the bounds on the R-parity violating couplings through the requirement of a phenomenologically realistic mass spectrum independent of the value for the branching ratio Br$`(bs\gamma )`$. However, in general, a realistic mass spectrum does not restrict the R-parity breaking parameters substantially. Moreover, the mass spectrum highly depends on the values of $`\mathrm{tan}\beta `$, $`A`$ and $`M(M_X)`$. Hence, such bounds would be strongly model dependent.
As reference model we chose $`\mathrm{tan}\beta =5`$, $`A=0`$ and $`M(M_X)=300`$ GeV. With vanishing R-parity violating couplings this leads to squark masses of $`600800`$ GeV, sleptons of $`260330`$ GeV, a lightest neutralino of 120 GeV and a lightest Higgs of about 100 GeV. Fig. 9 shows the behaviour of Br$`(bs\gamma )`$ in the neighbourhood of our reference model with, in addition, $`\lambda _{132}^{}=\lambda _{122}^{}=0.1`$. Interestingly, the value for the branching ratio is rather stable under a change of $`\mathrm{tan}\beta `$ and $`A`$.
As mentioned before, it is very difficult to isolate generic features of the different models. Let us make some comments:
* At least two of the $`\lambda `$’s must be non-zero to have an influence on the result. There are two exceptions: $`\lambda _{123}^{\prime \prime }`$ and $`\lambda _{223}^{\prime \prime }`$ alone will give a contribution due to the anti-symmetry of $`\lambda ^{\prime \prime }`$. However, their impact on the branching ratio is so small that no reasonable bounds can be found.
* As all the effective couplings depend on the inverse mass-squared of a heavy particle, it is clear that the influence of the new physics is bigger in models with a lower mass spectrum. To see the effect of smaller masses compare Figs. 10 (a) and (b) where for the latter $`M(M_X)=100`$ GeV is taken which results in, for instance, squark masses of about 250 GeV. Note, that our reference model leads to quite high masses.
* The squarks are always heavier than the sleptons. As a consequence, the influence of non-vanishing $`\lambda ^{\prime \prime }`$ on Br$`(bs\gamma )`$ is much smaller than for non-vanishing $`\lambda ^{}`$.
* The simplest non trivial situation consists of a single non-vanishing pair of R-parity violating couplings. We encounter the following scenarios:
+ The largest effect results from the pair $`(\lambda _{122}^{},\lambda _{132}^{})`$ (see Fig. 11). It is not reasonable to extract a bound for the product of these parameters because (through the RGE’s) the dependence on these couplings is much more complex. However, we draw the conclusion that significant effects result if both couplings exceed 0.1-0.2. Fig. 11 also shows that the branching ratio may strongly depend on the relative sign of the couplings. This happens if the new contributions are able to diminish the value of $`C_7`$.
+ A general feature of our model is that the influence of a pair of non-vanishing R-parity violating couplings on the branching ratio starts being significant if $`(\lambda ^{},\lambda ^{})0.20.4`$. For $`(\lambda ^{\prime \prime },\lambda ^{\prime \prime })`$, in most cases, the requirement of non-diverging RGE’s for these couplings gives more stringent bounds. There is one more aspect: Some of the contributing pairs have more than one distinct index. They appear only because of the mass differences of the different squark or slepton generations. This can be seen clearly in formulas (174): If the masses were independent of their index the $`\mathrm{\Gamma }`$’s would combine to an identity matrix leaving only those pairs with an identical sfermion index. The bounds on these pairs are much less stringent. This is exactly what happens in the case of $`\lambda _{123}^{\prime \prime }`$ and $`\lambda _{223}^{\prime \prime }`$ when being the only non-vanishing coupling. This is another reason why the results for this situation are not stringent.
+ A special situation is depicted in Fig. 12. The branching ratio seems to “explode” at $`\lambda _{312}^{}=\lambda _{313}^{}0.3`$. This is due to the mentioned decreasing mass squared of the lightest selectron, which appears in some denominators in the matching conditions. Hence, the specific position of the peak is highly model dependent.
* If several pairs of couplings are non-vanishing the picture can get more complex, as Fig. 13 shows. There, we combined the decreasing effect with the peak due to the small selectron mass.
* It is difficult to compare our results with existing constraints for the R-parity breaking couplings because the results are extremely model dependent. However, in the case of $`\lambda ^{}`$ our bounds are highly competitive. For comparison we put the current bounds in the captions of the respective figures.
## V Conclusions
To examine the influence of R-parity breaking on $`bs\gamma `$ one has to enlarge the operator basis substantially. At the leading-log level it consists of 28 operators, neglecting Higgs-lepton mixing avoiding this way the introduction of scalar operators. The corresponding $`\gamma `$-matrix can be found with the help of previously known results and diagonalized numerically. The matching conditions of the magnetic penguins $`O_7`$ and $`O_8`$ get new contributions. Their counterparts of opposite chirality, $`\stackrel{~}{O}_7`$ and $`\stackrel{~}{O}_8`$, also have to be considered. If one uses the semi-leptonic decay $`bu|ce\overline{\nu }`$ to cancel the large bottom mass dependence new contributions to this decay must be included.
R-parity breaking definitely has influence on the branching ratio of $`bs\gamma `$. However, its impact is highly model dependent because the (unknown) supersymmetric masses are mostly responsible for the size of the new contributions. In a cautious model the new couplings are able to change significantly the result if they are of order $`10^1`$. Moreover, 45 new (complex) Yukawa couplings offer an infinite number of possible scenarios. The simplest cases involve only one or two couplings present but also in these situations completely different evolutions of the branching ratio are possible. To give more concrete results we definitely need more informations on the parameters of our supersymmetric model.
###### Acknowledgements.
We wish to thank Francesca Borzumati for pointing out the importance of R-parity breaking, Christoph Greub, Tobias Hurth and especially Daniel Wyler for their support and many illuminating discussions. We also thank the authors of for pointing out to us some errors in the $`\gamma `$-matrix of the previous version of this work. Fortunately, the conclusions do not undergo significant changes.
## A Mixing matrices and interaction Lagrangian
### 1 mixing matrices
In this first appendix we present the mass mixing matrices for the relevant particles. They are needed for two reason: First, their eigenvalues correspond to the physical masses of the particles and second, the unitary matrices that diagonalize the mass matrices rotate the fields to their (physical) mass eigenstates.
#### Charginos
The charginos $`\chi _{1/2}^{\mathrm{ch}}`$ are a mixture of charged gauginos $`\lambda ^\pm `$ and Higgsinos $`h_1^{}`$ and $`h_2^+`$. Defining
$$\psi ^+=\left(\begin{array}{c}i\lambda ^+\\ h_2^+\end{array}\right)\psi ^{}=\left(\begin{array}{c}i\lambda ^{}\\ h_1^{}\end{array}\right)$$
(A1)
the mass terms are then
$$_m^{\mathrm{ch}}=\frac{1}{2}(\psi ^{+T}X^T\psi ^{}+\psi ^TX\psi ^+)+\text{h.c.},$$
(A2)
where
$$X=\left(\begin{array}{cc}M_2& g_2v_2\\ g_2v_1& \mu \end{array}\right).$$
(A3)
The two-component charginos $`\chi _i^\pm (i=1,2)`$ and the four-component charginos $`\chi _{1/2}^{\mathrm{ch}}`$ are then defined as
$$\begin{array}{cccccc}\chi _i^+\hfill & =\hfill & V_{ij}\psi _j^+\hfill & \chi _i^{}\hfill & =\hfill & U_{ij}\psi _j^{}\hfill \\ & & & & & \\ \chi _1^{\mathrm{ch}}\hfill & =\hfill & \left(\begin{array}{c}\chi _1^+\\ \overline{\chi _1^{}}\end{array}\right)\hfill & \chi _2^{\mathrm{ch}}\hfill & =\hfill & \left(\begin{array}{c}\chi _2^+\\ \overline{\chi _2^{}}\end{array}\right),\hfill \end{array}$$
(A4)
where the unitary matrices $`U`$ and $`V`$ diagonalize $`X`$:
$$M_D^{\mathrm{ch}}=U^{}XV^1=VX^{}U^T$$
(A5)
$`_m^{\mathrm{ch}}`$ then becomes
$$_m^{\mathrm{ch}}=M_{D\mathrm{\hspace{0.17em}11}}^{\mathrm{ch}}\overline{\chi _1^{\mathrm{ch}}}\chi _1^{\mathrm{ch}}M_{D\mathrm{\hspace{0.17em}22}}^{\mathrm{ch}}\overline{\chi _2^{\mathrm{ch}}}\chi _2^{\mathrm{ch}}.$$
(A6)
$`U`$ and $`V`$ can be found by observing that
$$M_D^{\mathrm{ch}\mathrm{\hspace{0.17em}2}}=VX^TXV^1=U^{}XX^TU^1.$$
(A7)
They are not fixed completely by these conditions. The freedom can be used to arrange the elements of $`M_D^{\mathrm{ch}}`$ to be positive: If the $`i^{\mathrm{th}}`$ eigenvalue of $`M_D^{\mathrm{ch}}`$ is negative simply multiply the $`i^{th}`$ row of $`V`$ with $`1`$.
#### Neutralinos
The neutralinos are linear combinations of the gauginos $`\lambda ^{}`$ and $`\lambda _3`$ and the neutral Higgsinos $`h_1^0`$ and $`h_2^0`$. If we define
$$\psi ^0=\left(\begin{array}{c}i\lambda ^{}\\ i\lambda _3\\ h_1^0\\ h_2^0\end{array}\right)$$
(A8)
the neutralino mass term reads
$$_m^0=\frac{1}{2}\psi ^{0T}Y\psi ^0+\text{h.c.},$$
(A9)
where
$$Y=\left(\begin{array}{cccc}M_1& 0& \frac{g_1v_1}{\sqrt{2}}& \frac{g_1v_2}{\sqrt{2}}\\ 0& M_2& \frac{g_2v_1}{\sqrt{2}}& \frac{g_2v_2}{\sqrt{2}}\\ \frac{g_1v_1}{\sqrt{2}}& \frac{g_2v_1}{\sqrt{2}}& 0& \mu \\ \frac{g_1v_2}{\sqrt{2}}& \frac{g_2v_2}{\sqrt{2}}& \mu & 0\end{array}\right)$$
(A10)
Two- and four-component neutralinos must be defined as
$$\begin{array}{c}\stackrel{~}{\chi }_i^0=N_{ij}\psi _j^0i=1,\mathrm{},4\hfill \\ \chi _i^0=\left(\begin{array}{c}\stackrel{~}{\chi }_i^0\\ \overline{\stackrel{~}{\chi }_i^0}\end{array}\right)\hfill \end{array}$$
(A11)
To diagonalize the mass matrix $`N`$ must obey
$$N_D=N^{}YN^1,$$
(A12)
where $`N_D`$ is a diagonal matrix. $`N`$ can be found using the property
$$N_D^2=NY^{}YN^1.$$
(A13)
The eigenvalues and eigenvectors are found numerically. Possible negative entries in $`N_D`$ are turned positive by multiplying the corresponding row of $`N`$ by a factor of $`i`$.
#### Quarks and Leptons
The situation in the quark and lepton sector is in almost complete analogy to the standard model. The quarks and leptons get their masses from the Yukawa potential when the Higgs bosons acquire a vacuum expectation value. We define the mass eigenstates by
$$\begin{array}{cc}u_{Li}^{(m)}=U_{ij}^Lu_{Lj}\hfill & u_{Ri}^{(m)}=U_{ij}^Ru_{Rj}\hfill \\ & \\ d_{Li}^{(m)}=D_{ij}^Ld_{Lj}\hfill & d_{Ri}^{(m)}=D_{ij}^Rd_{Rj}\hfill \\ & \\ e_{Li}^{(m)}=E_{ij}^Le_{Lj}\hfill & e_{Ri}^{(m)}=E_{ij}^Re_{Rj}.\hfill \end{array}$$
(A14)
The mixing matrices must satisfy
$$\begin{array}{cccccc}\hfill D^R\lambda ^{dT}D^L& =\hfill & \lambda _D^d\hfill & =\hfill & \text{diag}\left(\frac{m_{di}}{v_1}\right)\hfill & i=1,\mathrm{},3\hfill \\ & & & & & \\ \hfill U^R\lambda ^{uT}U^L& =\hfill & \lambda _D^u\hfill & =\hfill & \text{diag}\left(\frac{m_{ui}}{v_2}\right)\hfill & \\ & & & & & \\ \hfill E^R\lambda ^{eT}E^L& =\hfill & \lambda _D^e\hfill & =\hfill & \text{diag}\left(\frac{m_{ei}}{v_1}\right)\hfill & \end{array}$$
(A15)
$`v_1`$ $`=`$ $`\sqrt{2}{\displaystyle \frac{M_W}{g_2}}\mathrm{cos}\beta `$ (A16)
$`v_2`$ $`=`$ $`\sqrt{2}{\displaystyle \frac{M_W}{g_2}}\mathrm{sin}\beta .`$ (A17)
As one can see, the eigenvalues of $`\lambda ^u`$ and $`\lambda ^d`$ are fixed by the quark masses and the minimum of the Higgs potential. In the SM the only effect of the mixing which can be seen is the CKM-matrix $`K=U^LD^L`$ appearing in the flavour changing charged currents. Therefore it is possible and convenient to set
$$D^L=D^R=U^R=𝟙(𝕌^𝕃=𝕂)$$
(A18)
(To be more precise, one chooses $`\lambda ^d`$ and $`\lambda ^e`$ to be diagonal and $`\lambda ^u=K^T\text{diag}\left(m_{ui}/v_2\right)`$.) Although in our theory the mixing matrices appear in all kinds of combinations we adopt this convention here, emphasising that it is a *choice* made just for convenience. It is possible that one day an underlying theory fixes the values of $`\lambda ^u`$ and $`\lambda ^d`$ at some (high) scale.
Please note that in the text we neglect the superscript $`m`$ for the mass eigenstates.
#### Squarks and Sleptons
If supersymmetry were not broken squarks and sleptons would be rotated to their mass basis with the help of the same matrices as their fermionic partners. But since this situation is not realistic we need to introduce a further set of unitary rotation matrices. The notation must be set up carefully because the mass eigenstates of squarks and sleptons are linear combinations of the partners of *left-* and *right-*handed partners of the corresponding fermions. We define the $`6\times 3`$-matrices $`\mathrm{\Gamma }`$ in the following way:
$$\begin{array}{ccccc}\stackrel{~}{u}^{(m)}\hfill & =\hfill & \left(\mathrm{\Gamma }_L^u|\mathrm{\Gamma }_R^u\right)\left(\begin{array}{c}\stackrel{~}{u}_L\\ \stackrel{~}{u}_R\end{array}\right)\hfill & =\hfill & \mathrm{\Gamma }^u\stackrel{~}{u}\hfill \\ & & & & \\ \stackrel{~}{d}^{(m)}\hfill & =\hfill & \left(\mathrm{\Gamma }_L^d|\mathrm{\Gamma }_R^d\right)\left(\begin{array}{c}\stackrel{~}{d}_L\\ \stackrel{~}{d}_R\end{array}\right)\hfill & =\hfill & \mathrm{\Gamma }^d\stackrel{~}{d}\hfill \\ & & & & \\ \stackrel{~}{e}^{(m)}\hfill & =\hfill & \left(\mathrm{\Gamma }_L^e|\mathrm{\Gamma }_R^e\right)\left(\begin{array}{c}\stackrel{~}{e}_L\\ \stackrel{~}{e}_R\end{array}\right)\hfill & =\hfill & \mathrm{\Gamma }^e\stackrel{~}{e}\hfill \\ & & & & \\ \stackrel{~}{\nu }^{(m)}\hfill & =\hfill & \mathrm{\Gamma }^\nu \stackrel{~}{\nu }_L\hfill & & \end{array}$$
(A19)
To diagonalize the mass terms the mixing matrices have to satisfy
$`\mathrm{\Gamma }^um_{\mathrm{su}}^2\mathrm{\Gamma }^u`$ $`=`$ $`m_{\mathrm{su}D}^2`$ (A20)
$`\mathrm{\Gamma }^dm_{\mathrm{sd}}^2\mathrm{\Gamma }^d`$ $`=`$ $`m_{\mathrm{sd}D}^2`$ (A22)
$`\mathrm{\Gamma }^em_{\mathrm{se}}^2\mathrm{\Gamma }^e`$ $`=`$ $`m_{\mathrm{se}D}^2`$ (A24)
$`\mathrm{\Gamma }^\nu m_{\mathrm{s}\nu }^2\mathrm{\Gamma }^\nu `$ $`=`$ $`m_{\mathrm{s}\nu D}^2,`$ (A26)
where the matrices on the RHS are diagonal containing the masses squared of the phy-sical particles.
The mass matrices are (with the exception of the sneutrino) of the form $`\left(\begin{array}{cc}A& B\\ B^{}& C\end{array}\right)`$, where $`A`$, $`B`$ and $`C`$ are $`3\times 3`$-matrices. For the different fields they are (choosing $`\mu `$ real)
* up-squarks:
$`A`$ $`=`$ $`m_Q^2+v_2^2\lambda ^u\lambda ^u+M_Z^2\mathrm{cos}2\beta \left({\displaystyle \frac{1}{2}}{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _W\right)𝟙_\mathrm{𝟛}`$ (A27)
$`B`$ $`=`$ $`\left(h^u+\mu \mathrm{cot}\beta \lambda ^u\right)v_2`$ (A28)
$`C`$ $`=`$ $`m_U^2+v_2^2\lambda ^u\lambda ^u+{\displaystyle \frac{2}{3}}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W𝟙_\mathrm{𝟛}`$ (A29)
* down-squarks:
$`A`$ $`=`$ $`m_Q^2+v_1^2\lambda ^d\lambda ^d+M_Z^2\mathrm{cos}2\beta \left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{3}}\mathrm{sin}^2\theta _W\right)𝟙_\mathrm{𝟛}`$ (A30)
$`B`$ $`=`$ $`\left(h^d+\mu \mathrm{tan}\beta \lambda ^d\right)v_1`$ (A31)
$`C`$ $`=`$ $`m_D^2+v_1^2\lambda ^d\lambda ^d{\displaystyle \frac{1}{3}}M_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W𝟙_\mathrm{𝟛}`$ (A32)
* selectrons
$`A`$ $`=`$ $`m_L^2+v_1^2\lambda ^e\lambda ^e+M_Z^2\mathrm{cos}2\beta \left({\displaystyle \frac{1}{2}}+\mathrm{sin}^2\theta _W\right)𝟙_\mathrm{𝟛}`$ (A33)
$`B`$ $`=`$ $`\left(h^e+\mu \mathrm{tan}\beta \lambda ^e\right)v_1`$ (A34)
$`C`$ $`=`$ $`m_E^2+v_1^2\lambda ^e\lambda ^eM_Z^2\mathrm{cos}2\beta \mathrm{sin}^2\theta _W𝟙_\mathrm{𝟛}`$ (A35)
The sneutrinos have
$$m_{\mathrm{s}\nu }^2=m_{\stackrel{~}{e}}^2+\frac{M_Z}{2}\mathrm{cos}2\beta 𝟙_\mathrm{𝟛}$$
(A36)
#### Higgses
The Higgs sector consists of two $`SU(2)`$ doublets $`\stackrel{~}{h}_1`$ and $`\stackrel{~}{h}_2`$. The real and the imaginary part of the neutral components mix via the matrix
$$\left(\begin{array}{cc}\mu ^2+m_{H_1}^2+\frac{g_1^2+g_2^2}{4}(3v_1^2v_2^2)& \stackrel{~}{\mu }\frac{g_1^2+g_2^2}{2}v_1v_2\\ \stackrel{~}{\mu }\frac{g_1^2+g_2^2}{2}v_1v_2& \mu ^2+m_{H_2}^2\frac{g_1^2+g_2^2}{4}(v_1^23v_2^2)\end{array}\right)$$
(A37)
and
$$\left(\begin{array}{cc}\mu ^2+m_{H_1}^2+\frac{g_1^2+g_2^2}{4}(v_1^2v_2^2)& \stackrel{~}{\mu }\\ \stackrel{~}{\mu }& \mu ^2+m_{H_2}^2\frac{g_1^2+g_2^2}{4}(v_1^2v_2^2)\end{array}\right)$$
(A38)
respectively. The charged components give rise to a mass matrix of the form
$$\left(\begin{array}{cc}\mu ^2+m_{H_1}^2+\frac{g_1^2+g_2^2}{4}(v_1^2v_2^2)+\frac{g_2^2v_2^2}{2}& \stackrel{~}{\mu }+\frac{g_2^2}{2}v_1v_2\\ \stackrel{~}{\mu }+\frac{g_2^2}{2}v_1v_2& \mu ^2+m_{H_2}^2\frac{g_1^2+g_2^2}{4}(v_1^2v_2^2)+\frac{g_2^2v_1^2}{2}\end{array}\right)$$
(A39)
All the above expressions are valid at tree-level. The vacuum expectation values $`v_1`$ and $`v_2`$ are found by minimizing the effective Higgs potential. If one inserts the gained formulas into (A37) - (A39) one of the eigenvalues of (A38) and (A39) vanishes, indicating the eaten fields of the Higgs mechanism that takes place.
### 2 Interaction Lagrangian
For the evaluation of the matching conditions we need certain parts of the interaction Lagrangian. In addition to equations(72) - (80) these are
#### Squark-Quark-Chargino
$`_{\stackrel{~}{q}q\chi ^{\mathrm{ch}}}`$ $`=`$ $`\stackrel{~}{u}_j\overline{d_i}\left[A_{ij\mathrm{}}^dP_L+B_{ij\mathrm{}}^dP_R\right]\chi _{\mathrm{}}^{\mathrm{ch}c}+\stackrel{~}{u}_i^{}\overline{\chi _{\mathrm{}}^{\mathrm{ch}c}}\left[A_{ij\mathrm{}}^dP_R+B_{ij\mathrm{}}^dP_L\right]d_j`$ (A41)
$`+\stackrel{~}{d}_j\overline{u_i}\left[A_{ij\mathrm{}}^uP_L+B_{ij\mathrm{}}^uP_R\right]\chi _{\mathrm{}}^{\mathrm{ch}}+\stackrel{~}{d}_i^{}\overline{\chi _{\mathrm{}}^{\mathrm{ch}}}\left[A_{ij\mathrm{}}^uP_R+B_{ij\mathrm{}}^uP_L\right]u_j,`$
where
$`A_{ij\mathrm{}}^d`$ $`=`$ $`(\lambda _D^d\mathrm{\Gamma }_L^u)_{ij}U_\mathrm{}2^{}`$ (A42)
$`B_{ij\mathrm{}}^d`$ $`=`$ $`(K^{}\lambda _D^u\mathrm{\Gamma }_R^u)_{ij}V_\mathrm{}2g_2\mathrm{\Gamma }_{Lij}^uV_\mathrm{}1`$ (A44)
$`A_{ij\mathrm{}}^u`$ $`=`$ $`(\lambda _D^uK\mathrm{\Gamma }_L^d)_{ij}V_\mathrm{}2^{}`$ (A46)
$`B_{ij\mathrm{}}^u`$ $`=`$ $`(K\lambda _D^d\mathrm{\Gamma }_R^d)_{ij}U_\mathrm{}2g_2(K\mathrm{\Gamma }_L^d)_{ij}U_\mathrm{}1`$ (A48)
$`P_{L/R}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\gamma ^5)`$ (A50)
and $`\chi _{\mathrm{}}^{\mathrm{ch}c}`$ denotes the charge conjugated field.
#### Squark-Quark-Neutralino
$`_{\stackrel{~}{q}q\chi ^0}`$ $`=`$ $`\stackrel{~}{d}_j\overline{d_i}\left[C_{ij\mathrm{}}^dP_L+D_{ij\mathrm{}}^dP_R\right]\chi _{\mathrm{}}^0\stackrel{~}{d}_i^{}\overline{\chi _{\mathrm{}}^0}\left[C_{ij\mathrm{}}^dP_R+D_{ij\mathrm{}}^dP_L\right]d_j`$ (A52)
$`\stackrel{~}{u}_j\overline{u_i}\left[C_{ij\mathrm{}}^uP_L+D_{ij\mathrm{}}^uP_R\right]\chi _{\mathrm{}}^0\stackrel{~}{u}_i^{}\overline{\chi _{\mathrm{}}^0}\left[C_{ij\mathrm{}}^uP_R+D_{ij\mathrm{}}^uP_L\right]u_j,`$
where
$`C_{ij\mathrm{}}^d`$ $`=`$ $`(\lambda _D^d\mathrm{\Gamma }_L^d)_{ij}N_\mathrm{}3^{}\sqrt{2}g_1Q_d\mathrm{\Gamma }_{Rij}^dN_\mathrm{}1^{}`$ (A53)
$`D_{ij\mathrm{}}^d`$ $`=`$ $`(\lambda _D^d\mathrm{\Gamma }_R^d)_{ij}N_\mathrm{}3+{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Gamma }_{Lij}^d((2Q_d+1)g_1N_\mathrm{}1g_2N_\mathrm{}2)`$ (A55)
$`C_{ij\mathrm{}}^u`$ $`=`$ $`(\lambda _D^uK\mathrm{\Gamma }_L^u)_{ij}N_\mathrm{}4^{}\sqrt{2}Q_ug_1\mathrm{\Gamma }_{Rij}^uN_\mathrm{}1^{}`$ (A57)
$`D_{ij\mathrm{}}^u`$ $`=`$ $`(\lambda _D^u\mathrm{\Gamma }_R^u)_{ij}N_\mathrm{}4+{\displaystyle \frac{1}{\sqrt{2}}}(K\mathrm{\Gamma }_L^u)_{ij}((2Q_d+1)g_1N_\mathrm{}1+g_2N_\mathrm{}2).`$ (A59)
#### Squark-Quark-Gluino
$`_{\stackrel{~}{q}qg}`$ $`=`$ $`\sqrt{2}g_sT^{A\alpha \beta }[\stackrel{~}{u}_{i\alpha }^{}\overline{g_A}[(\mathrm{\Gamma }_L^uK^{})_{ij}P_L\mathrm{\Gamma }_{Rij}^uP_R]u_{j\beta }`$ (A60)
$`+`$ $`\stackrel{~}{d}_{i\alpha }^{}\overline{g_A}\left[\mathrm{\Gamma }_{Lij}^dP_L\mathrm{\Gamma }_{Rij}^dP_R\right]d_{j\beta }`$ (A61)
$`+`$ $`\stackrel{~}{u}_{i\alpha }\overline{u_{j\beta }}\left[(K\mathrm{\Gamma }_L^u)_{ij}P_R\mathrm{\Gamma }_{Rij}^uP_L\right]g_A`$ (A62)
$`+`$ $`\stackrel{~}{d}_{j\alpha }\overline{d_{i\beta }}[\mathrm{\Gamma }_{Lij}^dP_R\mathrm{\Gamma }_{Rij}^dP_L]g_A]`$ (A63)
#### Gluino-Gluino-Gluon
$$_{gg\stackrel{~}{g}}=\frac{i}{2}g_sf^{ABC}\overline{g_A}\gamma _\mu g_BA_C^\mu $$
(A64)
Note: There is a symmetry factor of two in the Feynman rule for this vertex.
## B Renormalization group equations
We present the full set of RGE’s for all the parameters of the MSSM including R-parity breaking terms. Our results are in complete agreement with , although we don’t restrict ourselves to couplings of the third generation. All the formulas can be derived from the expressions for the most general form of a softly broken SUSY . Let us begin with the parameters of the superpotential $`W`$ ($`t=\mathrm{ln}\mu `$).
$`16\pi ^2{\displaystyle \frac{d}{dt}}\mu `$ $`=`$ $`\text{Tr}\left(3\lambda ^d\lambda ^d+3\lambda ^u\lambda ^u+\lambda ^e\lambda ^e\right)\mu \left[g_1^2+3g_2^2\right]\mu `$ (B2)
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}^{}\right]\mu _{2i}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\mu _{2i}`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{ij}+\lambda _{i\mathrm{}m}\lambda _{j\mathrm{}m}^{}+3\lambda _{i\mathrm{}m}^{}\lambda _{j\mathrm{}m}^{}\right]\mu _{2j}`$ (B6)
$`+\left[3\lambda _{i\mathrm{}m}^{}\lambda _\mathrm{}m^d+\lambda _{i\mathrm{}m}\lambda _\mathrm{}m^e\right]\mu `$
$`+3\text{Tr}(\lambda ^u\lambda ^u)\mu _{2i}\left[g_1^2+3g_2^2\right]\mu _{2i}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\lambda _{ij}^u`$ $`=`$ $`\left[(\lambda ^u\lambda ^u)_{ik}+(\lambda ^d\lambda ^d)_{ik}+\lambda _{\mathrm{}km}^{}\lambda _{\mathrm{}im}^{}\right]\lambda _{kj}^u`$ (B10)
$`+\left[2(\lambda ^u\lambda ^u)_{kj}+\lambda _{k\mathrm{}m}^{\prime \prime }\lambda _{j\mathrm{}m}^{\prime \prime }\right]\lambda _{ik}^u`$
$`+3\text{Tr}(\lambda ^u\lambda ^u)\lambda _{ij}^u\left[{\displaystyle \frac{13}{9}}g_1^2+3g_2^2+{\displaystyle \frac{16}{3}}g_3^2\right]\lambda _{ij}^u`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\lambda _{ij}^d`$ $`=`$ $`\left[(\lambda ^u\lambda ^u)_{ik}+(\lambda ^d\lambda ^d)_{ik}+\lambda _{\mathrm{}km}^{}\lambda _{\mathrm{}im}^{}\right]\lambda _{kj}^d`$ (B16)
$`+\left[2(\lambda ^d\lambda ^d)_{kj}+2\lambda _{\mathrm{}mk}^{}\lambda _{\mathrm{}mj}^{}+2\lambda _{\mathrm{}km}^{\prime \prime }\lambda _{\mathrm{}jm}^{\prime \prime }\right]\lambda _{ik}^d`$
$`+\left[3\text{Tr}(\lambda ^d\lambda ^d)+\text{Tr}(\lambda ^e\lambda ^e)\right]\lambda _{ij}^d`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{k\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{k\mathrm{}m}^{}\right]\lambda _{kij}^{}`$
$`\left[{\displaystyle \frac{7}{9}}g_1^2+3g_2^2+{\displaystyle \frac{16}{3}}g_3^2\right]\lambda _{ij}^d`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\lambda _{ij}^e`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{ik}+\lambda _{k\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{k\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]\lambda _{kj}^e`$ (B22)
$`+\left[2(\lambda ^e\lambda ^e)_{kj}+\lambda _{\mathrm{}mk}^{}\lambda _{\mathrm{}mj}\right]\lambda _{ik}^e`$
$`+\left[3\text{Tr}(\lambda ^d\lambda ^d)+\text{Tr}(\lambda ^e\lambda ^e)\right]\lambda _{ij}^e`$
$`+{\displaystyle \frac{1}{2}}\left[3\lambda _\mathrm{}m^d\lambda _{k\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{k\mathrm{}m}^{}\right]\lambda _{kij}`$
$`\left[3g_1^2+3g_2^2\right]\lambda _{ij}^e`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\lambda _{ijk}`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{in}+\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]\lambda _{njk}`$ (B29)
$`+\left[(\lambda ^e\lambda ^e)_{jn}+\lambda _{n\mathrm{}m}^{}\lambda _{j\mathrm{}m}+3\lambda _{n\mathrm{}m}^{}\lambda _{j\mathrm{}m}^{}\right]\lambda _{ink}`$
$`+\left[2(\lambda ^e\lambda ^e)_{nk}+\lambda _{\mathrm{}mn}^{}\lambda _{\mathrm{}mk}\right]\lambda _{ijn}`$
$`+2\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]\lambda _{jk}^e`$
$`2\left[3\lambda _\mathrm{}m^d\lambda _{j\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{j\mathrm{}m}\right]\lambda _{ik}^e`$
$`\left[3g_1^2+3g_2^2\right]\lambda _{ijk}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\lambda _{ijk}^{}`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{in}+\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]\lambda _{njk}^{}`$ (B35)
$`+\left[(\lambda ^u\lambda ^u)_{jn}+(\lambda ^d\lambda ^d)_{jn}+\lambda _{\mathrm{}nm}^{}\lambda _{\mathrm{}jm}^{}\right]\lambda _{ink}^{}`$
$`+\left[2(\lambda ^d\lambda ^d)_{nk}+2\lambda _{\mathrm{}mn}^{}\lambda _{\mathrm{}mk}^{}+2\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{\mathrm{}km}^{\prime \prime }\right]\lambda _{ijn}^{}`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]\lambda _{jk}^d`$
$`\left[{\displaystyle \frac{7}{9}}g_1^2+3g_2^2+{\displaystyle \frac{16}{3}}g_3^2\right]\lambda _{ijk}^{}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\lambda _{ijk}^{\prime \prime }`$ $`=`$ $`\left[2(\lambda ^u\lambda ^u)_{ni}+\lambda _{n\mathrm{}m}^{\prime \prime }\lambda _{i\mathrm{}m}^{\prime \prime }\right]\lambda _{njk}^{\prime \prime }`$ (B40)
$`+\left[2(\lambda ^d\lambda ^d)_{nj}+2\lambda _{\mathrm{}nm}^{}\lambda _{\mathrm{}mj}^{}+2\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{\mathrm{}jm}^{\prime \prime }\right]\lambda _{ink}^{\prime \prime }`$
$`+\left[2(\lambda ^d\lambda ^d)_{nk}+2\lambda _{\mathrm{}nm}^{}\lambda _{\mathrm{}mk}^{}+2\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{\mathrm{}km}^{\prime \prime }\right]\lambda _{ijn}^{\prime \prime }`$
$`\left[{\displaystyle \frac{4}{3}}g_1^2+8g_3^2\right]\lambda _{ijk}^{\prime \prime }`$
The parameters of $`V_{\mathrm{soft}}`$ obey the following RGE’s:
$`16\pi ^2{\displaystyle \frac{d}{dt}}\stackrel{~}{\mu }`$ $`=`$ $`\text{Tr}\left(3\lambda ^d\lambda ^d+3\lambda ^u\lambda ^u+\lambda ^e\lambda ^e\right)\stackrel{~}{\mu }`$ (B45)
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}^{}\right]\stackrel{~}{\mu }_{2i}`$
$`+2\text{Tr}\left(3h^d\lambda ^d+3h^u\lambda ^u+h^e\lambda ^e\right)\mu `$
$`+2\left[3h_\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+h_\mathrm{}m^e\lambda _{i\mathrm{}m}^{}\right]\mu _{2i}`$
$`\left[g_1^2+3g_2^2\right]\stackrel{~}{\mu }+\left[2M_1g_1^2+6M_2g_2^2\right]\mu `$
$`16\pi ^2{\displaystyle \frac{d}{dt}}\stackrel{~}{\mu }_{2i}`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{ij}+\lambda _{j\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{j\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]\stackrel{~}{\mu }_{2j}+3\text{Tr}(\lambda ^u\lambda ^u)\stackrel{~}{\mu }_{2i}`$ (B51)
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]\stackrel{~}{\mu }`$
$`+2\left[(h^e\lambda ^e)_{ij}+C_{i\mathrm{}m}\lambda _{j\mathrm{}m}^{}+3C_{i\mathrm{}m}^{}\lambda _{j\mathrm{}m}^{}\right]\mu _{2j}`$
$`+6\text{Tr}(h^u\lambda ^u)\mu _{2i}+2\left[3\lambda _\mathrm{}m^dC_{i\mathrm{}m}^{}+\lambda _\mathrm{}m^eC_{i\mathrm{}m}\right]\mu `$
$`\left[g_1^2+3g_2^2\right]\stackrel{~}{\mu }_{2i}+\left[2M_1g_1^2+6M_2g_2^2\right]\mu _{2i}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}h_{ij}^u`$ $`=`$ $`5(\lambda ^u\lambda ^uh^u)_{ij}+(\lambda ^d\lambda ^dh^u)_{ij}+4(h^u\lambda ^u\lambda ^u)_{ij}`$ (B57)
$`+2(h^d\lambda ^d\lambda ^u)_{ij}+3\text{Tr}(\lambda ^u\lambda ^u)h_{ij}^u+6\text{Tr}(h^u\lambda ^u)\lambda _{ij}^u`$
$`+\lambda _{m\mathrm{}n}^{}\lambda _{min}^{}h_\mathrm{}j^u+\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{jnm}^{\prime \prime }h_i\mathrm{}^u+2\lambda _{mn\mathrm{}}^{}C_{mi\mathrm{}}^{}\lambda _{nj}^u`$
$`+2\lambda _{\mathrm{}mn}^{\prime \prime }C_{jmn}^{\prime \prime }\lambda _i\mathrm{}^u+{\displaystyle \frac{13}{9}}g_1^2(2M_1\lambda _{ij}^uh_{ij}^u)`$
$`+3g_2^2(2M_2\lambda _{ij}^uh_{ij}^u)+{\displaystyle \frac{16}{3}}g_3^2(2M_3\lambda _{ij}^uh_{ij}^u)`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}h_{ij}^d`$ $`=`$ $`5(\lambda ^d\lambda ^dh^d)_{ij}+(\lambda ^u\lambda ^uh^d)_{ij}+4(h^d\lambda ^d\lambda ^d)_{ij}`$ (B68)
$`+2(h^u\lambda ^u\lambda ^d)_{ij}+\text{Tr}(3\lambda ^d\lambda ^d+\lambda ^e\lambda ^e)h_{ij}^d`$
$`+\text{Tr}(6h^d\lambda ^d+2h^e\lambda ^e)\lambda _{ij}^d`$
$`+\left[\lambda _\mathrm{}m^eC_{kij}^{}+2h_\mathrm{}m^e\lambda _{kij}^{}\right]\lambda _{k\mathrm{}m}^{}`$
$`+3\left[\lambda _\mathrm{}m^dC_{kij}^{}+2h_\mathrm{}m^d\lambda _{kij}^{}\right]\lambda _{k\mathrm{}m}^{}`$
$`+\left[h_{kj}^d\lambda _{\mathrm{}im}^{}+2\lambda _{kj}^dC_{\mathrm{}im}^{}\right]\lambda _{\mathrm{}km}^{}`$
$`+2\left[h_{ik}^d\lambda _{\mathrm{}mj}^{}+2\lambda _{ik}^dC_{\mathrm{}mj}^{}\right]\lambda _{\mathrm{}mk}^{}`$
$`+2\left[h_{ik}^d\lambda _{\mathrm{}jm}^{\prime \prime }+2\lambda _{ik}^dC_{\mathrm{}jm}^{\prime \prime }\right]\lambda _{\mathrm{}km}^{\prime \prime }`$
$`+{\displaystyle \frac{7}{9}}g_1^2(2M_1\lambda _{ij}^dh_{ij}^d)+3g_2^2(2M_2\lambda _{ij}^dh_{ij}^d)`$
$`+{\displaystyle \frac{16}{3}}g_3^2(2M_3\lambda _{ij}^dh_{ij}^d)`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}h_{ij}^e`$ $`=`$ $`5(\lambda ^e\lambda ^eh^e)_{ij}+4(h^e\lambda ^e\lambda ^e)_{ij}`$ (B75)
$`+\text{Tr}(3\lambda ^d\lambda ^d+\lambda ^e\lambda ^e)h_{ij}^e+\text{Tr}(6h^d\lambda ^d+2h^e\lambda ^e)\lambda _{ij}^e`$
$`+\left[\lambda _\mathrm{}m^eC_{kij}+2h_\mathrm{}m^e\lambda _{kij}+h_{kj}^e\lambda _{i\mathrm{}m}+2\lambda _{kj}^eC_{i\mathrm{}m}\right]\lambda _{k\mathrm{}m}^{}`$
$`+3\left[\lambda _\mathrm{}m^dC_{kij}+2h_\mathrm{}m^d\lambda _{kij}+h_{kj}^e\lambda _{i\mathrm{}m}^{}+2\lambda _{kj}^eC_{i\mathrm{}m}^{}\right]\lambda _{k\mathrm{}m}^{}`$
$`+\left[h_{ik}^e\lambda _{\mathrm{}mj}+2\lambda _{ik}^eC_{\mathrm{}mj}\right]\lambda _{\mathrm{}mk}^{}`$
$`+3g_1^2(2M_1\lambda _{ij}^eh_{ij}^e)+3g_2^2(2M_2\lambda _{ij}^eh_{ij}^e)`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}C_{ijk}`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{in}+\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]C_{njk}`$ (B86)
$`+\left[(\lambda ^e\lambda ^e)_{jn}+\lambda _{n\mathrm{}m}^{}\lambda _{j\mathrm{}m}+3\lambda _{n\mathrm{}m}^{}\lambda _{j\mathrm{}m}^{}\right]C_{ink}`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]h_{jk}^e\left[3\lambda _\mathrm{}m^d\lambda _{j\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{j\mathrm{}m}\right]h_{ik}^e`$
$`+\left[2(\lambda ^e\lambda ^e)_{nk}+\lambda _{\mathrm{}mn}^{}\lambda _{\mathrm{}mk}\right]C_{ijn}`$
$`+\left[2(h^e\lambda ^e)_{in}+2\lambda _{n\mathrm{}m}^{}C_{i\mathrm{}m}+6\lambda _{n\mathrm{}m}^{}C_{i\mathrm{}m}^{}\right]\lambda _{njk}`$
$`+\left[2(h^e\lambda ^e)_{jn}+2\lambda _{n\mathrm{}m}^{}C_{j\mathrm{}m}+6\lambda _{n\mathrm{}m}^{}C_{j\mathrm{}m}^{}\right]\lambda _{ink}`$
$`+\left[6\lambda _\mathrm{}m^dC_{i\mathrm{}m}^{}+2\lambda _\mathrm{}m^eC_{i\mathrm{}m}\right]\lambda _{jk}^e`$
$`\left[6\lambda _\mathrm{}m^dC_{j\mathrm{}m}^{}+2\lambda _\mathrm{}m^eC_{j\mathrm{}m}\right]\lambda _{ik}^e`$
$`+\left[4(\lambda ^eh^e)_{nk}+2\lambda _{\mathrm{}mn}^{}C_{\mathrm{}mk}\right]\lambda _{ijn}`$
$`+3g_1^2(2M_1\lambda _{ijk}C_{ijk})+3g_2^2(2M_2\lambda _{ijk}C_{ijk})`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}C_{ijk}^{}`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{in}+\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{n\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]C_{njk}^{}`$ (B97)
$`+\left[(\lambda ^u\lambda ^u)_{jn}+(\lambda ^d\lambda ^d)_{jn}+\lambda _{\mathrm{}nm}^{}\lambda _{\mathrm{}jm}^{}\right]C_{ink}^{}`$
$`+\left[2(\lambda ^d\lambda ^d)_{nk}+2\lambda _{\mathrm{}mn}^{}\lambda _{\mathrm{}mk}^{}+2\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{\mathrm{}km}^{\prime \prime }\right]C_{ijn}^{}`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]h_{jk}^d`$
$`+\left[2(h^e\lambda ^e)_{in}+2\lambda _{n\mathrm{}m}^{}C_{i\mathrm{}m}+6\lambda _{n\mathrm{}m}^{}C_{i\mathrm{}m}^{}\right]\lambda _{njk}^{}`$
$`+\left[2(h^u\lambda ^u)_{jn}+2(h^d\lambda ^d)_{jn}+2\lambda _{\mathrm{}nm}^{}C_{\mathrm{}jm}^{}\right]\lambda _{ink}^{}`$
$`+\left[4(\lambda ^dh^d)_{nk}+4\lambda _{\mathrm{}mn}^{}C_{\mathrm{}mk}^{}+4\lambda _{\mathrm{}nm}^{\prime \prime }C_{\mathrm{}km}^{\prime \prime }\right]\lambda _{ijn}^{}`$
$`+\left[6\lambda _\mathrm{}m^dC_{i\mathrm{}m}^{}+2\lambda _\mathrm{}m^eC_{i\mathrm{}m}\right]\lambda _{jk}^d`$
$`+{\displaystyle \frac{7}{9}}g_1^2(2M_1\lambda _{ijk}^{}C_{ijk}^{})+3g_2^2(2M_2\lambda _{ijk}^{}C_{ijk}^{})`$
$`+{\displaystyle \frac{16}{3}}g_3^2(2M_3\lambda _{ijk}^{}C_{ijk}^{})`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}C_{ijk}^{\prime \prime }`$ $`=`$ $`\left[2(\lambda ^u\lambda ^u)_{ni}+\lambda _{n\mathrm{}m}^{\prime \prime }\lambda _{i\mathrm{}m}^{\prime \prime }\right]C_{njk}^{\prime \prime }`$ (B105)
$`+\left[2(\lambda ^d\lambda ^d)_{nj}+2\lambda _{\mathrm{}mn}^{}\lambda _{\mathrm{}mj}^{}+2\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{\mathrm{}jm}^{\prime \prime }\right]C_{ink}^{\prime \prime }`$
$`+\left[2(\lambda ^d\lambda ^d)_{nk}+2\lambda _{\mathrm{}mn}^{}\lambda _{\mathrm{}mk}^{}+2\lambda _{\mathrm{}nm}^{\prime \prime }\lambda _{\mathrm{}km}^{\prime \prime }\right]C_{ijn}^{\prime \prime }`$
$`+\left[4(\lambda ^uh^u)_{ni}+2\lambda _{n\mathrm{}m}^{\prime \prime }C_{i\mathrm{}m}^{\prime \prime }\right]\lambda _{njk}^{\prime \prime }`$
$`+\left[4(\lambda ^dh^d)_{nj}+4\lambda _{\mathrm{}mn}^{}C_{\mathrm{}mj}^{}+4\lambda _{\mathrm{}nm}^{\prime \prime }C_{\mathrm{}jm}^{\prime \prime }\right]\lambda _{ink}^{\prime \prime }`$
$`+\left[4(\lambda ^dh^d)_{nk}+4\lambda _{\mathrm{}mn}^{}C_{\mathrm{}mk}^{}+4\lambda _{\mathrm{}nm}^{\prime \prime }C_{\mathrm{}km}^{\prime \prime }\right]\lambda _{ijn}^{\prime \prime }`$
$`+{\displaystyle \frac{4}{3}}g_1^2(2M_1\lambda _{ijk}^{\prime \prime }C_{ijk}^{\prime \prime })+8g_3^2(2M_3\lambda _{ijk}^{\prime \prime }C_{ijk}^{\prime \prime })`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{U_iU_j}^2`$ $`=`$ $`\left[2(\lambda ^u\lambda ^u)_{ki}+\lambda _{k\mathrm{}m}^{\prime \prime }\lambda _{i\mathrm{}m}^{\prime \prime }\right]m_{U_kU_j}^2`$ (B110)
$`+\left[2(\lambda ^u\lambda ^u)_{jk}+\lambda _{j\mathrm{}m}^{\prime \prime }\lambda _{k\mathrm{}m}^{\prime \prime }\right]m_{U_iU_k}^2`$
$`+4\lambda _j\mathrm{}^u\lambda _{mi}^um_{Q_{\mathrm{}}Q_m}^2+4\lambda _j\mathrm{}^u\lambda _\mathrm{}i^um_{H_2H_2}^2+4\lambda _{j\mathrm{}m}^{\prime \prime }\lambda _{inm}^{\prime \prime }m_{D_{\mathrm{}}D_n}^2`$
$`+4(h^uh^u)_{ji}+2C_{j\mathrm{}m}^{\prime \prime }C_{i\mathrm{}m}^{\prime \prime }{\displaystyle \frac{32}{9}}g_1^2M_1^2\delta _{ij}{\displaystyle \frac{32}{3}}g_3^2M_3^2\delta _{ij}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{D_iD_j}^2`$ $`=`$ $`\left[2(\lambda ^d\lambda ^d)_{jk}+2\lambda _{\mathrm{}mj}^{}\lambda _{\mathrm{}mk}^{}+2\lambda _{\mathrm{}jm}^{\prime \prime }\lambda _{\mathrm{}km}^{\prime \prime }\right]m_{D_iD_k}^2`$ (B118)
$`+\left[2(\lambda ^d\lambda ^d)_{ki}+2\lambda _{\mathrm{}mk}^{}\lambda _{\mathrm{}mi}^{}+2\lambda _{\mathrm{}km}^{\prime \prime }\lambda _{\mathrm{}im}^{\prime \prime }\right]m_{D_kD_j}^2`$
$`+4\lambda _j\mathrm{}^d\lambda _\mathrm{}i^dm_{H_1H_1}^2+4\lambda _j\mathrm{}^d\lambda _{mi}^dm_{Q_{\mathrm{}}Q_m}^2+4\lambda _{\mathrm{}nj}^{}\lambda _{mni}^{}m_{L_{\mathrm{}}L_m}^2`$
$`+4\lambda _{n\mathrm{}j}^{}\lambda _{nmi}^{}m_{Q_{\mathrm{}}Q_m}^2+4\lambda _{n\mathrm{}j}^{\prime \prime }\lambda _{nmi}^{\prime \prime }m_{D_{\mathrm{}}D_m}^2+4\lambda _{\mathrm{}nj}^{\prime \prime }\lambda _{mni}^{\prime \prime }m_{U_{\mathrm{}}U_m}^2`$
$`+4\lambda _{knj}^{}\lambda _{ni}^dm_{H_1L_k}^2+4\lambda _{nj}^d\lambda _{kni}^{}m_{L_kH_1}^2`$
$`+4(h^dh^d)_{ji}+4C_{\mathrm{}mj}^{}C_{\mathrm{}mi}^{}+4C_{\mathrm{}jm}^{\prime \prime }C_{\mathrm{}im}^{\prime \prime }`$
$`{\displaystyle \frac{8}{9}}g_1^2M_1^2\delta _{ij}{\displaystyle \frac{32}{3}}g_3^2M_3^2\delta _{ij}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{Q_iQ_j}^2`$ $`=`$ $`\left[(\lambda ^u\lambda ^u)_{kj}+(\lambda ^d\lambda ^d)_{kj}+\lambda _{\mathrm{}jm}^{}\lambda _{\mathrm{}km}^{}\right]m_{Q_iQ_k}^2`$ (B126)
$`+\left[(\lambda ^u\lambda ^u)_{ik}+(\lambda ^d\lambda ^d)_{ik}+\lambda _{\mathrm{}km}^{}\lambda _{\mathrm{}im}^{}\right]m_{Q_kQ_j}^2`$
$`+2\lambda _\mathrm{}j^u\lambda _i\mathrm{}^um_{H_2H_2}^2+2\lambda _{mj}^u\lambda _i\mathrm{}^um_{U_mU_{\mathrm{}}}^2+2\lambda _\mathrm{}j^d\lambda _i\mathrm{}^dm_{H_1H_1}^2`$
$`+2\lambda _{mj}^d\lambda _i\mathrm{}^dm_{D_mD_{\mathrm{}}}^2+2\lambda _{mjn}^{}\lambda _{\mathrm{}in}^{}m_{L_mL_{\mathrm{}}}^2+2\lambda _{njm}^{}\lambda _{ni\mathrm{}}^{}m_{D_mD_{\mathrm{}}}^2`$
$`+2\lambda _{kjn}^{}\lambda _{in}^dm_{L_kH_1}^2+2\lambda _{jn}^d\lambda _{kin}^{}m_{H_1L_k}^2`$
$`+2(h^uh^u)_{ij}+2(h^dh^d)_{ij}+2C_{\mathrm{}jm}^{}C_{\mathrm{}im}^{}`$
$`{\displaystyle \frac{2}{9}}g_1^2M_1^2\delta _{ij}6g_2^2M_2^2\delta _{ij}{\displaystyle \frac{32}{3}}g_3^2M_3^2\delta _{ij}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{E_iE_j}^2`$ $`=`$ $`\left[2(\lambda ^e\lambda ^e)_{jk}+\lambda _{\mathrm{}mj}^{}\lambda _{\mathrm{}mk}\right]m_{E_iE_k}^2`$ (B132)
$`+\left[2(\lambda ^e\lambda ^e)_{ki}+\lambda _{\mathrm{}mk}^{}\lambda _{\mathrm{}mi}\right]m_{E_kE_j}^2`$
$`+4\lambda _{jk}^e\lambda _{ki}^em_{H_1H_1}^2+4\lambda _{jm}^e\lambda _\mathrm{}i^em_{L_mL_{\mathrm{}}}^2+4\lambda _{nmj}^{}\lambda _{n\mathrm{}i}m_{L_mL_{\mathrm{}}}^2`$
$`+4\lambda _{nkj}^{}\lambda _{ni}^em_{L_kH_1}^2+4\lambda _{nj}^e\lambda _{nki}m_{H_1L_k}^2`$
$`+4(h^eh^e)_{ij}+2C_{\mathrm{}mj}^{}C_{\mathrm{}mi}8g_1^2M_1^2\delta _{ij}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{L_iL_j}^2`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{ik}+\lambda _{k\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{k\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]m_{L_kL_j}^2`$ (B142)
$`+\left[(\lambda ^e\lambda ^e)_{kj}+\lambda _{j\mathrm{}m}^{}\lambda _{k\mathrm{}m}+3\lambda _{j\mathrm{}m}^{}\lambda _{k\mathrm{}m}^{}\right]m_{L_iL_k}^2`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]m_{H_1L_j}^2`$
$`+\left[3\lambda _{j\mathrm{}m}^{}\lambda _\mathrm{}m^d+\lambda _{j\mathrm{}m}^{}\lambda _\mathrm{}m^e\right]m_{L_iH_1}^2`$
$`+2\lambda _{kj}^e\lambda _{ik}^em_{H_1H_1}^2+2\lambda _{mj}^e\lambda _i\mathrm{}^em_{E_mE_{\mathrm{}}}^2+2\lambda _{mjn}^{}\lambda _{\mathrm{}in}m_{L_mL_{\mathrm{}}}^2`$
$`+2\lambda _{njm}^{}\lambda _{ni\mathrm{}}m_{E_mE_{\mathrm{}}}^2+6\lambda _{jmn}^{}\lambda _{i\mathrm{}n}^{}m_{Q_mQ_{\mathrm{}}}^2+6\lambda _{jnm}^{}\lambda _{in\mathrm{}}^{}m_{D_mD_{\mathrm{}}}^2`$
$`+2\lambda _{kjn}^{}\lambda _{in}^em_{L_kH_1}^2+2\lambda _{jn}^e\lambda _{kin}m_{H_1L_k}^2`$
$`+2(h^eh^e)_{ij}+2C_{j\mathrm{}m}^{}C_{i\mathrm{}m}+6C_{j\mathrm{}m}^{}C_{i\mathrm{}m}^{}`$
$`2g_1^2M_1^2\delta _{ij}6g_2^2M_2^2\delta _{ij}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{H_1H_1}^2`$ $`=`$ $`\text{Tr}(6\lambda ^d\lambda ^d+2\lambda ^e\lambda ^e)m_{H_1H_1}^2`$ (B150)
$`+\left[3\lambda _{i\mathrm{}m}^{}\lambda _\mathrm{}m^d+\lambda _{i\mathrm{}m}^{}\lambda _\mathrm{}m^e\right]m_{L_iH_1}^2`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]m_{H_1L_i}^2`$
$`+6\lambda _{nm}^d\lambda _\mathrm{}n^dm_{Q_mQ_{\mathrm{}}}^2+6\lambda _{mn}^d\lambda _n\mathrm{}^dm_{D_mD_{\mathrm{}}}^2`$
$`+2\lambda _{nm}^e\lambda _\mathrm{}n^em_{L_mL_{\mathrm{}}}^2+2\lambda _{mn}^e\lambda _n\mathrm{}^em_{E_mE_{\mathrm{}}}^2`$
$`+\text{Tr}(6h^dh^d+2h^eh^e)`$
$`2g_1^2M_1^26g_2^2M_2^2`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{L_iH_1}^2`$ $`=`$ $`\left[(\lambda ^e\lambda ^e)_{ik}+\lambda _{k\mathrm{}m}^{}\lambda _{i\mathrm{}m}+3\lambda _{k\mathrm{}m}^{}\lambda _{i\mathrm{}m}^{}\right]m_{L_kH_1}^2`$ (B158)
$`+\text{Tr}(3\lambda ^d\lambda ^d+\lambda ^e\lambda ^e)m_{L_iH_1}^2`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{i\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{i\mathrm{}m}\right]m_{H_1H_1}^2`$
$`+\left[3\lambda _\mathrm{}m^d\lambda _{k\mathrm{}m}^{}+\lambda _\mathrm{}m^e\lambda _{k\mathrm{}m}\right]m_{L_iL_k}^2`$
$`+6\lambda _{nm}^d\lambda _{in\mathrm{}}^{}m_{D_mD_{\mathrm{}}}^2+6\lambda _{mn}^d\lambda _{i\mathrm{}n}^{}m_{Q_mQ_{\mathrm{}}}^2`$
$`+2\lambda _{nm}^e\lambda _{ni\mathrm{}}m_{E_mE_{\mathrm{}}}^2+2\lambda _{mn}^e\lambda _{\mathrm{}in}m_{L_mL_{\mathrm{}}}^2`$
$`+6h_\mathrm{}m^dC_{i\mathrm{}m}^{}+2h_\mathrm{}m^eC_{i\mathrm{}m}`$
$`16\pi ^2{\displaystyle \frac{d}{dt}}m_{H_2H_2}^2`$ $`=`$ $`6\text{Tr}(\lambda ^u\lambda ^u)m_{H_2H_2}^2`$ (B162)
$`+6\lambda _{nm}^u\lambda _\mathrm{}n^um_{Q_mQ_{\mathrm{}}}^2+6\lambda _{mn}^u\lambda _n\mathrm{}^um_{U_mU_{\mathrm{}}}^2`$
$`+6\text{Tr}(h^uh^u)2g_1^2M_1^26g_2^2M_2^2`$
Gauge couplings and gaugino masses run in the following way
$$\begin{array}{cccccc}16\pi ^2\frac{d}{dt}g_1\hfill & =& 11g_1^3\hfill & 16\pi ^2\frac{d}{dt}M_1\hfill & =& 22g_1^2M_1\hfill \\ & & & & & \\ 16\pi ^2\frac{d}{dt}g_2\hfill & =& g_2^3\hfill & 16\pi ^2\frac{d}{dt}M_2\hfill & =& 2g_2^2M_2\hfill \\ & & & & & \\ 16\pi ^2\frac{d}{dt}g_3\hfill & =& 3g_3^3\hfill & 16\pi ^2\frac{d}{dt}M_3\hfill & =& 6g_3^2M_3\hfill \end{array}$$
(B163)
## C The $`\gamma `$-matrix
The QCD-mixing of our new operator basis leads to a $`28\times 28`$-matrix. How its elements are deduced is explained in section III B 2. Fortunately, many entries vanish giving us a chance to derive the eigenvalues and -vectors with the help of *Mathematica*. We split $`\gamma ^{\mathrm{eff}}`$ in the three block mentioned in section III B 2. The four-fermion operators give the block (we have included the ordering of the operators in the first row)
$$\left(\begin{array}{cccccccccccccccccccccccccccc}O_1& O_2& O_3& O_4& O_5& O_6& P_1& P_2& P_3& P_4& P_5& P_6& P_7& P_8& P_9& P_{10}& P_{11}& P_{12}& R_1& R_2& R_3& R_4& R_5& R_6& & & & \\ 2& 6& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 6& 2& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& \frac{22}{9}& \frac{22}{3}& \frac{4}{9}& \frac{4}{3}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& \frac{44}{9}& \frac{4}{3}& \frac{10}{9}& \frac{10}{3}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 2& 6& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& \frac{10}{9}& \frac{10}{3}& \frac{10}{9}& \frac{38}{3}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 0& 0& 0& 16& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \frac{22}{9}& \frac{22}{3}& \frac{4}{9}& \frac{4}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \frac{44}{9}& \frac{4}{3}& \frac{10}{9}& \frac{10}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 6& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \frac{10}{9}& \frac{10}{3}& \frac{10}{9}& \frac{38}{3}& 0& 0& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 6& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 6& 2& 0& 0& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 6& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 6& 2& 0& 0& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 2& 6& & & & \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& 0& 0& 0& 0& 6& 2& & & & \end{array}\right)$$
(C1)
The $`4\times 4`$-block of the magnetic penguins looks like
$$\left(\begin{array}{cccc}O_7& O_8& \stackrel{~}{O}_7& \stackrel{~}{O}_8\\ \frac{32}{3}& 0& 0& 0\\ \frac{32}{9}& \frac{28}{3}& 0& 0\\ 0& 0& \frac{32}{3}& 0\\ 0& 0& \frac{32}{9}& \frac{28}{3}\end{array}\right)$$
(C2)
The four last columns that mix four-fermion operators with magnetic penguins are given as rows. They are
$$\left(\begin{array}{ccccccccccccccc}& O_1& O_2& O_3& O_4& O_5& O_6& P_1& P_2& P_3& P_4& P_5& P_6& & \\ O_7\hfill & 0& \frac{832}{81}& \frac{928}{81}& \frac{272}{81}& \frac{64}{9}& \frac{592}{81}& 0& 0& 0& 0& 0& \frac{400}{81}& & \\ O_8\hfill & 6& \frac{140}{27}& \frac{1090}{27}& \frac{1024}{27}& \frac{118}{3}& \frac{1406}{27}& 0& 0& 0& 0& 0& \frac{238}{27}& & \\ \stackrel{~}{O}_7\hfill & 0& 0& 0& 0& 0& 0& \frac{896}{81}& \frac{896}{81}& \frac{400}{81}& \frac{400}{81}& \frac{400}{81}& 0& & \\ \stackrel{~}{O}_8\hfill & 0& 0& 0& 0& 0& 0& \frac{238}{27}& \frac{238}{27}& \frac{238}{27}& \frac{238}{27}& \frac{454}{27}& 0& & \\ & & & & & & & & & & & & & & \\ & P_7& P_8& P_9& P_{10}& P_{11}& P_{12}& R_1& R_2& R_3& R_4& R_5& R_6& & \\ & \frac{400}{81}& \frac{400}{81}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & \\ & \frac{238}{27}& \frac{454}{27}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& & \\ & 0& 0& \frac{928}{81}& \frac{272}{81}& \frac{64}{9}& \frac{592}{81}& 0& \frac{832}{81}& 0& \frac{832}{81}& 0& \frac{464}{81}& & \\ & 0& 0& \frac{1090}{27}& \frac{1024}{27}& \frac{118}{3}& \frac{1406}{27}& 6& \frac{140}{27}& 6& \frac{140}{27}& 6& \frac{140}{27}& & \end{array}\right)$$
(C3)
We emphasize that the matrix depicted here is $`\gamma ^{\mathrm{eff}}`$. In the HV-scheme it should coincide with the uncorrected $`\gamma `$. We have checked this explicitly for all the entries.
## D Definition of the functions $`F_1`$ \- $`F_4`$
These functions appear in all calculations of diagrams like those of Fig. 6.
$`F_1(x)`$ $`=`$ $`{\displaystyle \frac{1}{12(1x)^4}}\left[2+3x6x^2+x^3+6x\mathrm{log}x\right]`$ (D1)
$`F_2(x)`$ $`=`$ $`{\displaystyle \frac{1}{12(1x)^4}}\left[16x+3x^2+2x^36x^2\mathrm{log}x\right]`$ (D2)
$`F_3(x)`$ $`=`$ $`{\displaystyle \frac{1}{2(1x)^3}}\left[3+4xx^22\mathrm{log}x\right]`$ (D3)
$`F_4(x)`$ $`=`$ $`{\displaystyle \frac{1}{2(1x)^3}}\left[1x^2+2x\mathrm{log}x\right]`$ (D4)
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# The fractional - order controllers: Methods for their synthesis and application
## 1 Introduction
The real objects are generally fractional , however, for many of them the fractionality is very low. A typical example of a non-integer (fractional) order system is the voltage-current relation of a semi-infinite lossy $`RC`$ line or diffusion of the heat into a semi-infinite solid, where heat flow is equal to the half-derivative of temperature :
$$\frac{d^{0.5}T(t)}{dt^{0.5}}=q(t).$$
Their integer-order description can cause significant differences between mathematical model and the real system. The main reason for using integer-order models was the absence of solution methods for fractional-order differential equations. Recently, important achievements were obtained, which enable to take into account real order of dynamic systems.
$`PID`$ controllers belong to dominating industrial controllers and therefore are objects of steady effort for improvements of their quality and robustness. One of the possibilities to improve $`PID`$ controllers is to use fractional-order controllers with non-integer derivation and integration parts. For fractional-order systems the fractional controller $`CRONE`$ has been developed, $`PD^\delta `$ controller and the $`PI^\lambda D^\delta `$ controller has been suggested.
Before the actual design of controllers for dynamical systems it is necessary to identify these systems , and then determine their dynamical properties. The dynamical properties of the system under observation can be expressed mathematically and graphically in the form of various characteristics. These characteristics can be determined from the differential equation or transfer function via computation or experimentally by exciting the system from the equilibrium. Computation of the transfer characteristics of the fractional-order dynamical systems has been the subject of several publications, e.g. by numerical methods , as well as analytical methods . In the synthesis of a controller its parameters are determined according to the given requirements. These requirements are, for example, the stability measure, the accuracy of the regulation process, dynamical properties etc. The check whether the requirements have been met can be done with a simulation on the control circuit model. There are a large number of methods for the design of integer-order controllers, but the situation is worse in the case of fractional-order controllers where the methods are only being worked out. One of the methods being developed is the method (modification of roots locus method) of dominant roots , based on the given stability measure and the damping measure of the control circuit.
## 2 Basic mathematical tools for fractional calculus
The fractional calculus is a generalization of integration and derivation to non-integer order operators. The idea of fractional calculus has been known since the development of the regular calculus, with the first reference probably being associated with Leibniz and L’Hospital in 1695. At first, we generalize the differential and integral operators into one fundamental operator $`D_t^\alpha `$ which is known as fractional calculus:
$${}_{a}{}^{}D_{t}^{\alpha }=\{\begin{array}{cc}\frac{d^\alpha }{dt^\alpha }\hfill & \mathrm{}(\alpha )>0\text{,}\hfill \\ 1\hfill & \mathrm{}(\alpha )=0\text{,}\hfill \\ _a^t(d\tau )^\alpha \hfill & \mathrm{}(\alpha )<0\text{.}\hfill \end{array}$$
The two definitions used for the general fractional differintegral are the Grünwald definition and the Riemann-Liouville (RL) definition . The Grünwald definition is given here
$$_aD_t^\alpha f(t)=\underset{h0}{lim}\frac{1}{h^\alpha }\underset{j=0}{\overset{[\frac{ta}{h}]}{}}(1)^j\left(\genfrac{}{}{0pt}{}{\alpha }{j}\right)f(tjh),$$
(1)
where $`[x]`$ means the integer part of $`x`$. The RL definition is given as
$$_aD_t^\alpha f(t)=\frac{1}{\mathrm{\Gamma }(n\alpha )}\frac{d^n}{dt^n}_a^t\frac{f(\tau )}{(t\tau )^{\alpha n+1}}d\tau ,$$
(2)
for $`(n1<\alpha <n)`$ and where $`\mathrm{\Gamma }(x)`$ is the well known Euler’s gamma function.
The Laplace transform method is used for solving engineering problems. The formula for the Laplace transform of the RL fractional derivative (2) has the form :
$`{\displaystyle _0^{\mathrm{}}}e_0^{pt}D_t^\alpha f(t)𝑑t=`$
$`=p^\alpha F(p){\displaystyle \underset{k=0}{\overset{n1}{}}}p^k{}_{0}{}^{}D_{t}^{\alpha k1}f(t)|_{t=0},`$ (3)
for $`(n1<\alpha n)`$.
For numerical calculation of fractional-order derivation we can use the relation (4) derived from the Grünwald definition (1). This relation has the following form:
$$_{(tL)}D_t^\alpha f(t)h^\alpha \underset{j=0}{\overset{N(t)}{}}b_jf(tjh),$$
(4)
where $`L`$ is the ”memory length”, $`h`$ is the step size of the calculation,
$$N(t)=\text{min}\{\left[\frac{t}{h}\right],\left[\frac{L}{h}\right]\},$$
\[x\] is the integer part of $`x`$ and $`b_j`$ is the binomial coefficient:
$$b_0=1,b_j=\left(1\frac{1+\alpha }{j}\right)b_{j1}.$$
(5)
For the solution of the fractional-order differential equations (FODE) most effective and easy analytic methods were developed based on the formula of the Laplace transform method of the Mittag-Leffler function in two parameters . A two-parameter function of the Mittag-Leffler type is defined by the series expansion:
$$E_{\alpha ,\beta }(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(\alpha k+\beta )},(\alpha ,\beta >0).$$
(6)
The Mittag-Leffler function is a generalization of exponential function $`e^z`$ and the exponential function is a particular case of the Mittag-Leffler function. Here is the relationship given in :
$$E_{1,1}(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(k+1)}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{k!}=e^z.$$
## 3 Fractional - order closed <br>control loop
We will be studying feed-back control system with unit gain in the feed-back loop (Fig.1), where $`G_r(p)`$ is the controller transfer function, $`G_s(p)`$ is controlled system transfer function, $`W(p)`$ is an input, $`E(p)`$ is an error, $`U(p)`$ is output from controller and $`Y(p)`$ is output from system.
The transfer function of close feed-back control circuit (Fig.1) has the form:
$$G_c(p)=\frac{Y(p)}{W(p)}=\frac{G_r(p)G_s(p)}{1+G_r(p)G_s(p)}.$$
(7)
### 3.1 Fractional - order controlled <br>systems
The fractional-order controlled system will be represented with fractional model with transfer function given by the following expression :
$$G_s(p)=\frac{Y(p)}{U(p)}=\frac{1}{a_np^{\beta _n}+\mathrm{}+a_1p^{\beta _1}+a_0p^{\beta _0}},$$
(8)
where $`\beta _k`$, $`(k=0,1,2,\mathrm{},n)`$ are generally real numbers, $`\beta _n>\beta _{n1}>\mathrm{}>\beta _1>\beta _00`$ and $`a_k`$ $`(k=0,1,\mathrm{},n)`$ are arbitrary constants. In the time domain, the transfer function (8) corresponds to the $`n`$-term FODE with constant coefficients
$$a_nD_t^{\beta _n}y(t)+\mathrm{}+a_1D_t^{\beta _1}y(t)+a_0D_t^{\beta _0}y(t)=u(t).$$
(9)
Identification methods for determination of the coefficients $`a_k`$ and $`\beta _k`$, $`(k=0,1,\mathrm{},n)`$ were developed, based on minimization of the difference between the calculated ($`y^c`$) and experimentally measured ($`y^e`$) values
$$Q=\frac{1}{M+1}\underset{m=0}{\overset{M}{}}[y_m^ey_m^c]^2,$$
where $`M`$ is the number of measured values.
For the analytical solution of the $`n`$-term FODE (9) we can write formula in general form :
$`y(t)=`$
$`={\displaystyle \frac{1}{a_n}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^m}{m!}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{k_0+k_1+\mathrm{}+k_{n2}=m}{k_00;\mathrm{},k_{n2}0}}{}}(m;k_0,k_1,\mathrm{},k_{n2})`$
$`{\displaystyle \underset{i=0}{\overset{n2}{}}}\left({\displaystyle \frac{a_i}{a_n}}\right)^{k_i}t^{(\beta _n\beta _{n1})m+\beta _n+_{j=0}^{n2}(\beta _{n1}\beta _j)k_j1}`$
$`E_{\beta _n\beta _{n1},+\beta _n+_{j=0}^{n2}(\beta _{n1}\beta _j)k_j}^{(m)}\left({\displaystyle \frac{a_{n1}}{a_n}}t^{\beta _n\beta _{n1}}\right),`$
where $`E_{\lambda ,\mu }(z)`$ is the Mittag-Leffler function in two parameters (6),
$$E_{\lambda ,\mu }^{(n)}(y)\frac{d^n}{dy^n}E_{\lambda ,\mu }(y)=\underset{j=0}{\overset{\mathrm{}}{}}\frac{(j+n)!y^j}{j!\mathrm{\Gamma }(\lambda j+\lambda n+\mu )},$$
for $`(n=0,1,2,\mathrm{})`$.
For the numerical solution of the $`n`$-term FODE (9) we can write formula in general form:
$$y(k)=\frac{u(k)_{i=1}^n(a_ih^{\beta _i}_{j=1}^kb_jy(kj))}{_{i=0}^na_ih^{\beta _i}b_0},$$
for $`k=1,2,3,\mathrm{}`$, $`y(0)=0`$, $`y(1)=0`$, where $`u(k)`$ is a function on the right side of the differential equation.
### 3.2 Fractional - order controllers
The fractional-order controller will be represented by fractional-order $`PI^\lambda D^\delta `$ controller with transfer function given by the following expression :
$$G_r(p)=\frac{U(p)}{E(p)}=K+T_ip^\lambda +T_dp^\delta ,$$
(10)
where $`\lambda `$ and $`\delta `$ are an arbitrary real numbers $`(\lambda ,\delta 0)`$, $`K`$ is amplification (gain), $`T_i`$ is integration constant and $`T_d`$ is differentiation constant. In the time domain equation (10) has the form:
$$u(t)=Ke(t)+T_iD_t^\lambda e(t)+T_dD_t^\delta e(t).$$
(11)
Taking $`\lambda =1`$ and $`\delta =1`$, we obtain a classical $`PID`$ controller. If $`\lambda =0`$ $`(T_i=0)`$ we obtain a $`PD^\delta `$ controller, etc. All these types of controllers are particular cases of the $`PI^\lambda D^\delta `$ controller. The $`PI^\lambda D^\delta `$ controller (11) is more flexible and gives an opportunity to better adjust the dynamical properties of the fractional-order control system.
## 4 Synthesis of fractional - <br>order controllers
For the design of fractional-order controller a new method was suggested based on the dominant roots principle. This method is based on from poles distribution of the characteristic equation in the complex plane (Fig.2).
Values of dominant roots are designed for the quality requirement of the control circuit. Their significance is, that dominant roots are defined for stability measure $`S_t`$ and damping measure $`T_l`$. Under there conditions the complex conjugate roots satisfy the equation:
$$p_{1,2}=r\pm i\omega .$$
(12)
The parameters of controller were set up so that other poles were from dominant remote to the left side. The parameters design of the fractional-order controller can be divided into two stages:
1. Design of parameter $`K`$
Proportional parameter $`K`$ influence the value of static deviation $`E_t[\%]`$, control time $`T_r[s]`$, and overshoot $`P_r[\%]`$. Generally, with increased value of $`K`$, control time $`T_r[s]`$ is decreasing and static deviation $`E_t[\%]`$ is lower:
$$K(100/E_t)a_0.$$
2. Design of parameters $`T_d`$, $`\delta `$, $`T_i`$, $`\lambda `$
We define the required stability measure $`S_t=r`$ and damping measure $`T_l=\frac{r}{\omega }`$. This requirement is satisfied by the complex conjugate roots (12). We will similarly use characteristic equation as the classical root locus method . The characteristic equation of fractional-order control loop (7) has the form:
$$G_r(p)G_s(p)+1=0.$$
(13)
After substitution of the fractional-order controller transfer function (10) and the fractional-order controlled system transfer function (8) and after some manipulation we obtain characteristic equation (13) in the following form:
$$\underset{k=0}{\overset{n}{}}a_kp^{\beta _k}+(K+T_ip^\lambda +T_dp^\delta )=0.$$
(14)
This algebraic equation is in the general form for the $`n`$-term fractional-order controlled system and fractional-order $`PI^\lambda D^\delta `$ controller. Solution of this equation for known poles and parameter $`K`$ gives unknown parameters $`T_d`$, $`\delta `$, $`T_i`$, $`\lambda `$. The algebraic equation (14) is solved in the space of complex variable. Result are parameters of fractional-order controller for required stability measure and damping measure.
## 5 Control algorithm for <br>fractional - order controllers
The control algorithm was designed according to the control scheme in Fig.1. The position algorithm uses discrete time steps and consists of the following steps :
1. Filtering of required value:
$$w^{}(k)=w^{}(k1)+0.5(w(k)w^{}(k1)),$$
where $`w(k)`$ is required value.
2. Calculating the control error:
$$e(k)=w^{}(k)y(k),$$
for discrete time step $`(k=1,2,\mathrm{})`$, where $`e(k)`$ is regulation error and $`y(k)`$ is measured value.
3. Determination of control value:
$`u(k)=Ke(k)+{\displaystyle \frac{T_i}{T^\lambda }}{\displaystyle \underset{j=v}{\overset{k}{}}}q_je(kj)+`$
$`+{\displaystyle \frac{T_d}{T^\delta }}{\displaystyle \underset{j=v}{\overset{k}{}}}d_je(kj),`$
for discrete time step $`(k=1,2,\mathrm{})`$, where $`T`$ is the length of time step (sample period). The binomial coefficients $`d_j`$ and $`q_j`$ were calculated from the recurrent equation (5). The numerical algorithm (4) requires to store the whole history $`(v=0)`$. For improving their effectiveness we have used ”short memory” principle , where $`v=0`$ for $`k<(L/T)`$ or $`v=k(L/T)`$ for $`k>(L/T)`$. Besides the ”short memory” the control quality is influenced by time step $`T`$.
Fractional-order controllers can be realized as software or by passive or active electrical elements .
## 6 Example
We verify the above methods on an example from . Assume the system which we can describe by three-term $`(n=2)`$ differential equation with coefficients $`a_2=0.8,a_1=0.5,a_0=1,\beta _2=2.2,\beta _1=0.9,\beta _0=0`$. After its approximation with an integer-order system we have $`a_2=0.7414,a_1=0.2313,a_0=1,\beta _2=2,\beta _1=1,\beta _0=0`$. The integer-order $`PD`$ controller designed with the dominant roots method applied to the approximated system has the parameters $`K=20.5,T_d=2.7343,\delta =1`$. By applying the controller to the original system we do not achieve the required quality of the regulation process as with the approximated integer-order system. This was confirmed via simulation in the time domain and also by checking the stability measure and damping measure.
This proves the inadequacy of approximating the fractional-order system with an integer-order system for the purpose of controller design. It is suitable to consider fractional-order systems and also controllers, by means of which it is possible to obtain higher quality regulation and robustness. For the fractional-order $`PD^\delta `$ controller we then have $`K=20.5,T_d=5.79,\delta =0.95`$ and the required quality of regulation is ensured also in the original fractional-order system .
On the Fig.3 is showed the comparison of the unit-step responses the integer-order $`PD`$ controller applied on the fractional-order system (thin line) and the fractional-order $`PD^\delta `$ controller applied on the fractional-order system (thick line).
The dynamical properties of the closed loop with fractional-order controlled system and the fractional-order controller are better than the dynamical properties of the closed loop with the integer-order controller. The systems with the integer-order controller stabilizes slower and has larger surplus oscillations. We can see that use of the fractional-order controller leads to the improvement of the control of the fractional-order system.
## 7 Conclusion
The above methods make it possible to design fractional-order controllers with given measures of stability and damping. The results of previous works also show that fractional-order controllers are more robust , which means they are less sensitive to changes of the system parameters and controller parameters. This can even lead to qualitatively different dynamical phenomena in control circuits.
## Acknowledgement
The author thanks to Dr. Lubomir Dorcak for discussion of this problem and also to Dr. Ladislav Pivka for checking the language of this paper. This work was supported by grant VEGA 1/4333/97 from the Slovak Agency for Science.
Ivo Petras was born in Kosice, Slovak Republic, in 1973. He received the MSc. and Ph.D degrees in process control from the Department of Informatics & Process Control, Faculty B.E.R.G. of the Technical University of Kosice. His main research interests include fractional calculus in robust control and process control.
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# On the Expected Value of the Minimum Assignment
## 1 Problem Description and Background
Suppose that $`k`$, $`m`$ and $`n`$ are positive integers with $`kmn`$. A minimum $`k`$-assignment of an $`m\times n`$ matrix $`X`$ is a set of $`k`$ entries of $`X`$, no two of which belong to the same row or column, whose sum is as small as possible. We denote the value of this minimum sum by $`\mathrm{min}_k(X)`$.
We say that a random real number $`x`$ is exponentially distributed with rate $`a`$ if it is chosen according to the density $`ae^{ax}`$, $`x0`$. The mean value of a rate $`a`$ exponentially distributed quantity is $`1/a`$.
Suppose that we generate a random $`m\times n`$ matrix $`X`$ by choosing each entry independently from the exponential distribution with rate 1. In \[CS\] Coppersmith and Sorkin conjectured that the expected value of its minimum $`k`$-assignment is
###### Conjecture 1
$$E(\underset{k}{\mathrm{min}}(X))=\underset{i,j0,i+j<k}{}\frac{1}{(mi)(nj)}.$$
(1)
In \[AS\] Alm and Sorkin show that this conjecture is correct when $`k4`$, when $`k=m=5`$, and when $`k=m=n=6`$.
The conjecture of Coppersmith and Sorkin generalized a conjecture of Parisi \[P\] who considered the case $`k=m=n`$. In this case, as shown in \[CS\], (1) reduces to
$$E(\underset{k}{\mathrm{min}}(X))=\underset{i=1}{\overset{k}{}}\frac{1}{i^2}.$$
(2)
In this paper we describe our efforts to prove these conjectures. Our main result is Conjecture 2, which generalizes the Coppersmith-Sorkin conjecture.
We will say that a matrix $`X`$ is random exponential with rate matrix $`A=(a_{ij})`$ if each entry $`x_{ij}`$ is chosen independently according to the exponential distribution with rate $`a_{ij}`$. The expected value of the minimum $`k`$-assignment of such a matrix $`X`$ is then a function of the rate matrix $`A`$. We denote this function by $`E_k(A)`$.
We will show that $`E_k(A)`$ is a rational function of the rates $`a_{ij}`$ and give an explicit method for computing it, at least in principle. Then we will specialize to the case when the rate matrix has rank 1, for which we have the following explicit formula.
###### Conjecture 2
Suppose that $`r_1,\mathrm{},r_m`$ and $`c_1,\mathrm{},c_n`$ are positive real numbers and that $`a_{ij}=r_ic_j`$. Let $`X`$ be a random $`m\times n`$ matrix in which entry $`x_{ij}`$ is chosen independently from the exponential distribution with rate $`a_{ij}`$. Then the expected value of the minimum $`k`$-assignment of $`X`$ is
$$\underset{I,J}{}(1)^{k1|I||J|}\left(\genfrac{}{}{0pt}{}{m+n1|I||J|}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}.$$
Here the sum is over proper subsets $`I`$ of $`\{1,\mathrm{},m\}`$ and $`J`$ of $`\{1,\mathrm{},n\}`$ whose cardinalities $`|I|`$ and $`|J|`$ satisfy $`|I|+|J|<k`$.
###### Example 1
The expected value of the minimum 1-assignment of a random exponential matrix with rate matrix $`a_{ij}=r_ic_j`$ is
$$\frac{1}{(_ir_i)(_jc_j)}.$$
###### Example 2
The expected value of the minimum 2-assignment of a $`3\times 3`$ random exponential matrix with rate matrix $`a_{ij}=r_ic_j`$ is
$`\left({\displaystyle \frac{1}{r_2+r_3}}+{\displaystyle \frac{1}{r_1+r_3}}+{\displaystyle \frac{1}{r_1+r_2}}\right){\displaystyle \frac{1}{c_1+c_2+c_3}}`$
$`+`$ $`\left({\displaystyle \frac{1}{c_2+c_3}}+{\displaystyle \frac{1}{c_1+c_3}}+{\displaystyle \frac{1}{c_1+c_2}}\right){\displaystyle \frac{1}{r_1+r_2+r_3}}`$
$``$ $`{\displaystyle \frac{5}{(r_1+r_2+r_3)(c_1+c_2+c_3)}}`$
We will provide evidence in support of Conjecture 2.
We also have a stronger conjecture for which we will provide evidence, although perhaps this evidence is not as strong as that for Conjecture 2. A matrix can have several minimum $`k`$-assignments for some value of $`k`$. However, with probability 1, a random matrix has a single minimum $`k`$-assignment for each $`k`$. Suppose that $`M`$ is a $`(k1)\times (k1)`$ submatrix of $`X`$ and that $`\chi _M(X)`$ is the function with value 1 when $`M`$ contains a minimum $`(k1)`$-assignment of $`X`$ and 0 otherwise. Then we define the expected contribution of $`M`$ to the minimum $`k`$-assignment of $`X`$ as the expected value of the random variable $`\chi _M(X)\mathrm{min}_k(X)`$. It is clear that $`E(\mathrm{min}_k(X))`$ is the sum of the expected contributions of all the $`(k1)\times (k1)`$ submatrices. Our stronger conjecture gives a formula for the expected contribution of $`M`$ when the rate matrix has rank 1.
###### Conjecture 3
Suppose that $`A=(r_ic_j)`$ is a positive $`m\times n`$ matrix with rank 1 and that $`X`$ is a random exponential matrix with rate matrix $`A`$. Let $`I`$ be a set of $`k1`$ elements of $`\{1,\mathrm{},m\}`$, let $`J`$ be a set of $`k1`$ elements of $`\{1,\mathrm{},n\}`$, and let $`M`$ be the $`(k1)\times (k1)`$ submatrix of $`X`$ with rows indexed by $`I`$ and columns indexed by $`J`$. Then the expected value of $`\chi _M(X)\mathrm{min}_k(X)`$ is
$$\underset{i,j}{}\left(\underset{t=1}{\overset{k1}{}}\frac{r_{i_t}c_{j_t}}{\left(R_{s=1}^{t1}r_{i_s}\right)\left(C_{s=1}^{t1}c_{j_s}\right)}\right)\underset{t,u0,t+u<k}{}\frac{1}{(R_{s=1}^tr_{i_s})(C_{s=1}^uc_{j_s})}$$
where the outer sum is over permutations $`(i_1,\mathrm{},i_{k1})`$ and $`(j_1,\mathrm{},j_{k1})`$ of $`I`$ and $`J`$, respectively, and $`R`$ and $`C`$ denote the sums of all $`r_i`$’s and all $`c_j`$’s, respectively.
We shall see that Conjecture 3 implies Conjecture 2 and that Conjecture 2 in turn implies Conjecture 1.
Section 2 discusses what we know for the expected minimum assignment when the rate matrix is arbitrary.
In Section 3 we discuss the way we arrived at Conjecture 2 and give some equivalent formulations, one of which is directly implied by Conjecture 3.
We discuss the computational evidence for our conjectures in Section 4.
Section 5 gives additional evidence for Conjecture 2.
We would like to thank Jim Propp for bringing this problem to our attention.
## 2 Theory for a general rate matrix
### 2.1 Expected value for a general rate matrix
We begin by showing that the general formula for the expected value of the minimum assignment of a random exponential matrix is a rational function of the rates, with denominators factoring into linear terms of special form.
Recall that $`k,m,n`$ are positive integers with $`kmn`$ and that $`A=(a_{ij})`$ is a positive $`m\times n`$ matrix. We form a random matrix $`X`$ by choosing $`x_{ij}`$ independently from the exponential distribution with rate $`a_{ij}`$. The expected value of the minimum $`k`$-assignment of $`X`$ is then a function of $`A`$, which we will denote by $`E_k(A)`$.
By definition of expected value, $`E_k(A)`$ is given by the integral expression
$$E_k(A)=\left(\underset{i,j}{}a_{ij}\right)_{X0}\underset{k}{\mathrm{min}}(X)e^{AX}𝑑X$$
where the integral is taken over the space of all nonnegative matrices $`X`$. Here $`AX`$ denotes the dot product $`_{i,j}a_{ij}x_{ij}`$ and $`dX`$ denotes the product $`_{i,j}dx_{ij}`$.
We denote by $`𝒮_k`$ the set of all $`m\times n`$ matrices $`\sigma `$ such that all the entries of $`\sigma `$ are 0’s except for $`k`$ entries which are 1’s, no two in the same row or column. There are $`k!\left(\genfrac{}{}{0pt}{}{m}{k}\right)\left(\genfrac{}{}{0pt}{}{n}{k}\right)`$ such matrices in $`𝒮_k`$ and these we identify in the obvious way with the possible locations of the minimum $`k`$-assignment of $`X`$. In particular,
$$\underset{k}{\mathrm{min}}(X)=\underset{\sigma 𝒮_k}{\mathrm{min}}\left(\sigma X\right).$$
For each $`\sigma `$ we denote by $`P_\sigma `$ the set of nonnegative matrices $`X`$ for which $`\mathrm{min}_k(X)=\sigma X`$; that is, $`P_\sigma `$ is the set of nonnegative matrices $`X`$ for which the minimum $`k`$-assignment is $`\sigma `$. Thus, we have
$$E_k(A)=\left(\underset{i,j}{}a_{ij}\right)\underset{\sigma 𝒮_k}{}\left[_{XP_\sigma }(\sigma X)e^{AX}𝑑X\right].$$
(3)
Note that each of the sets $`P_\sigma `$ is a polyhedral cone determined by a finite set of homogeneous linear inequalities $`\sigma X\tau X`$ for all $`\tau 𝒮_k`$. As a consequence, each $`P_\sigma `$ can be decomposed into a finite collection $`𝒞_\sigma `$ of simplicial cones. It seems difficult to give an explicit description of $`𝒞_\sigma `$. Nevertheless, we can derive some useful properties of $`E_k(A)`$ from the fact that this decomposition exists. First we rewrite (3) as
$$E_k(A)=\left(\underset{i,j}{}a_{ij}\right)\underset{\sigma 𝒮_k}{}\underset{C𝒞_\sigma }{}\left[_C(\sigma X)e^{AX}𝑑X\right].$$
Each cone $`C`$ is the set of nonnegative linear combinations of a set of $`mn`$ linearly independent vectors $`V_i`$, $`i=1,\mathrm{},mn`$, where each $`V_i`$ is a nonnegative $`m\times n`$ matrix. For the part of the integral over $`C`$, we make the substitution $`X=_iu_iV_i`$, where $`U=(u_1,\mathrm{},u_{mn})`$ ranges over all nonnegative $`mn`$-tuples. We can then explicitly compute the integral over $`C`$ as
$`{\displaystyle _C}(\sigma X)e^{AX}𝑑X`$ $`=`$ $`|detV|{\displaystyle _{U0}}{\displaystyle \underset{i=1}{\overset{mn}{}}}u_i(\sigma V_i)e^{_{j=1}^{mn}u_jAV_j}dU`$
$`=`$ $`|detV|{\displaystyle \underset{i=1}{\overset{mn}{}}}\left[(\sigma V_i){\displaystyle _{U0}}u_ie^{_{j=1}^{mn}u_jAV_j}𝑑U\right]`$
$`=`$ $`|detV|\left({\displaystyle \underset{i=1}{\overset{mn}{}}}{\displaystyle \frac{\sigma V_i}{AV_i}}\right)\left({\displaystyle \underset{i=1}{\overset{mn}{}}}{\displaystyle \frac{1}{AV_i}}\right)`$
where $`|detV|`$ is the $`mn`$-volume of the parallelepiped determined by $`V_1,\mathrm{},V_{mn}`$. Thus, we obtain the expression
$$E_k(A)=\left(\underset{i,j}{}a_{ij}\right)\underset{\sigma 𝒮_k}{}\underset{C𝒞_\sigma }{}|detV|\left(\underset{i=1}{\overset{mn}{}}\frac{\sigma V_i}{AV_i}\right)\left(\underset{i=1}{\overset{mn}{}}\frac{1}{AV_i}\right).$$
(4)
Note that although the vectors $`V_i`$ depend on $`C`$ and $`\sigma `$, they do not depend on $`A`$. Thus, $`E_k(A)`$ is a rational function of the $`a_{ij}`$’s, homogeneous of degree $`1`$.
We can obtain more information about the rational function $`E_k(A)`$ by constructing, for each $`\sigma 𝒮_k`$, a finite set of generators of $`P_\sigma `$ in the sense that every element of $`P_\sigma `$ is a nonnegative linear combination of the generators.
For this purpose we define two classes of matrices. First, for any $`1im`$ and $`1jn`$ we define $`e_{ij}`$ to be the matrix that is all zero except for a single 1 at position $`(i,j)`$. Next, for any sets $`I\{1,\mathrm{},m\}`$ and $`J\{1,\mathrm{},n\}`$, we define $`V_{IJ}`$ to be the matrix obtained from the all 1’s matrix by zeroing out all entries in the rows indexed by $`I`$ and the columns indexed by $`J`$. It is easy to see that $`\mathrm{min}_k(V_{IJ})=\mathrm{max}(0,k|I||J|)`$. Thus $`V_{IJ}`$ is in $`P_\sigma `$ if and only if $`\sigma V_{IJ}=\mathrm{max}(0,k|I||J|)`$.
###### Theorem 1
Every element of $`P_\sigma `$ is a nonnegative linear combination of $`e_{ij}`$’s with $`e_{ij}\sigma =0`$ and $`V_{IJ}`$’s in $`P_\sigma `$ with $`|I|+|J|<k`$.
We prove Theorem 1 using a reduction procedure on the matrices of $`P_\sigma `$. Let $`X`$ be a matrix in $`P_\sigma `$ and suppose that $`\mathrm{min}_k(X)=s`$. We choose an arbitrary linear order for the $`e_{ij}`$’s and denote this ordered set by $`e_1,e_2,\mathrm{},e_{mn}`$. Then we choose a sequence of nonnegative real numbers $`\alpha _1,\alpha _2,\mathrm{},\alpha _{mn}`$ as follows. Once $`\alpha _1,\mathrm{},\alpha _{i1}`$ are chosen we select $`\alpha _i`$ as large as possible so that $`X(\alpha _1e_1+\alpha _2e_2+\mathrm{}+\alpha _ie_i)`$ is nonnegative and has minimum $`k`$-assignment with value $`s`$.
Set
$$Y=X(\alpha _1e_1+\alpha _2e_2+\mathrm{}+\alpha _{mn}e_{mn}).$$
Note that if $`e_i\sigma 0`$, then we will have $`\alpha _i=0`$, since otherwise $`\sigma (X(\alpha _1e_1+\alpha _2e_2+\mathrm{}+\alpha _ie_i))<s.`$ Thus $`X`$ is $`Y`$ plus a nonnegative linear combination of the $`e_{ij}`$’s given in Theorem 1.
We say that an entry $`y_{ij}`$ of a matrix $`Y`$ *participates* in a minimum $`k`$-assignment if there is a minimum $`k`$-assignment using the entry $`y_{ij}`$.
We say that a nonnegative matrix $`Y=(y_{ij})`$ is *$`k`$-reduced* if every nonzero entry of $`Y`$ participates in a minimum $`k`$-assignment.
It is straightforward to see that the matrix $`Y`$ resulting from our reduction process applied to $`XP_\sigma `$ is $`k`$-reduced and that $`YP_\sigma `$. It remains to show that every $`k`$-reduced matrix $`Y`$ with minimum $`k`$-assignment $`\sigma `$ is a nonnegative linear combination of the appropriate $`V_{IJ}`$. This will require a series of preliminary results.
First we need a simple combinatorial lemma.
###### Lemma 1
Suppose that $`T`$ is a matrix all of whose entries are 0, 1, or 2 and whose row and column sums are at most 2, and that the sum of all the entries in $`T`$ is $`2k`$. Then $`T=\sigma +\tau `$ for some $`k`$-assignments $`\sigma `$ and $`\tau `$.
Proof: We may assume there are no 2’s in $`T`$, since if there is a 2 we know that both $`\sigma `$ and $`\tau `$ must have a 1 there, and only 0’s everywhere else in its row and column. So we assume $`T`$ is a 0-1 matrix whose row and column sums are at most 2, such that the sum of all entries is $`2s`$ for some $`sk`$, and we want to find two $`s`$-assignments $`\sigma `$ and $`\tau `$ such that $`\sigma +\tau =T`$. Identify $`T`$ with a graph with vertices at each 1 of $`T`$, and edges between any two 1’s belonging to the same row or column. Clearly every vertex of $`T`$ has degree $`2`$, so every component of $`T`$ is a chain or a cycle. The vertices in each component can be alternately assigned to $`\sigma `$ and $`\tau `$. If there is an odd component (which must be a chain) there must be another odd component to balance it out (so one can have an extra $`\sigma `$ vertex, and the other can have an extra $`\tau `$ vertex), since the sum of all entries in $`T`$ is even. $`\mathrm{}`$
###### Lemma 2
Suppose that $`Y`$ is a $`k`$-reduced matrix and
$$S=\left[\begin{array}{cc}\hfill a& \hfill b\\ \hfill c& \hfill d\end{array}\right]$$
is a submatrix of $`Y`$. If $`a`$ and $`d`$ each participate in a minimum $`k`$-assignment of $`Y`$, then $`a+db+c`$. If also $`a+d=b+c`$, then $`b`$ and $`c`$ also each participate in a minimum $`k`$-assignment. These statements also hold for $`a`$ and $`d`$ switched with $`b`$ and $`c`$.
Proof: Suppose that $`a+d>b+c`$. Let $`\sigma _1`$ and $`\tau _1`$ be minimum assignments passing through $`a`$ and $`d`$ respectively. Form the matrix $`T_1=\sigma _1+\tau _1`$. Then form $`T`$ from $`T_1`$ by subtracting 1 at the positions of $`a`$ and $`d`$ and adding 1 at the positions of $`b`$ and $`c`$. The hypotheses of Lemma 1 apply to $`T`$ so that $`T=\sigma +\tau `$ for some $`k`$-assignments $`\sigma `$ and $`\tau `$. But since $`a+d>b+c`$ we must have
$$\sigma _1Y+\tau _1Y=T_1Y>TY=\sigma Y+\tau Y,$$
contradicting the minimality of the assignments $`\sigma _1`$ and $`\tau _1`$. Thus $`a+db+c`$.
Now suppose that $`b+c=a+d`$. Then the same construction yields
$$\sigma _1Y+\tau _1Y=T_1Y=TY=\sigma Y+\tau Y,$$
so that both $`\sigma `$ and $`\tau `$ are minimum $`k`$-assignments, and at least one includes $`b`$ and at least one includes $`c`$.
This proof obviously also holds with $`a`$ and $`d`$ switched with $`b`$ and $`c`$. $`\mathrm{}`$
###### Proposition 1
Suppose that $`Y=(y_{ij})`$ is a $`k`$-reduced $`m\times n`$ matrix. Then there exist $`\lambda _1,\mathrm{},\lambda _m`$ and $`\mu _1,\mathrm{},\mu _n`$ such that
$$y_{ij}=\mathrm{max}(0,\lambda _i+\mu _j)$$
(5)
and such that $`y_{ij}`$ participates in a minimum $`k`$-assignment precisely when $`y_{ij}=\mu _i+\lambda _j`$.
Proof: Let $`d=y_{tu}`$ denote the largest entry in $`Y`$. Take $`\lambda _i`$ to be the $`i^{th}`$ entry in the column of $`d`$, so that $`\lambda _i=y_{iu}`$. Let $`\mu _j`$ to be the $`j^{th}`$ entry in the row of $`d`$, decreased by $`d`$, so $`\mu _j=y_{tj}d`$.
First we prove (5).
When $`y_{ij}`$ is in the row or column of $`d`$, then we have $`y_{ij}=\lambda _i+\mu _j`$, so (5) is immediate.
Suppose that $`a`$ is in neither the row of $`d`$ nor the column of $`d`$. Let
$$S=\left[\begin{array}{cc}\hfill a& \hfill b\\ \hfill c& \hfill d\end{array}\right]$$
be the submatrix of $`Y`$ containing the rows and columns of $`a`$ and $`d`$ (where the order of the rows or columns in $`S`$ may be opposite to the order they occur in $`Y`$). Then we just need to show that $`a=\mathrm{max}(0,b+cd)`$.
First suppose that $`b+cd>0`$. Then $`b,c,d`$ must all be positive, because $`d`$ is maximum for the whole matrix. Since $`b,c`$ are both positive, they must both participate in a minimum $`k`$-assignment, so $`b+ca+d`$, by Lemma 2. But then, $`ab+cd`$, so $`a`$ is positive. Then $`a`$ and $`d`$ both participate in a minimum $`k`$-assignment, which implies $`a+db+c`$, so $`a=b+cd`$.
Next, suppose that $`b+cd0`$. If $`a>0`$, then $`a`$ and $`d`$ both participate in a minimum $`k`$-assignment, so $`d<a+db+c`$, a contradiction. Thus $`a=0`$.
Now we show that $`y_{ij}=\lambda _i+\mu _j`$ exactly when $`y_{ij}`$ participates in a minimum $`k`$-assignment.
Recall that for any $`y_{ij}`$ in either the row or column of $`d`$, we have $`y_{ij}=\lambda _i+\mu _j`$. So we need to show that all such entries participate in a minimum $`k`$-assignment. Let $`b`$ be any entry in the column of $`d`$. We already know that positive entries must participate in a minimum $`k`$-assignment, so we assume that $`b=0`$. If $`d=0`$, then our whole matrix is zero and our result is trivial, so we may assume that $`d>0`$ and therefore participates in a minimum $`k`$-assignment. If the minimum assignment using $`d`$ does not use the row of $`b`$, we can replace $`d`$ by $`b`$ and obtain a smaller assignment, a contradiction. So we can conclude that the minimum assignment using $`d`$ also uses an element $`a`$ from the row of $`b`$. Form the $`2\times 2`$ submatrix $`S`$ containing $`a`$ and $`d`$ as above. Since $`d`$ is the largest entry, $`cd`$, and we also have $`b=0a`$. Thus we can exchange $`a`$ and $`d`$ for $`b`$ and $`c`$, to obtain a minimum $`k`$-assignment in which $`b`$ participates. In the same way, we see that any entry in the row of $`d`$ must participate in a minimum $`k`$-assignment.
Finally, consider an element $`a`$ that is neither in the row nor the column of $`d`$ and form the $`2\times 2`$ submatrix $`S`$ containing $`a`$ and $`d`$ as above. We must show that $`a`$ participates in a minimum assignment exactly when $`a=b+cd`$. This is certainly true if $`a>0`$. So, let us assume that $`a=0`$. Also, since both $`b`$ and $`c`$ are in a row or a column of $`d`$, both participate in a minimum assignment, so that $`b+ca+d`$.
Now suppose that $`a`$ participates in a minimum $`k`$-assignment. Then $`a+db+c`$, so $`a=b+cd`$, as required.
Conversely, suppose that $`a=b+cd`$. Then, since $`b`$ and $`c`$ participate in minimum $`k`$-assignments, Lemma 2 shows that $`a`$ also participates. $`\mathrm{}`$
Proof of Theorem 1: Now we are ready to prove Theorem 1 by showing that any $`k`$-reduced matrix $`Y`$ in $`P_\sigma `$ is a nonnegative linear combination of a suitable collection of matrices $`V_{IJ}`$ from $`P_\sigma `$. Without loss of generality we can assume that the $`\lambda `$’s and $`\mu `$’s are weakly increasing. In this case the rows and columns of $`Y`$ are also weakly increasing.
If the matrix $`Y`$ is zero, there is nothing to prove. Otherwise, since every nonzero entry in $`Y`$ participates in a minimum $`k`$-assignment, we know that the minimum $`k`$-assignment is nonzero. Hence there is at least one nonzero entry among $`y_{1,k},y_{2,k1},\mathrm{},y_{k,1}`$. In particular, there is a pair $`(i,j)`$ such that $`y_{ij}>0`$ and $`i+jk+1`$. Now select such a pair $`(i,j)`$ to be minimal in the sense that if $`i1`$ then $`y_{i1,j}=0`$ and if $`j1`$ then $`y_{i,j1}=0`$. Let $`I=\{1,\mathrm{},i1\}`$ and let $`J=\{1,\mathrm{},j1\}`$. We will show that $`V_{IJ}`$ is in $`P_\sigma `$ and that $`Yy_{ij}V_{IJ}`$ is again in $`P_\sigma `$ and still $`k`$-reduced.
Suppose that $`1i^{}<i`$ and $`1j^{}<j`$. Then, since
$$\mathrm{max}(0,\lambda _i^{}+\mu _j)=y_{i^{}j}=0<y_{ij}=\mathrm{max}(0,\lambda _i+\mu _j),$$
we know that $`\lambda _i^{}<\lambda _i`$. Similarly $`\mu _j^{}<\mu _j`$. Since $`y_{i^{}j}=\mathrm{max}(0,\lambda _i^{}+\mu _j)=0`$, we have $`\lambda _i^{}+\mu _j0`$. But then
$$\lambda _i^{}+\mu _j^{}<\lambda _i^{}+\mu _j0.$$
If follows from Proposition 1 that none of the matrix entries $`y_{i^{}j^{}}`$ with $`i^{}<i`$ and $`j^{}<j`$ can participate in a minimum $`k`$-assignment.
To see that $`V_{IJ}P_\sigma `$, first note that any minimum $`k`$-assignment of $`Y`$ must use all of the first $`i1`$ rows. If not, since $`ik`$, there is some $`i_1>i`$ such that row $`i_1`$ participates. Then, since not all of the first $`i1`$ rows are used, we can replace the entry of the assignment in row $`i_1`$ with the entry in the same column of row $`i_0`$, for some $`i_0<i`$, to get an assignment with a value no larger. The entry being replaced could not come from a column preceding $`j`$, by the discussion above, so it must be positive. But also by the discussion above, $`\lambda _{i_0}<\lambda _{i_1}`$. Then the new assignment would be strictly smaller, a contradiction. Thus, any minimum assignment uses all of the first $`i1`$ rows and all of the first $`j1`$ columns, and does not use any entry which is in both the first $`i1`$ rows and the first $`j1`$ columns. So, if $`\tau `$ is any matrix representing a minimum $`k`$-assignment of $`Y`$, we must have $`V_{IJ}\tau =k|I||J|`$. In particular, $`V_{IJ}\sigma =k|I||J|`$, so $`V_{IJ}`$ is in $`P_\sigma `$ as claimed.
The preceding argument shows that if we replace $`Y`$ by $`Y^{}=YtV_{IJ}`$, for any $`t`$ satisfying $`0ty_{ij}`$, the effect on any minimum $`k`$-assignment is to subtract $`(k|I||J|)t`$ from its value. Thus, all assignments $`\tau `$ that are minimum for $`Y`$ will agree on $`Y^{}`$. We now show that each of these assignments $`\tau `$ is minimum for $`Y^{}`$ as well. Assume not. Then we could find $`t_1`$ and $`t_2`$ such that $`0t_1<t_2y_{ij}`$ and $`k`$-assignments $`\tau `$ and $`\varphi `$ such that $`\tau `$ is minimum for $`(YtV_{IJ})`$ when $`tt_1`$ but not when $`t_1<tt_2`$, and $`\varphi `$ is minimum for $`YtV_{IJ}`$ when $`t_1tt_2`$. Thus, we would have both $`\varphi `$ and $`\tau `$ minimum assignments for $`Yt_1V_{IJ}`$. But then our preceding argument applied to $`Yt_1V_{IJ}`$ tells us that $`\varphi `$ and $`\tau `$ must agree on $`YtV_{IJ}`$ for $`t_1tt_2`$, a contradiction.
Now we let $`Y^{}=Yy_{ij}V_{IJ}`$ and observe that $`Y^{}`$ is $`k`$-reduced. Indeed, all the assignments $`\tau `$ that were minimum for $`Y`$ are also minimum for $`Y^{}`$. Thus, any element that participated in a minimum $`k`$-assignment for $`Y`$ will also participate in a minimum $`k`$-assignment for $`Y^{}`$. Also, the replacement of $`Y`$ by $`Y^{}`$ creates no new nonzero elements, so the matrix $`Y^{}`$ will be $`k`$-reduced.
Since $`\sigma `$ in particular is a minimum assignment for $`Y`$, $`\sigma `$ will also be minimum for $`Y^{}`$, so that $`Y^{}`$ is also in $`P_\sigma `$.
Note that $`Y^{}`$ has at least one more zero entry than $`Y`$, namely the entry at $`(i,j)`$.
We can continue removing multiples of submatrices $`V_{IJ}`$, each time producing a matrix with at least one more zero entry. Thus, we eventually reach a matrix that is all zero. In effect, we have expressed $`Y`$ as a nonnegative linear combination of the generators as required. $`\mathrm{}`$
We have shown that every element of $`P_\sigma `$ is a nonnegative linear combination of certain $`e_{ij}`$’s and $`V_{IJ}`$’s in $`P_\sigma `$. We remark that the $`e_{ij}`$’s generate extreme rays of $`P_\sigma `$ but the $`V_{IJ}`$’s in general do not. It is not hard to show that the $`V_{IJ}`$’s which do generate extreme rays of $`P_\sigma `$ are those with $`|I|+|J|=k1`$.
We now observe that Theorem 1 allows us to make some conclusions about the rational function $`E_k(A)`$. A simplicial cone in a decomposition of $`P_\sigma `$ has some generators of the form $`e_{ij}`$ and some of the form $`V_{IJ}`$. For the generators of the form $`e_{ij}`$, we know that $`\sigma e_{ij}=0`$. Also the dot products $`e_{ij}V`$ in the denominator cancel with factors in the initial product of $`a_{ij}`$’s in (4). Thus, the denominator of the integral over a simplicial cone is a product of terms of the form $`AV_{IJ}`$. Finally, we can conclude:
###### Theorem 2
The expected minimum $`k`$-assignment of a random exponential matrix with rate matrix $`A`$ is a rational function of the entries of $`A`$. The denominator of the rational function is a product of sums, each being the sum of all entries in a submatrix of $`A`$ omitting $`i`$ rows and $`j`$ columns, where $`i+j<k`$.
###### Example 3
The minimum 2-assignment of a random $`2\times 2`$ exponential matrix with rate matrix $`A=(a_{ij})`$ is
$`{\displaystyle \frac{1}{a_{11}+a_{12}}}`$ $`+`$ $`{\displaystyle \frac{1}{a_{21}+a_{22}}}+{\displaystyle \frac{a_{11}a_{21}}{(a_{11}+a_{12})(a_{21}+a_{22})(a_{12}+a_{22})}}`$
$`+`$ $`{\displaystyle \frac{a_{12}a_{22}}{(a_{11}+a_{12})(a_{21}+a_{22})(a_{11}+a_{21})}}`$
The above formula is easily computed by the method we will sketch in Section 4.
In the case of a rank 1 rate matrix $`A=(r_ic_j)`$, Theorem 2 says that $`E_k(A)`$ is a rational function of the $`r_i`$’s and $`c_j`$’s and its denominator is a product of sums of subsets of the $`r_i`$’s omitting fewer than $`k`$ of the $`r_i`$’s and sums of subsets of the $`c_j`$’s omitting fewer than $`k`$ of the $`c_j`$’s. Since this is a consequence of Conjecture 2, it lends some support to the conjecture.
### 2.2 The nesting lemma
A real nonnegative $`m\times n`$ matrix $`X`$ has minimum $`k`$-assignments for each $`km`$. Generically there is only one of each but in some cases there are many minimum assignments of various sizes. It helps to know how these are related.
The following lemma is fundamental. Other proofs probably exist but we include ours here for completeness.
###### Lemma 3
Let $`k_1`$ and $`k_2`$ be two integers, with $`k_1k_2m`$.
Suppose that $`M_1`$ is a $`k_1\times k_1`$ submatrix of $`X`$ that contains a minimum $`k_1`$-assignment of $`X`$. Then there exists a $`k_2\times k_2`$ submatrix $`M_2`$ containing $`M_1`$ such that $`M_2`$ contains a minimum $`k_2`$-assignment of $`X`$.
Suppose that $`M_2`$ is a $`k_2\times k_2`$ submatrix of $`X`$ that contains a minimum $`k_2`$-assignment of $`X`$. Then there exists a $`k_1\times k_1`$ submatrix $`M_1`$ contained in $`M_2`$ such that $`M_1`$ contains a minimum $`k_1`$-assignment of $`X`$.
Proof: If $`k_1=k_2`$, there is nothing to prove. So assume $`k_1<k_2`$ and fix a minimum $`k_1`$-assignment and a minimum $`k_2`$-assignment.
Let $`G`$ be the graph on $`k_1`$ red vertices (representing the entries of the $`k_1`$-assignment) and $`k_2`$ blue vertices (representing the entries of the $`k_2`$-assignment), with edges between two vertices if the corresponding entries belong to the same row or column. (If the assignments share an entry then we have a red vertex and a blue vertex with two edges between them comprising a component which is a cycle of length 2.)
Then $`G`$ is bipartite and no vertex of $`G`$ has degree more than 2. Thus, every component of $`G`$ is a cycle or a chain in which the red and blue vertices alternate.
Suppose some component of $`G`$ has $`m_1`$ red vertices and $`m_2`$ blue vertices, and all of its red vertices have degree 2. Such a component is either a cycle, so that $`m_1=m_2`$, or a chain with blue vertices at each end, so that $`m_1+1=m_2`$. In either case we have $`m_1+1m_2`$.
Consider the submatrix $`M_1`$ spanned by the associated $`m_1`$ entries of the $`k_1`$-assignment and the submatrix $`M_2`$ spanned by the associated $`m_2`$ entries of the $`k_2`$-assignment. The component condition translates into the condition that the $`M_1`$ is contained in $`M_2`$. Also the remaining entries of the two assignments comprise a minimum $`(k_1m_1)`$-assignment and a minimum $`(k_2m_2)`$-assignment of the submatrix of $`X`$ complementary to $`M_2`$. Since $`k_1<k_2`$ and $`m_1+1m_2`$, we have $`k_1m_1k_2m_2`$, so the lemma follows by induction applied to the complementary submatrix.
Now assume that every component of $`G`$ has a red vertex of degree 1 or less. Such components are chains with one endpoint red. When the other endpoint is red, there are more red than blue vertices in the component. When the other endpoint is blue, the number of vertices of both colors is equal. In particular the number of red vertices is always at least as great as the number of blue vertices. It follows that, all together, there are at most $`k_1`$ blue vertices that are connected to some red vertex. Thus we can select $`k_2k_1`$ entries of the $`k_2`$-assignment that do not share any row or column with the $`k_1`$-assignment. We can now consider three sets of matrix entries, the set $`S_1`$ of $`k_1`$ entries of the $`k_1`$-assignment, the set $`S_2`$, just selected, of $`k_2k_1`$ of entries of the $`k_2`$-assignment and the set $`S_3`$ of the remaining $`k_1`$ entries of the $`k_2`$-assignment. Then the sets $`S_3`$ and $`S_1`$ are both in the submatrix complementary to that determined by $`S_2`$. Moreover the sums of the entries in $`S_1`$ and $`S_3`$ must be equal. For, if the sum of the $`S_1`$ entries were greater than that of $`S_3`$, then $`S_3`$ would be a smaller $`k_1`$-assignment than $`S_1`$, a contradiction. But if the sum of the $`S_1`$ entries were smaller than the sum of $`S_3`$ entries, then $`S_1S_2`$ would be give a smaller $`k_2`$-assignment than $`S_2S_3`$, which was a minimum $`k_2`$-assignment.
Thus $`S_3`$ is a minimum $`k_1`$-assignment contained in our minimum $`k_2`$-assignment and $`S_1S_2`$ is a minimum $`k_2`$-assignment containing our minimum $`k_1`$-assignment, which proves our lemma. $`\mathrm{}`$
### 2.3 A computational consequence of the nesting lemma
Lemma 3 suggests a way to compute the minimum $`k`$-assignment of a matrix $`X`$ when $`k<n`$. We proceed by finding the minimum $`k`$-assignments for $`X`$ for $`k=1,2,\mathrm{},`$ one at a time. Suppose we have found a minimum $`(k1)`$-assignment whose rows and columns determine a $`(k1)\times (k1)`$ submatrix $`M`$ of $`X`$. Then, when we search for a minimum $`k`$-assignment, we know that we can restrict our search to minimum $`k`$-assignments in one of the submatrices of $`X`$ obtained by appending a single new row and new column to $`M`$.
There are some simple properties that such an extension must have. Suppose the submatrix of a minimum $`k`$-assignment uses $`M`$ together with a new row $`i`$ and a new column $`j`$. Also, suppose that the minimum $`k`$-assignment uses an entry from column $`j`$ that is in row $`i^{}`$ of $`M`$. Then this entry in column $`j`$ must be the smallest entry in the part of row $`i^{}`$ outside of $`M`$. Similarly, if the minimum $`k`$-assignment uses an entry from row $`i`$ that is in column $`j^{}`$ of $`M`$, then this entry in row $`i`$ must be the smallest entry in the part of column $`j^{}`$ outside $`M`$. Finally, if the minimum $`k`$-assignment uses the entry $`x_{ij}`$, then this entry must be minimum in the submatrix of $`X`$ complementary to $`M`$.
The preceding discussion shows that the following strategy will work to construct the minimum $`k`$-assignment once we have found the minimum $`(k1)`$-assignment. We define the $`k\times k`$ *auxiliary matrix* $`\text{Aux}_M(X)`$ by appending to $`M`$ a new row and column as follows. To each row $`i`$ of $`M`$ we append a new entry which is the minimum of the entries in row $`i`$ of $`X`$ that are outside of $`M`$. To each column $`j`$ of $`M`$ we append a new entry which is the minimum of the entries in column $`j`$ of $`X`$ that are outside of $`M`$. At the intersection of the new row and new column we place the minimum entry of the submatrix of $`X`$ complementary to $`M`$.
We now find the minimum $`k`$-assignment of $`\text{Aux}_M(X)`$. This assignment will use an entry in the last row and an entry in the last column of $`\text{Aux}_M(X)`$, which can be the same entry if the assignment uses the entry in the last row and column. Each of these entries is a copy of some entry of $`X`$, which then tells us which row and column of $`X`$ need to be appended to $`M`$ to obtain the submatrix of the minimum $`k`$-assignment of $`X`$.
### 2.4 A consequence for the expected contribution
The discussion of the auxiliary matrix in the preceding section can be formalized to prove an interesting property of the expected contribution.
We will use the following simple facts about independent exponential random variables.
###### Proposition 2
Suppose that $`a_1,\mathrm{},a_m`$ are positive real numbers, and that $`x_1,\mathrm{},x_m`$ are independent random variables with $`x_i`$ chosen from the exponential distribution with rate $`a_i`$.
Let $`x`$ denote the random variable $`\mathrm{min}_ix_i`$. Then $`x`$ is distributed as an exponential random variable of rate $`a_1+\mathrm{}+a_m`$.
Let $`W`$ be the discrete random variable whose value is the least $`i`$ for which $`x=x_i`$. Then $`W`$ and $`x`$ are independent random variables, and the probability that $`W=i`$ is $`a_i/(a_1+\mathrm{}+a_m)`$.
Proof: Let $`c`$ be a positive real number. Then the probability that $`xc`$ is
$`a_1\mathrm{}a_m{\displaystyle _{x_1=c}^{\mathrm{}}}\mathrm{}{\displaystyle _{x_m=c}^{\mathrm{}}}e^{_{j=1}^ma_jx_j}𝑑x_1\mathrm{}𝑑x_m.`$
$`=`$ $`a_1\mathrm{}a_me^{c_{j=1}^ma_j}{\displaystyle _{u_1=0}^{\mathrm{}}}\mathrm{}{\displaystyle _{u_m=0}^{\mathrm{}}}e^{_{j=1}^ma_ju_j}𝑑u_1\mathrm{}𝑑u_m.`$
$`=`$ $`e^{c_{j=1}^ma_j},`$
where we have made the substitution $`x_i=u_i+c`$ on the second line. Taking the derivative with respect to $`c`$ we see that $`x`$ is distributed as an exponential random variable of rate $`a_1+\mathrm{}+a_m`$.
A similar substitution yields the part of the integral corresponding the event that $`x=x_1`$, as follows:
$`a_1\mathrm{}a_m{\displaystyle _{x_1=c}^{\mathrm{}}}{\displaystyle _{x_2=x_1}^{\mathrm{}}}\mathrm{}{\displaystyle _{x_m=x_1}^{\mathrm{}}}e^{_{j=1}^ma_jx_j}𝑑x_1\mathrm{}𝑑x_m.`$
$`=`$ $`a_1\mathrm{}a_m{\displaystyle _{u_1=0}^{\mathrm{}}}\mathrm{}{\displaystyle _{u_m=0}^{\mathrm{}}}e^{(a_1(c+u_1)+a_2(c+u_1+u_2)+\mathrm{}a_m(c+u_1+u_m))}𝑑u_1\mathrm{}𝑑u_m`$
$`=`$ $`{\displaystyle \frac{a_1}{_{j=1}^ma_j}}e^{c_{j=1}^ma_j}.`$
Thus, the conditional probability that $`x=x_i`$, given $`xc`$, is $`a_i/(a_1+\mathrm{}+a_m)`$. Since this is true for all $`c`$, the event $`x=x_i`$ is independent of the random variable $`x`$. $`\mathrm{}`$
Now let $`X`$ be any nonnegative $`m\times n`$ matrix and associate to $`X`$ the $`k\times k`$ matrix $`Y=\text{Aux}_M(X)`$ where $`M`$ is the upper left $`(k1)\times (k1)`$ submatrix of $`X`$. For any $`t<k`$, let $`M_t`$ denote the upper left $`t\times t`$ submatrix of $`X`$. By abuse of notation we will also let $`M_t`$ denote the upper left $`t\times t`$ submatrix of $`Y`$, since they are identical.
###### Lemma 4
For any $`t<k`$, let $`\sigma `$ be a $`t`$-assignment of $`M_t`$. Then $`\sigma `$ is a minimum $`t`$-assignment for $`X`$ if and only if it is a minimum $`t`$-assignment for $`Y`$.
Proof: Let $`\sigma `$ be a $`t`$-assignment of $`M_t`$.
If $`t=1`$ it is clear that the lemma holds. We proceed by induction on $`t`$.
Suppose that $`\sigma `$ is a minimum $`t`$-assignment of $`X`$. By induction on $`t`$, there is a minimum $`(t1)`$-assignment of $`Y`$ which lies in $`M_t`$, so by Lemma 3, there is some minimum $`t`$-assignment $`\tau `$ of $`Y`$ which uses at most one column and one row outside of $`M_t`$. By definition, each entry of $`Y`$ outside of $`M_t`$ equals some entry of $`X`$ outside of $`M_t`$. If we replace each of the entries of $`\tau `$ outside of $`M_t`$ by their equivalent entries in $`X`$, then we get a $`t`$-assignment $`\tau ^{}`$ of $`X`$. Furthermore, $`\tau ^{}X\sigma X`$, since $`\sigma `$ is minimum. But $`\tau ^{}X=\tau Y`$ and $`\sigma X=\sigma Y`$, so $`\tau Y\sigma Y`$. Therefore, $`\sigma `$ is a minimum $`t`$-assignment of $`Y`$.
Now suppose that $`\sigma `$ is a minimum $`t`$-assignment of $`Y`$. By induction on $`t`$, there is a minimum $`(t1)`$-assignment of $`X`$ which lies in $`M_t`$, so by Lemma 3, there is some minimum $`t`$-assignment $`\tau `$ of $`X`$ which uses at most one column and one row outside of $`M_t`$. We can replace each entry $`x_{ij}`$ of $`\tau `$ by the entry $`y_{\mathrm{min}(i,k),\mathrm{min}(j,k)}`$ to get a $`t`$-assignment $`\tau ^{}`$ of $`Y`$. Furthermore, $`\tau ^{}Y\sigma Y`$, since $`\sigma `$ is minimum. But $`\tau X\tau ^{}Y`$ and $`\sigma X=\sigma Y`$, so $`\tau X\sigma X`$. Therefore, $`\sigma `$ is a minimum $`t`$-assignment of $`X`$. $`\mathrm{}`$
###### Lemma 5
The minimum $`(k1)`$-assignment of $`X`$ uses its first $`k1`$ rows and columns exactly when the minimum $`(k1)`$-assignment of $`Y`$ uses its first $`k1`$ rows and columns. In this case, the value of the minimum $`k`$-assignment of $`X`$ is the same as the value of the minimum $`k`$-assignment of $`Y`$.
Moreover in this case there is a minimum $`k`$-assignment $`\tau `$ of $`X`$ and a minimum $`k`$-assignment $`\tau ^{}`$ of $`Y`$ which correspond entry by entry; i.e., the entries of $`\tau `$ in $`M`$ are located in the same positions as the entries of $`\tau ^{}`$ in $`M`$, and for any entry $`x_{ij}`$ of $`\tau `$ outside of $`M`$ there is a corresponding entry $`y_{\mathrm{min}(i,k),\mathrm{min}(j,k)}`$ of $`\tau ^{}`$ outside of $`M`$.
Proof: The first paragraph is the case $`t=k1`$ in Lemma 4. The second paragraph follows from our discussion of auxiliary matrix in the preceding section. $`\mathrm{}`$
We will use the more detailed statement in the second paragraph to prove Lemma 6 in Section 3.2.
Now let $`X`$ be a random exponential $`m\times n`$ matrix with rate matrix $`A`$, and associate to $`A`$ the $`k\times k`$ matrix $`B=(b_{ij})`$ with
$$b_{ij}=a_{ij}\text{when }1i,jk1$$
and
$$b_{ik}=\underset{j^{}k}{}a_{ij^{}},i=1,\mathrm{},k1$$
and
$$b_{kj}=\underset{i^{}k}{}a_{i^{}j},j=1,\mathrm{},k1$$
and
$$b_{kk}=\underset{i^{},j^{}k}{}a_{i^{}j^{}}.$$
Note that $`Y=\text{Aux}_M(X)`$ is a random exponential matrix with rate matrix $`B`$, by Proposition 2.
The statement in the first paragraph of the preceding Lemma shows that the expected contribution of the submatrix of $`X`$ consisting of its first $`k1`$ rows and first $`k1`$ columns to the minimum $`k`$-assignment of $`X`$ is the same as the expected contribution of the submatrix of $`Y`$ consisting of its first $`k1`$ rows and first $`k1`$ columns to the minimum $`k`$-assignment of $`Y`$.
In the rank 1 rate matrix case, if $`C(k,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n)`$ denotes the expected contribution of the submatrix of the first $`k1`$ rows and columns when the rate matrix is $`A=(r_ic_j)`$, then
$$C(k,r_1,\mathrm{},r_{k1},r_k+\mathrm{}+r_m,c_1,\mathrm{},c_{k1},c_k+\mathrm{}+c_n)$$
is the expected contribution when the rate matrix is $`B`$. These two functions must agree, so the entries of $`A`$ outside of the first $`k1`$ rows and columns enter into the contribution function only via the sums defining the entries in the $`k^{th}`$ row and column of $`B`$.
Now note that Conjecture 3 has the feature that it predicts the same contribution in the two cases above. So this is a bit of evidence in favor of Conjecture 3.
Thus, we can summarize our discussion in the rank 1 rate matrix case as follows:
###### Theorem 3
Let $`k`$, $`m`$ and $`n`$ be integers $`kmn`$. Then Conjecture 3 holds for $`k`$-assignments in an $`m\times n`$ matrix if and only if it holds for $`k`$-assignments in a $`k\times k`$ matrix.
Thus, if one could prove only the cases $`k=m=n`$ of Conjecture 3, that would prove the general case of that conjecture as well as Conjecture 2 and Conjecture 1.
## 3 Rank 1 rate matrices
In this section we restrict our attention to random exponential matrices for which the rate matrix has rank 1. We describe how we arrived at Conjecture 2 and then give equivalent formulations, which provide different kinds of confirmation. Finally, we prove a probability result about the locations of the minimum $`\mathrm{}`$-assignments for $`1\mathrm{}k`$.
We discovered Conjecture 2 while experimenting with the computations described in Section 4. Finding that Mathematica had trouble carrying out the computation when the rate matrix consisted of $`mn`$ indeterminates, we decided to try a simpler case, with rate matrix of the form $`a_{ij}=a_i`$ for all $`i`$ and $`j`$. We noticed in this case that the answer has a surprisingly simple form—in particular, it can be written as a linear combination of the reciprocals of sums of the $`a_i`$’s. Next we found that, when the rate matrix has rank 1, so that $`a_{ij}=r_ic_j`$, the expected value seems to be a linear combination of terms of the form
$$\frac{1}{(_{iI}r_i)(_{jJ}c_j)}$$
(6)
with $`I`$ a proper subset of $`\{1,\mathrm{},m\}`$, and $`J`$ a proper subset of $`\{1,\mathrm{},n\}`$. It is easily shown that the rational functions $`\frac{1}{(_{iI}r_i)(_{jJ}c_j)}`$ are linearly independent over the real numbers. Thus the coefficients in such a linear combination are uniquely determined.
Making the assumption that the expected value is indeed a linear combination of terms of the form (6), we arrived at Conjecture 2 by considering certain limiting conditions on the expected value. We will describe these limiting conditions in Section 5.
### 3.1 Equivalent formulations of Conjecture 2
From now on we use the shorthand notation $`[m]`$ for the set $`\{1,\mathrm{},m\}`$.
Let us introduce the notation
$$F(k,r,c)=\underset{I,J}{}(1)^{k1|I||J|}\left(\genfrac{}{}{0pt}{}{m+n1|I||J|}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}$$
(7)
for the formula given in Conjecture 2. Here recall that $`r=(r_1,\mathrm{},r_m)`$ is an $`m`$-tuple of positive real numbers, $`c=(c_1,\mathrm{},c_n)`$ is an $`n`$-tuple of positive real numbers, and the sum is over proper subsets $`I[m]`$ and $`J[n]`$. The binomial coefficient enforces the condition $`|I|+|J|<k`$. In what follows we will often not mention such constraints explicitly.
In this section we derive alternative ways to write (7) and conclude that $`F(k,r,c)`$ is positive, Conjecture 2 implies Conjecture 1, and Conjecture 3 implies Conjecture 2.
Note that $`F(k,r,c)`$ can be written more succinctly as
$$F(k,r,c)=\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k1mn}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)},$$
(8)
using binomial coefficients with negative numerator.
###### Proposition 3
$$F(k,r,c)=\underset{|I^{}|+|J^{}|<k,II^{},JJ^{}}{}(1)^{|I^{}||I|+|J^{}||J|}\frac{1}{(_{iI}r_i)(_{jJ}c_j)}.$$
(9)
Here $`I^{}`$ and $`I`$ are proper subsets of $`[m]`$ and $`J^{}`$ and $`J`$ are proper subsets of $`[n]`$.
Proof: Comparison with (8) shows that, for a fixed $`I`$ and $`J`$ with $`|I|+|J|<k`$, we need to evaluate the sum
$$\underset{II^{},JJ^{},|I^{}|+|J^{}|<k}{}(1)^{|I^{}||I|+|J^{}||J|}.$$
If we denote by $`i`$ and $`j`$ the cardinalities of $`I`$ and $`J`$ and by $`t`$ and $`u`$ the cardinalities of $`I^{}`$ and $`J^{}`$, we can rewrite this sum as
$`{\displaystyle \underset{ti,uj,t+u<k}{}}(1)^{ti+uj}\left({\displaystyle \genfrac{}{}{0pt}{}{mi}{ti}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{uj}}\right)`$
$`=`$ $`{\displaystyle \underset{t0,u0,t+u<kij}{}}(1)^{t+u}\left({\displaystyle \genfrac{}{}{0pt}{}{mi}{t}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{u}}\right)`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{k1ij}{}}}(1)^l{\displaystyle \underset{t0,u0,t+u=l}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{mi}{t}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{u}}\right)`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{k1ij}{}}}(1)^l\left({\displaystyle \genfrac{}{}{0pt}{}{m+nij}{l}}\right)`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{k1ij}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{i+jmn+l1}{l}}\right)`$
$`=`$ $`{\displaystyle \underset{l=0}{\overset{k1ij}{}}}\left(\left({\displaystyle \genfrac{}{}{0pt}{}{i+jmn+l}{l}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{i+jmn+l1}{l1}}\right)\right)`$
$`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{k1mn}{k1ij}}\right)`$
which agrees with (8). $`\mathrm{}`$
We can rewrite (9) as a double sum with the inner sum over $`I`$ and $`J`$ and the outer sum over $`I^{}`$ and $`J^{}`$. Then, for fixed $`I^{}`$ and $`J^{}`$, the inner sum factors as
$$\left(\underset{II^{}}{}(1)^{|I^{}||I|}\frac{1}{_{iI}r_i}\right)\left(\underset{JJ^{}}{}(1)^{|J^{}||J|}\frac{1}{_{jJ}c_j}\right)$$
(10)
Now we show that each factor has an interesting probabilistic interpretation.
Suppose that an urn contains $`m`$ balls labeled $`1,2,\mathrm{},m`$ and for each $`i`$, ball $`i`$ has weight $`r_i`$. We select balls one at a time without replacement, at each time selecting a ball with probability proportional to the weights of those balls still in the urn. Let $`\mathrm{Pr}(r,I^{})`$ denote the probability that the set of balls in $`I^{}`$ are the first $`t`$ balls to be chosen, where $`t`$ is the cardinality of $`I^{}`$. Then
$$\mathrm{Pr}(r,I^{})=\underset{\pi }{}\underset{i=1}{\overset{t}{}}\frac{r_{\pi _i}}{R_{j=1}^{i1}r_{\pi _i}}$$
(11)
where $`R=_{i=1}^mr_i`$ and the outer sum is over all $`t!`$ orderings $`(\pi _1,\mathrm{},\pi _t)`$ of $`I^{}`$.
We can calculate $`\mathrm{Pr}(r,I^{})`$ in a different way as follows. Suppose we draw all $`m`$ balls from the urn. If we fix any subset $`U`$ of balls, then the probability that a particular ball $`u`$ from $`U`$ is chosen before any other ball from $`U`$ is the weight of $`u`$ divided by the sum of the weights of the balls in $`U`$.
Now, for $`iI^{}`$, let $`E_i`$ denote the event that the first time a ball is drawn from the set consisting of $`i`$ together with the complement of $`I^{}`$, the ball chosen is from the complement of $`I^{}`$. Then $`E_i`$ has probability
$$\frac{_{jI^{}}r_j}{r_i+_{jI^{}}r_j}.$$
In order for our set $`I^{}`$ to be the set of the first $`t`$ balls chosen, it is necessary and sufficient that none of the events $`E_i`$ occur. For any subset $`I`$ of $`I^{}`$ the probability that all of the events $`E_i`$, $`iI`$, occur is
$$\frac{_{jI^{}}r_j}{_{i(I^{}I)}r_i}$$
So, by the Inclusion-Exclusion Principle,
$$\mathrm{Pr}(r,I^{})=\underset{II^{}}{}(1)^{|I^{}||I|}\frac{_{iI^{}}r_i}{_{iI}r_i}$$
(12)
which is $`(_{iI^{}}r_i)`$ times the first factor in (10). The analogous result holds for the second factor in (10). We conclude that
###### Proposition 4
$$F(k,r,c)=\underset{I,J,|I|+|J|<k}{}\frac{\mathrm{Pr}(r,I)\mathrm{Pr}(c,J)}{(_{iI}r_i)(_{jJ}c_j)},$$
(13)
where the sum is over proper subsets $`I`$ of $`[m]`$ and $`J`$ of $`[n]`$. Hence $`F(k,r,c)`$ is always positive.
$`\mathrm{}`$
Now we rewrite $`F(k,r,c)`$ to show that Conjecture 2 implies Conjecture 1 and Conjecture 3 implies Conjecture 2.
Let $`\mathrm{Pr}(r,(i_1,\mathrm{},i_t))`$ denote the probability that the first $`t`$ selections from our urn are $`i_1,\mathrm{},i_t`$ in that order. Then, from (13) and (11), we get
$$F(k,r,c)=\underset{t,u0,t+u<k}{}\underset{i,j}{}\frac{\mathrm{Pr}(r,(i_1,\mathrm{},i_t))\mathrm{Pr}(c,(j_1,\mathrm{},j_u))}{(R_{s=1}^tr_{i_s})(C_{s=1}^uc_{j_s})},$$
(14)
where the sum is over all sequences $`i=(i_1,\mathrm{},i_t)`$ of distinct integers in $`[m]`$ and $`j=(j_1,\mathrm{},j_u)`$ of distinct integers in $`[n]`$. But, certainly
$$\mathrm{Pr}(r,(i_1,\mathrm{},i_t))=\underset{i}{}\mathrm{Pr}(r,(i_1,\mathrm{},i_t,\mathrm{},i_{k1}))$$
where the sum is over all extensions of $`(i_1,\mathrm{},i_t)`$ to a $`(k1)`$-long sequence $`i`$ of distinct integers in $`[m]`$. Thus, we can rewrite (14) and obtain:
###### Proposition 5
$$F(k,r,c)=\underset{i,j}{}\underset{t,u0,t+u<k}{}\frac{\mathrm{Pr}(r,(i_1,\mathrm{},i_{k1}))\mathrm{Pr}(c,(j_1,\mathrm{},j_{k1}))}{(R_{s=1}^tr_{i_s})(C_{s=1}^uc_{j_s})},$$
(15)
where the outer sum in (15) is over pairs of ordered sequences of $`k1`$ distinct integers from $`[m]`$ and $`[n]`$.
Note that each term in the above sum corresponds to a flag of submatrices of sizes $`1\times 1,\mathrm{},(k1)\times (k1)`$. In this form, specializing to the case that all the $`r`$’s and $`c`$’s are 1, it is easy to see that Conjecture 2 implies Conjecture 1.
We now group the terms in the outer sum according to the (unordered) sets $`I=\{i_1,\mathrm{},i_{k1}\}`$ and $`J=\{j_1,\mathrm{},j_{k1}\}`$. It then becomes
$$F(k,r,c)=\underset{I,J}{}\underset{i,j}{}\underset{t,u0,t+u<k}{}\frac{\mathrm{Pr}(r,(i_1,\mathrm{},i_{k1}))\mathrm{Pr}(c,(j_1,\mathrm{},j_{k1}))}{(R_{s=1}^tr_{i_s})(C_{s=1}^uc_{j_s})}$$
(16)
where the outer sum is over sets $`I`$ and $`J`$ of size $`k1`$ and the inner sum is over permutations $`(i_1,\mathrm{},i_{k1})`$ of $`I`$ and permutations $`(j_1,\mathrm{},j_{k1})`$ of $`J`$. In this form we can see that the term of the outer sum corresponding to the sets $`I`$ and $`J`$ is the expected value of the contribution of the submatrix with row indices $`I`$ and column indices $`J`$ predicted by Conjecture 3. Since the sum of the expected contributions of all submatrices is the expected minimum $`k`$-assignment, we now see that Conjecture 3 implies Conjecture 2.
Finally, for any $`TI[m]`$, let $`\mathrm{Pr}(r,T,I)`$ denote the probability that the first $`|T|`$ balls drawn from the urn comprise the set $`T`$ and that the first $`|I|`$ balls drawn comprise the set $`I`$. Then we can rewrite our formula for the expected contribution of a submatrix with rows $`I`$ and columns $`J`$ as
$$\underset{TI,UJ,|T|+|U|<k}{}\frac{\mathrm{Pr}(r,T,I)\mathrm{Pr}(c,U,J)}{(_{tT}r_t)(_{uU}c_u)}$$
(17)
### 3.2 Flag probabilities
In this section, we prove ap probability result in the special case that the rate matrix has rank 1. This result may be the reason that simple formulas exist for $`E_k(A)`$ when the matrix $`A`$ has rank 1.
It is possible for a matrix to have many minimum $`k`$-assignments for some $`k`$. However, with probability 1, a random matrix $`X`$ has a unique minimum $`k`$-assignment for each $`k`$. So, if we let $`M_k`$ denote a $`k\times k`$ submatrix of $`X`$ containing a minimum $`k`$-assignment of $`X`$, then, with probability 1, $`M_k`$ is unique. By Lemma 3, the submatrices $`M_k`$ are nested: $`M_1M_2\mathrm{}M_k.`$ We will call this the *flag of submatrices of $`X`$*. This flag can also be described by the list $`i_1,i_2,\mathrm{},i_k`$ of appended rows and the list $`j_1,j_2,\mathrm{},j_k`$ of appended columns. Thus, $`M_l`$ is the submatrix with rows $`i_1,\mathrm{},i_l`$ and columns $`j_1,\mathrm{},j_l`$.
It is natural to ask for the probability that a random matrix has a given flag of submatrices. We know of no formula for this probability for general rate matrices. However, for rate matrices of rank 1, we can prove a simple formula for the probability of each flag. Moreover, this formula will involve the probabilities $`\mathrm{Pr}(r,(i_1,\mathrm{},i_k))`$ calculated in the previous section.
We first need the following
###### Lemma 6
Let $`X`$ be an exponential random matrix with rank 1 rate matrix $`A=(r_ic_j)`$. If the minimum $`(k1)`$-assignment of $`X`$ uses the first $`k1`$ rows and first $`k1`$ columns, then the minimum $`k`$-assignment uses an additional row and column. The probability of row $`i^{}`$ being the additional row is $`\frac{r_i^{}}{_{i=k}^mr_i}`$, the probability of column $`j^{}`$ being the additional column is $`\frac{c_j^{}}{_{j=k}^nc_j}`$, and these events are independent.
Proof: Let $`M`$ denote the upper left $`(k1)\times (k1)`$ submatrix of $`X`$, and $`Y`$ the $`k\times k`$ matrix $`\text{Aux}_M(X)`$ defined in Section 2.3. Then the upper left submatrix of $`Y`$ is identical to $`M`$, and, by abuse of notation, we will let $`M`$ denote that submatrix of $`Y`$ also. If the minimum $`(k1)`$-assignment of $`X`$ lies in $`M`$, then by Lemma 4, the minimum $`(k1)`$-assignment of $`Y`$ lies in $`M`$, and by Lemma 5, the minimum $`k`$-assignments of $`X`$ and $`Y`$ correspond entry by entry.
There are two cases to consider.
In the first case, for some $`sk1`$ and $`tk1`$, the minimum $`k`$-assignment of $`Y`$ uses the entries $`y_{sk}`$ and $`y_{kt}`$. These entries correspond to the minimum entry in row $`s`$ of $`X`$ outside of $`M`$ and the minimum entry in column $`t`$ of $`X`$ outside of $`M`$. In the second case, the minimum $`k`$-assignment of $`Y`$ uses entry $`y_{kk}`$. This entry corresponds to the minimum entry in the submatrix of $`X`$ complementary to $`M`$. In both cases, by Proposition 2, the locations of these minima in $`X`$ are independent of the random variables making up the entries of $`Y`$, and thus independent of the events that the minimum $`(k1)`$-assignment of $`Y`$ lies in $`M`$ and the minimum $`k`$-assignment of $`Y`$ uses particular entries outside of $`M`$.
Thus in the first case, the probability that the minimum entry in the part of row $`s`$ outside of $`M`$ comes from column $`j^{}`$ is $`\frac{r_sc_j^{}}{_{j=k}^nr_sc_j}=\frac{c_j^{}}{c_k+\mathrm{}+c_n}`$, and the probability that the minimum entry in the part of column $`t`$ outside of $`M`$ comes from row $`i^{}`$ is $`\frac{r_i^{}c_t}{_{i=k}^mr_ic_t}=\frac{r_i^{}}{r_k+\mathrm{}r_m}`$. Moreover, the locations of these minima within row $`s`$ and column $`t`$ are independent events since the parts of row $`s`$ and column $`t`$ outside of $`M`$ are disjoint.
In the second case, the probability of the minimum entry in the submatrix of $`X`$ complementary to $`M`$ coming from row $`i^{}`$ and column $`j^{}`$ is $`\frac{r_i^{}c_j^{}}{{\scriptscriptstyle r_ic_j}}`$ where the sum in the denominator is over all locations $`(i,j)`$ in the submatrix of $`X`$ complementary to $`M`$. Thus, the probability that the minimum entry comes from row $`i^{}`$ is $`\frac{r_i^{}}{_{i=k}^mr_i}`$ and the probability that the minimum entry comes from column $`j^{}`$ is $`\frac{c_j^{}}{_{j=k}^nc_j}`$. $`\mathrm{}`$
From this lemma we can immediately conclude the following theorem, which imparts further meaning to the formal $`\mathrm{Pr}(r,I)`$ and $`\mathrm{Pr}(r,(i_1,\mathrm{},i_k))`$ functions used in Section 3.1:
###### Theorem 4
Suppose that $`A=(r_ic_j)`$ is a rank 1 rate matrix and $`X`$ an exponential random matrix with rate matrix $`A`$. Let $`(i_1,\mathrm{},i_k)`$ be a sequence of distinct elements of $`[m]`$ and $`(j_1,\mathrm{},j_k)`$ a sequence of distinct elements from $`[n]`$. Then the probability that $`X`$ has the associated flag of submatrices is $`\mathrm{Pr}(r,(i_1,\mathrm{},i_k))\mathrm{Pr}(c,(j_1,\mathrm{},j_k))`$. Furthermore, if $`I[m]`$ and $`J[n]`$ are sets of size $`k`$, then the probability that the minimum $`k`$-assignment of $`X`$ uses the rows indexed by $`I`$ is $`\mathrm{Pr}(r,I)`$ and the probability that it uses the columns indexed by $`J`$ is $`\mathrm{Pr}(c,J)`$, and these events are independent.
$`\mathrm{}`$
We will see that this formula enters in an essential way into the proof of Theorem 5.
## 4 Computational evidence for our conjectures
By Proposition 2, it is an easy matter to compute the expected value of the minimum 1-assignment for an arbitrary rate matrix.
###### Example 4
The expected value of the minimum 1-assignment when the rate matrix is $`A=(a_{ij})`$ is
$$\frac{1}{_{ij}a_{ij}}.$$
For $`k2`$ the computation is more complicated. In \[AS\] and \[CS\] the authors calculate the expected value of the minimum assignment for a random exponential matrix when the rates are all 1 and $`k`$ is small.
The method in \[AS\] applies just as well to the case of arbitrary rate matrices. The essence of their idea is to introduce a slightly more general expectation problem in which they choose all the entries of the random matrix $`X`$ as before, except that there is a set $`Z`$ of fixed zeroes in $`X`$. Let us denote the expected value of the minimum assignment in this case by $`E(A,Z)`$.
It is then sometimes possible to establish a recursive calculation of $`E(A,Z)`$. The base of the recursion occurs when there exist $`k`$ zeroes in $`Z`$, no two in the same row or column. In this case we know that the expected value of the minimum assignment is zero. For the inductive part of the calculation we can sometimes express an expected value $`E(A,Z)`$ as a constant plus a linear combination of $`E(A,Z^{})`$ where $`Z^{}`$ is obtained from $`Z`$ by adjoining one more position to $`Z`$.
This arises as follows. Suppose that $`X`$ is a random exponential matrix except for a set $`Z`$ of positions in $`X`$ where the entries are fixed zeroes. Suppose further that we have a set $`S`$ of positions in $`X`$, disjoint from $`Z`$, such that any minimum $`k`$-assignment of $`X`$ meets the set $`S`$ in exactly $`r`$ positions. (In other words, every nonnegative matrix with zero set $`Z`$ has the property that its minimum $`k`$-assignments all meet $`S`$ in exactly $`r`$ positions.) Abusing notation, we also let $`S`$ denote the matrix which is 1 at the positions in the set $`S`$ and zero otherwise.
We will derive the following formula:
$$E(A,Z)=\frac{r}{AS}+\underset{(i,j)S}{}\frac{a_{ij}}{AS}E(A,Z\{(i,j)\}).$$
(18)
Indeed, the integral for $`E(A,Z)`$, which involves only the variables $`x_{ij}`$ for $`(i,j)Z`$, is given by
$$E(A,Z)=\left(\underset{(i,j)Z}{}a_{ij}\right)_X\underset{k}{\mathrm{min}}(X)e^{AX}𝑑X.$$
We can derive (18) by breaking up this integral into $`|S|`$ parts, each corresponding to a position in $`S`$ containing the minimum entry among all positions in $`S`$. For the part of the integral where $`x_{i_0j_0}`$ is the minimum entry in $`S`$, we make a change of variables with Jacobian 1, as follows. We express the $`x_{ij}`$ in terms of new variables $`y_{ij}`$ by setting $`x_{ij}=y_{i_0j_0}+y_{ij}`$ when $`(i,j)S\{(i_0,j_0)\}`$ and $`x_{ij}=y_{ij}`$ otherwise. $`X`$ can then be written as $`Y+y_{i_0j_0}S`$ where $`Y`$ is a nonnegative matrix with fixed zeroes at $`Z\{(i_0,j_0)\}`$. From our hypothesis about $`S`$, we have
$$\underset{k}{\mathrm{min}}(X)=\underset{k}{\mathrm{min}}(Y)+ry_{i_0j_0}.$$
(Otherwise, there would be a non-minimum $`k`$-assignment of $`X`$, meeting $`S`$ in fewer than $`r`$ positions, that becomes a minimum $`k`$-assignment of a matrix $`XtS`$ for some $`t<y_{i_0j_0}`$. But, the matrix $`XtS`$ still has zero set $`Z`$, so our hypothesis on $`S`$ would be contradicted.) Thus, this part of the integral becomes
$$\left(\underset{(i,j)Z}{}a_{ij}\right)_{Y,y_{i_0j_0}}(\underset{k}{\mathrm{min}}(Y)+ry_{i_0j_0})e^{A(Y+y_{i_0j_0}S)}𝑑y_{i_0j_0}𝑑Y.$$
This can be computed as the sum of two integrals in the obvious way. The first is $`a_{i_0j_0}E(A,Z\{(i_0,j_0)\})/(AS)`$ and the second is $`ra_{i_0j_0}/(AS)^2`$. When we sum these expressions over all $`(i_0,j_0)S`$ we obtain (18).
When $`k=m=n4`$, it is easy to see that, when we are not in the base case, there always exists a set $`S`$ of positions in $`X`$ and disjoint from $`Z`$ such that every minimum $`k`$-assignment of $`X`$ meets $`S`$ in the same number of positions.
To illustrate the method we now discuss the case $`k=m=n=4`$. First note that if any row or column of $`X`$ has no fixed zeroes, then we can take that row or column to be the set $`S`$. So we can suppose that every row or column has at least one fixed zero.
Now suppose that there is a $`3\times 3`$ submatrix $`S`$ of $`X`$ that has no fixed zeroes. Without loss of generality, we may take this to be the upper left $`3\times 3`$ submatrix, so the matrix $`X`$ has the form
$$X=\left[\begin{array}{cccc}\hfill & \hfill & \hfill & \hfill 0\\ & \hfill & \hfill & \hfill 0\\ & \hfill & \hfill & \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill .\end{array}\right]$$
(19)
where $``$ means that entry is positive and . means that nothing is known about that entry. Then any minimum 4-assignment must use either two or three entries from $`S`$. If it uses three, then it must use $`x_{44}`$ and some entry $`x_{ij}`$, $`i,j<4`$. But then we can decrease the value of the assignment by replacing $`x_{ij}`$ and $`x_{44}`$ with $`x_{i4}`$ and $`x_{4j}`$, both of which are zero. This contradicts the minimality of the 4-assignment we started with. Thus, any minimum 4-assignment must use exactly two entries from $`S`$.
If every $`3\times 3`$ submatrix of $`X`$ has a fixed zero, and we are not in the base case, then the Hall marriage theorem implies that there is a $`2\times 3`$ or $`3\times 2`$ submatrix $`S`$ that has no fixed zeroes. Suppose the former, which we can take to be the upper left $`2\times 3`$ submatrix of $`X`$. Each of the first three columns has at least one fixed zero. The fixed zeros in those columns must be in more than one row, since every $`3\times 3`$ submatrix has a fixed zero. Thus, we may assume the matrix $`X`$ has the form
$$X=\left[\begin{array}{cccc}\hfill & \hfill & \hfill & \hfill 0\\ & \hfill & \hfill & \hfill 0\\ \hfill 0& \hfill 0& \hfill .& \hfill .\\ \hfill .& \hfill .& \hfill 0& \hfill .\end{array}\right]$$
(20)
Any minimum $`4`$-assignment must use one or two entries from $`S`$. Suppose a minimum $`4`$-assignment uses two entries from $`S`$. It cannot use $`x_{11}`$, since then there would be a smaller $`4`$-assignment consisting of $`x_{11}`$,$`x_{24}`$,$`x_{32}`$,and $`x_{43}`$. Similarly it cannot use $`x_{12}`$, $`x_{21}`$, or $`x_{22}`$. But, it can only use one of $`x_{13}`$ and $`x_{23}`$, so it must use exactly one entry from $`S`$.
Thus, we can always find a suitable set $`S`$ to continue the recursive calculation.
We have used this method to compute the expected minimum $`k`$-assignment for various small cases. For the case $`km=n3`$, this is easily carried out by Mathematica and confirms our conjecture.
When $`k=m=n=4`$ and $`k=m=n=5`$ we were not patient enough to wait for Mathematica to simplify the complete rational expression, even when the rate matrix has rank 1. However, we were able to check that we obtained the correct answer for many random choices or $`r_i`$’s and $`c_j`$’s. For this purpose, we used an ordinary C program, but, instead of using exact rational arithmetic, we carried out our calculations modulo a large prime. Even so, the evidence seems to be overwhelming that our conjecture is correct in these cases.
It is possible, although somewhat more complicated, to compute the expected contribution of a $`(k1)\times (k1)`$ submatrix to the expected minimum assignment of a $`k\times k`$ matrix when $`2k4`$. In the cases $`k=2`$ and $`k=3`$ we were able to check directly with Mathematica that Conjecture 3 was valid, which proves Conjecture 1 and Conjecture 2 whenever $`k3`$.
When $`k=4`$ we obtained computational evidence for the validity of Conjecture 3, checking its validity in a large number of random cases modulo a prime. This provides confirmation of the other conjectures when $`k=4`$ and $`m`$ and $`n`$ are arbitrary.
## 5 Additional evidence for the main conjecture
Let $`A=(r_ic_j)`$ as usual and denote $`E_k(A)`$ by $`E(k,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n)`$, or simply $`E(k,r,c)`$. Recall from Section 3.1 the notation $`F(k,r,c)`$ for the formula in Conjecture 2. In this section we will show that $`E`$ and $`F`$ share several properties.
Let us consider how $`E`$ behaves if we let a collection of the $`r_i`$’s approach 0. For simplicity we assume that $`r_1,r_2,\mathrm{},r_l`$ approach 0. Then the random matrices will have very large entries in the first $`l`$ rows. When $`kml`$, there are assignments which avoid the first $`l`$ rows, so in the limit that $`r_1,r_2,\mathrm{},r_l0`$, the minimum assignment will avoid those rows and become equal to $`E(k,r_{l+1},\mathrm{},r_m,c).`$
Now suppose that $`k>ml`$. Then a $`k`$-assignment must use at least $`k(ml)=k+lm`$ of the first $`l`$ rows. But, in our limiting case, these rows will be very large, so the minimum $`k`$-assignment will use as few as possible, or exactly $`k+lm`$ of them. The contribution of the entries from these rows to the minimum $`k`$-assignment will dominate the minimum $`k`$-assignment, so in the minimum $`k`$-assignment this contribution will be as small as possible. In particular, in the limit as $`r_1,r_2,\mathrm{},r_l0`$, this part of the minimum $`k`$-assignment will be $`E(k+lm,r_1,\mathrm{},r_l,c_1,\mathrm{},c_n)`$. By Theorem 4 we know that a set $`K`$ of $`k+lm`$ columns will be used by the part of the assignment in the first $`l`$ rows with probability $`\mathrm{Pr}(c,K)`$. When this happens the expected contribution from the remaining rows is $`E(ml,r_{l+1},\mathrm{},r_m,c^{}(K))`$ where by $`c^{}(K)`$ denotes the $`c_j`$’s corresponding to columns not in $`K`$. Thus, we should have the following
###### Theorem 5
When $`kml`$,
$$\underset{r_1,\mathrm{},r_l0}{lim}E(k,r,c)=E(k,r_{l+1},\mathrm{},r_m,c).$$
(21)
When $`k>ml>0`$,
$`\underset{r_1,\mathrm{},r_l0}{lim}(E(k,r,c)`$ $`E(k+lm,r_1,\mathrm{},r_l,c))`$
$`={\displaystyle \underset{K}{}}\mathrm{Pr}(c,K)E(ml,r_{l+1},\mathrm{},r_m,c^{}(K))`$ (22)
where the sum is over $`K[n]`$ such that $`|K|=k+lm`$.
Proof: Let $`Z`$ be a random exponential $`m\times n`$ matrix with all entries of mean 1. Then define $`X=Z/A`$ to be the term by term quotient of the random matrix $`Z`$ by the fixed rate matrix $`A`$, where $`a_{ij}=r_ic_j`$. Then $`X`$ is a random exponential matrix with rate matrix $`A`$. In particular,
$$E(k,r,c)=E_k(A)=E(\underset{k}{\mathrm{min}}(Z/A)).$$
Let $`A_u`$ denote the $`l\times n`$ matrix that is comprised of the first $`l`$ rows of $`A`$, and let $`A_d`$ denote $`(ml)\times n`$ matrix consisting of the last $`ml`$ rows of $`A`$. Furthermore, let $`Z_u`$ and $`Z_d`$ denote random exponential matrices of the same corresponding shapes, again with rate 1.
If $`kml`$, then it is easy to see that
$$\underset{r_1,\mathrm{},r_l0}{lim}\underset{k}{\mathrm{min}}(Z/A)=\underset{k}{\mathrm{min}}(Z_d/A_d)$$
pointwise almost everywhere (i.e., *almost surely*, as random variables). Since $`\mathrm{min}_k(Z/A)\mathrm{min}_k(Z_d/A_d)`$, and the expectation of $`\mathrm{min}_k(Z_d/A_d)`$ is finite, we can apply the dominated convergence theorem to show that the limit of the expectation is equal to the expectation of the limit. Thus,
$$\underset{r_1,\mathrm{},r_l0}{lim}E(\underset{k}{\mathrm{min}}(Z/A))=E(\underset{k}{\mathrm{min}}(Z_d/A_d))$$
and therefore
$$\underset{r_1,\mathrm{},r_l0}{lim}E(k,r,c)=E(k,r_{l+1},\mathrm{},r_m,c).$$
Now suppose $`k>ml>0`$. Consider the nonnegative random variable $`R=\mathrm{min}_k(Z/A)\mathrm{min}_{k+lm}(Z_u/A_u)`$. (It is nonnegative because the minimum $`k`$-assignment of $`Z/A`$ must use at least $`k+lm`$ elements from $`Z_u/A_u`$.) Furthermore,
$$\underset{k}{\mathrm{min}}(Z/A)\underset{k+lm}{\mathrm{min}}(Z_u/A_u)+\underset{ml}{\mathrm{max}}(Z_d/A_d)$$
where $`\mathrm{max}_kX`$ denotes the *maximum* $`k`$-assignment of a matrix $`X`$. Then the second summand is a nonnegative random variable of finite expectation, independent of $`r_1,\mathrm{},r_l`$, dominating $`R`$. Let $`\chi _K`$ denote the random variable that is 1 or 0 depending upon whether the minimum $`(k+lm)`$-assignment of the matrix $`Z_u/A_u`$ uses precisely the columns from the set $`K`$ or does not. Then, by Theorem 4, we know that $`E(\chi _K)=\mathrm{Pr}(c,K)`$, independent of $`r`$. Let $`C_K(Z_d/A_d)`$ denote the submatrix of $`Z_d/A_d`$ obtained when the columns indexed by $`K`$ are removed. We can then show
$$\underset{r_1,\mathrm{},r_l0}{lim}\left(R\underset{K}{}\chi _K\underset{ml}{\mathrm{min}}(C_K(Z_d/A_d))\right)=0$$
with convergence pointwise almost everywhere, where the sum is over $`K[n]`$ such that $`|K|=k+lm`$. Furthermore, we can bound the finite sum by a random variable independent of $`r_1,\mathrm{},r_l`$. (The random variable $`_K\mathrm{min}_{ml}(C_K(Z_d/A_d))`$ will do.) Thus, by the dominated convergence theorem, we can take the limit of the expectations and obtain
$`\underset{r_1,\mathrm{},r_l0}{lim}E(R)`$ $`=`$ $`\underset{r_1,\mathrm{},r_l0}{lim}{\displaystyle \underset{K}{}}E(\chi _K\underset{ml}{\mathrm{min}}(C_K(Z_d/A_d)))`$
$`=`$ $`\underset{r_1,\mathrm{},r_l0}{lim}{\displaystyle \underset{K}{}}E(\chi _K)E(\underset{ml}{\mathrm{min}}(C_K(Z_d/A_d)))`$
$`=`$ $`{\displaystyle \underset{K}{}}\mathrm{Pr}(c,K)E(ml,r_{l+1},\mathrm{},r_m,c^{}(K))`$
But, $`E(R)=E(k,r,c)E(k+lm,r_1,\mathrm{},r_l,c)`$, so we are done.$`\mathrm{}`$
We stated Theorem 5 in terms of limits as the first $`l`$ $`r`$’s approach 0, in order to simplify the notation. However, since $`E(k,r,c)`$ is symmetric in the $`r`$’s and $`c`$’s, the analogous results hold for any set of $`l`$ $`r`$’s approaching 0.
From Example 1, we have that
$$E(1,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n)=\frac{1}{(_ir_i)(_jc_j)}.$$
(23)
Also when $`l=1`$, taking symmetry into account, Theorem 5 reduces to the following.
Suppose that $`m>1`$ and $`1im`$. Then, if $`k<m`$,
$$\underset{r_i0}{lim}E(k,r,c)=E(k,r_1,\mathrm{},\widehat{r}_i,\mathrm{},r_m,c)$$
(24)
while if $`k=m`$,
$`\underset{r_i0}{lim}(`$ $`E(k,r,c){\displaystyle \frac{1}{r_i_jc_j}})`$
$`={\displaystyle \frac{_jc_jE(k1,r_1,\mathrm{},\widehat{r}_i,\mathrm{},r_k,c_1,\mathrm{},\widehat{c}_j,\mathrm{},c_n)}{_jc_j}}.`$ (25)
###### Proposition 6
There is at most one set of functions
$$G(k,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n),$$
each a linear combination of terms of the form (6), that satisfy the equations (23), (24) and (5) with $`E`$ replaced by $`G`$.
Proof: Let $`H(k,m,n)=H(k,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n)`$ denote the difference between two sets of functions satisfying the conditions of the proposition. We will show by induction on $`k`$ and $`m`$ that $`H(k,m,n)=0`$.
It is clear that $`H(1,m,n)=0`$ for all $`m`$ and $`n`$. Given values of $`k`$ and $`m`$, we may suppose that $`H(k^{},m^{},n)=0`$ and $`H(k,m^{},n)=0`$ for $`k^{}<k`$, $`m^{}<m`$, and arbitrary $`n`$. Then, by equations (24) and (5) and induction, we have, for any $`i`$, $`lim_{r_i0}H(k,m,n)=0`$. Suppose that
$$H(k,m,n)=\underset{I}{}\frac{h_I}{_{iI}r_i},$$
where $`I`$ runs over nonempty subsets of $`[m]`$ and the $`h_I`$’s are rational functions of the $`c`$’s. Suppose that $`h_I0`$ for some $`I`$. Let $`I_0`$ be a minimal $`I`$ such that $`h_I0`$ and let $`tI_0`$. Since $`lim_{r_t0}H(m,n,k)`$ exists, we must have $`h_{\{t\}}=0`$, so $`I_0\{t\}`$. Also
$$\underset{r_t0}{lim}H(k,m,n)=\underset{I}{}\frac{h_I+h_{I\{t\}}}{_{iI\{t\}}r_i}$$
where the sum is over all $`I`$ strictly containing $`\{t\}`$. Since the terms $`\frac{1}{_{iI}r_i}`$ are linearly independent, we must have $`h_I+h_{I\{t\}}=0`$ for all $`I`$ strictly containing $`\{t\}`$. In the case $`I=I_0`$, this contradicts the minimality of $`I_0`$. Thus $`h_I=0`$ for all $`I`$ and $`H(k,m,n)=0`$. $`\mathrm{}`$
Now we show that the rational functions $`F(k,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n)`$ satisfy the same limit conditions that are proved about $`E(k,r_1,\mathrm{},r_m,c_1,\mathrm{},c_n)`$ in Theorem 5.
In particular the functions $`F(k,r,c)`$ satisfy the conditions of Proposition 6. Thus they are the only possible linear combinations of terms of the form (6) that could equal $`E(k,r,c)`$.
The fact that the $`F`$’s satisfy all the limit conditions proved about $`E`$ provides additional evidence for Conjecture 2.
###### Theorem 6
When $`kml`$,
$$\underset{r_1,\mathrm{},r_l0}{lim}F(k,r,c)=F(k,r_{l+1},\mathrm{},r_m,c).$$
(26)
When $`k>ml>0`$,
$`\underset{r_1,\mathrm{},r_l0}{lim}(F(k,r,c)`$ $`F(k+lm,r_1,\mathrm{},r_l,c))`$
$`={\displaystyle \underset{K}{}}\mathrm{Pr}(c,K)F(ml,r_{l+1},\mathrm{},r_m,c^{}(K))`$ (27)
where the sum is over $`K[n]`$ such that $`|K|=k+lm`$.
Proof: We use the alternate form of $`F(k,r,c)`$ given by (8):
$$F(k,r,c)=\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k1mn}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)},$$
where the binomial coefficient enforces the condition $`|I|+|J|<k`$.
Now, on the left side of (6), before passing to the limit, the first term is
$$\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k1mn}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}$$
where the sum is over subsets $`I[m]`$ and $`J[n]`$ and the second term is
$$\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k+lm1ln}{k+lm1|I||J|}\right)\frac{1}{(_{iI,il}r_i)(_{jJ}c_j)}$$
where the sum is over subsets $`I[l]`$ and $`J[n]`$. By substituting $`I\{l+1,\mathrm{},m\}`$ for $`I`$ in the second term we get
$$\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k1mn}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}$$
where the sum is over subsets $`I[m]`$ and $`J[n]`$ such that $`\{l+1,\mathrm{},m\}I`$. Thus, before taking the limit, the left side of (6) equals
$$\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k1mn}{k1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}.$$
(28)
where the sum is over all subsets $`J[n]`$ and those subsets $`I[m]`$ that do not contain $`\{l+1,\mathrm{},m\}`$.
The last condition on $`I`$ implies that $`_{iI}r_i`$ is nonzero if we set $`r_1,\mathrm{},r_l=0`$. Thus, we can obtain the limit on the left of (6) simply by replacing $`r_1,\mathrm{},r_l`$ by zero in (28). When we do this, the effect is that we combine terms with $`I`$ having a fixed intersection with $`\{l+1,\mathrm{},m\}`$. Note that this intersection is never all of $`\{l+1,\mathrm{},m\}`$.
Suppose then that $`I`$ is a set strictly contained in $`\{l+1,\mathrm{},m\}`$. For each $`i`$ there are $`\left(\genfrac{}{}{0pt}{}{l}{i}\right)`$ ways of extending $`I`$ to a $`(|I|+i)`$-element subset of $`[m]`$ whose intersection with $`\{l+1,\mathrm{},m\}`$ is $`I`$. Thus, after taking the limit on the left of (6), we obtain
$`{\displaystyle \underset{I,J}{}}\left({\displaystyle \underset{i=0}{\overset{l}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{l}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{k1mn}{k1|I||J|i}}\right)\right){\displaystyle \frac{1}{(_{iI,i>l}r_i)(_{jJ}c_j)}}`$
$`=`$ $`{\displaystyle \underset{I,J}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{k+l1mn}{k1|I||J|}}\right){\displaystyle \frac{1}{(_{iI,i>l}r_i)(_{jJ}c_j)}}`$
where the sum is over proper subsets $`I\{l+1,\mathrm{},m\}`$ and all subsets $`J[n]`$. Note that in the case that $`kml`$, this expression is precisely $`F(k,r_{l+1},\mathrm{},r_m,c)`$ so we have proved (26).
We continue with the proof of (6). We obtain a slightly more convenient expression if we replace $`I`$ by $`I[l]`$ in the preceding expression. Then the left side becomes
$$\underset{I,J}{}\left(\genfrac{}{}{0pt}{}{k+l1mn}{k+l1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}$$
(29)
where the sum is over sets $`I`$ strictly contained in $`[m]`$ and containing $`[l]`$ and $`J[n]`$.
Now we turn to the right side of (6).
The expression (8) gives $`F(ml,r_{l+1},\mathrm{},r_m,c^{}(K))`$ as a sum over certain subsets $`I`$ of $`\{l+1,\mathrm{},m\}`$ and $`J`$ of $`c^{}(K)`$. But this expression is simpler if we replace $`I`$ by $`I[l]`$ and $`J`$ by $`JK`$. Then the right side can be written
$$\underset{I[m],KJ[n]}{}\mathrm{Pr}(c,K)\left(\genfrac{}{}{0pt}{}{k+l1mn}{k+l1|I||J|}\right)\frac{1}{(_{iI}r_i)(_{jJ}c_j)}$$
where the sum is over proper subsets $`I[m]`$ containing $`[l]`$, subsets $`KJ[n]`$ such that $`|K|=k+lm`$.
Now we have shown that both the left and right sides of (6) are linear combinations of the same reciprocal sums $`1/(_{iI}r_i)`$, so to prove (6) it will suffice to prove that the coefficients of the same reciprocal sums are equal on both sides.
For a given $`I[m]`$, the coefficient on the left and right depend only on the cardinality of $`I`$. We introduce the abbreviations $`H=k+l1|I|`$ and $`L=k+lm`$. Then $`0<LH<k`$. After using these abbreviations and equating coefficients we are reduced to proving
$$\underset{J[n]}{}\left(\genfrac{}{}{0pt}{}{Ln1}{H|J|}\right)\frac{1}{_{jJ}c_j}=\underset{KJ[n]}{}\mathrm{Pr}(c,K)\left(\genfrac{}{}{0pt}{}{Ln1}{H|J|}\right)\frac{1}{_{jJ}c_j}$$
(30)
where in the sum on the right the subset $`K`$ must have cardinality $`L`$.
Now use the expression (12) for $`\mathrm{Pr}(c,K)`$ to rewrite the right side of (30) as
$`{\displaystyle \underset{AKJ[n]}{}}\left((1)^{L|A|}{\displaystyle \frac{_{jK}c_j}{_{jA}c_j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|J|}}\right){\displaystyle \frac{1}{_{jJ}c_j}}`$
where we still require that $`|K|=L`$. We sum this first over $`K`$. In the term $`_{jK}c_j`$, the number of $`K`$’s for which a given $`c_j`$ occurs depends only on whether $`j`$ belongs to $`J`$. Thus, we can rewrite the right side of (30) as
$`{\displaystyle \underset{AJ}{}}(1)^{L|A|}\left({\displaystyle \frac{\left(\genfrac{}{}{0pt}{}{|J||A|}{L|A|}\right)_{jJ}c_j+\left(\genfrac{}{}{0pt}{}{|J||A|1}{L|A|}\right)_{jJA}c_j}{\left(_{jA}c_j\right)\left(_{jJ}c_j\right)}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|J|}}\right).`$
The numerator of the fraction can be rewritten as
$`\left({\displaystyle \genfrac{}{}{0pt}{}{|J||A|}{L|A|}}\right){\displaystyle \underset{jJ}{}}c_j+\left({\displaystyle \genfrac{}{}{0pt}{}{|J||A|1}{L|A|}}\right){\displaystyle \underset{jA}{}}c_j\left({\displaystyle \genfrac{}{}{0pt}{}{|J||A|1}{L|A|}}\right){\displaystyle \underset{jJ}{}}c_j`$
$`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{|J||A|1}{L|A|1}}\right){\displaystyle \underset{jJ}{}}c_j+\left({\displaystyle \genfrac{}{}{0pt}{}{|J||A|1}{L|A|}}\right){\displaystyle \underset{jA}{}}c_j`$
so the right side of (30) can be rewritten as a sum of two terms:
$$\underset{A}{}(1)^{L|A|}\left(\underset{JA}{}\left(\genfrac{}{}{0pt}{}{|J||A|1}{L|A|1}\right)\left(\genfrac{}{}{0pt}{}{Ln1}{H|J|}\right)\right)\frac{1}{_{jA}c_j}$$
(31)
and
$$\underset{J}{}\left(\underset{AJ}{}(1)^{L|A|}\left(\genfrac{}{}{0pt}{}{|J||A|1}{L|A|}\right)\right)\left(\genfrac{}{}{0pt}{}{Ln1}{H|J|}\right)\frac{1}{_{jJ}c_j}.$$
(32)
Now, comparing with the left side of (30), it suffices to show that the inner sum in (31) equals $`(1)^{L|A|}\left(\genfrac{}{}{0pt}{}{Ln1}{H|A|}\right)`$ when $`|A|<L`$ and 0 otherwise, and the inner sum in (32) equals 1.
It is easy to see that the inner sum of (31) equals 0 when $`|A|=L`$, since the first binomial coefficient has a negative lower term in that case. When $`|A|<L`$, because $`|J|L`$, the inner sum is over $`J`$ strictly larger than $`A`$. Collecting terms according to the cardinality $`j`$ of $`J`$, we obtain
$`{\displaystyle \underset{j=|A|+1}{\overset{H}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n|A|}{j|A|}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{j|A|1}{L|A|1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{Hj}}\right)`$
$`={\displaystyle \underset{i=1}{\overset{H|A|}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n|A|}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{i1}{L|A|1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|A|i}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{n|A|}{0}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{01}{L|A|1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|A|0}}\right)`$
$`+{\displaystyle \underset{i=0}{\overset{H|A|}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n|A|}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{i1}{L|A|1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|A|i}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{1}{L|A|1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|A|}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{nL+H|A|}{H|A|}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{L1H}}\right)`$
$`=(1)^{L|A|}\left({\displaystyle \genfrac{}{}{0pt}{}{Ln1}{H|A|}}\right)`$
where the third equality holds by substituting into the identity (\[R\], p.16)
$$\left(\genfrac{}{}{0pt}{}{m}{p}\right)\left(\genfrac{}{}{0pt}{}{n}{q}\right)=\underset{i=0}{\overset{p}{}}\left(\genfrac{}{}{0pt}{}{n+i}{p+q}\right)\left(\genfrac{}{}{0pt}{}{mn+q}{i}\right)\left(\genfrac{}{}{0pt}{}{nm+p}{pi}\right)$$
and the fourth equality holds because $`L1H<0.`$
In the inner sum of (32) we can collect terms according to the cardinality $`a`$ of $`A`$. Recalling that $`A`$ must be contained in $`J`$, we obtain
$`{\displaystyle \underset{a}{}}(1)^{La}\left({\displaystyle \genfrac{}{}{0pt}{}{|J|}{a}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{|J|a1}{La}}\right)={\displaystyle \underset{a}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{|J|}{a}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{L|J|}{La}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{L}{L}}\right)=1.`$
This proves Theorem 6. $`\mathrm{}`$
Finally we prove that $`E`$ and $`F`$ have another property in common.
###### Theorem 7
Both $`E(k,r,c)`$ and $`F(k,r,c)`$ are monotonically decreasing functions of $`r`$ and $`c`$. In particular, if $`r_1<r_1^{}`$, then
$$E(k,r_1,r_2,\mathrm{},r_m,c)>E(k,r_1^{},r_2,\mathrm{},r_m,c)$$
and
$$F(k,r_1,r_2,\mathrm{},r_m,c)>F(k,r_1^{},r_2,\mathrm{},r_m,c).$$
Furthermore, $`F`$ is differentiable to any degree $`\mathrm{}0`$ in each $`r_i`$ and $`c_j`$, and $`(1)^{\mathrm{}}\frac{^{\mathrm{}}F}{r_i^{\mathrm{}}}>0`$.
Proof: Without loss of generality, because of symmetry, we can restrict ourselves to considering the behavior of $`E`$ and $`F`$ as functions of $`r_1`$.
Recall from the proof of Theorem 5 that $`E(k,r,c)`$ is the expectation of the random variable $`\mathrm{min}_k(Z/A)`$, where $`Z`$ is an $`m\times n`$-matrix-valued random variable with exponentially distributed independent entries of mean 1, where $`A=(r_ic_j)`$ is the rank 1 rate matrix, and where $`Z/A`$ denotes the element by element quotient. Suppose $`r_1<r_1^{}`$ and $`r_i=r_i^{}`$ for $`i=2,\mathrm{},m`$. Define $`A^{}=(r_i^{}c_j)`$. Then $`AA^{}`$ term by term, so $`Z/AZ/A^{}`$, and $`\mathrm{min}_k(Z/A)\mathrm{min}_k(Z/A^{})`$. Hence,
$$E(k,r,c)=E(\underset{k}{\mathrm{min}}(Z/A))E(\underset{k}{\mathrm{min}}(Z/A^{}))=E(k,r^{},c).$$
Since there is a nonzero probability that the minimum $`k`$-assignment of $`Z/A`$ uses the first row, the inequality is actually strict.
We start with the following formula for $`F`$, from (9) and (10):
$$F(k,r,c)=\underset{|I^{}|+|J^{}|<k}{}\left(\underset{II^{}}{}(1)^{|I^{}||I|}\frac{1}{_{iI}r_i}\right)\left(\underset{JJ^{}}{}(1)^{|J^{}||J|}\frac{1}{_{jJ}c_j}\right)$$
For $`\mathrm{}1`$, $`I^{}[m]`$, and $`J^{}[n]`$, we define the functions
$`f(\mathrm{},r,I^{})`$ $`=`$ $`{\displaystyle \underset{II^{}}{}}(1)^{|I^{}||I|}{\displaystyle \frac{1}{(_{iI}r_i)^{\mathrm{}}}}`$
$`g(\mathrm{},c,J^{})`$ $`=`$ $`{\displaystyle \underset{JJ^{}}{}}(1)^{|J^{}||J|}{\displaystyle \frac{1}{(_{jJ}c_j)^{\mathrm{}}}},`$
so that
$$F(k,r,c)=\underset{|I^{}|+|J^{}|<k}{}f(1,r,I^{})g(1,c,J^{}).$$
(33)
First we prove that $`f(\mathrm{},r,I^{})>0`$. We define the partial sum $`R=_{iI^{}}^mr_i`$. Then
$`0`$ $`<`$ $`{\displaystyle _0^{\mathrm{}}}t^\mathrm{}1e^{Rt}\left({\displaystyle \underset{iI^{}}{}}(1e^{r_it})\right)𝑑t`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}t^\mathrm{}1{\displaystyle \underset{II^{}}{}}(1)^{|I|}\mathrm{exp}(t(R{\displaystyle \underset{iI}{}}r_i))dt`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}t^\mathrm{}1{\displaystyle \underset{II^{}}{}}(1)^{|I^{}||I|}\mathrm{exp}(t({\displaystyle \underset{iI}{}}r_i))dt`$
$`=`$ $`{\displaystyle \underset{II^{}}{}}(1)^{|I^{}||I|}{\displaystyle _0^{\mathrm{}}}t^\mathrm{}1\mathrm{exp}(t({\displaystyle \underset{iI}{}}r_i))𝑑t`$
$`=`$ $`(\mathrm{}1)!{\displaystyle \underset{II^{}}{}}(1)^{|I^{}||I|}{\displaystyle \frac{1}{(_{iI}r_i)^{\mathrm{}}}}`$
$`=`$ $`(\mathrm{}1)!f(\mathrm{},r,I^{}),`$
so $`f(\mathrm{},r,I^{})>0`$ as claimed. The proof that $`g(\mathrm{},c,J^{})>0`$ is essentially the same.
To finish the proof of the theorem it suffices to show that for $`\mathrm{}1`$,
$$(\frac{}{r_1})^{\mathrm{}}F(k,r,c)=\underset{|I^{}|+|J^{}|=k1,1I^{}}{}\mathrm{}!f(\mathrm{}+1,r,I^{})g(1,c,J^{}),$$
(34)
because all terms on the right hand side are positive. First, rewrite (33) as follows
$`F(k,r,c)`$ $`=`$ $`{\displaystyle \underset{|I^{}|+|J^{}|<k1,1I^{}}{}}(f(1,r,I^{})+f(1,r,I^{}\{1\}))g(1,c,J^{})`$
$`+`$ $`{\displaystyle \underset{|I^{}|+|J^{}|=k1,1I^{}}{}}f(1,r,I^{})g(1,c,J^{}).`$
The functions $`f(1,r,I^{})+f(1,r,I^{}\{1\})`$ and $`g(1,c,J^{})`$ are independent of $`r_1`$, so those partial derivatives vanish. Meanwhile, if $`iI^{}`$, then all the denominators in $`f(1,r,I^{})`$ involve $`r_1`$, so we have
$$(\frac{}{r_1})^{\mathrm{}}f(1,r,I^{})=\mathrm{}!f(\mathrm{}+1,r,I^{}).$$
This proves (34) and Theorem 7. $`\mathrm{}`$
The functions $`E(k,r,c)`$ and $`F(k,r,c)`$ share many properties. Both functions are rational, homogeneous of degree $`1`$, and symmetric in the $`r_i`$’s and the $`c_j`$’s. Both rational functions have denominators that factor into linear factors which are either sums of $`r_i`$’s or sums of $`c_j`$’s. Both are positive, and both are monotonically decreasing in each $`r_i`$ and $`c_j`$. Finally, $`E`$ and $`F`$ share various limit properties with $`r_i`$ or $`c_j`$ tending to 0. We could also consider limits as $`m`$ or $`n`$ go to $`\mathrm{}`$. However, it is conceivable that the properties we have already found are sufficient to guarantee that such a function is unique. In any case, one plan to prove $`E=F`$ would be to extend the list of common properties until equality is forced.
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# Cooper pair dispersion relation in two dimensions
## Abstract
The Cooper pair binding energy vs. center-of-mass-momentum dispersion relation for Bose-Einstein condensation studies of superconductivity is found in two dimensions for a renormalized attractive delta interaction. It crosses over smoothly from a linear to a quadratic form as coupling varies from weak to strong.
For the attractive interfermion potential
$$V(r)=v_0\delta (𝐫),$$
(1)
where $`v_00`$ is the interaction strength, one can apply the Lippmann-Schwinger as well as the Cooper-pair (CP) equation in two dimensions (2D) for two fermions of mass $`m`$ with momenta wavevectors $`𝐤_1`$ and $`𝐤_2`$ in free space (i.e., vacuum) and in the momentum space above the filled Fermi sea, respectively. Combining these two equations so as to eliminate (the regularized, infinitesimally small ) $`v_0`$, one obtains the renormalized CP equation
$`{\displaystyle \underset{k}{}}{\displaystyle \frac{1}{B_2+\mathrm{}^2k^2/m}}`$
$`{\displaystyle \underset{k}{}}{}_{}{}^{}{\displaystyle \frac{1}{\mathrm{}^2k^2/m+\mathrm{\Delta }_K2E_F+\mathrm{}^2K^2/4m}}=0,`$ (2)
where $`𝐤(𝐤_1𝐤_2)/2`$ is the relative, $`𝐊𝐤_1+𝐤_2`$ the center-of-mass momentum (CMM), usually taken as zero in BCS theory, $`E_F`$ the Fermi energy, $`B_20`$ the (single-bound-state) pair binding energy in vacuum, and $`2E_F\mathrm{\Delta }_K`$ the total energy of the CP with $`\mathrm{\Delta }_K0`$ its binding energy now simply a function of $`B_2`$. The prime on the second summation denotes the restriction $`|𝐤\pm 𝐊/2|>k_F`$. In 2D one finds that $`\mathrm{\Delta }_K=B_2`$ exactly, but only for $`K0`$. Expanding Eq. (2) for small but nonzero $`K`$ and subtracting from the expression for $`K=0`$ gives a small-CMM expansion valid for any dimensionless coupling $`B_2/E_F=\mathrm{\Delta }_0/E_F`$, namely
$`\epsilon _K(\mathrm{\Delta }_0\mathrm{\Delta }_K)={\displaystyle \frac{2}{\pi }}\mathrm{}v_FK+`$
$`\left[1\left\{2\left({\displaystyle \frac{4}{\pi }}\right)^2\right\}{\displaystyle \frac{E_F}{B_2}}\right]{\displaystyle \frac{\mathrm{}^2K^2}{2(2m)}}+O(K^3),`$ (3)
where a nonnegative CP excitation energy $`\epsilon _K`$ has been defined, and the Fermi velocity $`v_F`$ is given by $`E_F/k_F=\mathrm{}v_F/2`$. It is this excitation energy that must be inserted into the Bose-Einstein (BE) distribution function to determine the critical temperature in a picture of superconductivity (or of superfluidity in, e.g., liquid <sup>3</sup>He) as a BE condensation (BEC) of CPs .
The leading term in (3) is linear in the CMM, followed by a quadratic term; similar results hold in 3D. In the strong coupling limit ($`B_2=\mathrm{\Delta }_0E_F`$) the quadratic term is exactly the CM kinetic energy of what was originally a CP (becoming what is sometimes called a “local pair”), namely,
$$\underset{\mathrm{\Delta }_0E_F}{lim}\epsilon {}_{K}{}^{}=\frac{\mathrm{}^2K^2}{2(2m)},$$
(4)
the familiar nonrelativistic kinetic energy of the composite pair of mass $`2m`$ and CMM $`K`$ in vacuum. It is this dispersion relation that has been assumed in virtually all BEC studies in 3D of superconductivity (see, e.g., among others)—but which in 2D is well-known to give zero transition BEC temperature $`T_c`$. However, recent BCS-Bose crossover picture root-mean-square radii calculations compared with empirical coherence lengths of several typical 2D-like cuprates suggest these superconductors to be well within the BCS or weak-coupling regime, implying that the (nearly) linear relation is appropriate for them rather than the quadratic one.
Fig. 1 displays exact numerical results obtained from (2) for different couplings of a dimensionless CP excitation energy $`\epsilon {}_{K}{}^{}/\mathrm{\Delta }_0`$ as function of $`K/k_F`$. For weak enough coupling (and/or high enough density) the exact dispersion relation is practically linear over almost the entire interval of $`K/k_F`$ up to breakup—despite the divergence in (3) of the quadratic and possibly higher-order terms as $`B_20`$. For stronger coupling the quadratic dispersion relation (4) begins to dominate, i.e., as $`B_2/E_F\mathrm{}`$, meaning that either $`B_2\mathrm{}`$ or $`E_F0`$, the latter implying that the Fermi sea vanishes and the vacuum is recovered as a limit.
In conclusion, the CP problem with nonzero CMM evolves with weak to strong zero-range pairwise interaction to give a linear dispersion relation in the CMM that gradually crosses over to a quadratic relation with increasing coupling. These results will play a critical role in a model of superconductivity (or superfluidity) based on BE condensation of CPs, particularly in 2D where only the linear relation can give rise to nonzero $`T_c`$’s.
Partial support from UNAM-DGAPA-PAPIIT (México) # IN102198, CONACyT (México) # 27828 E, DGES (Spain) # PB95-0492 and FAPESP (Brazil) is gratefully acknowledged.
Figure Caption
1. Dimensionless exact CP excitation energy $`\epsilon _K/\mathrm{\Delta }_0`$ vs $`K/k_F`$ in 2D (full curves) calculated from (2) for different $`B_2/E_F`$. Dot-dashed line is linear approximation (virtually coincident with the exact curve for all $`B_2/E_F\stackrel{<}{}\mathrm{\hspace{0.33em}0.1}`$) while dashed curve is exact finite range result for a potential more general than (1) and whose double Fourier transform is $`V_{pq}(v_0/L^2)[(1+p^2/p_0^2)(1+q^2/p_0^2)]^{1/2}`$ (see Nozières and Schmitt-Rink ), which for $`p_0\mathrm{}`$ becomes (1). Dots mark values of CMM wavenumber where the CP breaks up, i.e., where $`\mathrm{\Delta }_K`$ turns negative.
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# Decoherence and thermalization dynamics of a quantum oscillator
## 1 Introduction
Recently, a significant interest to the decoherence processes in quantum mechanics is observed, in particular, due to the problem of stability of quantum superpositions (frequently modeled by some kinds of the ‘Schrödinger cats’ -) under the influence of the environment -. However, despite that the qualitative picture of the phenomenon seems more or less clear, there are no unique quantitative measures of the rate of decoherence or the rate of thermalization. The decoherence implies a degradation of the quantum interference effects (manifesting themselves in the existence of quantum superpositions) due to the interactions with the ‘external world’. Since these effects are inherent to the pure quantum states, while they disappear in quantum mixtures, it seems natural, on the face of it, to identify the ‘degree of decoherence’ with the degree of ‘impurity’ of the quantum state, expressed in terms of the ‘canonical entropy’ $`S=\mathrm{Tr}\left(\widehat{\rho }\mathrm{ln}\widehat{\rho }\right)`$ or in terms of the ‘linear entropy’ $`s=1\mathrm{Tr}\widehat{\rho }^2`$, which is more simple for calculations. However, a deeper analysis shows that such an identification leads to certain difficulties, especially in the low temperature case.
Indeed, let us consider the evolution of an initial pure ($`s(0)=S(0)=0`$) quantum state due to a weak interaction with a large reservoir at low temperature. For $`t>0`$, $`s(t)`$ and $`S(t)`$ assume positive values, so the rate of increase of $`s(t)`$ or $`S(t)`$ at $`t0`$ can provide us some hints to the time scale of the initial phase of the decoherence process . However, tracing the evolution of the entropies for the long time interval, we discover that for a small enough temperature of the environment, the entropies, after reaching some maxima, finally decrease to very small values which tend to zero when $`T0`$ (because the stationary thermal mixed state is very close in this case to the ground pure state). Then, identifying the measure of quantum impurity with the measure of decoherence, one should accept a strange result that the degree of decoherence of the final equilibrium state is almost the same (close to zero) as it was initially, despite that the thermal states are usually believed to be the most ‘incoherent’
This example shows that at low temperatures the entropies can serve as the measures of decoherence only at the initial stage of the decoherence process. Thus, several questions arise. The first one: is it possible to find some other measures which could be used in the whole interval $`0t<\mathrm{}`$ and for any temperature of the environment? Another question is: under which conditions (at which time scale) the usage of the entropies as the measures of the decoherence can be justified?
In the present paper, we answer both questions, introducing the new parameter $`𝒞`$ whose connection with the degree of coherence is indubitable (see section 2). This parameter equals one for pure quantum states and $`0`$ for the thermodynamical equilibrium states, for any temperature $`T>0`$. In section 3 we introduce another parameter $`𝒟`$, which can be considered as the ‘measure of thermalization’, since it equals zero for any pure state and $`1`$ for the thermodynamically equilibrium state of any quantum system with an equidistant energy spectrum, for any temperature $`T>0`$.
Following the time evolution of the parameters $`𝒞`$ and $`𝒟`$ in the process of the thermal relaxation of various initial states (Fock’s, coherent, squeezed, ‘cat’) of the harmonic oscillator, described in the framework of the ‘standard master equation’ (sections 4 and 5), we find the conditions, under which the linear entropy can serve as a reasonable ‘measure of decoherence’. Moreover, in section 6 we demonstrate that the decoherence of highly excited initial states at low temperatures goes through three distinct stages, characterized not by some unique ‘decoherence time’, but at least by two times with quite different dependences on the initial energy and the temperature.
The first time $`t_1`$ (which is usually identified with the time of decoherence) is, roughly speaking, inversly proportional to the product of the energy of quantum fluctuations by the number of photons per mode of the reservoir. During this short time interval the parameter $`𝒞`$ rapidly decreases from $`1`$ to some finite value which depends on the initial state. Then $`𝒞(t)`$ remains at a more or less constant level (for the ‘cat’ states) or even can increase with time (for the ‘squeezed’ states). And only after the ‘ultimate decoherence time’ $`t_d`$, which increases (logarithmically) with the increase of the initial energy, the coherence coefficient goes monotonously to the final zero value.
The evolution of the ‘thermalization parameter’ $`𝒟(t)`$ is similar to certain extent to the behaviour of $`1𝒞(t)`$: the initial rapid increase from zero to some intermediate value, then some stabilization or even decrease, and the final rapid transition to the equilibrium unit value after the ‘thermalization time’ $`t_T`$, which also depends logarithmically on the initial energy. The difference between $`t_d`$ and $`t_T`$ consists in their temperature dependences: $`t_d`$ has a finite limit value when the temperature $`T`$ tends to zero, whereas $`t_T`$ is inversly proportional to the absolute temperature at $`T0`$, in accordance with the third law of thermodynamics (the inattainability of the absolute zero of temperature implies that the rate of the relaxation processes must go to zero as $`T0`$).
## 2 The measure of ‘coherence’
The controversies of the identification the decoherence measure with the von Neumann’ or linear entropies take their origin in the invariance of these entropies with respect to the choice of the basis in the Hilbert space: the entropies do not distinguish the equilibrium state (or other stationary states in the case of ‘colored’ or ‘squeezed’ reservoirs) from any other mixed one. But such a symmetry with respect to the choice of the basis in the Hilbert space is obviously broken in the relaxation processses, when all possible initial states tend to the unique equilibrium state, whose density matrix is diagonal in the distinct basis formed by the energy eigenstate of the Hamiltonian (or by some other distinct basis in the case of more sophisticated artificial reservoirs considered recently). Thus it seems natural to suppose that the measure of decoherence must depend explicitly on this distinct basis $`|nn|`$ (the concept of the broken symmetry of the Hilbert space was used as a basis for introducing the polarized distance between different quantum states in ). In some special cases, when it is known exactly to which specific family of quantum states (e.g., coherent state, even/odd coherent state, etc.) the initial quantum state belongs, the prefered basis may be different from the energy one, so that some special approaches can be used, as well. We shall discuss such a situation later on. However, in the generic case, when the type of the quantum state is not known beforehand, the only available information is contained in the set of matrix elements of the statistical operator with respect to the physically distinguished energy basis: $`\widehat{\rho }(t)=_{mn}\rho _{mn}|mn|`$. Therefore, it seems reasonable to define the measure of coherence in terms of the coefficients $`\rho _{mn}`$. Since the decoherence is usually identified with the disappearence of the off-diagonal elements of the density matrix $`\rho _{mn}`$, it is natural to define the normalized measure of coherence as
$$𝒞(t)=\underset{mn}{}\left|\rho _{mn}\right|^2(t)/\underset{mn}{}\left|\rho _{mn}\right|^2(0).$$
(1)
Then $`𝒞(0)=1`$, while $`𝒞0`$ for any ‘completely incoherent’ state without off-diagonal matrix elements in the energy basis (provided initially at least one off-diagonal element was different from zero). Introducing the ‘diagonal part’ of the operator $`\widehat{\rho }`$
$$\widehat{\rho }_d=\underset{n}{}p_n|nn|,p_nn|\widehat{\rho }|n,$$
(2)
and taking into account the property $`\text{Tr}\left(\widehat{\rho }\widehat{\rho }_d\right)=\text{Tr}\left(\widehat{\rho }_d^2\right)`$ one can rewrite (1) in the form
$$𝒞=\frac{\text{Tr}\left[\widehat{\rho }(t)\widehat{\rho }_d(t)\right]^2}{\text{Tr}\left[\widehat{\rho }(0)\widehat{\rho }_d(0)\right]^2}\frac{\mu (t)\lambda (t)}{\mu (0)\lambda (0)}$$
(3)
$$\mu \text{Tr}\widehat{\rho }^2,\lambda \text{Tr}\widehat{\rho }_d^2=\underset{n}{}p_n^2$$
(4)
We shall call $`\mu `$ the ‘total purity’ and $`\lambda `$ the ‘diagonal purity’. In many cases of practical interest both ‘purities’ can be calculated rather easily. For example, if one knows the Wigner function (we assume $`\mathrm{}1`$)
$$W(q,p)=𝑑ve^{ipv}q\frac{v}{2}\left|\widehat{\rho }\right|q+\frac{v}{2}$$
(5)
$$\mathrm{Tr}\widehat{\rho }=W(q,p)𝑑q𝑑p/(2\pi )=1$$
then
$$\mu =W^2(q,p)𝑑q𝑑p/(2\pi ).$$
(6)
As to the ‘diagonal purity’, it can be calculated either by means of a direct summation of the series in (4), or, equivalently, it can be expressed as the integral
$$\lambda (t)=_0^{2\pi }\frac{d\phi }{2\pi }\left|G(e^{i\phi };t)\right|^2,$$
(7)
where
$$G(z;t)\underset{n=0}{\overset{\mathrm{}}{}}p_n(t)z^n$$
(8)
is the diagonal generating function. Below we use both methods.
## 3 The measure of thermalization
A qualitative measure of thermalization can be introduced in the following way. The analysis of the low temperature behaviour of the entropies shows that the troubles mentioned in the introduction arise due to the double nature of the ground state, described by the density operator $`\widehat{\rho }_0|00|`$. On one hand, this state is pure, with $`\mathrm{Tr}\widehat{\rho }_0^2=1`$. On the other hand, it is the limit of the equilibrium states, which are conceived to be completely decoherent. Therefore it seems reasonable to exclude the state $`\widehat{\rho }_0`$ in some way. One of the possibilities is to take a simple expression for the linear entropy and to divide it by a proper time-dependent factor which would ensure a nonzero limit at $`t\mathrm{}`$. This goal can be achieved, for instance, if one chooses as the normalizing factor the Hilbert-Schmidt distance between the states $`\widehat{\rho }(t)`$ and $`\widehat{\rho }_0`$. If the system under study has a finite number of energy levels (e.g., spin systems), then there are some grounds to treat the state with the maximal energy $`\widehat{\rho }_f|E_{max}E_{max}|`$ on the same footing as $`\widehat{\rho }_0`$. Thus we arrive at the parameter (introduced for the first time in , but identified erroneously with the measure of decoherence)
$`𝒟`$ $`=`$ $`{\displaystyle \frac{1\mathrm{Tr}\widehat{\rho }^2}{\left[\mathrm{Tr}\left(\widehat{\rho }\widehat{\rho }_0\right)^2\mathrm{Tr}\left(\widehat{\rho }\widehat{\rho }_f\right)^2\right]^{1/2}}}`$ (9)
$`=`$ $`{\displaystyle \frac{1\mu }{\left[\left(1+\mu p_f\right)\left(1+\mu p_0\right)\right]^{1/2}}},`$
where $`p_0\mathrm{Tr}\left(\widehat{\rho }\widehat{\rho }_0\right)=0|\widehat{\rho }|0`$ is the ground state occupation probability, while $`p_f\mathrm{Tr}\left(\widehat{\rho }\widehat{\rho }_f\right)`$ is the occupation probability of the level with the maximal energy (evidently, for quantum systems with infinite dimensional Hilbert spaces, such as a harmonic oscillator, $`p_f0`$ for any physical state possessing finite energy).
For systems with equidistant spectra, $`E_{n+1}E_n=\mathrm{\Delta }E=const`$, the equilibrium occupation probabilities read
$$p_n=\xi ^n(1\xi )/\left(1\xi ^M\right)$$
where $`M`$ is the total number of levels, $`n=0,1,\mathrm{},M1`$ and $`\xi =\mathrm{exp}(\beta \mathrm{\Delta }E)<1`$ is the Boltzmann factor. Then
$$\mu _{eq}=(1\xi )\left(1+\xi ^M\right)/\left[(1+\xi )\left(1\xi ^M\right)\right],$$
$$p_0^{(eq)}=(1\xi )/\left(1\xi ^M\right)$$
$$p_f^{(eq)}=\xi ^{M1}(1\xi )/\left(1\xi ^M\right)$$
and we see that $`𝒟_{eq}1`$ for any value $`0<\xi <1`$, while $`𝒟0`$ for all pure states. For this reason, we may consider $`𝒟`$ as ‘the measure of thermalization’. Note that we have an indeterminacy in (9) if $`\xi =T=0`$, thus the case $`T=0`$ must be excluded. But as we know the limit of exact zero temperature is an idealization, thus we do not have to worry on this issue.
For systems with nonequidistant spectra the value of $`𝒟_{eq}`$ depends on temperature, nonetheless the limits at $`T0`$ and $`T\mathrm{}`$ still equal $`1`$. For $`T\mathrm{}`$ we have $`p_0=p_1=\mathrm{}=p_f=1/M`$, consequently $`\mu =M(1/M)^2=1/M`$ and $`𝒟=1`$. In the low temperature case $`T0`$, the equilibrium statistical operator is close to $`p_0|00|+p_1|11|`$ with $`p_11`$ (where $`|1`$ is the first excited state), while the contribution of other states can be neglected (we consider the systems with discrete energy spectra). Then $`p_0+p_1=1`$, $`\mu =p_0^2+p_1^2`$, $`1+\mu 2p_0=2p_1^2`$, $`1+\mu 2p_f=2`$, and $`1\mu =2p_1`$ (up to higher order terms). As a result, we have $`𝒟=1`$ again.
## 4 Decoherence dynamics
### 4.1 Time evolution of the quantum state
We confine ourselves to the analysis of the time dependence of the ‘coherence coefficient’ $`𝒞`$ (3) in the process of thermal relaxation of the harmonic oscillator described in the framework of the standard master equation (more general models were considered, e.g., in )
$`d\widehat{\rho }/dt=\gamma \left(1+\nu \right)\left(2\widehat{a}\widehat{\rho }\widehat{a}^{}\widehat{a}^{}\widehat{a}\widehat{\rho }\widehat{\rho }\widehat{a}^{}\widehat{a}\right)`$
$`+\gamma \nu \left(2\widehat{a}^{}\widehat{\rho }\widehat{a}\widehat{a}\widehat{a}^{}\widehat{\rho }\widehat{\rho }\widehat{a}\widehat{a}^{}\right)i[\widehat{a}^{}\widehat{a},\widehat{\rho }].`$ (10)
Here $`\widehat{a}`$ and $`\widehat{a}^{}`$ are the usual bosonic annihilation and creation operators, $`\nu `$ is the equilibrium mean number of quanta in the reservoir corresponding to the given mode, and $`\gamma >0`$ is a damping coefficient ($`\mathrm{}=\omega =1`$).
An immediate consequence of equation (10) is the universal expression for the purity loss rate in an initial pure state $`\widehat{\rho }^2=\widehat{\rho }=|\psi \psi |`$ (cf. )
$$\dot{\mu }|_{t=0}=2\mathrm{T}\mathrm{r}\left(\widehat{\rho }\widehat{\dot{\rho }}\right)|_{t=0}=4\gamma \left[\nu +(1+2\nu )\sigma _a\right],$$
(11)
where $`\sigma _a\psi |\widehat{a}^{}\widehat{a}|\psi \left|\psi |\widehat{a}|\psi \right|^2`$. The ‘primary’ purity loss rate is minimal for the coherent states with $`\sigma _a0`$. In the generic case it is roughly proportional to the average number of thermal photons in the reservoir and to the ‘reduced’ energy of quantum fluctuations in the initial pure state,
$$_0=\frac{1}{2}\left[\sigma _p^{(0)}+\sigma _q^{(0)}\right]\sigma _a+\frac{1}{2},$$
(12)
where $`\sigma _q^{(0)}`$ and $`\sigma _p^{(0)}`$ are the variances of the quadrature components $`\widehat{q}=\left(\widehat{a}+\widehat{a}^{}\right)/\sqrt{2}`$ and $`\widehat{p}=i\left(\widehat{a}^{}\widehat{a}\right)/\sqrt{2}`$ in the initial pure state $`|\psi `$.
Due to equation (11) the initial evolution of the ‘purity’ has the linear dependence on time $`\mu (t)=1t/t_1+\mathrm{}`$, where the ‘primary purity loss time’ equals
$$t_1=(4\gamma )^1\left[\nu +(1+2\nu )\sigma _a\right]^1.$$
(13)
Note that some ‘microscopic’ models, based on an explicit coupling of the system under study with a large reservoir, result in a quadratic time dependence $`\mu (t)`$ at $`t0`$ . This apparent contradiction is explained by the fact that the ‘microscopic’ models and the phenomenological master equations describe the evolution of the subsystem in different time scales. Actually, the master equation describes a ‘coarse-grained’ evolution averaged over many periods of the fast oscillation, so the physical meaning of the limit $`t0`$ in the case of the master equation is quite different from the same (formally) limit in the ‘microscopic’ models.
To calculate the time dependence of the ‘purity’ $`\mu `$ in the whole interval $`0<t<\mathrm{}`$ with the aid of formula (6) we need the time dependent Wigner function $`W(q,p,t)`$. It obeys the Fokker–Planck equation which follows immediately from (10):
$`{\displaystyle \frac{W}{t}}`$ $`=`$ $`{\displaystyle \frac{}{q}}\left([\gamma qp]W\right)+{\displaystyle \frac{}{p}}\left([\gamma p+q]W\right)`$ (14)
$`+\gamma \left(\nu +{\displaystyle \frac{1}{2}}\right)\left({\displaystyle \frac{^2W}{q^2}}+{\displaystyle \frac{^2W}{p^2}}\right).`$
The solution to equation (14) can be written as
$$W(q,p;t)=𝒦(q,p;t|q^{},p^{},0)W(q^{},p^{};0)𝑑q^{}𝑑p^{}.$$
(15)
The propagator $`𝒦(q,p;t|q^{},p^{},0)`$ was calculated by means of different methods in ; the explicit expressions in the case of most general multidimensional time-dependent quadratic operator (with respect to $`q`$, $`p`$, $`/q`$, $`/p`$) in the right-hand side of the Fokker–Planck equation were given in . In the case involved the general form found in is reduced to (see appendix)
$`𝒦(q,p;t|q^{},p^{},0)=\left(2\pi \sigma _{}u\right)^1\mathrm{exp}\{\left(2\sigma _{}u\right)^1[q_t^2+p_t^2`$
$`+e^{2\gamma t}(q^2+p^2)2e^{\gamma t}(q^{}q_t+p^{}p_t)]\},`$ (16)
where $`\sigma _{}\nu +\frac{1}{2}`$,
$$q_t=q\mathrm{cos}tp\mathrm{sin}t,p_t=p\mathrm{cos}t+q\mathrm{sin}t,$$
and the ‘compact time’ $`u`$ is given by
$$u(t)1e^{2\gamma t}.$$
(17)
The consequence of the master equation (10) is the closed set of equations for the diagonal elements of the density matrix in the Fock (energy) basis
$`\dot{p}_n`$ $`=`$ $`2\gamma (1+\nu )\left[(n+1)p_{n+1}np_n\right]`$ (18)
$`+2\gamma \nu \left[np_{n1}(n+1)p_n\right].`$
These equations, in turn, are equivalent to the simple first order partial differential equation for the diagonal generating function (8)
$$\frac{G}{t}=2\gamma (1z)[1+\nu (1z)]\frac{G}{z}2\gamma \nu (1z)G$$
(19)
The solution to (19) reads
$$G(z,u)=\frac{1}{1+\nu u(1z)}G_0\left(\frac{z+u(1+\nu )(1z)}{1+\nu u(1z)}\right)$$
(20)
where $`G_0(z)G(z,0)`$. Putting $`z=1`$ in (20) we verify the normalization condition $`G(1,t)1`$.
### 4.2 Initial coherent states
As the first example we consider the evolution of the initial coherent state $`|\alpha `$, $`\alpha \sqrt{a}\mathrm{exp}(i\varphi )`$. Applying the propagator (16) to the initial Wigner function
$$W^{(coh)}(0)=2\mathrm{exp}\left[(q\sqrt{2a}\mathrm{cos}\varphi )^2(p\sqrt{2a}\mathrm{sin}\varphi )^2\right]$$
we obtain
$`W^{(coh)}(q,p,t)`$ $`=`$ $`2\xi _\nu \mathrm{exp}\{\xi _\nu ([q\sqrt{2a}e^{\gamma t}\mathrm{cos}(\varphi t)]^2`$ (21)
$`[p\sqrt{2a}e^{\gamma t}\mathrm{sin}(\varphi t)]^2)\},`$
where
$$\xi _\nu (u)(1+2u\nu )^1$$
(22)
coincides with the ‘total purity’: $`\mu ^{(coh)}=\xi _\nu (u)`$.
To calculate the ‘diagonal purity’ we use the explicit expression for the time-dependent diagonal matrix elements in terms of the Laguerre polynomials
$$p_n^{(coh)}=\frac{(u\nu )^n}{(1+u\nu )^{n+1}}\mathrm{exp}\left[\frac{a(u1)}{1+u\nu }\right]L_n\left[\frac{a(u1)}{u\nu (1+u\nu )}\right]$$
(23)
Then the sum in (4) is reduced to the known series
$$\underset{n=0}{\overset{\mathrm{}}{}}L_n(x)L_n(y)z^n=(1z)^1\mathrm{exp}\left[z\frac{x+y}{z1}\right]I_0\left[2\frac{\sqrt{xyz}}{1z}\right]$$
(24)
($`I_0(z)`$ is the modified Bessel function), so we obtain
$$\lambda ^{(coh)}=\xi _\nu (u)\mathrm{exp}\left(\eta \right)I_0\left(\eta \right),$$
(25)
$$\eta =2a(1u)\xi _\nu (u).$$
(26)
For $`a1`$ it is sufficient to take into account the first terms of the Taylor expansion of the function $`e^\eta I_0(\eta )`$ to obtain
$$𝒞^{(coh)}(t)\xi _\nu ^2(u)(1u)=\mu ^2(t)e^{2\gamma t}.$$
(27)
In this case the time dependence of the ‘purity’ has a little in common with the time dependence of the ‘coherence’; the same is true even for $`a1`$: see figure 1.
For the highly excited initial states with $`a1`$ the asymptotics of the modified Bessel function, $`I_0(x)(2\pi x)^{1/2}e^x`$ for $`x1`$, yields
$$𝒞^{(coh)}(t)\mu ^{(coh)}(t)\left[1\frac{1}{2\sqrt{a\pi }}\left(\sqrt{1+2u\nu }e^{\gamma t}1\right)\right].$$
In this case, the contribution of the diagonal elements to the total purity is small compared with the contribution of the off-diagonal terms, therefore the sum over $`mn`$ in (1) can be replaced by the sum over all values of $`m`$ and $`n`$, and the correlation coefficient can be approximated by the purity $`\mu `$. This is just the case considered in most of the papers devoted to the decoherence of initial ‘macroscopic’ quantum states. However, $`𝒞\mu `$ only under the condition $`\lambda \mu `$. Consequently, the linear entropy $`1\mu `$ can be considered as a measure of decoherence only for $`a1`$, and under the additional restriction $`a\xi _\nu (u)\mathrm{exp}(2\gamma t)1`$, i.e., at the time scale
$$tt_{}(2\gamma )^1\mathrm{ln}[(a+2\nu )/(1+2\nu )].$$
(28)
For $`tt_{}`$ the identification of the ‘purity’ with the ‘coherence’ leads to incorrect results as mentioned in the introduction.
The ‘final decoherence time’ $`t_d`$ can be defined by means of the equation
$$𝒞(t_d)=\beta \mu _{eq}$$
(29)
where $`\beta <1`$ is some small number whose choice is a matter of convention (say, $`\beta =0.1`$), and $`\mu _{eq}=(1+2\nu )^1`$ is the equilibrium value of the ‘purity’ (we take into account that $`𝒞`$ is proportional to $`\xi _\nu =\mu ^{(coh)}`$, according to equation (25)). To solve equation (29) for sufficiently small $`\beta `$ we can use the asymptotical form of the coherence function at $`\gamma t1`$ (when $`u1`$)
$$𝒞^{(coh)}(t)2a(1+2\nu )^2e^{2\gamma t},a1,\gamma t1.$$
Thus we obtain the estimation
$$t_d^{(coh)}(2\gamma )^1\mathrm{ln}[2a\mu _{eq}/\beta ]$$
(30)
which holds for $`a\mu _{eq}1`$. The evolution of the ‘total purity’ and the ‘coherence coefficient’ for highly excited initial coherent states at different temperatures is shown in figure 2.
### 4.3 Initial ‘cat’ states
Now let us consider the family of the initial ‘Schrödinger cat’ states
$$|\alpha ;\phi =𝒩\left(|\alpha +e^{i\phi }|\alpha \right),$$
(31)
$$𝒩=\left(2\left[1+\mathrm{cos}\phi \mathrm{exp}(2a)\right]\right)^{1/2},a|\alpha |^2.$$
(32)
The special cases of this family are even ($`\phi =0`$) and odd ($`\phi =\pi `$) coherent states , and the Yurke-Stoler states ($`\phi =\pi /2`$) . The Wigner function of the state (31) reads (hereafter we assume $`\alpha =\sqrt{a}`$ to be real)
$`W^{(cat)}(q,p;0)=4𝒩^2\mathrm{exp}\left(q^2p^2\right)`$
$`\times \left[e^{2a}\mathrm{cosh}(\sqrt{8a}q)+\mathrm{cos}(\sqrt{8a}p+\phi )\right].`$ (33)
Applying the propagator (16) to this function we obtain (see also )
$`W^{(cat)}(q,p;t)=4𝒩^2\xi _\nu \mathrm{exp}\left[\xi _\nu \left(q^2+p^2\right)\right]`$
$`\times \{\mathrm{exp}[2a(1u)\xi _\nu ]\mathrm{cosh}\left[\xi _\nu \sqrt{8a(1u)}q_t\right]`$
$`+\mathrm{exp}[2a(1+2\nu )u\xi _\nu ]\mathrm{cos}[\xi _\nu \sqrt{8a(1u)}p_t+\phi ]\}`$ (34)
where the function $`\xi _\nu (u)`$ was defined in equation (22). Calculating the integral in (6) we find the ‘total purity’
$`\mu =2𝒩^4\xi _\nu (u)[1+4\mathrm{cos}\phi e^{2a}+\mathrm{cos}(2\phi )e^{4a}`$
$`+\mathrm{exp}[4a(1u)\xi _\nu ]+\mathrm{exp}[4au(1+2\nu )\xi _\nu ]]`$ (35)
The photon distribution function can be written as
$$p_n^{(cat)}=2𝒩^2\left[p_n^{(coh)}(a)+\mathrm{cos}\phi e^{2a}p_n^{(coh)}(a)\right],$$
(36)
where $`p_n^{(coh)}(a)`$ is given by (23). Calculating again the sum $`p_n^2`$ with the aid of formula (24) we find
$`\lambda ^{(cat)}=4𝒩^4\xi _\nu (u)\{I_0(\eta )[e^\eta +\mathrm{cos}^2\phi e^{\eta 4a}]`$
$`+2\mathrm{cos}\phi e^{2a}J_0(\eta )\}`$ (37)
where $`J_0(z)`$ means the usual Bessel function, and $`\eta `$ was defined in equation (26). If $`a1`$, then we have $`\mu \xi _\nu (u)`$ and the same formula (27) for $`\lambda `$.
For highly excited (‘macroscopic’: $`a1`$) initial cat states the phase $`\phi `$ becomes unimportant, and $`𝒩^21/2`$. Until $`a(1u)a\mathrm{exp}(2\gamma t)1`$, the ‘diagonal purity’ is small, similarly to the case of coherent states: $`\lambda ^{(cat)}\mathrm{exp}(\gamma t)/\sqrt{a}`$. Then the coherence coefficient can be replaced by the ‘total purity’ $`\mu `$, which rapidly decreases from $`1`$ to the value $`\frac{1}{2}\xi _\nu (u)`$:
$$\mu \frac{1}{2}\xi _\nu (u)\left(1+\mathrm{exp}\left[4au(1+2\nu )\xi _\nu (u)\right]\right),$$
staying at this level until $`a(1u)`$ becomes smaller than $`1`$. In particular, in the low temperature case $`\nu 1`$ we observe the ‘plateau’ $`\mu 𝒞\frac{1}{2}`$: see figure (3).
At the final stage of the evolution $`\mu `$ goes to the equilibrium value $`\mu _{eq}=(1+2\nu )^1`$ as
$$\mu \frac{1}{2}\xi _\nu (u)\left(1+\mathrm{exp}[4a(1u)\xi _\nu (u)]\right),$$
but now it becomes compatible with the ‘diagonal purity’ $`\lambda \xi _\nu (u)\mathrm{exp}[2a(1u)\xi _\nu (u)]`$, so
$$𝒞^{(cat)}\frac{1}{2}\xi _\nu (u)\left(1\mathrm{exp}[2a(1u)\xi _\nu (u)]\right)^2.$$
For $`a(1u)1`$ we have $`𝒞^{(cat)}2a^2\mu _{eq}^3\mathrm{exp}(4\gamma t)`$. Then equation (29) yields the ‘ultimate decoherence time’ $`t_d^{(cat)}(2\gamma )^1\mathrm{ln}\left[a\mu _{eq}\sqrt{2/\beta }\right]`$ (if $`a\mu _{eq}1`$), which is only slightly less than the similar time $`t_d^{(coh)}`$ (30).
### 4.4 Decoherence in the accompanying basis
In the special case of the ‘cat’ states there exists another natural choice of the ‘diagonal’ part of the statistical operator, different from (2). Namely, one can define the ‘accompanying’ diagonal operator as $`\widehat{\rho }_{mix}=\frac{1}{2}\left(|\alpha \alpha |+|\alpha \alpha |\right)`$. The Wigner function of this quantum mixture is the sum of two coherent Wigner functions (21) with opposite values of parameter $`\alpha `$:
$`W_{mix}(q,p;t)={\displaystyle \frac{1}{2}}\left[W_\alpha +W_\alpha \right]`$
$`=2\xi _\nu (u)\mathrm{exp}\left\{\xi _\nu (u)\left[q^2+p^2+2a(1u)\right]\right\}`$
$`\times \mathrm{cosh}\left[\xi _\nu (u)\sqrt{8a(1u)}q_t\right].`$ (38)
Then the ‘accompanying’ (normalized) coherence coefficient can be defined as
$`={\displaystyle \left[W^{(cat)}(t)W_{mix}(t)\right]^2𝑑q𝑑p}`$
$`\times \left\{{\displaystyle \left[W^{(cat)}(0)W_{mix}(0)\right]^2𝑑q𝑑p}\right\}^1`$ (39)
Performing the calculations we obtain
$`={\displaystyle \frac{\xi _\nu (u)\left(1\mathrm{exp}\left[4a(1u)\xi _\nu \right]\right)}{\left(1\mathrm{cos}^2\phi e^{4a}\right)\left(1e^{4a}\right)}}`$
$`\times \left(\mathrm{exp}\left[4au(1+2\nu )\xi _\nu \right]\mathrm{cos}^2\phi e^{4a}\right)`$ (40)
The time evolution of this parameter is essentially different from the evolution of the coefficient $`𝒞`$: see figure (4). This is not surprising, since the phenomena observed from the moving (‘accompanying’) frame of reference in many cases look very different, compared with their appearance in the fixed frame.
### 4.5 Initial squeezed states
Now let us consider the states possessing the Gaussian Wigner functions
$`W(q,p)=d^{1/2}\mathrm{exp}\{{\displaystyle \frac{1}{2d}}[\sigma _p(q\overline{q})^2`$
$`2\sigma _{qp}(q\overline{q})(p\overline{p})+\sigma _q(p\overline{p})^2]\},`$ (41)
where $`\sigma _q`$, $`\sigma _p`$ and $`\sigma _{qp}=\sigma _{pq}`$ are the (co)variances of the quadrature components, whereas $`\overline{q}`$ and $`\overline{p}`$ stand for the average values of these quadratures. The parameter $`d\sigma _p\sigma _q\sigma _{pq}^2`$ must satisfy the Schrödinger–Robertson uncertainty relation $`d1/4`$. It is related to the ‘purity’ of the state as $`\mu =(4d)^{1/2}`$. In the thermal state, $`\sigma _q^{(eq)}=\sigma _p^{(eq)}=\sigma _{}\frac{1}{2}+\nu `$, $`\sigma _{pq}=\overline{q}=\overline{p}=0`$.
The evolution of the five parameters determining the Gaussian state is governed by the set of equations following from the master equation (10)
$$d\overline{q}/dt=\overline{p}\gamma \overline{q},d\overline{p}/dt=\overline{q}\gamma \overline{p}$$
(42)
$`\dot{\sigma }_q`$ $`=`$ $`2\sigma _{qp}2\gamma \sigma _q+\gamma (1+2\nu )`$ (43)
$`\dot{\sigma }_p`$ $`=`$ $`2\sigma _{qp}2\gamma \sigma _p+\gamma (1+2\nu )`$ (44)
$`\dot{\sigma }_{qp}`$ $`=`$ $`\sigma _p\sigma _q2\gamma \sigma _{qp}`$ (45)
The solutions read
$`\overline{q}(t)`$ $`=`$ $`e^{\gamma t}\left[\overline{q}_0\mathrm{cos}t+\overline{p}_0\mathrm{sin}t\right]`$ (46)
$`\overline{p}(t)`$ $`=`$ $`e^{\gamma t}\left[\overline{p}_0\mathrm{cos}t\overline{q}_0\mathrm{sin}t\right]`$ (47)
$`\sigma _q(t)=\sigma _{}+e^{2\gamma t}[(\sigma _q^{(0)}\sigma _{})\mathrm{cos}^2t`$
$`+(\sigma _p^{(0)}\sigma _{})\mathrm{sin}^2t+\sigma _{pq}^{(0)}\mathrm{sin}(2t)]`$ (48)
$`\sigma _p(t)=\sigma _{}+e^{2\gamma t}[(\sigma _p^{(0)}\sigma _{})\mathrm{cos}^2t`$
$`+(\sigma _q^{(0)}\sigma _{})\mathrm{sin}^2t\sigma _{pq}^{(0)}\mathrm{sin}(2t)]`$ (49)
$$\sigma _{qp}(t)=e^{2\gamma t}\left[\sigma _{qp}^{(0)}\mathrm{cos}(2t)+\frac{1}{2}\left(\sigma _p^{(0)}\sigma _q^{(0)}\right)\mathrm{sin}(2t)\right]$$
(50)
Any initial pure Gaussian state is unitarily equivalent to the squeezed state, defined as an eigenstate of the canonically transformed operator $`\widehat{b}=\mathrm{cosh}\rho \widehat{a}+\mathrm{sinh}\rho \widehat{a}^{}`$ with a complex eigenvalue $`\alpha \sqrt{a}\mathrm{exp}(i\varphi )`$ and a real ‘squeezing parameter’ $`\rho `$ (sometimes it is called also a ‘two-photon state’ ). Therefore we parametrize the initial variances and average values as
$$\sigma _q^{(0)}=\frac{1}{2}e^{2\rho },\sigma _p^{(0)}=\frac{1}{2}e^{2\rho },\sigma _{pq}^{(0)}=0,$$
$$\overline{q}_0=\sqrt{2a}e^\rho \mathrm{cos}\varphi ,\overline{p}_0=\sqrt{2a}e^\rho \mathrm{sin}\varphi .$$
The ‘total purity’ $`\mu `$ does not depend on the first order moments of the coordinates:
$$\mu =\left[(1+2u\nu )^2+4u(1u)(1+2\nu )\mathrm{sinh}^2\rho \right]^{1/2}.$$
(51)
On the contrary, the ‘diagonal purity’ $`\lambda `$ depends on all 5 parameters of the one-dimensional Gaussian Wigner function.
The generic diagonal generating function found in can be expressed as
$$G(z)=[𝒢(z)]^{1/2}\mathrm{exp}\left(\frac{1}{D}\left[\frac{zg_1z^2g_2}{𝒢(z)}g_0\right]\right)$$
(52)
where
$$𝒢(z)=\frac{1}{4}\left[(1+z)^2+4d(1z)^2+2(\sigma _q+\sigma _p)\left(1z^2\right)\right]$$
$$D=1+2(\sigma _p+\sigma _q)+4d$$
$$g_0=\overline{p}^2(2\sigma _q+1)+\overline{q}^2(2\sigma _p+1)4\overline{p}\overline{q}\sigma _{pq}$$
$`g_1`$ $`=`$ $`2\overline{p}^2\left[\sigma _q^2+\sigma _{pq}^2+\sigma _q+{\displaystyle \frac{1}{4}}\right]+2\overline{q}^2\left[\sigma _p^2+\sigma _{pq}^2+\sigma _p+{\displaystyle \frac{1}{4}}\right]`$
$`4\overline{p}\overline{q}\sigma _{pq}(\sigma _q+\sigma _p+1)`$
$`g_2`$ $`=`$ $`2\overline{p}^2\left(\sigma _q^2+\sigma _{pq}^2{\displaystyle \frac{1}{4}}\right)+2\overline{q}^2\left(\sigma _p^2+\sigma _{pq}^2{\displaystyle \frac{1}{4}}\right)`$
$`4\overline{p}\overline{q}\sigma _{pq}(\sigma _q+\sigma _p)`$
For the initial pure squeezed states $`G(z,t)`$ can be written as
$`G(z;t)=\left(fbz+cz^2\right)^{1/2}`$
$`\times \mathrm{exp}\left[a(1u){\displaystyle \frac{FBz+Cz^2}{fbz+cz^2}}\right]`$ (53)
where
$$f=(1+u\nu )^2+(1u)(1+u+2u\nu )\mathrm{sinh}^2\rho ,$$
$$b=2u\nu (1+u\nu )+2u(1+2\nu )(1u)\mathrm{sinh}^2\rho ,$$
$$c=(u\nu )^2(1u)(1u2u\nu )\mathrm{sinh}^2\rho ,$$
$$F=\frac{1}{2}\left[1u+R(1+u+2u\nu )\right]$$
$$B=1u+Ru(1+2\nu )$$
$$C=\frac{1}{2}\left[1uR(1u2u\nu )\right]$$
$$R=\mathrm{cosh}(2\rho )\mathrm{sinh}(2\rho )\mathrm{cos}(2\varphi )$$
Using formula (7) and the relations
$$fb+c=1,FB+C=0$$
we obtain after some algebra the following integral representation for the time dependent ‘diagonal purity’ of a generic initially squeezed state:
$$\lambda =_0^{\frac{\pi }{2}}\frac{2d\gamma }{\pi \sqrt{\mathrm{\Phi }(\gamma )}}\mathrm{exp}\left[4a(1u)\left(V+Y\mathrm{sin}^2\gamma \right)\frac{\mathrm{sin}^2\gamma }{\mathrm{\Phi }(\gamma )}\right]$$
(54)
where
$$\mathrm{\Phi }(\gamma )=\left(1+2b\mathrm{sin}^2\gamma \right)^2+(1u)^2\mathrm{sinh}^2(2\rho )\mathrm{sin}^2(2\gamma )$$
$$V=R(1+2u\nu )+(1u)\left(R1+4R\mathrm{sinh}^2\rho \right)$$
$`Y`$ $`=`$ $`4uB\left[\nu (1+u\nu )+(1u)(1+2\nu )\mathrm{sinh}^2\rho \right]`$
$`+2(1u)\left(1R2R\mathrm{sinh}^2\rho \right)`$
The evolution of the ‘purity’ and the ‘coherence coefficient’ for the highly squeezed ($`\rho >1`$) initial state is illustrated in figure (5). In this case $`\mu 𝒞`$ up to the values of the dimensionless time $`\tau =2\gamma t1`$. In contradistinction to the cases of coherent or ‘cat’ states, the coherence coefficient is not monotonous function of time, but it tries to follow the increase of the ‘purity’ at $`\tau >1`$, before going finally to zero.
The integral (54) can be easily calculated in the long-time limit $`1u=\mathrm{exp}(2\gamma t)\epsilon 1`$ at zero temperature ($`\nu =0`$): $`\lambda =12\epsilon \left(\mathrm{sinh}^2\rho +aR\right)+𝒪(\epsilon ^2)`$. Comparing this expression with the similar expansion of the ‘total purity’ $`\mu `$ (51) we obtain $`𝒞\mu \lambda 2\epsilon aR=2\epsilon _{cl}`$, where
$$_{cl}\frac{1}{2}\left(\overline{q}_0^2+\overline{p}_0^2\right)=aR$$
(55)
is the initial ‘classical’ energy (the total energy without the contribution of the vacuum fluctuations). Then equation (29) yields the ‘ultimate decoherence time’ $`t_d^{(sq)}(2\gamma )^1\mathrm{ln}\left(2_{cl}/\beta \right)`$ which has the same order of magnitude as the time $`t_d^{(coh)}`$ (30) for the coherent state with the same ‘classical energy’ $`|\alpha |^2`$ (at zero temperature $`\mu _{eq}=1`$). If $`a=0`$ (the initial squeezed vacuum state), then one should calculate $`\mu `$ and $`\lambda `$ up to the second order terms with respect to $`\epsilon `$. In this case we obtain $`𝒞\mu \lambda \frac{1}{4}\epsilon ^2\mathrm{sinh}^2(2\rho )`$, and $`t_d^{(sq)}(2\gamma )^1\mathrm{ln}\left[\mathrm{sinh}(2\rho )/2\sqrt{\beta }\right]`$. Since we consider the case $`\rho 1`$, we can replace $`\mathrm{sinh}(2\rho )/2`$ by $`\mathrm{sinh}^2(\rho )=E\frac{1}{2}E`$, where $`E`$ is the total energy in the case discussed.
## 5 Thermalization dynamics
The concrete evolution of the thermalization coefficient depends on a peculiar ‘competition’ between the ‘total purity’ $`\mu (u)`$ (which was calculated in the preceding sections) and the ground state probability $`p_0(u)`$, which can can be easily found from the ‘diagonal generating function’
$$p_0(t)=\frac{1}{1+\nu u(t)}G_0\left(\frac{u(t)(1+\nu )}{1+\nu u(t)}\right).$$
(56)
In the case of initial coherent state we have ($`a|\alpha |^2`$) $`G_0^{(coh)}(z)=\mathrm{exp}\left[a(z1)\right]`$. Consequently,
$$p_0^{(coh)}=\frac{1}{1+u\nu }\mathrm{exp}\left[\frac{a(1u)}{1+u\nu }\right],$$
(57)
so the ‘thermalization coefficient’ reads
$$𝒟_a^{(coh)}(u)=\left\{1+\frac{1+2u\nu }{(u\nu )^2}\left(1\mathrm{exp}\left[\frac{a(1u)}{1+u\nu }\right]\right)\right\}^{1/2}$$
Evidently, the case $`\nu =0`$ should be excluded in this expression. The dependence on the displacement parameter $`a`$ disappears for $`a(1u)1`$, when all the functions $`𝒟_a^{(coh)}(u)`$ merge to $`𝒟_{\mathrm{}}^{(coh)}(u)u\nu /(1+u\nu )`$. This slow evolution is transformed into a fast transition to the equilibrium value if $`a(1u)1`$:
$$𝒟_a^{(coh)}(1\epsilon )1\frac{1+2\nu }{2(1+\nu )}\frac{a\epsilon }{\nu ^2}+𝒪\left(\epsilon ^2\right).$$
(58)
For the ‘cat’ states equations (23) and (36) yield
$`p_0^{(cat)}={\displaystyle \frac{2𝒩^2}{1+u\nu }}\mathrm{exp}\left[{\displaystyle \frac{a(1u)}{1+u\nu }}\right]`$
$`\times \left\{1+\mathrm{cos}\phi \mathrm{exp}\left[{\displaystyle \frac{2au(1+\nu )}{1+u\nu }}\right]\right\}.`$ (59)
At low temperatures ($`\nu 1`$) all the exponential functions ‘die out’ if $`a(1u)/(1+\nu )1`$ and $`au/(1+\nu )1`$, and we observe the ‘plateau’ $`𝒟^{(cat)}(u)\frac{1}{3}`$ (see figure 6).
For the squeezed state, formula (53) yields
$$p_0^{(sqz)}=G(0,t)=f^{1/2}\mathrm{exp}\left[a(1u)F/f\right].$$
(60)
For large values of the squeezing parameter $`\rho `$ and $`\nu 1`$, the ‘purity’ (which does not depend on $`a`$) is given by $`\mu \left(2\mathrm{sinh}\rho \sqrt{u(1u)}\right)^11`$, unless $`u`$ is close enough to $`0`$ or $`1`$. If $`a(1u)1`$, then $`p_01`$, and we observe a universal (independent of $`\alpha `$) behavior of the thermalization coefficient
$$𝒟_{\mathrm{}}^{(sq)}(u)\frac{1\mu }{1+\mu }\frac{4u(1u)\mathrm{sinh}^2\rho }{\left(\sqrt{1+4u(1u)\mathrm{sinh}^2\rho }+1\right)^2}$$
The $`𝒟`$-factor rapidly increases for a small time interval $`t<t_1(\gamma \mathrm{sinh}^2\rho )^1`$ (note that $`\mathrm{sinh}^2\rho `$ is just the ‘reduced’ energy of fluctuations in the initial squeezed state, $`_0\frac{1}{2}`$). For $`t>t_1`$ we observe some ‘plateau’, whose extension corresponds approximately to the interval $`0.07<u<0.93`$ (inside this interval, the values of the function $`u(1u)`$ are not less than a half of the maximum value at $`u=0.5`$). For the values of $`u`$ close to $`1`$, the $`𝒟`$-factor may decrease, following the decreasing linear entropy, but finally the term $`p_0`$ enters the game and prevents the thermalization coefficient from falling down to zero. This final stage of evolution seems very fast in terms of the ‘compact time’ $`u`$, but it is not so dramatic with respect to the scaled time $`\tau =2\gamma t`$: see figure 7.
It is interesting to consider also the thermalization of the initial $`M`$-photon Fock state $`|M`$. In this case the off-diagonal elements of the statistical operator in the Fock basis are equal to zero identically for any time $`t0`$, so the ‘total purity’ $`\mu `$ coincides with the diagonal one $`\lambda `$. The initial diagonal generating function equals $`G_0(z)=z^M`$. Consequently,
$$p_0(t)=\frac{\left[u(t)(1+\nu )\right]^M}{\left[1+\nu u(t)\right]^{M+1}},$$
whereas the integral (7) can be transformed to the form
$$\lambda (t)=_0^{2\pi }\frac{d\phi }{2\pi }\frac{\left[a+b\mathrm{cos}\phi \right]^M}{\left[cd\mathrm{cos}\phi \right]^{M+1}}$$
(61)
$$a=u^2(1+\nu )^2+(1uu\nu )^2,d=2u\nu (1+u\nu )$$
$$b=2u(1uu\nu ),c=\nu ^2u^2+(1+u\nu )^2$$
To calculate the integral (61) we designate it as $`I_M`$ and introduce a new generating function
$$Q(y)=\underset{n=0}{\overset{\mathrm{}}{}}I_ny^n=_0^{2\pi }\frac{d\phi }{2\pi [cya(d+yb)\mathrm{cos}\phi ]}$$
The last integral is given by the expression
$$Q(y)=\left[(cya)^2(d+yb)^2\right]^{1/2}$$
which has the same structure as the known generating function of the Legendre polynomials, so after some algebra we obtain
$$\lambda =\frac{\left|12u(1+\nu )\right|^M}{(1+2u\nu )^{M+1}}P_M\left(\frac{(1u)^2+u^2(1+2\nu )^2}{(1+2u\nu )\left|12u(1+\nu )\right|}\right)$$
The typical dependences $`𝒟(u)`$ for the coherent and Fock states with different initial energies are given in figure 8.
## 6 Three stages of decoherence and thermalization
We see that at low temperatures the decoherence and thermalization of highly excited initial states go through three distinct stages. The first one is rather short, its characteristic time being determined completely by the initial energy of quantum fluctuations, $`t_1(\gamma )^1`$. However, the coefficients $`𝒞`$ and $`𝒟`$ do not assume their equilibrium values ($`0`$ and $`1`$, respectively) at the end of this stage, but they remain approximately constant for a rather long period of time. The total destruction of coherence is observed only after the time $`t_d(2\gamma )^1\mathrm{ln}\left(E\right)t_1`$, where $`E`$ is either the total energy or its ‘classical’ part (depending on the initial state). This time tends to a finite limit when the temperature $`T`$ goes to zero.
The disappearance of the off-diagonal matrix elements of the statistical operator (decoherence) does not mean that the energy level populations reach their equilibrium values. This happens only after the ‘thermalization time’ $`t_T`$, which can be evaluated from the asymptotical behavior of the thermalization parameter $`𝒟`$ at $`t\mathrm{}`$ in the form of the Taylor expansion with respect to the small variable $`\epsilon =1u=\mathrm{exp}(2\gamma t)`$. For example, for a generic squeezed state with nonzero mean values of the quadrature components we have (see also (58))
$$𝒟(1\epsilon )=1\frac{(1+2\nu )_{cl}}{2\nu ^2(1+\nu )}\epsilon +𝒪\left(\epsilon ^2\right)$$
(62)
where the ‘classical energy’ is given by (55). Assuming (for $`\nu 1`$) $`\left(_{cl}/\nu ^2\right)\mathrm{exp}(2\gamma t)1`$ we obtain the estimation $`t_T(2\gamma )^1\mathrm{ln}\left(_{cl}/\nu ^2\right)`$ which shows that the ‘thermalization time’ may exceed essentially not only the decay time $`\gamma ^1`$ but the ‘ultimate decoherence time’ $`t_d`$, too. In particular, for the initial coherent state ($`\rho =0`$) we have $`t_T\gamma ^1\mathrm{ln}|\alpha /\nu |`$.
The situation resembles the classical theory of magnetic relaxation, where we have also two characteristic times: the time of transverse relaxation (dephasing) $`T_2`$ (analog of $`t_d`$) and the time of longitudinal relaxation $`T_1`$ (analog of $`t_T`$). The difference is that in our case both times depend not only on the properties of the environment (through the constants $`\gamma `$ and $`\nu `$), but also on the initial state (through its energy). Besides, there exists the third time – the ‘primary decoherence time’ $`t_1`$.
For the states with zero mean values of the quadratures, the expansion of $`1𝒟`$ begins with the quadratic term $`\epsilon ^2=\mathrm{exp}(4\gamma t)`$. For the initial vacuum squeezed state ($`\alpha =0`$) we have
$$𝒟=1\left[\frac{\mathrm{sinh}(2\rho )}{4\nu (1+\nu )}\epsilon \right]^2+𝒪\left(\epsilon ^3\right).$$
(63)
If $`\rho >1`$ and $`\nu 1`$, one can rewrite (63) as
$$𝒟1\left[E\epsilon /(2\nu )\right]^2$$
(64)
where $`E=\mathrm{sinh}^2\rho +\frac{1}{2}\frac{1}{4}\mathrm{exp}(2\rho )`$ is the total energy of the initial state (it coincides with the energy of fluctuations in the case involved). Consequently, $`t_T^{sqvac}(2\gamma )^1\mathrm{ln}\left(E/\nu \right)`$.
A similar behavior of $`𝒟(u)`$ at $`1u1`$ is observed for the ‘cat’ states. If $`a1`$, then the dependence on the phase $`\phi `$ becomes unimportant, and we obtain
$$𝒟=1\frac{(a\epsilon )^2(1+2\nu )}{4\nu ^2(1+\nu )^2}+𝒪\left(\epsilon ^3\right).$$
In this case the total energy $`Ea1`$, and we arrive again at the equation (64) (if $`\nu 1`$) which yields $`t_T^{cat}(2\gamma )^1\mathrm{ln}\left(E/\nu \right)`$, similarly to the case of the vacuum squeezed state.
For the Fock states we obtain
$$𝒟(1\epsilon )=1\frac{M(M+1)}{4\nu (1+\nu )^2}\epsilon ^2+\mathrm{}.$$
If $`M1`$, then $`EM`$, and we have $`t_T^{Fock}\gamma ^1\mathrm{ln}(E/\sqrt{\nu })`$. We see that the ‘thermalization time’ $`t_T`$ depends logarithmically on the initial energy. Besides, it has the strong temperature dependence, growing as $`T^1`$ at $`T0`$ (remember that $`\nu =\left[\mathrm{exp}(\mathrm{}\omega /k_BT)1\right]^1`$, so $`\nu \mathrm{exp}(\mathrm{}\omega /k_BT)`$ at $`T0`$), in a complete agreement with the third law of thermodynamics. The dependence of the ‘thermalization time’ on the mean equilibrium photon number $`\nu `$ enables ordering different families of quantum states with respect to their robustness against the thermalization (while the ‘primary time’ $`t_1`$ is the same for all states with equal values of the energy of quantum fluctuations). The coherent states are the most robust ones, then follow squeezed and ‘cat’ states, whereas the Fock states, being ‘the most unclassical states’, are thermalized much faster than all the others.
Another interesting feature of the decoherence and thermalization process is the existence of ‘plateaus’ in the dependences $`𝒞(u)`$ and $`𝒟(u)`$ for several different types of states (excluding the coherent states) possessing high initial energy (provided the temperature is low enough). In the cases of the squeezed and Fock states the altitudes of ‘plateaus’ tend to $`1`$ for $`𝒟(u)`$ and to $`0`$ for $`𝒞(u)`$ when the initial energy increases. But for the coherent ‘cat’ states, the metastable values of the coherence and thermalization coefficients remain finite even for $`E\mathrm{}`$: $`𝒞_{plt}\frac{1}{2}`$ and $`𝒟_{plt}\frac{1}{3}`$. Consequently, some degree of coherence (with respect to the fixed energy basis) survives in the ‘cat’ states for a long period of time $`t<t_d(2\gamma )^1\mathrm{ln}a`$. Perhaps, this fact could be important for applications.
## 7 Conclusion
We may conclude that the new quantitative measures of decoherence and thermalization shed new light on the details of the decoherence process accompanying the ‘standard’ thermal relaxation of a quantum harmonic oscillator, showing that this process has three distinct stages in the case of highly excited initial pure states and low temperatures. In particular, we have shown that at low temperatures the ‘ultimate decoherence’ is achieved after rather long interval of time, which is essentially greater than the relaxation time and the ‘primary decoherence time’ which was the central subject of previous studies. Our analysis permits to find the conditions under which the ‘purity’ or the ‘linear entropy’ can serve as reasonable measures of (de)coherence: the initial energy of the quantum state must be much greater then the mean energy of the reservoir oscillators, $`E_01+2\nu `$. However, even under this condition the ‘purity’ can be used only to describe the initial stage of the relaxation process, but it cannot replace the true measures of ‘coherence’ for the whole time interval.
## Acknowledgements
ALS thanks CAPES (Brasil) for support. SSM thanks CNPq (Brasil) for partial financial support.
## Appendix A Propagator of the Fokker-Planck equation
Under certain conditions the process of relaxation of linear $`r`$-dimensional quantum systems (such as a system of coupled oscillators or a charged particle in a homogeneous electromagnetic field and in a confining parabolic potential) can be described in the framework of the Fokker-Planck equation for the Wigner function , -
$$\frac{W}{t}=\frac{}{y_i}\left[\left(\mathrm{𝐀𝐲}+𝐊\right)_iW\right]+D_{ij}\frac{^2W}{y_iy_j}$$
(A.1)
where $`i,j=1,2,\mathrm{},2r`$; the $`2r`$-dimensional vector $`𝐲`$ consists of the linear combinations of the Cartesian coordinates $`q_i`$ and canonically conjugated momenta $`p_i`$ (in the simplest case $`𝐲=(𝐪,𝐩)`$). The drift matrix $`𝐀`$ and vector $`𝐊`$ do not depend on the phase space vector variable $`𝐲`$, although they may have, in general, arbitrary dependences on time. However, the diffusion symmetric matrix $`𝐃D_{ij}`$ cannot be arbitrary, since the physically acceptable solutions to equation (A.1) must satisfy the condition of the positive semidefiniteness of the corresponding statistical operator. This condition is fulfilled provided the matrix $`𝐃_{}=𝐃+\frac{i\mathrm{}}{4}\left(𝐀\mathrm{\Sigma }+\mathrm{\Sigma }\stackrel{~}{𝐀}\right)`$ is positively semidefinite . The elements of the antisymmetric c-number matrix $`\mathrm{\Sigma }=\mathrm{\Sigma }_{jk}`$ are the commutators $`\mathrm{\Sigma }_{jk}=\frac{i}{\mathrm{}}[\widehat{y}_j,\widehat{y}_k]`$. In the case of the single space coordinate the matrix condition $`𝐃_{}0`$ is equivalent to three scalar conditions
$$D_{pp}D_{qq}D_{pq}^2det𝐃\frac{\mathrm{}^2}{16}(\text{Tr}A)^2$$
(A.2)
$$D_{pp}0,D_{qq}0,𝐃=\begin{array}{cc}D_{pp}& D_{pq}\\ D_{pq}& D_{qq}\end{array}.$$
Since (A.1) can be considered as the Schrödinger equation with an effective quadratic (although non-Hermitian) Hamiltonian, the propagator $`G(𝐲,𝐲^{},t)`$,
$$W(𝐲,t)=𝑑𝐲^{^{}}G(𝐲,𝐲^{},t)W(𝐲^{^{}},0),$$
can be found with the aid of the method of quantum time-dependent invariants given in . However, to find its explicit expression it is sufficient to know that this propagator is a Gaussian, so, as any Gaussian Wigner function it can be written as
$`G(𝐲,𝐲^{},t)=(2\pi )^r\left[det_{}(t)\right]^{1/2}`$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{2}}\left[𝐲𝐲_{}(t)\right]_{}^1\left[𝐲𝐲_{}(t)\right]\right\}`$ (A.3)
where $`𝐲_{}(𝐲^{},t)`$ is the mean value of the phase space vector $`𝐲`$ and $`_{}(t)`$ is the variance matrix. The explicit form of $`𝐲_{}`$ and $`_{}`$ can be obtained by solving the equations (which are immediate consequences of the Fokker-Planck equation (A.1))
$$\dot{}_{}=𝐀_{}+_{}\stackrel{~}{𝐀}+2𝐃$$
(A.4)
$$\dot{𝐲}=\mathrm{𝐀𝐲}+𝐊$$
(A.5)
with the initial conditions $`_{}(0)=0`$ and $`𝐲_{}(𝐲^{},0)=𝐲^{}`$, which are equivalent to the property $`G(𝐲,𝐲^{},0)=\delta \left(𝐲𝐲^{}\right)`$ distinguishing the propagator from all other Gaussians.
In the case under study equations (A.4) coincide with the set (42)-(45). Putting their solutions (46)-(50) to the right-hand side of (A.3) we obtain the propagator (16).
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# Andreev Conductance of Chaotic and Integrable Quantum Dots
## I Introduction
Though the actual process of Andreev reflection is simple to describe – an electron in a normal metal incident on a superconductor is reflected back as a hole –, it serves as the fundamental basis for some of the most striking effects known in mesoscopic physics . In particular, Andreev reflection may be viewed as the phenomenon underlying the proximity effect, in which a superconductor is able to strongly influence the properties of a nearby normal metal. Andreev reflection is also the process responsible for the unique conductance properties of normal–superconducting (NS) junctions, such as the conductance enhancement observed in NS point contacts.
Attention has recently turned to so-called “Andreev billiards” as ideal systems in which to study the proximity effect . Such systems consist of an isolated normal metal region, small enough that electrons remain phase coherent within it (i.e. a quantum dot), coupled weakly through a point contact to a superconducting electrode. Calculations of the density of states in such structures have shown, rather remarkably, that the proximity effect is sensitive to the nature of the classical dynamics– a gap of size $`\stackrel{~}{\mathrm{\Delta }}`$ in the spectrum is present in the case of a chaotic billiard, whereas in the integrable case the density of states vanishes linearly at the Fermi energy .
In the present work we also focus on the Andreev billiard, but now add a second point contact leading to a normal electrode and investigate the Andreev conductance of the resulting structure. The Andreev conductance is the differential conductance $`dI/dV`$ at voltages smaller than the superconducting gap in the S electrode, where Andreev reflection at the interface between the normal metal dot and the superconductor is the only current carrying mechanism. Unlike previous studies , we calculate the full voltage ($`V`$) and magnetic field ($`B`$) dependence of the conductance.
Two questions are of particular interest in this study. First, does the sensitivity of the proximity effect to chaotic versus integrable dynamics, as seen in the density of states, also manifest itself in the Andreev conductance? We find that indeed it does. In the case where the contact to the superconducting electrode is much wider than that to the normal electrode, the $`V`$ and $`B`$ dependent conductance of a chaotic dot is consistent with assuming the dot has itself become superconducting. The conductance is doubled with respect to the normal state, and remains bias-independent until $`eV`$ reaches the induced gap $`\stackrel{~}{\mathrm{\Delta }}`$ in the Andreev billiard. This is in sharp contrast to what is found in exact calculations for integrable dots. For both rectangular and circular dot geometries, the conductance drops off linearly with voltage without any plateau near $`V=0`$. We argue that this difference is ubiquitous for integrable versus chaotic systems.
The second motivation for the present study is the question of re-entrance effects– is the behaviour of the conductance simply monotonic in $`V`$ and $`B`$? In the case of a diffusive NS junction, it is well known that this is not the case . At zero voltage and also at voltages large enough to break electron-hole degeneracy, the conductance of the junction is the same as in the normal state. However, a conductance enhancement does occur at intermediate voltages. We examine re-entrance effects in the conductance of chaotic Andreev billiards, both in the case of ballistic contacts, and in the case where both point contacts contain opaque tunnel barriers.
The remainder of this paper is organized as follows. In Section II, we present our model for the chaotic dot and introduce the technique used to calculate the voltage and field-dependent Andreev conductance. In Section III, we discuss results in the case where the contact to the superconductor dominates the contact to the normal electrode, and contrast the results for a chaotic billiard with exact quantum mechanical calculations for two integrable systems. In Section IV, we discuss re-entrance effects. We conclude with a synopsis of our key results in Section V.
## II Model and Technical Details
### A Formulation of the Problem
We consider a chaotic quantum dot coupled via point contacts to a normal metal and a superconductor, each having respectively $`N_N`$ and $`N_S`$ propagating modes at the Fermi energy $`E_F`$ (see Fig. II A). We assume that the ergodic time is much shorter than other relevant timescales of the dot (i.e. the inverse superconducting gap $`\mathrm{}/\mathrm{\Delta }`$ and the dwell time), so that random matrix theory (RMT) may be used to describe its transport and spectral properties. In RMT, the Hamiltonian of dot is represented by a $`M\times M`$ Hermitian matrix $`H`$ which, at zero magnetic field, is real symmetric and a member of the Gaussian orthogonal ensemble:
$$P(H)=\mathrm{exp}(\frac{1}{4}M\lambda ^2\text{tr}H^2).$$
(1)
The matrix size $`M`$ is sent to infinity at the end of the calculation. The energy scale $`\lambda `$ is related to the mean level spacing $`2\delta `$ by $`\lambda =\frac{2M\delta }{\pi }`$. More specifically, $`2\delta `$ is the mean level spacing for particle-like excitations in the absence of a coupling to the superconductor; with the superconductor, the relevant excitations are of mixed particle-hole type, and have a level spacing $`\delta `$ for excitation energies much larger than the gap energy.
In the case of a non-zero magnetic field, the Hermitian matrix $`H`$ is a member of the Pandey-Mehta distribution :
$$P(H)e^{\frac{M(1+\gamma ^2)}{4\lambda ^2}\mathrm{\Sigma }_{i,j=1}^M\left[(\text{Re}H_{ij})^2+\gamma ^2(\text{Im}H_{ij})^2\right]}$$
(2)
As $`\gamma `$ increases from $`0`$ to $`1`$, the distribution evolves from one with complete time-reversal symmetry (Gaussian orthogonal ensemble) to one with no time-reversal symmetry (Gaussian unitary ensemble). The parameter $`\gamma `$ may be related to the magnetic flux $`\mathrm{\Phi }`$ through a two-dimensional dot having area $`A`$ :
$$M\gamma ^2=C\left(\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\right)^2\frac{\mathrm{}v_F}{\sqrt{A}\delta },$$
(3)
where $`\mathrm{\Phi }_0=hc/e`$ is the flux quantum, $`v_F`$ is the Fermi velocity, and $`C`$ is a constant of order unity.
The normal $`(N_N+N_S)\times (N_N+N_S)`$ scattering matrix $`S(\epsilon )`$ of the system at an energy $`\epsilon `$ above $`E_F`$ can be expressed in terms of the matrix $`H`$ :
$`S(\epsilon )`$ $`=`$ $`12\pi iW^{}\left(\epsilon Hi\pi WW^{}\right)^1W`$ (4)
$`=`$ $`\left(\begin{array}{cc}r_{NN}(\epsilon )& t_{NS}(\epsilon )\\ t_{SN}(\epsilon )& r_{SS}(\epsilon )\end{array}\right),`$ (7)
where $`W`$ is an $`M\times (N_N+N_S)`$ matrix representing the coupling between the point contacts and the dot, having elements:
$`W_{mn}`$ $`=`$ $`\delta _{mn}{\displaystyle \frac{1}{\pi }}\left(2M\delta \right)^{1/2}\left({\displaystyle \frac{2T_n2\sqrt{1T_n}}{T_n}}\right)^{1/2}`$ (8)
$`=`$ $`\delta _{mn}\left({\displaystyle \frac{\lambda }{\pi }}Z_n\right)^{1/2}`$ (9)
The $`T_n`$ are the transmission probabilities for each mode, which we take simply as $`T_N`$ for modes coupled to the normal electrode, and $`T_S`$ for those coupled to the superconducting electrode (this in turn defines $`Z_N`$ and $`Z_S`$).
For voltages below the excitation gap $`\mathrm{\Delta }`$ of the superconductor, electrons and holes incident on the dot-superconductor interface may be Andreev reflected. In this case scattering from the dot, as seen from the normal metal contact, can be represented by a $`2N_N\times 2N_N`$ scattering matrix $`𝒮`$:
$$𝒮=\left(\begin{array}{cc}r_{ee}(\epsilon )& r_{eh}(\epsilon )\\ r_{he}(\epsilon )& r_{hh}(\epsilon )\end{array}\right).$$
(10)
Here, $`r_{he}(\epsilon )`$ is the $`N_N\times N_N`$ matrix describing the Andreev reflection of an incoming electron in the N point contact to an outgoing hole in the same lead; $`r_{eh}(\epsilon )`$, $`r_{ee}(\epsilon )`$ and $`r_{hh}(\epsilon )`$ are defined analogously. These matrices may be written in terms of the sub-matrices of the normal scattering matrices $`S(\epsilon )`$ . For $`r_{he}(\epsilon )`$ one has:
$$r_{he}(\epsilon )=t_{NS}^{}(\epsilon )M_{ee}(\epsilon )\alpha (\epsilon )t_{SN}(\epsilon ),$$
(11)
with
$`M_{ee}(\epsilon )`$ $`=`$ $`[(1\alpha (\epsilon )^2r_{SS}^{}(\epsilon )r_{SS}^{}(\epsilon )]^1`$ (12)
$`\alpha (\epsilon )`$ $`=`$ $`\mathrm{exp}(i\mathrm{arccos}(\epsilon /\mathrm{\Delta }))`$ (13)
For voltages below the gap $`\mathrm{\Delta }`$ in the S electrode, the zero temperature conductance of the system is given by the Tabikane-Ebisawa formula :
$$G(eV)=\frac{4e^2}{h}\text{tr }r_{he}^{}(eV)r_{he}^{}(eV).$$
(14)
We wish to calculate the ensemble-averaged Andreev conductance of the dot for arbitrary values of voltage $`V`$ and magnetic field $`B`$; previous studies focused exclusively on the cases where $`V,B`$ were $`0`$ or large enough to completely break the symmetry between electrons and holes. To this end, we first rewrite Eq. (10) in a manner which is formally equivalent to the normal-state expression (4). Letting $`\mathrm{\Omega }_S=\pi W_SW_S^{}`$ and $`\mathrm{\Omega }_N=\pi W_NW_N^{}`$, we define the $`2M\times 2M`$ effective particle-hole Hamiltonian :
$$=\left(\begin{array}{cc}H& 0\\ 0& H^{}\end{array}\right),$$
(15)
and the self-energies from the leads:
$`\mathrm{\Sigma }_N^0(\epsilon )`$ $`=`$ $`i\left(\begin{array}{cc}\mathrm{\Omega }_N& 0\\ 0& \mathrm{\Omega }_N\end{array}\right)`$ (18)
$`\mathrm{\Sigma }_S^0(\epsilon )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{\mathrm{\Delta }^2\epsilon ^2}}}\left(\begin{array}{cc}(\epsilon /\mathrm{\Delta })\mathrm{\Omega }_S& \mathrm{\Omega }_S\\ \mathrm{\Omega }_S& (\epsilon /\mathrm{\Delta })\mathrm{\Omega }_S\end{array}\right)`$ (21)
With these definitions in place, a direct algebraic manipulation shows that the particle-hole scattering matrix $`\mathrm{SS}`$ in (10) can be expressed in terms of the effective retarded Green function $`𝒢_R`$,
$$𝒢_R(\epsilon )=\left[\epsilon (\epsilon )\mathrm{\Sigma }_N^0(\epsilon )\mathrm{\Sigma }_S^0(\epsilon )\right]^1,$$
(22)
by
$`\mathrm{SS}(\epsilon )`$ $`=`$ $`12\pi i\left(\begin{array}{cc}W_N^{}& 0\\ 0& W_N^{}\end{array}\right)𝒢_R(\epsilon )\left(\begin{array}{cc}W_N& 0\\ 0& W_N\end{array}\right)`$ (27)
The major simplification offered by Eq. (27) compared to Eq. (10) is that it allows one to compute the conductance in a manner analogous to that used in the normal state, as we now demonstrate. From this point onward, we focus on the case $`\mathrm{\Delta }\epsilon ,eV`$. In this regime, the properties of the system become independent of the specific details of the superconductor. Calculations retaining a finite $`\mathrm{\Delta }`$ will be presented elsewhere .
### B Details of the Calculation - $`B=0`$
We proceed to average Eq. (14) for the Andreev conductance over the Gaussian orthogonal ensemble defined by Eq. (1), which is the appropriate ensemble for zero magnetic field. We make use of the relation
$$H_{ij}H_{kl}=\left(\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}\right)\frac{\lambda ^2}{M}.$$
(28)
By expressing the scattering matrix in terms of the Green function $`𝒢`$, the trace appearing in the conductance formula may be represented as a standard “conductance bubble” diagram (see Fig. II B). Further, at this stage of the calculation $`\lambda ^2/M`$ is taken to be a small parameter, meaning that the usual diagrammatic technique for impurity averaging may be used. To obtain the leading-order result in $`1/N_N`$,$`1/N_S`$, we need only sum diagrams which have no crossed lines (lines indicate the averaging of matrix elements of $`H`$ using the rule (28); see Fig. II B). Note that $`\lambda `$ will be sent to infinity only at the very end of the calculation; this procedure corresponds to first taking the matrix size to infinity while keeping a finite bandwidth, and then independently taking the bandwidth to infinity. This ordering of limits is necessary to generate a well-defined perturbative expansion.
The first step in this framework is to calculate the averaged matrix Green function $`𝒢`$. As in , we find the following self-consistent Dyson equation:
$$𝒢(\epsilon )=\left[\epsilon \mathrm{\Sigma }_N^0(\epsilon )\mathrm{\Sigma }_S^0(\epsilon )\mathrm{\Sigma }(\epsilon )\right]^1,$$
(30)
where the self-energy from averaging $`\mathrm{\Sigma }`$ is given by
$`\mathrm{\Sigma }`$ $`=`$ $`1_M\left(\begin{array}{cc}\mathrm{\Sigma }_{ee}& \mathrm{\Sigma }_{eh}\\ \mathrm{\Sigma }_{he}& \mathrm{\Sigma }_{hh}\end{array}\right)`$ (33)
$`=`$ $`1_M{\displaystyle \frac{\lambda ^2}{M}}\left(\begin{array}{cc}\text{tr }𝒢_{ee}& \text{tr }𝒢_{eh}\\ \text{tr }𝒢_{he}& \text{tr }𝒢_{hh}\end{array}\right).`$ (36)
In this equation, $`1_M`$ denotes the $`M\times M`$ unit matrix, and $``$ denotes a direct product; it is pictured diagrammatically in Fig. (II B). Having performed the summation of diagrams, we now let $`M\mathrm{}`$ keeping $`\epsilon `$ , $`\delta `$ , $`N_S`$ and $`N_N`$ fixed.
Three of the equations thus obtained relate the components of the self energy to one another:
$$\mathrm{\Sigma }_{ee}=\mathrm{\Sigma }_{hh}\text{ , }\mathrm{\Sigma }_{eh}=\mathrm{\Sigma }_{he},$$
(38)
$$\mathrm{\Sigma }_{eh}^2\mathrm{\Sigma }_{ee}^2=\lambda ^2.$$
(39)
The last equation allows us to parameterize $`\mathrm{\Sigma }`$ in terms of a pairing angle $`\theta (\epsilon )`$:
$`\lambda \mathrm{sin}\theta (\epsilon )`$ $`=`$ $`\mathrm{\Sigma }_{eh}(\epsilon ),`$ (41)
$`\lambda \mathrm{cos}\theta (\epsilon )`$ $`=`$ $`i\mathrm{\Sigma }_{ee}(\epsilon ).`$ (42)
The remaining self-energy equations now take the form:
$`\mathrm{tan}(\theta (\epsilon ))`$ $`=`$ $`{\displaystyle \frac{N_SQ_S}{N_NQ_Ni\frac{\pi \epsilon }{2\delta }}},`$ (44)
$`Q_N`$ $`=`$ $`{\displaystyle \frac{T_N}{2}}\left(2T_N\left[1\mathrm{cos}\theta (\epsilon )\right]\right)^1,`$ (45)
$`Q_S`$ $`=`$ $`{\displaystyle \frac{T_S}{2}}\left(2T_S\left[1+\mathrm{sin}\theta (\epsilon )\right]\right)^1.`$ (46)
To obtain a unique solution of the self-energy equations we have imposed the boundary condition $`\mathrm{\Sigma }_{ee}(\epsilon )i\lambda `$ as $`\epsilon \mathrm{}`$; this represents the physical condition that we recover a normal metal with a constant density of states at large energies. The pairing angle $`\theta `$ can is related the density of states of the dot:
$$\rho (\epsilon )=\frac{\text{Im}}{\pi }\text{tr }𝒢(\epsilon +0^+)=\frac{\text{Re}[\mathrm{cos}(\theta (\epsilon ))]}{\delta };$$
(47)
$`\theta =\pi /2`$ corresponds to a fully superconducting state, while $`\theta =0`$ corresponds to the normal state. It is interesting to note the analogy between $`\theta (\epsilon )`$ and the pairing angle used in the quasi-classical theory of dirty normal-superconducting interfaces and in the circuit theory of Andreev conductance . In the present case, the normalization condition (39) does not need to be externally imposed, but is a direct consequence of the averaging procedure.
While the leading order solution of the self energy $`\mathrm{\Sigma }`$ is sufficient if one is only interested in the density of states (as in ref. ), the conductance calculation requires that $`N_S/M`$ , $`N_N/M`$ and $`\epsilon /M\delta `$ corrections to Eq. (39) be calculated. This correction may be expressed in terms of the leading order self energy solution:
$`{\displaystyle \frac{\mathrm{\Sigma }_{eh}^2\mathrm{\Sigma }_{ee}^2}{\lambda ^2}}1`$ $`=`$ $`{\displaystyle \frac{1}{M}}(N_SQ_S(\mathrm{sin}(\theta )+Z_S)`$ (49)
$`+N_NQ_N(\mathrm{cos}(\theta )+Z_N)+i{\displaystyle \frac{\pi \epsilon }{2\delta }})`$
Having computed the self-energy from averaging $`\mathrm{\Sigma }`$ and thus the average Green function $`𝒢`$, we can proceed to sum diagrams for the conductance. In analogy to the usual impurity technique , two sets of contributions arise (see Fig. II Ba) to leading order: $`G=(4e^2/h)(g_{\mathrm{Dir}}+g_{\mathrm{Diff}})`$. The first is a direct contribution which is completely determined by $`𝒢`$, that is, by the ensemble-averaged probability amplitude $`r_{he}`$:
$`g_{\mathrm{Dir}}`$ $`=`$ $`(2\pi )^2\text{tr }W_N^{}𝒢W_N^{}W_N^{}𝒢^{}W_N^{}`$ (50)
$`=`$ $`\text{tr }r_{he}r_{he}^{}.`$ (51)
This contribution may be interpreted as arising from Andreev reflections which effectively occur at the interface of the normal lead and the cavity. The second contribution is due to the fluctuations of $`r_{he}`$:
$$g_{\mathrm{Diff}}=\left(r_{eh}r_{eh}\right)\left(r_{eh}r_{eh}\right)^2.$$
(52)
It describes the current carried by quasiparticles in the dot which are Andreev reflected at the interface of the dot and the superconductor. Diagrammatically, it is equivalent to a diffusion ladder, where the averaging links the upper and lower branches of the conductance bubble (see Fig. II Bb-d) . We ignore here quantum corrections, which are formally smaller by a factor of $`\mathrm{max}(1/N_N,1/N_S)`$. Note that in computing the diffusion sum, each individual graph is of order $`1/M`$; this means that terms of order $`1/M`$ must be retained when computing the $`4\times 4`$ matrix inverse arising from summing the series. As with the calculation of the average Green function $`𝒢`$, we only let $`M`$ tend to infinity after having performed the partial summation of diagrams.
To present the results of the conductance calculation, we first define the following kernel functions:
$`\mathrm{\Lambda }(\epsilon )`$ $`=`$ $`2N_N\left(1\left|{\displaystyle \frac{Q_N}{Z_N}}\right|^2\left|1+Z_N\mathrm{cos}(\theta )\right|^2\right)+2N_S\left(1\left|{\displaystyle \frac{Q_S}{Z_S}}\right|^2\left|1+Z_S\mathrm{sin}(\theta )\right|^2\right)+`$ (55)
$`2\text{Re}\left[N_NQ_N\left(Z_N+\mathrm{cos}(\theta )\right)+N_SQ_S\left(Z_S+\mathrm{sin}(\theta )\right)\right]+{\displaystyle \frac{\pi \epsilon }{\delta }}\text{Im }[\mathrm{cos}(\theta )],`$
$`\mathrm{\Omega }(\epsilon )`$ $`=`$ $`2N_N\left|Q_N\mathrm{sin}(\theta )\right|^2+2N_S\left|Q_S\mathrm{cos}(\theta )\right|^2,`$ (56)
$`\mathrm{\Pi }_N(\epsilon )`$ $`=`$ $`2N_N|Q_N|^2\left\{\left(1+6Z_N^2+Z_N^4\right)\left(1+|\mathrm{cos}(\theta )|^2\right)+8Z_N\left(1+Z_N^2\right)\text{Re}\left[\mathrm{cos}(\theta )\right]\right\},`$ (57)
$`\mathrm{\Pi }_S(\epsilon )`$ $`=`$ $`2N_S|Q_S|^2\left\{\left(1+6Z_S^2+Z_S^4\right)\left(1+|\mathrm{cos}(\theta )|^2\right)+8Z_S\left(1+Z_S^2\right)\text{Re}\left[\mathrm{cos}(\theta )\right]\right\},`$ (58)
where $`Z_N`$ and $`Z_S`$ were defined below Eq. (8). With these definitions, we find:
$$g_{\mathrm{Dir}}=2N_N\frac{T_N^2}{\left|22T_N\mathrm{sin}^2(\theta /2)\right|^2}\left|\mathrm{sin}(\theta )\right|^2,$$
(60)
$`g_{\mathrm{Diff}}`$ $`=`$ $`8N_N^2|{\displaystyle \frac{Q_N^2}{Z_N}}|^2(\mathrm{\Lambda }^2\mathrm{\Omega }^2)^2(\mathrm{\Pi }_N|\mathrm{sin}(\theta )|^2`$ (62)
$`+(1Z_N^2)^2\mathrm{\Lambda }|\mathrm{sin}(\theta )|^2+\mathrm{\Pi }_S|\mathrm{cos}(\theta )|^2).`$
Equations (55) - (62), along with Eq.(II B) for the self energy $`\mathrm{\Sigma }`$, determine the ensemble-averaged $`B=0`$ Andreev conductance through the dot to leading order in $`1/N_N`$ , $`1/N_S`$ for all voltages such that $`eV\mathrm{\Delta }`$. The extension of these formulas to non-zero $`B`$ are presented in the Appendix.
## III Probing Induced Superconductivity via Conductance
### A Chaotic Dot
In this section, we consider the situation where the coupling between the dot and the superconductor is much stronger than the coupling between the dot and the normal lead, $`N_NT_NN_ST_S`$. In this limit, the normal metal will only weakly perturb the properties of the dot-superconductor system; thus, we may view the conductance through the structure as being a probe of the induced superconductivity in the dot. As mentioned in the Introduction, previous studies have examined the density of states of such Andreev billiards in the absence of a normal lead . It was found for the case of a chaotic dot that an energy gap $`\stackrel{~}{\mathrm{\Delta }}`$ opens up on a scale set by the inverse of the time needed for a particle to escape from the dot to the superconductor:
$`\stackrel{~}{\mathrm{\Delta }}`$ $`=`$ $`c(E_S)`$ (63)
$`E_S`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{}}{\tau _{\mathrm{S}.\mathrm{esc}}}}\right)=\left({\displaystyle \frac{N_ST_S\delta }{2\pi }}\right)`$ (64)
The parameter c is of order unity and a monotonic function of $`T_S`$; it varies from $`0.6`$ in the case of no tunnel barrier, to $`1`$ in the case of an opaque tunnel barrier. It was also found that the shape of the density of states above $`\stackrel{~}{\mathrm{\Delta }}`$ was vastly different in these two limits: in the tunnel regime $`T_S1`$, the density of states was BCS-like, having a square-root singularity at $`\stackrel{~}{\mathrm{\Delta }}`$, whereas in the case of no tunnel barrier, the density of states gradually increased above $`\stackrel{~}{\mathrm{\Delta }}`$. Note that the gap $`\stackrel{~}{\mathrm{\Delta }}`$ does not depend on the gap $`\mathrm{\Delta }`$ in the bulk superconductor, which was taken to infinity in the calculations.
The question that naturally arises here is how the induced superconductivity seen in the density of states manifests itself in the conductance. Some insight may be obtained by considering equations (60) and (62) to lowest order in the small parameter $`\eta =N_NT_N/N_ST_S`$. Without the normal lead, we have a gap in the dot density of states up to an energy $`\stackrel{~}{\mathrm{\Delta }}`$. As the density of states is proportional to $`\text{Re}\left[\mathrm{cos}(\theta (\epsilon ))\right]`$, this implies that $`\text{Re}\left[\theta (\epsilon )\right]=\pi /2+O(\eta )`$ for $`\epsilon <\stackrel{~}{\mathrm{\Delta }}`$ . Using this, the “sub-gap” (i.e. $`\epsilon <\stackrel{~}{\mathrm{\Delta }}`$) direct contribution to the conductance takes the form:
$`g=g_{\mathrm{Dir}}(\epsilon )`$ $`=`$ $`{\displaystyle \frac{2N_NT_N^2}{(2T_N)^24(1T_N)(\text{Im}[\mathrm{sin}(\theta _0(\epsilon ))])^2}}`$ (66)
$`+O(\eta ),`$
while the diffusion contribution $`g_{\mathrm{Diff}}`$ is of order $`\eta ^2`$ and thus negligible. As the energy is increased above the gap, the density of states returns to its normal-state value, and consequently $`\theta (\epsilon )0`$ to leading order. For $`\epsilon \stackrel{~}{\mathrm{\Delta }}`$ the direct term is negligible (being proportional to $`\mathrm{sin}(\theta )`$), while the diffusion term gives the normal-state conductance:
$$g=g_{\mathrm{Diff}}(\epsilon )=N_NT_N+O(\eta ).$$
(67)
The above considerations become extremely suggestive when one considers the case $`T_N=1`$. There is a perfect conductance doubling for voltages below the effective gap $`\stackrel{~}{\mathrm{\Delta }}`$, while at higher voltages the conductance drops to its normal state value:
$$g=\{\begin{array}{cc}2N_N,\hfill & \text{for }\epsilon \stackrel{~}{\mathrm{\Delta }}\hfill \\ N_N,\hfill & \text{for }\epsilon \stackrel{~}{\mathrm{\Delta }}\text{.}\hfill \end{array}$$
(68)
This is precisely what would be expected if the normal point contact were in perfect contact with a bulk superconductor having a gap $`\stackrel{~}{\mathrm{\Delta }}`$– conductance doubling would be expected below the gap due to Andreev reflection, whereas above the gap the normal-state conductance would be recovered as quasiparticles would now carry the current. In this respect, note that below the induced gap $`\stackrel{~}{\mathrm{\Delta }}`$ both the averaged reflection probability $`|r_{eh}|^2`$ and the averaged reflection amplitude $`\left|r_{eh}\right|`$ are non-zero; for $`eV\stackrel{~}{\mathrm{\Delta }}`$, the latter vanishes. Nonetheless, in both cases (i.e. for voltages above and below $`\stackrel{~}{\mathrm{\Delta }}`$) Andreev reflection is the only current carrying process. For energies below $`\stackrel{~}{\mathrm{\Delta }}`$, the dot itself appears as if it were superconducting and Andreev reflection effectively occurs at the normal lead–dot interface. For energies far above $`\stackrel{~}{\mathrm{\Delta }}`$, current through the dot is effectively carried by quasiparticles in the dot (which are either electron or hole-like) which Andreev reflect at the dot-superconductor interface.
The quantum dot continues to act as if it were a superconductor in the case where the contact to the normal lead is no longer perfect (i.e. $`T_N1`$); one need only note that if $`\text{Im}[\mathrm{sin}(\theta _0(\epsilon ))]`$ is replaced by $`\epsilon /\stackrel{~}{\mathrm{\Delta }}`$, Equation (66) becomes identical to the Blonder-Tinkham-Klapwijk (BTK) formula for the sub-gap conductance of a normal point contact - superconductor junction, where there is a tunnel barrier at the interface having transmission $`T_N`$ . The BTK formula is derived assuming there is no spatial separation between the sites of Andreev scattering and normal scattering; the only energy dependence thus comes from the Andreev reflection phase $`\alpha (\epsilon )`$, not from the lack of electron hole degeneracy at finite voltages. The replacement of $`\epsilon /\stackrel{~}{\mathrm{\Delta }}`$ by $`\text{Im}[\mathrm{sin}(\theta _0(\epsilon ))]`$ in Eq. (66) means that the Andreev reflection phase $`\stackrel{~}{\alpha }`$ in the present case is not necessarily the usual $`\alpha (\epsilon )`$ given in Eq. (11), but is rather defined through:
$$\text{Im}[\mathrm{sin}(\theta _0(\epsilon ))]^2=\frac{1+\text{Re}[\stackrel{~}{\alpha }(\epsilon )^2]}{2}.$$
(69)
This notwithstanding, the overall implication is still that for $`eV\stackrel{~}{\mathrm{\Delta }}`$, we can effectively consider all Andreev reflections as occurring at the N-dot interface, as opposed to at the dot-S interface– the dot is indeed acting as though it were itself a superconductor.
We have numerically solved Equations (44)-(46) for the pairing angle $`\theta (\epsilon )`$, and used this to compute the Andreev conductance versus voltage. In Figures III A and III A, we plot the calculated conductance vs. voltage curves for various values of $`T_N`$, and compare to what would be expected from the BTK theory for a simple N-S interface. We find an excellent agreement in the case of $`T_S1`$ (see Fig. III A) even for voltages above the gap $`\stackrel{~}{\mathrm{\Delta }}`$ in the dot. Such an agreement is not unlikely, as for $`T_S1`$ the dot density of states is BCS-like, and consequently the effective Andreev phase $`\stackrel{~}{\alpha }`$ is just equal to the usual Andreev phase $`\alpha `$. In the opposite case of a transparent contact to the superconductor (Fig. III A), clear deviations from the BTK lineshapes are seen. These deviations result completely from the fact $`\stackrel{~}{\alpha }\alpha `$, which is to be expected as the dot density of states in this case is quite different from the BCS form.
The “induced superconductivity” effect in the dot also manifests itself in other manners. A straightforward calculation in the limit $`T_N0`$ shows that the Andreev conductance becomes proportional to the dot density of states (i.e.$`\text{Re}[\mathrm{cos}(\theta )]`$); thus, the Andreev conductance becomes equivalent to a conventional superconductor tunneling-density of states measurement. We have also calculated the magnetic field dependence of the Andreev conductance in the limit of small $`\eta `$ (see Appendix A); here too, our results are consistent with a picture in which the dot itself acts as a superconductor. The conductance enhancement at $`0`$ field remains constant until a critical flux $`\mathrm{\Phi }_C`$ given by:
$$\left(\frac{\mathrm{\Phi }_C}{\mathrm{\Phi }_0}\right)=C\sqrt{\frac{\stackrel{~}{\mathrm{\Delta }}}{E_{\mathrm{erg}}}}=C\sqrt{\frac{\stackrel{~}{\mathrm{\Delta }}\tau _{\mathrm{erg}}}{\mathrm{}}}$$
(70)
where $`\tau _{\mathrm{erg}}`$ is the ergodic time, and C is a geometry-dependent constant of order unity. This is the same critical field required to close the gap $`\stackrel{~}{\mathrm{\Delta }}`$. Note that unlike a conventional BCS superconductor, where the critical field is proportional to the gap, in the present case, Eq. (70) implies that the critical field is proportional to $`\sqrt{\stackrel{~}{\mathrm{\Delta }}}`$.
The conclusions reached here are markedly different from what would be expected from a naive trajectory-based semi-classical analysis . Consider the situation of a ballistic dot with no tunnel barriers, where $`N_SN_N`$ . In the semiclassical picture, electrons entering the dot from the normal lead typically bounce off the walls of the dot several times before hitting the superconductor, where they Andreev reflect. At $`V=0`$, holes are the time-reversed partners of electrons; thus when an Andreev reflection occurs, the hole will retrace the path of the incoming electron and cancel its acquired phase. Andreev trajectories will thus interfere constructively, leading to a large conductance enhancement at $`V=0`$ and a non-vanishing average reflection amplitude $`r_{he}`$. This enhancement should be lost however as the voltage is increased– at finite $`V`$, electrons and holes are no longer degenerate, and the phase acquired by the hole will not precisely cancel that acquired by the electron. The presence of a residual phase $`\delta \varphi =eVL/(\mathrm{}v_F)`$ (where $`L`$ is the length of the trajectory) will lead to destructive interference, and a consequent decrease in the conductance. In this picture, even a small voltage should impair the conductance enhancement seen at zero voltage.
To make this picture quantitative, we may use the fact that our exact random matrix theory indicates that the conductance below $`\stackrel{~}{\mathrm{\Delta }}`$ is proportional to $`|r_{he}|^2`$ (see Eq. (66)). Using the semiclassical Andreev phase $`\delta \varphi =eVL/(\mathrm{}v_F)`$ and the fact that path lengths in chaotic systems have an exponential distribution function $`P(L)=\mathrm{exp}(L/\overline{L})/\overline{L}`$ , we estimate the average of the semiclassical reflection amplitude as:
$$r_{he}_{\mathrm{S}.\mathrm{C}.}𝑑Le^{i\frac{eVL}{\mathrm{}v_F}}P(L)=\frac{1}{1i\frac{eV\overline{L}}{\mathrm{}v_F}},$$
(71)
where $`\overline{L}`$ is the mean path length to the superconductor; for a ballistic dot, we have $`\overline{L}=v_F\tau _{S.esc}`$. Hence:
$$\left|r_{he}\right|^2=\left(1+\left(\frac{eV}{c\stackrel{~}{\mathrm{\Delta }}}\right)^2\right)^1,$$
(72)
where we have used $`\stackrel{~}{\mathrm{\Delta }}=c\mathrm{}/\tau _{S.esc.}`$ (see Eq. (63)) with $`c0.6`$ in the case of a perfect contact to the superconductor. The semiclassical approach thus predicts that $`\left|r_{he}\right|^2`$ (and hence the Andreev conductance) will fall off with voltage as a Lorentzian for $`eV<\stackrel{~}{\mathrm{\Delta }}`$. This is clearly at odds with our fully quantum mechanic calculation at $`T_N=1`$, which shows that $`\left|r_{he}\right|^2=1`$ for all voltages below $`\stackrel{~}{\mathrm{\Delta }}`$. We return to the discrepancy between the semiclassical and quantum mechanical calculations in the next subsection.
### B Integrable Dot
In this subsection, we examine the conductance of an N-dot-S system where the classical dynamics of the dot are integrable. Previous studies have indicated that the proximity effect in the density of states is very different for chaotic and integrable billiards; in the latter, there is no induced gap, but rather the density of states tends to zero linearly at $`E_F`$. We argue here that the proximity effect’s strong sensitivity to chaos versus integrability also manifests itself in the Andreev conductance.
We consider two different integrable systems: a rectangular dot and a circular dot. In each case, the dot is coupled via narrow leads to both a normal metal and a superconductor; there are no tunnel barriers. As in the previous subsection, we consider the case where the width of the normal contact is much smaller than the width of the superconducting contact (see inset of Fig. III B), so that the normal contact serves as a probe of the proximity effect in the quantum dot. The scattering matrix of the system is computed by numerically matching wavefunctions across the structure, and the conductance then follows from equation (14). The number of modes in the normal lead was fixed at $`N_N=\text{int}\left(k_FW/\pi \right)`$ ; we averaged results over small variations in $`k_F`$ which did not change $`N_N`$. Our results are displayed in Fig. III B; similar results are obtained if one changes the position of the two point contacts.
The results for the two integrable systems are similar to one another, and differ significantly from what was found for a chaotic dot– as opposed to flat sub-gap region followed by a rapid drop-off, we have instead a gradual linear reduction of the conductance with voltage. We argue that the difference exhibited here is generic for integrable systems. The reason is the same as that given in to explain the difference in density of states, namely, that the distribution of path lengths is very different for integrable and chaotic dynamics . In the first case, there is a power law distribution of path lengths, meaning that there is an appreciable probability to find long paths. Even for small voltages, these long paths will quickly acquire a random phase, leading to destructive interference and a reduction of the conductance. In contrast, chaotic systems have an exponential distribution of path lengths– long paths are quite rare. In this case, small voltages will not be able to cause any significant phase-randomization, and thus there will be no resulting destructive interference of Andreev trajectories.
Though the semiclassical reasoning in terms of path lengths presented here provides a qualitative account for the difference between the chaotic and integrable Andreev conductance lineshapes, attempts to translate it into a quantitative theory have not been successful. As was demonstrated in the previous section, a simple semiclassical theory for the chaotic case fails to recover the correct lineshape. A similar problem is encountered when one tries to do the analogous calculation for the integrable case– such a calculation predicts the Andreev conductance should fall off quadratically with voltage, albeit at a faster rate than in the chaotic case. This behaviour is clearly at odds with our exact calculation showing a linear dependence on voltage. In many ways, the failure of a semiclassical approach is not surprising. It is known that semiclassical approximations are unreliable in superconducting systems, as the usual diagonal approximation is worse than for normal systems .
Finally, we note that there is a striking connection between the Andreev conductance’s sensitivity to the nature of the dot’s classical dynamics and weak localization. It has been shown both theoretically and experimentally that the magnetic field dependence of the weak localization correction of a quantum dot with two normal leads is very different for chaotic versus integrable dots. In the chaotic case, one finds a smooth Lorentzian field dependence, while in the integrable case a much sharper profile is found, with a cusp at zero magnetic field. The similarity between the Andreev conductance effect and that in weak localization is not coincidental; both effects rely on the interference of time reversed paths.
Despite this strong similarity, it is worth noting that the effect in the Andreev conductance is much more pronounced. Here, the difference between the chaotic and integrable lineshapes is more severe than in the weak localization effect (i.e. the chaotic case is much flatter than a Lorentzian). The magnitude of the effect is also much larger in the Andreev case; the signature of chaos versus integrability is the entire conductance itself, not a quantum correction like weak localization. For this reason, the effect in the Andreev conductance should be observable in a single sample, whereas an ensemble average is required in the weak localization case.
## IV Re-entrance Effects
We shift focus in this section, and examine so-called “re-entrance” phenomena in the Andreev conductance of a chaotic quantum dot. These effects are loosely defined by non-monotonic behaviour of the conductance in either voltage or magnetic field. They are well known in the case of diffusive NS systems , where the word “re-entrant” is used because the Andreev conductance is the same as the normal conductance at zero voltage and magnetic field, and at high voltage or field, but not for intermediate values. The theory developed here allows us to address this behaviour using a scattering approach, whereas previous approaches almost exclusively made use of the quasi-classical Green function technique. In what follows, we discuss the cases of ballistic contacts ($`T_N=T_S=1`$) and tunnel contacts ($`T_N,T_S1`$) separately.
### A Ballistic Contacts
In the absence of tunnel junctions, the equations determining the conductance simplify considerably. We find:
$$g_{\mathrm{Dir}}(\epsilon )=2N_N\mathrm{tan}^2\left(\frac{\theta (\epsilon )}{2}\right),$$
(73)
$`g_{\mathrm{Diff}}(\epsilon )`$ $`=`$ $`{\displaystyle \frac{N_N^2}{2}}|\mathrm{cos}\left({\displaystyle \frac{\theta (\epsilon )}{2}}\right)|^4\times `$ (75)
$`{\displaystyle \frac{N_N\left|\mathrm{tan}\left(\frac{\theta (\epsilon )}{2}\right)\right|^2+N_S\left|\frac{\mathrm{cos}(\theta (\epsilon ))}{1+\mathrm{sin}(\theta (\epsilon ))}\right|^2}{\mathrm{\Lambda }(\epsilon )^2\mathrm{\Omega }(\epsilon )^2}}.`$
As discussed in the previous section, these two contributions to the Andreev conductance can be interpreted as representing two distinct physical consequences of the proximity effect. The direct term $`g_{\mathrm{Dir}}`$ represents processes in which the dot mimics a bulk superconductor, with Andreev reflections effectively occurring locally at the N-dot interface. It decreases monotonically with voltage and magnetic field, going to zero at large voltages or magnetic fields. This reflects the fact that the induced superconductivity effect is sensitive to the averaged $`\mathrm{𝑎𝑚𝑝𝑙𝑖𝑡𝑢𝑑𝑒}`$ for Andreev reflection, which is largest at $`V=B=0`$. On the other hand, the diffusion term $`g_{\mathrm{Diff}}`$ increases monotonically with $`V`$ and $`B`$, tending to the classical result for two conductances in series:
$$g=g_{\mathrm{Diff}}=\frac{(2N_S)(N_N)}{2N_S+N_N}.$$
(76)
It represents a contribution from Andreev quasiparticles in the dot, and is thus sensitive to the dot’s density of states. In the present case of ballistic contacts, the density of states does not have any BCS-type peak and thus the diffusion term rises steadily with applied $`V`$ or $`B`$.
Given that the direct and diffusion contributions react in opposite fashions to an increase in $`V`$ or $`B`$, it is not surprising that a non-monotonic $`V`$ or $`B`$ dependence of the total conductance can be found if the relative strengths of these two terms are varied. The latter can be achieved by tuning the ratio of $`N_N/N_S`$. If $`N_NN_S`$, the direct term will be dominant at $`V,B=0`$, and we expect a monotonic decrease of $`G`$ as a magnetic field or finite bias are applied. In the opposite limit $`N_NN_S`$, the diffusion term dominates at $`V,B=0`$, and $`G`$ is expected to increase with $`V`$ or $`B`$. A non-monotonic $`V`$ or $`B`$ dependence can thus be anticipated in the intermediate regime where $`N_N`$ and $`N_S`$ are comparable.
To quantify the competition between the direct and diffusion contributions, we examine these terms at $`V=B=0`$. This is done by solving the self-energy equations (38)-(46) to determine $`\theta (\epsilon )`$, and then substitute this into Eqs. (73) and (75). Letting $`N_{\mathrm{Sum}}=\sqrt{N_N^2+6N_NN_S+N_S^2}`$, we obtain:
$`g_{\mathrm{Dir}}`$ $`=`$ $`(N_{\mathrm{Sum}}N_SN_N){\displaystyle \frac{N_{\mathrm{Sum}}}{N_N}}2N_S`$ (78)
$`g_{\mathrm{Diff}}`$ $`=`$ $`(N_{\mathrm{Sum}}N_S3N_N){\displaystyle \frac{N_S}{N_N}}+{\displaystyle \frac{4N_NN_S}{N_{\mathrm{Sum}}}}`$ (79)
with the total conductance given by
$$g=g_{\mathrm{Dir}}+g_{\mathrm{Diff}}=(N_S+N_N)\left(1\frac{N_N+N_S}{N_{\mathrm{Sum}}}\right)$$
(80)
Plotted in Figure IV A as a function of $`N_N/N_S`$ are the zero field values of $`g_{\mathrm{Dir}}`$ , $`g_{\mathrm{Diff}}`$ and $`g`$, normalized by the high field conductance given by the classical formula (76). As $`N_N/N_S`$ is increased from zero, the total conductance initially decreases with $`g_{\mathrm{Dir}}`$ as the local Andreev reflection effect is suppressed, while at larger values it starts to increase with $`g_{\mathrm{Diff}}`$ as the density of states at $`E_F`$ returns to its normal state value. At $`N_N\frac{1}{2}N_S`$, we find that the total conductance at zero field is the same as the classical result. Not surprisingly, we find that re-entrance effects in both magnetic field and voltage are maximized here (see Figure IV A).
It is instructive to make a comparison at this point to the non-monotonic re-entrance behaviour of diffusive NS systems (i.e. a diffusive normal metal in good contact with a superconductor). In the diffusive case, the $`0`$ field and high field conductances are the same, being equal to the normal state conductance. The lack of any change due to superconductivity at $`V=B=0`$ is usually explained as an exact cancellation of the conductance doubling effect of Andreev reflection by a suppression of the density of states at the Fermi energy. A conductance enhancement is however found at finite voltage, roughly at $`eVE_T=\mathrm{}D/L^2`$, where $`D`$ is the diffusion constant and $`L`$ is the length of the normal metal. Quasi-classical calculations find maximum conductance enhancements on the order of $`10\%`$ of the normal state conductance .
In the present system, we can tune the relative significance of the Andreev reflection enhancement and density of states suppression terms by varying $`N_N/N_S`$. We find that non-monotonic effects are maximized when we adjust this ratio to mimic the diffusive system, by insisting that the $`0`$ field and high field conductances are the same; this occurs at $`N_N1/2N_S`$. The conductance maximum in voltage is smaller however, being on the order of $`0.01G_{Classical}`$, and occurs roughly at $`eV=E_S=N_S\delta /2\pi `$. This is the inverse of the time needed to escape to the superconductor, and is thus the analog of $`E_T`$ in the diffusive system, which represents the inverse of the time needed for a particle to diffuse across the normal metal and reach the superconductor. Note also that if the normal lead was removed, the size of the induced gap $`\stackrel{~}{\mathrm{\Delta }}`$ in the dot density of states is $`E_S`$.
Finally, we also find pronounced non-monotonic behaviour in magnetic field for this range of $`N_N/N_S`$, with the magnitude of the effect being larger than that in voltage (Fig. IV A). The maximum conductance occurs roughly at a flux $`\mathrm{\Phi }_C`$ given by Eq. (70), the same flux that would be required to close the gap in the density of states in the absence of the normal lead.
### B Tunnel Regime
We turn now to the case where both point contacts contain opaque tunnel barriers ($`T_S,T_N1`$). In this regime, it is sufficient to consider the conductance to lowest non-vanishing order in $`T_N`$ and $`T_S`$. Of particular interest here is the well-known “reflectionless tunneling” effect – the $`V=0,B=0`$ Andreev conductance of a N-I-N-I-S structure (where $`I`$ is an insulating region having transmission $`T`$) is found to be proportional to $`T`$, not to $`T^2`$ as one has for a single barrier. It is as though the Andreev reflected hole is not reflected at all by the tunnel barriers. A similar effect occurs for a N-I-S system where the normal region is sufficiently disordered. The origin of this striking behaviour is now understood to result from the constructive interference of trajectories which reflect many times between the two barriers, leading to a fraction $`T`$ of the conductance channels being open (i.e. having a transmission probability close to unity) .
In the present case, we are able to examine the effects of finite voltage and magnetic field on reflectionless tunneling when a quantum dot separates the barriers. At large voltages or magnetic fields, the pairing angle $`\theta (\epsilon )`$ tends to its normal-state value of 0, and we find that the conductance is given to leading order by the classical series addition formula:
$$g_{\mathrm{class}}=\frac{(\frac{1}{2}N_ST_S^2)(N_NT_N)}{\frac{1}{2}N_ST_S^2+N_NT_N}$$
(81)
For $`N_ST_S^2N_NT_N`$, this simplifies to $`g_{\mathrm{class}}=\frac{1}{2}N_ST_S^2`$. Unlike the case at $`V=0`$,$`B=0`$, there is no order $`T`$ reflectionless tunneling contribution here, as the necessary constructive interference is lost when electron-hole degeneracy or time-reversal symmetry is broken. Below we show that reflectionless tunneling does survive at small values of $`B`$ and $`V`$, and describe how this contribution evolves as $`B`$ and $`V`$ are increased.
We begin by solving Eqs. (44)-(46) for the self energy to lowest order in $`T_N,T_S`$ for $`B=0`$; we find the pairing angle is given simply by:
$$\theta (\epsilon )=\mathrm{arctan}\left(\frac{E_S}{E_Ni\epsilon }\right).$$
(82)
$`E_S`$ and $`E_N`$ are the inverse escape times to the superconductor and normal metal leads respectively, and are defined by:
$$E_N=\frac{N_NT_N\delta }{2\pi },E_S=\frac{N_ST_S\delta }{2\pi }.$$
(83)
Using this result for $`\theta (\epsilon )`$, we next write Eqs. (60)-(62) for the conductance to lowest order in $`T_S,T_N`$. The direct contribution $`g_{\mathrm{Dir}}`$ corresponds to Andreev reflection at the $`N`$-dot interface and is order $`T_N^2`$. The only order $`T`$ contribution is found in the diffusion term $`g_{\mathrm{Diff}}`$, which yields:
$$g_{\mathrm{Diff}}(\epsilon )=N_NT_N\frac{E_NE_S^2}{\stackrel{~}{E}^2(\epsilon )}\sqrt{\frac{2}{\stackrel{~}{E}^2(\epsilon )+E_N^2+E_S^2\epsilon ^2}}+O(T^2),$$
(84)
where we have defined:
$$\stackrel{~}{E}(\epsilon )=\left[\left(E_S^2E_N^2\epsilon ^2\right)^2+4\left(E_SE_N\right)^2\right]^{\frac{1}{4}}.$$
(85)
Eq. (84) gives the complete voltage dependence of the reflectionless tunneling effect. At $`V=0`$, it reduces to:
$$g_{\mathrm{Diff}}(0)=\frac{\left(N_NT_N\right)^2\left(N_ST_S\right)^2}{\left[\left(N_NT_N\right)^2+\left(N_ST_S\right)^2\right]^{\frac{3}{2}}},$$
(86)
which is similar to the formula found in , generalized to the case where the point contacts have different widths. Depending on the relative magnitudes of $`E_S`$ and $`E_N`$, the conductance drops monotonically with voltage, or shows a maximum around $`eVE_S`$ (see Fig. IV B). The location of the maximum is at $`eV=\sqrt{7/6}E_S`$ if $`E_SE_N`$. Similar behaviour is found for the magnetic field dependence of the conductance (see inset of Fig. IV B).
It is easy to understand the origin of this non-monotonic behaviour within our theory. When $`N_S`$ becomes larger than $`N_N`$, the effect of the superconductor on the dot density of states becomes significant. As we have seen in Section II, the induced density of states will have a sharp peak at $`E_S`$ when $`T_S1`$; it is this peak that is manifesting itself in the conductance.
It is interesting to note that Eq. (84) has the same form as what was found for a diffusive N-I-N-I-S system using a quasi-classical Green’s function approach . In that system, the quantum dot is replaced by a diffusive normal wire of length $`d`$, width $`W`$ and mean free path $`l`$, and the normal and superconducting leads are not attached via point contacts, but are also wires of the same width. In the limit $`T_N,T_Sl/d`$ the resistances of the tunnel barriers dominate, and an expression identical to Eq. (84) was obtained, but now the energies $`E_N,E_S`$ are given by:
$$E_N=T_N\frac{\mathrm{}v_F}{4d},E_S=T_S\frac{\mathrm{}v_F}{4d}$$
(87)
where $`v_F`$ is the Fermi velocity in the normal inter-barrier region. This is identical to the definition of $`E_N`$ and $`E_S`$ in Eq. (83) if one takes $`\delta `$ to be the level spacing in the wire. Note that this correspondence is not surprising; previous studies have also found that adding strong tunnel barriers makes a diffusive normal wire with many channels equivalent to a quantum dot.
## V Conclusions
We have studied the voltage ($`V`$) dependence of the Andreev conductance of chaotic and integrable quantum dots; we have also examined the magnetic field ($`B`$) dependence in the chaotic case. In the regime where the contact to the superconductor dominates, we find that the voltage dependence of the Andreev conductance is extremely sensitive to the nature of the dot’s classical dynamics– in the chaotic case, the dot itself mimics a superconductor and the conductance is initially flat in voltage, whereas in the integrable case the conductance falls off linearly with voltage. This effect is large in that the entire conductance forms the signature of chaos vs. integrability; also, it does not require any ensemble averaging to be done. Both these facts make it particularly amenable to experiment. Also important in this regard is that one does not require an extremely clean contact between the dot and the superconductor as long as the contact is large.
We have also studied non-monotonic “re-entrance” phenomenon in the $`V`$ and $`B`$ dependence of the Andreev conductance of a chaotic dot. We find that such behaviour is ubiquitous, and is the result of two competing processes– Andreev reflection at the dot-normal lead interface, which decreases as $`V`$ or $`B`$ increases, versus quasiparticles being injected into the dot, which increases with $`V`$ and $`B`$.
## VI Acknowledgements
We thank C. W. J. Beenakker for a useful discussion. A.C. acknowledges the support of the Olin Foundation and the Cornell Center for Materials Research. Work supported in part by the NSF under grant DMR-9805613.
## A Calculating the Andreev Conductance for $`B0`$
In this appendix, we outline the method used for calculations at non-zero magnetic field $`B`$. Formally, the magnetic field enters the model in a different fashion than a voltage difference. The latter is dealt with by the fact that the Andreev scattering matrix $`\mathrm{SS}`$ defined in Eq. (27) is an explicit function of energy. The field dependence however does not appear directly in the expression for $`\mathrm{SS}`$, but rather only emerges in the averaging procedure– as we use the Pandey-Mehta distribution defined in Eq. (2), the ensemble of random matrices is itself a function of field.
As in the calculation at $`B=0`$, the first step in obtaining the conductance is to calculate the averaged matrix Green function $`𝒢`$ defined in Eq. (30). The diagrams used are the same as the $`B=0`$ case, but now use of the distribution (2) leads to a different self-energy $`\mathrm{\Sigma }`$ in the Dyson equation:
$`\mathrm{\Sigma }(\epsilon ,\gamma )`$ $`=`$ $`1_M\left(\begin{array}{cc}\mathrm{\Sigma }_{ee}& (12\gamma )\mathrm{\Sigma }_{eh}\\ (12\gamma )\mathrm{\Sigma }_{he}& \mathrm{\Sigma }_{hh}\end{array}\right)`$ (A3)
$`=`$ $`1_M{\displaystyle \frac{\lambda ^2}{M}}\left(\begin{array}{cc}\text{tr }𝒢_{ee}& (2\gamma 1)\text{tr }𝒢_{eh}\\ (2\gamma 1)\text{tr }𝒢_{he}& \text{tr }𝒢_{hh}\end{array}\right),`$ (A6)
where $`\gamma `$ is a function of magnetic field (see Eq. 3). The additional factors of $`(12\gamma )`$ here reflect the fact that breaking time-reversal symmetry suppresses off-diagonal superconducting correlations. Solving the Dyson equation (30) to leading order in $`1/M`$, we find that relations (38) and (39) continue to hold, meaning that we may still parameterize the self-energies in terms of a pairing angle $`\theta (\epsilon ,\gamma )`$ through Eq. (41). The equation determining $`\theta (\epsilon ,\gamma )`$ now takes the form:
$$\mathrm{tan}(\theta (\epsilon ,\gamma ))=\frac{N_SQ_S2\gamma \mathrm{sin}(\theta )}{N_NQ_Ni\frac{\pi \epsilon }{2\delta }},$$
(A7)
where $`Q_N`$ and $`Q_S`$ are functions of $`\theta `$ defined in Eqs. (45) and (46). The $`1/M`$ corrections to the self-energy in the presence of a magnetic field read:
$`{\displaystyle \frac{\mathrm{\Sigma }_{eh}^2\mathrm{\Sigma }_{ee}^2}{\lambda ^2}}1={\displaystyle \frac{1}{M}}\left(N_SQ_S\left(\mathrm{sin}(\theta )+Z_S\right)+N_NQ_N\left(\mathrm{cos}(\theta )+Z_N\right)+i{\displaystyle \frac{\pi \epsilon }{2\delta }}+2\gamma \mathrm{sin}^2(\theta )\right).`$ (A8)
The next step in the calculation is to sum diagrams for the conductance. The necessary diagrams are the same as those retained in the $`B=0`$ calculation (i.e. direct and diffusion terms), although their evaluation is different. The direct term is still given by Eq. (60), but with $`\theta (\epsilon ,\gamma )`$ now determined from Eq. (A7).
The diffusion term acquires a form different from Eq. (62), as factors of $`(12\gamma )`$ now appear in graphs where the particle-hole indices of upper and lower branches do not match. These factors lead to $`1/M`$ corrections both to the matrix inverse that arises when summing the diffusion ladder, and to the matrix prefactors of the ladder. As discussed earlier, such corrections are important when calculating the conductance. The result is:
$`g_{\mathrm{Diff}}=8N_N^2\left|{\displaystyle \frac{Q_N^2}{Z_N}}\right|^2{\displaystyle \frac{\left(1Z_N^2\right)^2\stackrel{~}{\mathrm{\Lambda }}|\mathrm{sin}(\theta )|^2+(\mathrm{\Pi }_N+\gamma \mathrm{\Pi }_B)|\mathrm{sin}(\theta )|^2+\mathrm{\Pi }_S|\mathrm{cos}(\theta )|^2}{\stackrel{~}{\mathrm{\Lambda }}^2\stackrel{~}{\mathrm{\Omega }}^2}},`$ (A9)
where
$`\mathrm{\Pi }_B(\theta )`$ $`=`$ $`2\left((1+Z_N)^2\left(1+2\text{Re}[\mathrm{sin}^2(\theta )]\right)+4Z_N(1+Z_N)\text{Re}[\mathrm{cos}(\theta )]+4Z_N^2\left(\left|\mathrm{cos}(\theta )\right|^22\text{Re}[\mathrm{sin}^2(\theta )]\right)\right),`$ (A10)
$`\stackrel{~}{\mathrm{\Lambda }}(\theta ,\gamma )`$ $`=`$ $`\mathrm{\Lambda }(\theta )12\gamma \mathrm{sin}^2(\theta ),`$ (A11)
$`\stackrel{~}{\mathrm{\Omega }}(\theta ,\gamma )`$ $`=`$ $`\mathrm{\Omega }(\theta )+4\gamma \left(1\left|\mathrm{cos}(\theta )\right|^2+4\text{Re}[\mathrm{sin}^2(\theta )]\right)\stackrel{~}{\mathrm{\Lambda }}(\theta ,\gamma )+8\gamma \left(N_N\left|{\displaystyle \frac{Q_N}{Z_N}}\right|^2Y_N(\theta )+N_S\left|{\displaystyle \frac{Q_S}{Z_S}}\right|^2Y_S(\theta )+\gamma Y_2(\theta )\right),`$ (A12)
$`Y_N(\theta )`$ $`=`$ $`\left(\left|{\displaystyle \frac{Z_N}{Q_N}}\right|^21\right)\left(1\left|\mathrm{cos}(\theta )\right|^2+4\text{Re}[\mathrm{sin}^2(\theta )]\right)2Z_N\left(5\text{Re}[\mathrm{sin}^2(\theta )]\text{Re}[\mathrm{cos}(\theta )]+\text{Im}[\mathrm{sin}^2(\theta )]\text{Im}[\mathrm{cos}(\theta )]\right)`$ (A14)
$`+Z_N^2\left(1\left|\mathrm{cos}(\theta )\right|^22\text{Re}[\mathrm{sin}^2(\theta )]\left(1+2\left|\mathrm{cos}(\theta )\right|^2\right)\right),`$
$`Y_S(\theta )`$ $`=`$ $`\left(\left|{\displaystyle \frac{Z_N}{Q_N}}\right|^2\left|1Z_S\mathrm{sin}(\theta )\right|^2\right)\left(1\left|\mathrm{cos}(\theta )\right|^2+4\text{Re}[\mathrm{sin}^2(\theta )]\right)Z_S^2\left|\mathrm{sin}(\theta )\mathrm{cos}(\theta )\right|^2,`$ (A15)
$`Y_2(\theta )`$ $`=`$ $`1|\mathrm{cos}(\theta )|^2(1+4\text{Re}\left(\mathrm{sin}(\theta )^2\right))+\text{Re}[3\mathrm{sin}^2(\theta )+4\mathrm{sin}^4(\theta )]+4|\mathrm{sin}(\theta )|^4.`$ (A16)
Here $`\mathrm{\Lambda }(\theta )`$, $`\mathrm{\Omega }(\theta )`$, $`\mathrm{\Pi }_N(\theta )`$ and $`\mathrm{\Pi }_S(\theta )`$ are given by Eqs. (55) - (58). The above equations, together with Eq. (A7) for $`\theta (\epsilon ,\gamma )`$, determine the Andreev conductance for arbitrary voltage and magnetic field.
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# 1 Reconstructed particle trajectories in CHAOS for the 𝜋⁺_𝑖→𝜋⁺𝜋⁻𝑝 reaction on ¹²𝐶, the geometrical disposition of the wire chambers (WC), the first level trigger hardware (CFT) and the magnet return yokes in the corners. Two CFT segments are removed to permit the particle beam (𝜋_𝑖) to traverse the spectrometer. The CFT segments which are hit by particles are marked with crosses, and the energy deposited in Δ𝐸1 and Δ𝐸2 is indicated by boxes. The proton has a momentum slightly above the CHAOS threshold (185 MeV/c), and its energy is fully deposited in Δ𝐸1.
INTERACTING PION PAIRS
IN NUCLEAR MATTER
N. GRION
Istituto Nazionale di Fisica Nucleare, 34127 Trieste, Italy
The CHAOS Collaboration
Abstract
The pion-production $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reactions were studied on nuclei $`{}_{}{}^{2}H`$, $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$ at a pion energy of $`T_{\pi ^+}`$=283 MeV using the CHAOS spectrometer. The experimental results are reduced to differential cross sections and compared to both theoretical predictions and reaction phase space. Near the $`2m_\pi `$ threshold pion pairs couple to $`(\pi \pi )_{I=J=0}`$ when produced in the $`\pi ^+\pi ^+\pi ^{}`$ reaction channel. The $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ ratio between the $`\pi ^+\pi ^\pm `$ invariant masses on nuclei and on the nucleon is also presented. The marked near-threshold enhancement of $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ is consistent with theoretical predictions addressing the partial restoration of chiral symmetry. On the opposite, nuclear matter weakly influences $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^+}{}^{}`$.
1 Introduction
The influence of the nuclear medium on the $`\pi \pi `$ interaction was studied at TRIUMF by means of the pion induced pion-production reaction $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ (henceforth labelled $`\pi 2\pi `$). The initial study was directed to the deuterium, to understand the $`\pi 2\pi `$ behaviour on both a neutron and a proton then on complex nuclei $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$ to derive possible $`\pi \pi `$ medium modifications by direct comparison of the $`\pi 2\pi `$ data.
Early $`\pi 2\pi `$ measurements found that the near-threshold behaviour of the $`\pi ^+\pi ^{}`$ invariant mass $`M_{\pi ^+\pi ^{}}^A`$ increasingly peak toward the $`2m_\pi `$ threshold as the nucleus mass number increases. Such a behaviour was explained by a theoretical approach of which considered a $`\pi \pi `$ pair a strongly interacting system when coupled to the I=J=0 quantum numbers. The theory studies the $`(\pi \pi )_{I=J=0}`$ properties in nuclear matter by dressing the single-pion propagator to account for the $`P`$wave coupling of pions to $`ph`$ and $`\mathrm{\Delta }h`$ configurations. The model is able to explain the general features of the $`M_{\pi ^+\pi ^{}}^A`$ distributions, which are predicted to increasingly accumulate strength near the $`2m_\pi `$ threshold for $`\rho `$, the nuclear medium density, approaching $`\rho _0`$, the saturation density. Within the same theoretical framework, however, the near-threshold strength of the $`\pi \pi `$ $`T`$matrix is considerably reduced when the $`\pi \pi `$ interaction is constrained to be chiral symmetric, which may indicate that effects other than the in-medium $`(\pi \pi )_{I=J=0}`$ interaction contribute to the observed strength. Conversely, the absence of any in-medium modification of the $`(\pi \pi )_{I=J=0}`$ interaction leads to $`M_{\pi ^+\pi ^{}}^A`$ distributions which lack of strength at threshold. In recent theoretical works on the $`(\pi \pi )_{I=J=0}`$ interaction in nuclear matter, the effects of standard many-body correlations are combined with those deriving from the restoration of chiral symmetry in nuclear matter. As a result, the $`M_{\pi ^+\pi ^{}}^A`$ distributions are shown to regain strength near the $`2m_\pi `$ threshold as $`A`$ (thus the average $`\rho `$) increases. Such a property was earlier outlined in some theoretical works, which demonstrated that the $`M_{\pi ^+\pi ^{}}^A`$ enhancement near threshold is a distinct consequence of the partial restoration of the chiral symmetry at $`\rho <\rho _0`$ . The purpose of this contribution is to present a set of $`\pi 2\pi `$ data and to discuss it to the light of the most recent theoretical findings.
2 The experiment
The experiment was carried out at TRIUMF. Incident pions were produced by the collision of 480 MeV protons on a 10 mm thick graphite target. The M11 pion beam line transported the 282.7 MeV pions to the final focus. Pion pairs were detected in coincidence to ensure an unique identification of the pion production process.
CHAOS is a magnetic spectrometer which was designed for the detection of multi-particle events in the medium-energy range . The magnetic field is generated by a dipole whose pole tip has 95 cm in diameter. The magnet is capable of producing a field intensity up to 1.6 T with an uniformity of about 1%. The top side has a 12 cm bore at the centre for the insertion of targets. Fig. 1 illustrates
a reconstructed $`\pi _i^+\pi ^+\pi ^{}p`$ events on $`{}_{}{}^{12}C`$ and the geometrical disposition of the wire chambers (WC), the CHAOS first level trigger hardware (CFT) and the magnet return yokes in the corners. WC1 and WC2 are multiwire proportional chambers which are capable of handling rates exceeding 5$`\times 10^6`$ particles/s for extended periods of time with high efficiency ($``$95%). They have a cylindrical shape with a diameter of 22.8 cm and 45.8 cm, respectively. WC3 is a cylindrical drift chamber designed to operate in a magnetic field, the chamber diameter is 68.6 cm. The outermost chamber WC4 is a vector drift chamber 122.6 cm in diameter, which operates outside of the magnetic field of CHAOS. The segments of WC3 and WC4 which were crossed by the incident particle beam were turned off. The CFT hardware consists of three coaxial cylindrical layers of fast-counting detectors. The first two layers are NE110 plastic scintillators 0.3 cm and 1.2 cm thick, respectively. $`\mathrm{\Delta }E1`$ is 72 cm far from the magnet centre and spans a zenith angle of $`\pm 7^{}`$; thus, it defines the geometrical solid angle of CHAOS $`\mathrm{\Omega }`$=1.5 sr. The third layer is a SF5 lead-glass 12.5 cm thick, about 5 radiation lengths, which is used as a Cerenkov counter. The three layers were segmented in order to provide an efficient triggering system to multi-particle events. Each segment covered an azimuthal angle (i.e., in-the-reaction plane) of 18.
3 Analysis
The data reduction was based on fully reconstructed $`\pi ^+\pi ^+\pi ^{}`$ and $`\pi ^+\pi ^+\pi ^+`$ events. In order to form differential cross sections these events were binned with their weights, which include the $`\pi ^+\pi ^\pm `$ decay rates inside CHAOS. The capability of the CHAOS spectrometer of measuring the kinetic energies ($`T`$) and laboratory angles ($`\theta `$) of each $`\pi \pi _1\pi _2`$ event permits the determination of the five-fold differential cross section $`^5\sigma /(T\mathrm{\Theta })_{\pi _1}(T\mathrm{\Theta })_{\pi _2}\mathrm{\Phi }_{\pi _1\pi _2}`$, where $`\mathrm{\Phi }_{\pi _1\pi _2}`$ is the zenithal angle between $`\pi _1`$ and $`\pi _2`$, which CHAOS enabled measurement at both $`180^{}\pm 7^{}`$ and $`0^{}\pm 7^{}`$. Four-fold differential cross sections were then obtained by integrating out the $`\mathrm{\Phi }_{\pi \pi }`$ dependence, which was performed by using a linear function joining the two measured data points. Such an assessment of the $`\mathrm{\Phi }_{\pi \pi }`$ dependence reflected in a systematic uncertainty of 8% ($`\sigma `$) for the deuterium and 6-7% ($`\sigma `$) for nuclei. Furthermore, the data were reduced to single differential cross sections $`d\sigma /d𝒪_\pi `$, where $`𝒪_\pi `$ represents $`(T_\pi `$ or $`\mathrm{\Theta }_\pi )_{1,2}`$ or a combination of them. The cross section was then related to measured quantities $`\frac{d\sigma }{d𝒪_\pi }`$= $`f_e\frac{N(𝒪_\pi )}{\mathrm{\Delta }𝒪_\pi }`$, where $`f_e`$ is a parameter which is determined by the experimental conditions.
4 Results of the $`\pi \pi \pi `$ reaction in nuclei
A general property of the $`\pi 2\pi `$ process on nuclei in the low-energy $`M_{\pi \pi }`$ regime was outlined by previous experimental works: it is a quasi-free process both when it occurs on deuterium and on complex nuclei. Furthermore, a common reaction mechanism underlies the process whether it occurs on a nucleon or a nucleus. Thus the study of the $`\pi ^+`$$`{}_{}{}^{2}H\pi ^+\pi ^\pm NN`$ reaction is dynamically equivalent to studying the elementrary $`\pi ^+n\pi ^+\pi ^{}p`$ and $`\pi ^+p\pi ^+\pi ^+n`$ reactions separately.
In the present measurement, the $`\pi ^+A\pi ^+\pi ^\pm A^{}`$ reactions were studied under the same experimental conditions. Thus for a given observable the distributions are directly comparable. In addition, the moderate out-of-the-reaction plane angular acceptance of CHAOS may condition the intrinsic shape of the distributions. They are not corrected for it. The error bars explicitly reported on the spectra are the statistical uncertainties.
Fig. 2 shows the single differential cross sections (diamonds) as a function of the $`\pi \pi `$ invariant mass ($`M_{\pi \pi }`$, MeV) for the two $`\pi ^+\pi ^+\pi ^{}`$ and $`\pi ^+\pi ^+\pi ^+`$ reaction channels. The horizontal error bars are not indicated since they lie within symbols. The distributions span the energy interval available to the $`\pi 2\pi `$ reaction which ranges from $`2m_\pi `$, the low-energy threshold, up the 420 MeV, the maximum allowed by the reaction. The $`\pi A\pi \pi N[A1]`$ phase space simulations (dotted histograms) are also provided and are normalized to the area subtended by the experimental distributions.
Regardless of the nucleus mass number, the invariant mass for the $`\pi ^+\pi ^+\pi ^+`$ distributions closely follow phase space and the energy maximum increases with the increase of $`A`$, that is, with the increase of the nuclear Fermi momentum. The $`\pi ^+\pi ^+\pi ^{}`$ channel discloses a different behaviour; as compared to phase space, the $`{}_{}{}^{2}H`$ invariant mass displays little strength from $`2m_\pi `$ to 310 MeV while, on the same energy interval, the $`{}_{}{}^{12}C`$, $`{}_{}{}^{40}Ca`$ and $`{}_{}{}^{208}Pb`$ $`\pi ^+\pi ^{}`$ invariant mass distributions increasingly peak as $`A`$ increases.
In order to explain the nature of the reaction mechanism contributing to the peak structure, it is useful to examine the cos$`\mathrm{\Theta }_{\pi \pi }^{CM}`$ distribution in the invariant mass interval of the peak, where $`\mathrm{\Theta }_{\pi \pi }^{CM}`$ is the angle between the direction of a final pion and the direction of the incoming pion beam in the $`\pi ^+\pi ^{}`$ rest frame. Fig. 3 shows the cos$`\mathrm{\Theta }_{\pi \pi }^{CM}`$ distributions (diamonds) for $`2m_\pi M_{\pi ^+\pi ^{}}310`$ MeV and $`310<M_{\pi ^+\pi ^{}}420`$ MeV, the latter being shown for comparison. The vertical error bars are the overall uncertainties, which sum in quadrature the systematic and the statistic uncertainties. The differential cross sections are best-fitted (solid line) with a partial wave expansion limited to the three lowest waves, i.e. S, P and D. For all the studied nuclei is $`\chi _\nu ^21`$ which indicates that a proper number of waves was used in the expansion. In the case of heavier nuclei, the $`\pi ^+\pi ^{}`$ system predominantly couples $`S`$wave $``$ 95% and a remaining 5% is spent in a $`D`$wave state. Furthermore, within the sensitivity of the $`\chi _\nu ^2`$method any $`P`$wave coupling of the two pions is excluded.
Two recent theoretical works have modelled the $`\pi 2\pi `$ reaction on nuclei whose results are reported in Fig. 4. The (short and long) dashed lines denote the calculations of while the results of the predictions are shown with the full lines . In the case of $`Ca`$, $`R1`$ ($`R2`$) indicates the predictions for $`\rho `$=0.7$`\rho _n`$ ($`\rho `$=0.5$`\rho _n`$), where $`\rho _n`$ is the nuclear saturation density, while $`V1`$ is the result of the calculations for a mean $`\rho `$=0.24$`\rho _n`$. For both $`{}_{}{}^{2}H`$ and $`{}_{}{}^{40}Ca`$ the curves are normalized to the experimental data. For the $`\pi ^+\pi ^+\pi ^+`$ channel, the predictions agree with the invariant mass distributions. The $`R1`$ and $`R2`$ distributions slightly differ from each other: $`R1`$ presents a broader shape which is due to the larger nuclear Fermi momentum as a consequence of the higher nuclear density used. The lower $`\rho `$ used for $`R1`$ seems to better fit the measured distribution, which is a trend also supported by $`V1`$. Therefore, both models and indicate that the average nuclear density of $`{}_{}{}^{40}Ca`$ for the production reaction to take place cannot exceed 0.5$`\rho _n`$.
In the $`\pi ^+\pi ^+\pi ^{}`$ channel the distribution predicted by for $`{}_{}{}^{2}H`$ is able to describe the data, although it tends to overestimate the low-energy $`M_{\pi \pi }`$ yield. The present version of the model surely improves the previous versions, thus allowing for the construction of nuclear medium effects on a reliable ground. The approach of includes several medium effects: Fermi motion, pion absorption, pion quasi-elastic scattering and $`(\pi \pi )_{I=J=0}`$ medium modifications, which are able to reproduce only a moderate $`M_{\pi \pi }`$ strength in the near-threshold region. In addition, an increase of $`\rho `$ from 0.24$`\rho _n`$ to 0.5$`\rho _n`$ and to 0.7$`\rho _n`$ is unlikely to improve the agreement with the experimental cross section, see also Fig. 8 of Ref.. In the case of, the $`R2`$ prediction (0.5$`\rho _n`$) seems to better reproduce the $`M_{\pi \pi }`$ distribution, although some of the near threshold yield already derives from the the low-energy $`\pi ^+`$ $`{}_{}{}^{2}H\pi ^+\pi ^{}pp`$ production. Therefore, $`V1`$, $`R1`$ and $`R2`$ suggest that the $`M_{\pi ^+\pi ^{}}`$ missing strength near the 2m<sub>π</sub> threshold should be searched in a stronger $`\rho `$-dependence of the $`(\pi \pi )_{I=J=0}`$ interaction, rather than requiring an unlikely high-density nuclear environment for the pion-production process to occur.
The observable $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ is presented in comparison with recent theoretical predictions. $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ is defined as the composite ratio $`\frac{M_{\pi \pi }^A}{\sigma _T^A}/\frac{M_{\pi \pi }^N}{\sigma _T^N}`$, where $`\sigma _T^A`$ ($`\sigma _T^N`$) is the measured total cross section of the $`\pi 2\pi `$ process in nuclei (nucleon). This observable has the property of yielding the net effect of nuclear matter on the $`(\pi \pi )_{I=J=0}`$ interacting system regardless of the $`\pi 2\pi `$ reaction mechanism used to produce the pion pair . Therefore, $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ can be compared with the predictions which explicitly calculate both $`M_{\pi \pi }^{Ca}`$ and $`M_{\pi \pi }^{{}_{}{}^{2}H}`$, but also with the theories described in because they calculate the mass distribution of an interacting $`(\pi \pi )_{I=J=0}`$ system both in vacuum and in nuclear matter. Since the above calculations are reported either in arbitrary units or in units which are complex to scale , theoretical predictions are normalized to the experimental distributions at $`M_{\pi \pi }`$=350$`\pm `$10 MeV, where $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ presents a flat behaviour Fig. 5.
For both reaction channels, the full and dotted curves in Fig. 5 are obtained by simply dividing $`M_{\pi \pi }^{Ca}`$/ $`M_{\pi \pi }^{{}_{}{}^{2}H}`$. Furthermore, for both approaches the underlying medium effect is the $`P`$wave coupling of $`\pi `$’s to $`ph`$ and $`\mathrm{\Delta }h`$ configurations, which accounts for the near-threshold enhancement. When applied to the $`𝒞`$$`{}_{}{}^{Ca}{}_{\pi \pi }{}^{}`$, both and predict the same result; in fact, they well describe the behaviour of $`𝒞`$$`{}_{}{}^{Ca}{}_{++}{}^{}`$ throughout the $`M_{\pi \pi }`$ energy range, while for $`𝒞`$$`{}_{}{}^{Ca}{}_{+}{}^{}`$ only part of the near-threshold strength is reproduced. The models of Refs. examine the medium modifications on the scalar-isoscalar meson, the $`\sigma `$meson. Nuclear matter is assumed to partially restore chiral symmetry and consequently $`m_\sigma `$ to vary with $`\rho `$, the variation being parametrised as $`1p\frac{\rho }{\rho _n}`$, where $`p`$ can range in the interval 0.1$`p`$0.3 for and 0.2$`p`$0.3 for , and for both is $`\rho `$=$`\rho _n`$. Both models are capable of yielding large strength near the $`2m_\pi `$ threshold, therefore the capability of predicting $`𝒞`$$`{}_{}{}^{A}{}_{+}{}^{}`$ is compared for a common minimum value of the parameter $`p`$=0.2. In Fig. 5 the predictions of and are reported with dash-dotted line and dashed line, respectively. The model of provides a larger near-threshold strength, which is due to the combined contributions of the in-medium $`P`$wave coupling of pions to $`ph`$ and $`\mathrm{\Delta }h`$ configurations and to the partial restoration of chiral symmetry in nuclear matter. This model, however, is still too schematic for a conclusive comparison to the present data, therefore full theoretical calculations are called for.
5 Conclusions
$`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ was found to yield the net effect of nuclear matter on the $`\pi \pi `$ system regardless of the $`\pi 2\pi `$ reaction mechanism used to produce the pion pair. These distributions display a marked dependence on the charge state of the final pions: (i) the $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ distributions peak at the $`2m_\pi `$ threshold and the yield increases as $`A`$ increases thus denoting that pion pairs form a strongly interacting system; furthermore, the $`\pi \pi `$ system couples to the $`I=J=0`$ channel, the $`\sigma `$meson channel. (ii) In the $`\pi ^+\pi ^+\pi ^+`$ channel, the $`𝒞`$$`{}_{}{}^{A}{}_{\pi \pi }{}^{}`$ behaviour barely depends on both $`A`$ and $`T`$ thus indicating that nuclear matter weakly affect the $`(\pi \pi )_{I,J=2,0}`$ interaction. The $`𝒞`$$`{}_{}{}^{A}{}_{\pi ^+\pi ^{}}{}^{}`$ observable was compared with theories studying the $`(\pi \pi )_{I=J=0}`$ in-medium modifications associated to the partial restoration of chiral symmetry in nuclear matter, and with model calculations which only include standard many-body correlations, i.e. the $`P`$wave coupling of $`\pi ^{}s`$ to $`ph`$ and $`\mathrm{\Delta }h`$ configurations. It was found that both mechanisms are necessary to interpret the data, although chiral symmetry restoration yields the larger near-threshold contribution. Whether this conclusion is correct, the $`\pi 2\pi `$ CHAOS data would indicate an example of a distinct $`QCD`$ effect in low-energy nuclear physics.
Montecarlo simulations of the $`\pi ^+A\pi ^+\pi ^\pm N[A1]`$ reaction phase space revealed useful to interpret some of the $`\pi 2\pi `$ data. In the case of $`M_{\pi ^+\pi ^+}`$, $`\pi ^+\pi ^+`$ pairs distribute according to phase space. In addition, simulations are able to describe the high-energy part of the distributions which are sensitive to the nuclear Fermi momentum of the interacting $`\pi ^+p[A1]\pi ^+\pi ^+n[A1]^{}`$ proton. For the $`M_{\pi ^+\pi ^{}}`$ distributions the $`\pi \pi `$ dynamics overwhelms the dipion kinematics: unlike phase space, the near-threshold $`\pi ^+\pi ^{}`$ yield is suppressed in the elementary production reaction, $`\pi ^+`$ $`{}_{}{}^{2}H\pi ^+\pi ^{}pp`$ in the present work, while in the same energy range medium modifications strongly enhance $`M_{\pi ^+\pi ^{}}`$. A guideline to the interpretation of the $`M_{\pi ^+\pi ^\pm }^A`$ behaviour should combine the effects of the chiral symmetry restoration in nuclear matter and standard many-body correlations. Such an approach would exclude high-density nuclear matter for both the production reaction to take place and the $`\pi \pi `$ system to undergo medium modification.
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# A periodic elastic medium in which periodicity is relevant
## I Introduction
Periodic elastic media arise in a surprising array of problems, including spin or charge density waves; flux line lattices; and random magnets. A model frequently used to describe a manifold, defined by the single valued height variable $`h(\stackrel{}{r})`$ in a periodic elastic medium is
$$_{pem}=𝑑\stackrel{}{r}\left\{\frac{\gamma }{2}[h(\stackrel{}{r})]^2+\eta [h(\stackrel{}{r})]+V_p[h(\stackrel{}{r})]\right\},$$
(1)
where $`V_p`$ is a periodic potential in the height direction and the random potential $`\eta `$ is not periodic. This is directly analogous to the model used to study lattice effects in thermal roughening and in field theoretic studies of commensurate phases in Ising magnets with competing interactions. In the model (1), the periodic potential is non-random and tends to pin the interface while the quenched random pinning $`\eta [h(\stackrel{}{r})]`$ prefers to make the interface wander. The surface tension term $`\frac{\gamma }{2}[h(\stackrel{}{r})]^2`$ seeks a flat interface and also competes with the quenched random pinning. Field theoretic calculations suggest that at long distances, for (1+1)- and (2+1)- dimensional interfaces, the periodic pinning potential is irrelevant, and hence that the interface scaling behavior is in the random-bond-Ising universality class where width $`w^2=h^2h^2L^{2\zeta }`$ with the roughness exponent $`\zeta =2/3`$ in $`(1+1)`$, and $`\zeta 0.21(4d)`$ in $`(d+1),d2`$ . Note that lattice calculations are strongly affected by a lattice pinning potential and have a flat phase even for large lattice sizes .
Another problem which has been heavily studied is the random substrate problem. This was introduced to model the effect of a random substrate on layers of absorbed atoms, and also serves as a model for the effect of a p-fold random field on the XY-model. There is now a consensus that there is a disorder dominated glassy phase in this model (in 2 substrate dimensions) at low temperatures that is reflected in long distance correlations which behave as $`C(r)\mathrm{ln}^2|r|`$ (in contrast to thermally rough correlations in dimension $`(2+1)`$, which grow as as $`C(r)\mathrm{ln}|r|`$). There has been some uncertainty about whether the leading order correlations found by Cardy and Ostlund (CO) are correct, with functional renormalization group calculations agreeing with CO , and variational calculations disagreeing. The substrate roughness is randomly drawn from the interval $`(0,1)`$ (in lattice units). This corresponds to a different sort of periodic elastic medium than that described in (1) above. Here, the random substrate leads to a periodically repeated disorder seen by an interface lying above the random substrate. This arises due to the fact that the first, third, fifth, etc. atoms deposited at the same position on the random substrate see exactly the same disorder when they land. This corresponds to a random-bond Ising magnet in which the disorder is repeated with period $`\lambda =2`$ along the growth direction. In general, the disorder may range over a scale $`(0,\lambda 1)`$, and this leads to a periodic variation in the disorder on length scale $`\lambda `$. The continuum model for this system is simply,
$$_p=𝑑\stackrel{}{r}\left\{\frac{\gamma }{2}[h(\stackrel{}{r})]^2+\eta [h(\stackrel{}{r})]\right\},$$
(2)
but where $`\eta `$ is periodic in $`h(\stackrel{}{r})`$, so that we require $`\eta [h(\stackrel{}{r})+\lambda ]=\eta [h(\stackrel{}{r})]`$. There has been considerable study of the random substrate ($`\lambda =2`$) problem, with the early controversy now being resolved in favor of a “super-rough” “Bragg-glass” phase in $`(2+1)`$-dimensions in which $`w\mathrm{ln}(L)`$. Exact ground-state calculations have been very useful in resolving this controversy . It is quite easy to see (see Section III) that in $`(1+1)`$-dimensions, the random substrate problem behaves as a random walk (RW), so that $`wL^{1/2}`$. Note however it has been recently argued that although typical dislocations do not destroy the “Bragg-glass” ground state, optimal dislocations have negative energy, and hence are expected to destroy the Bragg glass in $`(2+1)`$-dimensions.
In this paper we study the Hamiltonian (2) as a function of the periodicity $`\lambda `$ of the disorder. We show that at long length scales in $`(1+1)`$\- and $`(2+1)`$-dimensions, the periodicity is relevant and that the random substrate universality class holds. The paper is arranged as follows: Section II sets up the model and describes the way in which we calculate the exact positions of interfaces in random Ising magnets. The scaling theory describing the behavior of these interfaces is developed and tested in Section III. We give a brief conclusion in Section IV.
## II Discrete model and exact algorithm
The model which we use to analyze the effect of periodic disorder on interface properties is a spin-half Ising system with random bonds (RB) on square and cubic lattices. The Hamiltonian is given by,
$$_{RB}=\underset{ij}{}J_{ij}S_iS_j,$$
(3)
where $`J_{ij}>0`$ are coupling constants and the spin variables $`S_i`$ take the values $`\pm 1`$. The spins on two opposite boundaries of the lattices, $`z=1`$ and $`z=L`$, are fixed and have opposite signs so that an interface must exist in the lattice. Our calculations are at zero temperature and we find the ground-state interface properties for interfaces whose average normals lie in the $`\{10\}`$ or $`\{11\}`$ directions of square lattices and in the $`\{100\}`$ or $`\{111\}`$ directions of cubic lattices. The coupling constants are random in a slab of size $`L^{d1}\times \lambda `$ and then periodically repeated $`L/\lambda `$ times along a chosen direction. The distributions used for the $`J_{ij}`$’s vary here from case to case but are always chosen so that the interfaces are rough even for small lattices sizes, and even in the $`\{100\}`$ orientation cubic systems. In Fig. 1, we illustrate the way in which the periodic disorder is implemented for the $`\{10\}`$ and $`\{11\}`$ directions of a square lattice. As is now well known , the ground state interface of the system (3) can be found exactly using the maximum flow algorithm. We have a custom implementation of the push-relabel algorithm for this problem and using it we are able to find the exact ground state interface in Ising systems of size one millions sites in about one minute of CPU time on a high end workstation.
## III Scaling theory and numerical results
Consider the ground-state interface of a square lattice in which the bond disorder has period two in the $`\{11\}`$ orientation (e.g., Fig. 1(a)). It is obvious that the interface is highly degenerate, as the ground state interface may start in any of $`L/2`$ equivalent positions. Consider now starting to create a ground state interface from the left side of Fig. 1(a). To minimize the interface energy one chooses the weakest bond. Having chosen this weakest bond, the interface crosses this weakest bond and chooses the weakest bond in the next column. This process of choosing the weakest bond continues across the sample and, for period two, the random walk so generated gives the exact ground state. The reason this ground state is exact is that at each step, all of the possible random bonds in each column are tested (there are only two!). Thus in this limit, $`wL^{1/2}`$ as for a random walk. In contrast, if the period diverges, the model returns to the random bond Ising universality class (or equivalently the directed polymer (DP) in a random medium) for which $`wL^{2/3}`$. For finite $`\lambda `$, we expect that the interface will seek to optimize its global wandering until the roughness reaches the wavelength of the periodicity . After that it has exhausted all possibilities and then returns to a random walk behavior. We thus have,
$$w(L,\lambda )\{\begin{array}{ccc}L^{2/3},\hfill & \text{ }\hfill & w\lambda ,\hfill \\ L^{1/2},\hfill & & w\lambda .\hfill \end{array}$$
(4)
A natural scaling form based on these limiting behaviors is,
$$w(L,\lambda )L^{2/3}f\left(\frac{L}{\lambda ^{3/2}}\right).$$
(5)
where the scaling function $`f(z)`$ for the roughness has the asymptotic behavior,
$$f(z)\{\begin{array}{ccc}\text{const}\hfill & \text{ }\hfill & z1,\hfill \\ z^{1/6}\hfill & & z1.\hfill \end{array}$$
(6)
Tests of the asymptotic behaviors (4) and the scaling function (5) and the results are presented in Figs. 2 and 3 for the $`\{10\}`$ orientation. It is seen that the predictions of the scaling theory are nicely confirmed. Similar results were found for the $`\{11\}`$ orientation, too.
We turn now to the behavior of random surfaces in $`(2+1)`$ dimensions. There, renormalization group (RG) techniques have been applied to the random-phase sine-Gordon model , to random bond interfaces and to fairly general models of periodic elastic media. Numerically, exact maximum-flow-minimum-cut and minimum-cost-matching algorithms and Monte Carlo methods ) have been used. In the random substrate problem, there is a low temperature “super-rough” phase where $`w^2\mathrm{ln}^2(L)`$, while in the random manifold problem, the surface roughness is found to behave as $`wL^{\zeta _{RB}}`$, where $`\zeta _{RB}=0.42\pm 0.01`$. The qualitative reasoning expressed in the first paragraph of this section also applies to higher dimensions, so that we expect the behavior of $`_p`$ to be in the random substrate universality classes at long length scales $`w>\lambda `$, while the random manifold universality class is dominant at short length scales $`w<\lambda `$. The limiting behaviors in dimension (2+1) are then,
$$w(L,\lambda )\{\begin{array}{ccc}L^{\zeta _{RB}},\hfill & \text{ }\hfill & w\lambda ,\hfill \\ \mathrm{ln}L,\hfill & & w\lambda .\hfill \end{array}$$
(7)
We thus expect,
$$w(L,\lambda )L^{\zeta _{RB}}f\left(\frac{L}{\lambda ^{1/\zeta _{RB}}}\right),$$
(8)
and that the scaling function in $`(2+1)`$-dimensions is
$$f(z)\{\begin{array}{ccc}\text{const}\hfill & \text{ }\hfill & z1,\hfill \\ \mathrm{ln}z/z^{\zeta _{RB}}\hfill & & z1,\hfill \end{array}$$
(9)
with the scaling parameter $`z=L/\lambda ^{1/\zeta _{RB}}`$. The asymptotic behaviors of Eq. (7) are illustrated in Figs. 4(a) and (b) for interfaces in the $`\{100\}`$ orientation. The logarithmic asymptotic behavior is clearly confirmed in Fig. 4(a), but the random manifold behavior is still strongly effected by finite size effects. This is understandable as large system sizes are necessary to see the asymptotic random manifold behavior, even in the $`\lambda \mathrm{}`$ limit . Though finite size effects are clearly evident in the scaling plot of Fig. 4(c), the data collapse at large $`\lambda `$ is quite satisfying. It is clear that the random substrate (Bragg glass) universality class is dominant at large enough length scales. We have tested the behavior in the $`\{111\}`$ orientations and find that $`\{111\}`$ interfaces behave in a similar manner.
## IV Conclusions
We have studied the scaling behavior of an elastic manifold in the presence of a periodically repeated “strong” bond disorder. We find that in $`(1+1)`$\- and in $`(2+1)`$\- dimensions, and at long distances, the periodicity is relevant so these interfaces are in the random substrate universality class. This is to be contrasted with an interface in a system with a periodic potential and with random disorder. In the latter problem the periodic potential is claimed to be irrelevant on long length scales in $`(1+1)`$\- and $`(2+1)`$-dimensions for any disorder, though at weak disorder numerical work on $`\{100\}`$ orientation cubic lattices indicate a strong tendency to order due to lattice effects.
ETS and MJA thank the Academy of Finland for financial support. PMD thanks the DOE under contract DOE-FG02-090-ER45418 for support.
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# References
Equivalence of Descriptions of Gravity in Both Curved and Flat Space-time
Mei Xiaochun
( Institute of Theoretical Physics in Fuzhou, No.303, Building 2, Yinghu Garden,
Xihong Road, Fuzhou, 350025, P.R.Chian, E-mail: fzbgk@pub3.fz.fj.cn )
## Abstract
It is proved in the manuscript that as long as the proper coordinate transformation is introduced, the equations of geodetic lines described in curved space-time can be transformed into the dynamic equations in flat space-time. That is to say, the Einstain theory of gravity and other gravitational theories based on the curved space-time can be identically transformed into flat space-time to describe. As an example, the Schwarzschild solution of the spherical symmetry gravitational field is transformed into flat space-time to study. The results show that there exists no any singularity in the all processes and the whole space-time including at the point r=0. So it seems more rational to discuss the problems of gravitation in flat space-time.
PACS number 0400
Introduction
The general theory of relativity based on the curved space-time has got great success and becomes main current theory now. However, there still exist some foundational problems in it just as the definition of gravitational field’s energy, the quantization of gravitation and the problem of singularity and so on. So it is always an attractive idea to re-establish gravitational theory in flat space-time. Since the 1940’s, many theories based on flat space-time were put forward. Though all of those theories are coincident with the gravitational theory of Einstain under the conditions of weak fields, it can not be proved that they are better than the theory of Einstain by experiments at present. So according to the current viewpoint, the space-time of gravitation field should be non-Euclidean one. The flat space-time is always regarded as the boundary condition where the gravitational field is far away.
In the paper, the author does not try to establish an independent theory in flat space-time. But it can be proved that as long as the proper coordinate transformations are introduced and, the equations of geodetic lines in curved space-time can be transformed into the dynamic equations in flat space-time. It means that the Einstein’s theory of gravitation and other theories based on the curved space-time can also be identically transformed into the flat space-time to describe. Then, the method is used to discuss the Schwarzschild solution of the spherical symmetry gravitational field. The results show that there exists no any singularity again in the all processes and the whole space-time including the point . So it seams more rational to discuss gravitational problems in the flat space-time.
The paper includes three chapters. The first chapter discusses how to transform the Schwarzschild solution of Einstein’s theory in the spherical symmetry gravitation field into the dynamic equations in flat space-time. The second chapter provides a general proof to transform all gravitational theories based on curved space-time into the theories based on flat space-time. The third chapter discusses some problems of foundational concepts, for example, weather space-time is curved or not when exit gravitational fields exist.
1. The Transformation of the Schwarzschild Solution
According to the general theory of relativity, the Schwarzschild metric of the spherical symmetry gravitation field is
$$ds^2=c^2(1\frac{\alpha }{r})dt^2(1\frac{\alpha }{r})^1dr^2r^2(\mathrm{sin}^2\theta d\phi ^2+d\theta ^2)$$
(1)
In the formula, we take $`\alpha =2GM/c^2`$ . Let $`\theta =\pi /2`$ and put Eq.(1) into the geodetic line equation, according to the familiar results in the general theory of relativity, we have the integrals
$$c(1\frac{\alpha }{r})\frac{dt}{ds}=\epsilon $$
(2)
$$r^2\frac{d\phi }{ds}=\frac{L}{c}$$
(3)
Here $`\epsilon `$ and $`L`$ are constants. From above two formulas, the linear element $`ds`$ can be eliminated and we can get
$$r^2(1\frac{\alpha }{r})^1\frac{d\phi }{dt}=\frac{L}{\epsilon }$$
(4)
Defining
$$d\tau =(1\frac{\alpha }{r})dt$$
(5)
and regarding $`\tau `$ as the eigen time , $`t`$ as the coordinate time and taking $`\epsilon =1`$,we can write Eq.(2)as
$$cd\tau =ds$$
(6)
Eq.(4) becomes
$$r^2\frac{d\phi }{d\tau }=L$$
(7)
Here $`L`$ is the angel momentum of unit mass. Eq.(7) is just the conservation formula of angel momentum.
Let’s first discuss the motions of particles with static masses. By using Eq.(6), Eq.(1) can be written as
$$(1\frac{\alpha }{r})(\frac{dt}{d\tau })^2\frac{1}{c^2}(1\frac{\alpha }{r})^1(\frac{dr}{d\tau })^2\frac{r^2}{c^2}(\frac{d\phi }{d\tau })^2=1$$
(8)
By using Eq.(5) and (7), we get
$$(\frac{dr}{d\tau })^2=\frac{c^2\alpha }{r}(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{c^2r^2})$$
(9)
Taking the differential about $`d\tau `$ in the formula above, we get
$$\frac{d^2r}{d\tau ^2}\frac{L^2}{r^3}=\frac{c^2\alpha }{2r^2}(1+\frac{3L^2}{c^2r^2})$$
(10)
It should be noted that each quantity in Eq.(10) is defined in caved space-time. In order to express the equation in flat space-time, the further transformation is needed. Let $`r_0`$, $`\phi _0`$ and $`t_0`$ represent the space-time coordinates in flat space-time, because of the invariability of the 4-diamention interval $`ds^2`$ , we have
$$ds^2=c^2dt_0^2dr_0^2r_0^2d\phi _0^2=c^2(1\frac{\alpha }{r})dt^2(1\frac{\alpha }{r})^1dr^2r^2d\phi ^2$$
(11)
Let $`r_0=r`$,$`\phi _0=\phi `$ we get the transformation relation between $`t_0`$ and $`t`$
$$c^2dt_0^2=c^2(1\frac{\alpha }{r})dt^2+[1(1\frac{\alpha }{r})^1]dr^2$$
(12)
Considering Eq.(5) and (9), we get
$$dr=c(1\frac{\alpha }{r})\sqrt{\frac{\alpha }{r}(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{c^2r^2})}dt$$
(13)
Put it into Eq. 12 , we have
$$dt_0=\sqrt{(1\frac{\alpha }{r})[1\frac{\alpha ^2}{r^2}(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{c^2r^2})]}dt$$
(14)
Comparing it with Eq.(5), we get
$$d\tau =(1\frac{\alpha }{r})^{\frac{1}{2}}[1\frac{\alpha ^2}{r^2}(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{c^2r^2})]^{\frac{1}{2}}dt_0$$
(15)
Combining Eq.(7)with (10) and let $`r_0r`$ the results of the Einstein’s theory can be expressed in the similar formula of the Newtonian gravitation in flat space-time
$$\frac{d^2\stackrel{}{r}}{d\tau ^2}=GM(1+\frac{3L^2}{c^2r^2})\frac{\stackrel{}{r}}{r^3}$$
(16)
Let $`u=1/r`$ and by considering Eq.(7), Eq.(16) can be transformed into
$$\frac{d^2u}{d\phi ^2}+u=\frac{c^2\alpha }{2L^2}(1+\frac{3L^2}{c^2}u^2)$$
(17)
The formula can describe the perihelion precession of the Mercury.
On the other hand, We have used the eigen time $`\tau `$ in Eq.(16). It can be proved that the effect of special relativity has been considered in the formula. The square of a particle’s speed in the center gravitational field is
$$V^2=V_r^2+V_\phi ^2$$
(18)
From Eq.(7),(9) and (15), we have
$$V_r^2=(\frac{dr}{dt_0})^2=(\frac{dr}{d\tau }\frac{d\tau }{dt_0})^2=\frac{c^2\alpha }{r}(1\frac{\alpha }{r})(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{c^2r^2})[1\frac{\alpha ^2}{r^2}(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{c^2r^2})]^1$$
$$V_\phi ^2=(r\frac{d\phi }{dt_0})^2=(r\frac{d\phi }{d\tau }\frac{d\tau }{dt_0})^2=\frac{L^2}{r^2}(1\frac{\alpha }{r})[1\frac{\alpha ^2}{r^2}(1\frac{L^2}{\alpha c^2r}+\frac{L^2}{cr^2})]^1$$
(19)
Therefore, we have
$$V^2=V_r^2+V_\phi ^2=\frac{c^2\alpha }{r}(1\frac{\alpha }{r})(1+\frac{L^2}{c^2r^2})[1\frac{\alpha ^2}{r^2}(1\frac{L^2}{\alpha ^2c^2r}+\frac{L^2}{c^2r^2})]^1$$
$$1\frac{V^2}{c^2}=(1\frac{\alpha }{r})[1\frac{\alpha ^2}{r^2}(1\frac{L^2}{\alpha ^2c^2r}+\frac{L^2}{c^2r^2})]^1$$
(20)
Comparing Eq.(20) with Eq.(15), we obtain
$$d\tau =\sqrt{1\frac{v^2}{c^2}}dt_0$$
(21)
It is completely the same as the formula of time delay in the special theory of relativity. Therefore, Eq.(16) can be written as
$$\frac{d\stackrel{}{p}}{dt}=GMm(1+\frac{3L^2}{c^2r^2})\sqrt{1\frac{V^2}{c^2}}\frac{\stackrel{}{r}}{r^3}$$
(22)
Here $`m`$ is the mass of particle, and $`d\tau `$ is determined by Eq.(15). Because of $`\stackrel{}{L}\stackrel{}{V}\times \stackrel{}{r}`$, it can be seen that there exist the extra two items relative to $`V^2/c^2`$ comparing with The Newtonian theory.
The problem of energy conservation is discussed below. For simplification, we only discuss the situation with $`L=0`$ . In this case, the particle moves along the radium direction. By using Eq.(20), Eq.(22) becomes
$$\frac{d\stackrel{}{p}}{dt}=\frac{GMm}{\sqrt{1+\alpha /r}}\frac{\stackrel{}{r}}{r^3}$$
(23)
By producing $`d\stackrel{}{r}`$ on the two sides of Eq.(23) and taking the integral, we have
$$\frac{d\stackrel{}{p}}{dt}𝑑\stackrel{}{r}=\frac{GMm}{\sqrt{1+\alpha /r}}𝑑\stackrel{}{r}$$
(24)
The left side of Eq.(24) can be written as
$$\frac{d\stackrel{}{p}}{dt}𝑑\stackrel{}{r}=\frac{d\stackrel{}{p}}{dt}\frac{d\stackrel{}{r}}{dt}𝑑t=\stackrel{}{V}𝑑\stackrel{}{p}=\stackrel{}{V}𝑑\frac{m\stackrel{}{V}}{\sqrt{1V^2/c^2}}=\frac{mV^2}{\sqrt{1V^2/c^2}}+mc^2\sqrt{1\frac{V^2}{c^2}}+E_1$$
(25)
Here $`E_1`$ is a constant. So Eq.(24)can be written as
$$\frac{mV^2}{\sqrt{1V^2/c^2}}+mc^2\sqrt{1\frac{V^2}{c^2}}=mc^2\sqrt{1+\frac{\alpha }{r}}+E$$
(26)
Here $`E`$ is a constant. Supposes $`V0`$ when $`r\mathrm{}`$,we have $`E=0`$. Eq.(26) can be written as
$$\frac{mV^2}{\sqrt{1V^2/c^2}}+mc^2(\sqrt{1\frac{V^2}{c^2}}1)+mc^2(1\sqrt{1+\frac{\alpha }{r}})=0$$
(27)
Let $`T`$ represents the kinetic energy of the particle, $`U`$ represents the potential energy of the particle. We define
$$T=\frac{mV^2}{\sqrt{1V^2/c^2}}+mc^2(\sqrt{1\frac{V^2}{c^2}}1)$$
(28)
$$U=mc^2(1\sqrt{1+\frac{\alpha }{r}})$$
(29)
Eq.(27)is just the formula of energy conservation $`T+U=E=0`$ . When $`\alpha /r<<1`$ , $`V<<c`$ from Eq.(27 we get the result of the Newtonian theory.
$$\frac{mV^2}{2}\frac{GMm}{r}=0$$
(30)
The motion equation of photon in the center gravitational field is discussed as follows. For photons,$`ds=0`$ so $`ds`$ can not be used as the parameter of the equation of geodetic line. In this case, we take $`d\tau `$ to replace $`ds`$ and get the same results by solving the equations of gravitational field
$$(1\frac{\alpha }{r})\frac{dt}{d\tau }=\epsilon $$
(31)
$$r^2\frac{d\phi }{d\tau }=L$$
(32)
Let $`\epsilon =1`$ we have
$$d\tau =(1\frac{\alpha }{r})dt$$
(33)
Because
$$ds^2=c^2(1\frac{\alpha }{r})dt^2(1\frac{\alpha }{r})^1dr^2r^2d\phi ^2=0$$
(34)
We get
$$c^2(1\frac{\alpha }{r})(\frac{dt}{d\tau })^2(1\frac{\alpha }{r})^1(\frac{dr}{d\tau })^2(r\frac{d\phi }{d\tau })^2=0$$
(35)
From the formula above we have
$$(\frac{dr}{d\tau })^2=c^2[1(1\frac{\alpha }{r})\frac{L^2}{c^2r^2}]$$
(36)
Taking the differential about $`d\tau `$, we get
$$\frac{d^2r}{d\tau ^2}\frac{L^2}{r^3}=\frac{3\alpha L^2}{2r^4}$$
(37)
By using Eq.(33) and (36), we get
$$(\frac{dr}{dt})^2=c^2(1\frac{\alpha }{r})^2[1(1\frac{\alpha }{r})\frac{L^2}{c^2r^2}]$$
(38)
Suppose the speed of photon in the gravitational field is $`V`$, from Eq. 32 33 and 38 we have
$$V=\sqrt{(\frac{dr}{dt})^2+(r\frac{d\phi }{dt})^2}=c(1\frac{\alpha }{r})\sqrt{1+\frac{\alpha L^2}{c^2r^3}}$$
(39)
It is obvious that $`V`$ constant, so the speed of light would change with $`r`$ in gravitational field. Then let’s discuss how to transform the results into the flat reference system. For light’s motion, if we write the metric in the flat reference system as
$$ds^2=c^2dt_0^2dr_0^2dr_0^2r_0^2d\phi _0^2=0$$
(40)
the result shows that light move in a uniform speed $`c`$ in the gravitational field. However, this is improper for it contradicts Eq.(39). Suppose light’s speed is $`V_0`$ in the flat space-time, the metric should be written as
$$ds^2=u_0^2dt_0^2dr_0^2r_0^2d\phi _0^2=0$$
(41)
From Eq.(34) and (41)we have
$$V_0^2dt_0^2dr_0^2r_0^2d\phi _0^2=c^2(1\frac{\alpha }{r})dt^2(1\frac{\alpha }{r})^1dr^2r^2d\phi ^2$$
(42)
Let $`r_0=r`$,$`\phi _0=\phi `$ we get from Eq.(42)
$$V_0^2dt_0^2=c^2(1\frac{\alpha }{r})dt^2+[1(1\frac{\alpha }{r})^1]dr^2$$
(43)
By using Eq.(38), we have
$$V_0^2dt_0^2=c^2(1\frac{\alpha }{r})^2(1+\frac{\alpha L^2}{c^2r^3})dt^2$$
(44)
There exists one degree of freedom to choose the relations between $`t_0`$ and $`t`$ here. If taking $`t_0=t`$, we get
$$V_0=c(1\frac{\alpha }{r})\sqrt{1+\frac{\alpha L^2}{c^2r^3}}$$
(45)
Comparing Eq.(45) with Eq.(39) we obtain $`V_0=V`$, that is to say, the speeds of lights are the same in the both curved and flat space-times. Therefore, by the relation $`\alpha =dV_0/dt_0=dV/dt`$, the accelerations and forces are the same, so that the forms of motion equations are also the same. So by connecting Eq.(32) with (37), we can directly write the motion equation of photons in the vector’s form in flat space-time as
$$\frac{d^2\stackrel{}{r}}{d\tau ^2}=\frac{3\alpha L^2\stackrel{}{r}}{2r^5}$$
(46)
Let $`u=1/r`$ , the formula can be transformed into
$$\frac{d^2u}{d\phi ^2}+u=\frac{3\alpha }{2}u^2$$
(47)
The formula can be used to describe the deviation of light in the solar gravitational field. As for the time delay experiments of radar waves in the solar gravitational field, by considering $`t_0=t`$ in Eq.(38) under the condition of week field, we can get
$$cdt_0=(1+\frac{\alpha }{r})(1\frac{L^2}{r^2})^{\frac{1}{2}}(1\frac{\alpha L^2}{c^2r^3})dr$$
(48)
Suppose radar waves just swept over the surface of the sun with the radium $`r_0`$ and the speed of radar waves is the speed of light in vacuum, we have $`L=cr_0`$ in the light of angular conservation. The integral of Eq.(48) is
$$ct_0=\sqrt{r^2r_0^2}+\alpha ln\frac{\sqrt{r^2r_0^2}+r}{r_0}\alpha \frac{\sqrt{r^2r_0^2}}{2r}$$
(49)
The same result can be reached from the formula <sup>(2)</sup> .
However, it can be seen from Eq.(38) that the photon would move in the speed over light’ speed in vacuum when $`L=0`$ and $`r<\alpha /2`$. This is unacceptable (The problem will be discussed again later.). So we re-define the transform relations between $`t_0`$ , $`t`$ and $`\tau `$ as
$$dt_0=(1\frac{\alpha ^2}{r^2})dt=(1+\frac{\alpha }{r})d\tau $$
(50)
Put it into Eq.(36), we have
$$\frac{dr}{dt_0}=c(1+\frac{\alpha }{r})^1\sqrt{1(1\frac{\alpha }{r})\frac{L^2}{c^2r^2}}$$
(51)
From the formula, Eq.(49) can also be reached under the condition of weak field. It is obvious that Eq.(50) is the simplest form to obtain the formula (49).from (36). In this way, we have
$$V_0=c\sqrt{1+\frac{\alpha L^2}{c^2r^3}}(1+\frac{\alpha }{r})^1$$
(52)
Because $`\stackrel{}{L}=\stackrel{}{V_0}\times \stackrel{}{r}`$ ,when $`r0`$ , $`V_00`$ , there is no the motion of over light’s speed again according to Eq.(52).
It should be noted that Eq. 46 is not the dynamic equation of the photon in the gravitational field. Because when $`L=0`$ , $`d^2\stackrel{}{r}/d\tau ^2=0`$ , it seems that the photon is not acted by force. However, the photon has acceleration in the gravitational field. So there should be a force acting on the photon. Therefore, Eq. 46 can only be regarded as the equation of kinematics from which the velocity and acceleration of the phone can be obtained. But it can not be regarded as the equation of dynamics of the photon from which the force can be obtained.
Now let’s discuss how to obtain the dynamic equation of the photon based on Eq.(46). Let photon’s speed $`V_0V`$ $`(`$as well as $`r_0r`$ , $`t_0t`$ $`)`$ , when $`r\mathrm{}`$, $`Vc`$ . When $`r<\mathrm{}`$ , $`V<c`$ , showing that the photon is acted by repulsion and does retarded motion. In order to obtain the dynamic equation of photon in the gravitational field, we suppose to have an imaginary particle with speed $`\stackrel{}{V^{}}`$ and
$$\stackrel{}{V}^{}=\stackrel{}{c}\stackrel{}{V}$$
(53)
Here $`V`$ is the speed of photon in the gravitational field, $`c`$ is the speed of photon in vacuum. The directions of $`V`$ and $`c`$ are supposed always the same. When photon’s initial speed $`V=c`$ , the speed of imaginary particle is $`V^{}=0`$ . When the photon falls down in the gravitational field with $`V<c`$ , we have $`V^{}>0`$ . When $`V=0`$ , we have $`V^{}=c`$ . So it is obvious that the imaginary particle does acceleration motion in the gravitational field similar to the general particles with static masses. Therefore, the force acted on the imaginary particle in the gravitational field of spherical symmetry can be supposed to be
$$\stackrel{}{F}^{}=GMm(1+\frac{3L^2}{c^2r^2})\sqrt{1\frac{V^2}{c^2}}\frac{\stackrel{}{r}}{r^3}$$
(54)
Here $`m`$ is the static mass of the imaginary particle. The dynamic equation of imaginary particle is
$$\frac{d\stackrel{}{p}^{}}{dt}=GMm(1+\frac{3L^2}{c^2r^2})\sqrt{1\frac{(\stackrel{}{c}\stackrel{}{V})^2}{c^2}}\frac{\stackrel{}{r}}{r^3}$$
(55)
On the other hand, similar to the particles with static masses, the relativity momentum of a photon in the gravitational field can be defined as
$$\stackrel{}{p}=\frac{m\stackrel{}{V}}{R}$$
(56)
Here $`m`$ is so-called photon’s static mass and $`R=R(r,\theta )`$ is the function remained to be decided. Then, we define the imaginary particle’s momentum $`\stackrel{}{p}^{}`$ as
$$\stackrel{}{p}^{}=\stackrel{}{p}_c\stackrel{}{p}=m\stackrel{}{c}\frac{mV}{R}$$
(57)
Here $`\stackrel{}{p}_c`$ is photon’s momentum in vacuum. Put Eq.(57) into Eq.(55), we get the dynamic equation of the photon in the gravitational field of spherical symmetry
$$\frac{d\stackrel{}{p}}{dt}=GMm(1+\frac{3L^2}{c^2r^2})\sqrt{\frac{2cVV^2}{c^2}}\frac{\stackrel{}{r}}{r^3}=\stackrel{}{F}$$
(58)
However, it can be seen that it is actually unnecessary for us to introduce imaginary particle. In fact, we can directly suppose that the dynamic equation of photon in the central gravitational field is just Eq.(58), so long as from it we can reach the identical results comparing with the Einstein’s theory. The concrete form of function $`R`$ is discussed as follows. We have from Eq.(56)
$$\frac{d\stackrel{}{p}}{dt}=\frac{d}{dt}\frac{m}{R}\frac{d\stackrel{}{r}}{dt}=\frac{m}{R}\frac{d^2\stackrel{}{r}}{dt^2}+m\stackrel{}{V}\frac{d}{dt}\frac{1}{R}$$
(59)
Put it into Eq.(58), we get photon’s acceleration in the gravitational field
$$\frac{d^2\stackrel{}{r}}{dt^2}=R\frac{\stackrel{}{F}}{m}R\stackrel{}{V}\frac{d}{dt}\frac{1}{R}$$
(60)
On the other hand, from Eq.(46) and (50), we can get the result of the Einstein’s theory in the flat space-time
$$\frac{d^2\stackrel{}{r}}{dt^2}=\frac{3\alpha L^2\stackrel{}{r}}{2r^5}(1+\frac{\alpha }{r})^2+\frac{\alpha V_r\stackrel{}{V}}{r^2}(1+\frac{\alpha }{r})^1$$
(61)
Comparing Eq.(60) with (61), we get
$$R(\frac{\stackrel{}{F}}{m}\stackrel{}{V}\frac{d}{dt}\frac{1}{R})=\frac{3\alpha L^2\stackrel{}{r}}{2^2r^5}(1+\frac{\alpha }{r})^2+\frac{\alpha V_r\stackrel{}{V}}{r^2}(1+\frac{\alpha }{r})^1$$
(62)
Decomposing the formula in the both directions of $`\stackrel{}{e}_r`$ and $`\stackrel{}{e}_\phi `$ ,we have
$$R(\frac{F}{m}V_r\frac{d}{dt}\frac{1}{R})=\frac{3\alpha L^2}{2^2r^4}(1+\frac{\alpha }{r})^2+\frac{\alpha V_r^2}{r^2}(1+\frac{\alpha }{r})^1$$
(63)
$$RV_\phi \frac{d}{dt}\frac{1}{R}=\frac{\alpha V_rV_\phi }{r^2}(1+\frac{\alpha }{r})^1$$
(64)
Putting Eq.(64)into(63) and using Eq.(52), we get
$$R=\frac{3L^2}{c^2r^2}(1+\frac{\alpha }{r}+\frac{3L^2}{r^2}+\frac{3\alpha L^2}{C^2r^3})^1(1+\frac{2\alpha }{r}+\frac{\alpha L^2}{c^2r^3}+\frac{2\alpha ^2L^2}{c^2r^4})^{\frac{1}{2}}$$
(65)
When $`L=0`$ we have $`V=V_r`$ ,$`V_\phi =0`$ . In this case, Eq.(64) does not exist. But from Eq.(63)directly, we have
$$\frac{dR}{dr}=\frac{\alpha }{2r^2}\sqrt{1+\frac{2\alpha }{r}}(1+\frac{\alpha }{r})R^2+\frac{\alpha }{r^2}R$$
(66)
This is the quasi-one order Bernoulii equation the solution is
$$R=[e^x\sqrt{1+2x}(1+x)e^x𝑑x+C]^1$$
(67)
Here $`x=\alpha /r`$ When $`x`$ , $`R=1`$ the integral constant C can be determined.
After the function $`R`$ is determined by the method above, Eq.(58) can be regarded as the dynamic equation of photon in the gravitational field of spherical symmetry in flat space-time. It coincides with the results of the Einstein’s theory and can be used to explain the deviation of light as well as the time delay experiments of radar waves in the solar gravitational field. Besides, Eq.(58) can also be used to explain the gravitational red shift of spectral line. Let’s first establish the energy conservation equation of photon in the gravitational field. Similar to Eq.(24), we multiply $`d\stackrel{}{r}`$ on the two sides of Eq.(58) and take the integral
$$\frac{d\stackrel{}{p}}{dt}𝑑\stackrel{}{r}=GMm(1+\frac{3L^2}{c^2r^2})\sqrt{\frac{2cVV^2}{c^2}}\frac{\stackrel{}{r}}{r^3}𝑑\stackrel{}{r}$$
(68)
We only considering the situation with $`L=0`$, by using Eq.(52), the two sides of Eq.(68) can be written as
$$\frac{mc^2}{R}(1+\frac{\alpha }{r})^2\frac{mc^2\alpha }{R^2r}(1+\frac{\alpha }{r})^3𝑑r=\frac{mc^2\alpha }{2r^2}(1+\frac{\alpha }{r})^1\sqrt{1+\frac{2\alpha }{r}}𝑑r$$
(69)
The integral of the right side is
$$\frac{mc^2\alpha }{2r^2}(1+\frac{\alpha }{r})^1\sqrt{1+\frac{2\alpha }{r}}𝑑r=mc^2(\sqrt{1+\frac{2\alpha }{r}}\alpha rctg\sqrt{1+\frac{2\alpha }{g}})+C_2$$
(70)
When $`r\mathrm{}`$ , we have
$$mc^2(\sqrt{1+\frac{2\alpha }{r}}\alpha rctg\sqrt{1+\frac{2\alpha }{g}})_r\mathrm{}=mc^2(1\frac{\pi }{4})$$
(71)
So when $`L=0`$ we can define the potential energy of photon in the central gravitational field as
$$U(r)=mc^2(\sqrt{1+\frac{2\alpha }{r}}\alpha rctg\sqrt{1+\frac{2\alpha }{g}}1+\frac{\pi }{4})$$
(72)
When $`r\mathrm{}`$ we have $`U(r)0`$. Let
$$\frac{mc^2}{R}(1+\frac{\alpha }{r})^2\frac{mc^2\alpha }{Rr^2}(1+\frac{\alpha }{r})^3𝑑r=K(r)+C_1$$
(73)
Eq.(69) can be written as
$$K(r)+mc^2(1\frac{\pi }{4})+mc^2(\sqrt{1+\frac{2\alpha }{r}}\alpha rctg\sqrt{1+\frac{2\alpha }{r}}1+\frac{\pi }{4})=C$$
(74)
Here $`C`$ is a constant. If we choose another proper constant $`b`$ and let $`V_0`$ represent the light’s frequency in vacuum, and define the total energy of photon as $`E=C+b=hv_0=mc^2`$ , the kinetic energy of photon can be defined as
$$T(r)=K(r)+mc^2(1\frac{\pi }{4})+b$$
(75)
In this way, when $`L=0`$ , the formula of energy conservation of photon in the gravitational field of spherical symmetry can be written as
$$T(r)+U(r)=E=hv_0$$
(76)
On the other hand, in the general theory of relativity, the red shift of spectral line is considered caused by time delay of gravitational field. In the gravitational theory of curved space-time, the photon is actually considered to be free one and moves along the curved geodetic line without potential energy or no force acted on it. Therefore, the total energy of a photon is equal to its kinetic energy in the gravitational field. But if the formula $`E=hv`$ is considered tenable in any point in the gravitational field, because is a variable, is also a variable, that is to say, that the energy of photon in the gravitational field is not conservative. This is just the price we have to pay in the gravitational theory based on the curved space-time to explain the red shift of spectral line, thought now people seem to avoid this problem. However, this is unacceptable when we describer gravitational force in flat space-time. In order to keep energy conservation for photon in the description of flat space-time, the rational way is to suppose that the frequency of photon is only relative to its kinetic energy, and has nothing to do with its potential energy with $`T=hv`$ . Here $`v`$ is the frequency of photon in the point in the field. So the formula of energy conservation of photon in the gravitational field is $`hv+U(r)=hv_0=mc^2`$ . In this way, the formula of red shift becomes
$$\frac{v}{v_0}=\frac{vv_0}{v_0}=(\sqrt{1+\frac{2\alpha }{r}}\alpha rctg\sqrt{1+\frac{2\alpha }{r}}1+\frac{\pi }{4})$$
(77)
Under the condition of weak field with $`\alpha /r<<1`$ by the developing of the Taylor series, we have
$$\frac{v}{v_0}=\frac{vv_0}{v_0}=(\frac{\alpha }{r}\frac{\alpha }{2r})=\frac{GM}{r}$$
(78)
Under the same condition, the result of the general theory of relativity is
$$\frac{v}{v_0}=\frac{vv_0}{v_0}=\sqrt{1\frac{\alpha }{r}}1=\frac{GM}{r}$$
(79)
The both are the same. But in the strong field, they are not the same, especially when $`\alpha >r`$, Eq.(79) becomes meaningless, but Eq.(77)is still meaningful. In fact Eq.(77) can be used to explain the big red shifts of the quasi-stellar objects. Taking $`\alpha /r=7.5`$, we have the red shift value $`Z=v/v_0=2.47`$ . Suppose the mass of the quasi-stellar object is $`10^{40}Kg`$, it can be calculated that the radium of he quasi-stellar object is $`r=1.98\times 10^{12}m`$ and its mean density is $`\rho =3.08\times 10^2Kg/m^3`$ . It is only $`0.22`$ times comparing with the mean density of the sun with $`\rho =1.40\times 10^3Kg/m^3`$ . Taking $`\alpha /r=12`$ we have $`Z=3.42`$ , $`r=1.24\times 10^{12}m`$ , $`\rho =1.26\times 10^2Kg/m^3`$ . The even density is similar to the sun’s density. According to the general theory of relativity, if $`\alpha /r>1`$ $`r`$ is in the inside of black hole. But according to Eq.(77), only when $`r0`$ , we have infinite red shift with $`Z\mathrm{}`$ . For the common stars, $`r0`$ so according to Eq.(77), the black holes do not exist actually.
What is mainly shown above is that the results are the same under the condition of weak field when the gravitational theory is described in the both flat and curved space-time. Besides the red shift of spectral line, the following discussion will further show their differences in the strong fields owing to the fact that the coordinate times are different in the both situations. We also only discuss particle’s motions along the radial direction of the spherical symmetry gravitational field with $`L=0`$ . For a particle with static mass, when $`L=0`$ , the speed is from Eq.(9) and (5)
$$V=\frac{dr}{dt}=\pm c\sqrt{\frac{\alpha }{r}}(1\frac{\alpha }{r})$$
(80)
In the formula it has been supposed that $`V=0`$ when $`r\mathrm{}`$ . Definite the direction of velocity is positive along the radius vector’s positive direction. The direction of particle’s velocity is negative when the particle falls down the gravitational field. The acceleration is
$$\alpha =\frac{dV}{dt}=\frac{1}{2}\frac{c^2\alpha }{r^2}(1\frac{\alpha }{r})(1\frac{3\alpha }{r})$$
(81)
Take the integral of Eq. 80 and suppose $`r=r_0`$ when $`t=0`$ , we have
$$ct=\pm \frac{1}{\sqrt{\alpha }}[\frac{2}{3}(r^{\frac{3}{2}}r_0^{\frac{3}{2}})+2\alpha (\sqrt{r}\sqrt{r_0})+\alpha ^{\frac{3}{2}}ln\frac{(\sqrt{r}\sqrt{\alpha })(\sqrt{r_0}+\sqrt{\alpha })}{(\sqrt{r}+\sqrt{\alpha })(\sqrt{r_0}\sqrt{\alpha })}]$$
(82)
(82) Because $`t>0`$ the negative sign is taken when a particle falls down along the direction of radius vector. Conversely, it takes positive sign.
Let’s first discuss the particle’s motion in the area $`r\alpha `$ . When , $`r>a`$, $`V`$ is a negative number. When $`r=\alpha `$ , $`V=0`$ the particle arrives at the Schwarzschild event horizon. It can be known from Eq.(81) that when $`r=3\alpha `$ and $`r=\alpha `$ , acceleration becomes zero. When $`r>3\alpha `$ ,$`\alpha <0`$ the particle is acted by gravitation and is accelerated downward. When $`\alpha <r<3\alpha `$ , $`\alpha >0`$ , the particle is acted by repulsive force and is retarded Gravitation becomes repulsion, it seems unimaginable. .When $`r=\alpha `$ , the particle is not acted by force and at rest on the surface of event horizon. It can be known from Eq.(82) that when $`r=\alpha `$ , $`t=\mathrm{}`$ , that is to say, it takes the particle an infinite time to reach the event horizon.
Then let’s discuss the particle’s motion beneath the event horizon $`r<\alpha `$ . It can be known from Eq.(82) that the time has no definition inside the event horizon because the logarithm of a negative number has no definition(At present, some peasants think that it means space and time to be exchanged each other in the black holes. This is absurd. How one dimension’s time can be transformed into three dimension’s space?). So speaking strictly, it is meaningless to talk about particle’s speed, acceleration and motion in the area $`r<\alpha `$ . If this problem is neglected temporary, we can write Eq.(80)as
$$V=\pm c\sqrt{\frac{\alpha }{r}}(1\frac{\alpha }{r})+V_0$$
(83)
When $`r=\alpha `$ , $`V=0`$ so $`V_0=0`$ . Therefore, the velocity and acceleration can still be expressed by Eq.(80) and (81). It is obvious that we have $`\alpha <0`$ inside the event horizon, meaning that the particles only acted by gravitation. Suppose a particle at the point $`r`$ has a velocity upward, it would be retarded until $`a=V=0`$ when it arrives at the event horizon and stays there at last. If a particle has a velocity downward, is would be accelerated and reach light’s speed at a certain place. After that, the particle would move in the speed over light’s speed and reaches an infinite speed at the point $`r=0`$. The results is the same as that analyzed in the current theory using the method of light cone, except that the particle’s speed would be over light’s speed.
As for a photon’s motion, when $`L=0`$, from Eq.(38) we have
$$V=\frac{dr}{dt}=\pm c(1\frac{\alpha }{r})$$
(84)
$$\alpha =\frac{dV}{dt}=\frac{c^2\alpha }{r^2}(1\frac{\alpha }{r})$$
(85)
Let $`r=r_0`$ when $`t=0`$ , we have integral from Eq.(84)
$$ct=\pm (rr_0+\alpha ln\frac{r\alpha }{r_0\alpha })$$
(86)
When a photon falls down in the area $`r>\alpha `$ , because $`\alpha >0`$ they are retarded by repulsion. When the photon arrives at the event horizon $`r=\alpha `$ , the speeds and accelerations are equal to zero and the infinite time is need. In the area $`r<\alpha `$ , time has no definition for the same reason. Despite of this problem, we have $`\alpha <0`$ in the area $`r<\alpha `$. Suppose a photon has a velocity upwards inside the event horizon, it would be retarded by gravitation and has $`\alpha =V=0`$ when it arrives at the surface of event horizon. If the photon has a velocity downwards, it would be accelerated. When it arrives at the point $`r=\alpha /2`$ , its speed would reach light’s speed in vacuum again. After that time, the photon would move in the speed over light’s speed in vacuum. When the photon arrives at the point $`r=0`$ , its speed becomes infinite <sup>(3)</sup>.
It is obvious that there exist some things irrational, especially particles would move in the speeds over light’s speed in vacuum. In fact in the current theory of black holes, the motions with the speeds over light’s speed can not be avoided during the processes in which material collapses toward the center singularities of gravitational fields. At present, those problems are attributed to the improper selections of coordinates. In order to eliminate those defects, people now transform the problems into other coordinate systems to discuss, for example, the Eddington and the Kruskal coordinate system. In the new coordinate system, though the singularities on the surfaces of event horizons can be eliminated, they can not yet be eliminated at the point $`r=0`$. Hawking even proved that it was impossible to eliminate all singularities in the general theory of relativity <sup>(3)</sup>.
Now let’s discuss the problems in flat space-time. When $`t_0=0`$ let $`r=r_0`$ . According to Eq.(19), when a particle falls free down the gravitational field, we have
$$V=c\sqrt{\frac{\alpha }{r}}(1+\frac{\alpha }{r})^{\frac{1}{2}}$$
(87)
$$\alpha =\frac{1}{2}\frac{c^2\alpha }{r^2}(1+\frac{\alpha }{r})^2$$
(88)
$$ct_0=\frac{2}{3\sqrt{\alpha }}[(r_0+\alpha )^{\frac{3}{2}}(r+\alpha )^{\frac{3}{2}}]$$
(89)
It is obvious that every thing is normal in the area $`r>\alpha `$ . The particle is monotonously accelerated by gravitation. There is no any singularity in the whole space-time and in all physical quantities. When the particle arrives at the point $`r=0`$ , we have
$$V=\underset{x\mathrm{}}{lim}\frac{c\sqrt{x}}{\sqrt{1+x}}c(x=\frac{\alpha }{r})$$
(90)
$$\alpha =\underset{x\mathrm{}}{lim}\frac{c^2x^2}{2\alpha (x+1)^2}\frac{c^2}{2\alpha }$$
(91)
It can be seen that the particle’s speed tends to light’s speed but can not yet reach it. Besides, acceleration and time are finite. When a particle moves along the positive direction of radius vector, as long as it has a speed at the point $`r`$
$$Vc\sqrt{\frac{\alpha }{r}}(1+\frac{\alpha }{r})^{\frac{1}{2}}$$
(92)
the particle can escape the gravitational field and has a speed $`V0`$ when it reach the point $`r\mathrm{}`$ .
As for a photon, when it falls down the gravitational field in flat space-time, according to Eq.(51), its velocity and acceleration are
$$V=\frac{dr}{dt_0}=c(1+\frac{\alpha }{r})^1$$
(93)
$$\alpha =\frac{dV}{dt_0}=\frac{c^2\alpha }{r^2}(1+\frac{\alpha }{r})^3$$
(94)
It can be seen that the photon is acted by repulsion and dose the retarded motion. When $`t=0`$ let $`r=r_0`$ ,we have
$$ct_0=rr_0+\alpha ln\frac{r_0}{r}$$
(95)
There is no any singularity in the area $`r>\alpha `$ . When the photon arrives at the point $`r=0`$, its speed $`V=0`$ acceleration is also finite
$$\alpha =\underset{x\mathrm{}}{lim}\frac{c^2x^2}{\alpha (1+x)^3}0(x=\frac{\alpha }{r})$$
(96)
But it takes photons an infinite time to reach the point $`r=0`$ .
Therefore after the Schwazschild solution is transformed into flat space-time to describe, all original singularities disappear (In flat space-time, singularity appears in the form of over light speed’s motion.). So it is obvious that the singularities in the general theory of relativity are actually caused by describing the theory in the curved space-time. The gravitational field itself has no singularities. Meanwhile, it is known a photon can escape from the gravitational field when it moves along the direction of radius by the action of repulsion as long as it is not at the point $`r=0`$ . In this way, the black holes, at least the singular black holes with infinite densities and infinite small volumes, do not exist.
If observers are in the reference system which falls free down the gravitational field, in this case, the time is $`\tau `$ and the distance between the observers and the center mass is $`r`$ according to Eq. 9 , we have
$$V=\frac{dr}{d\tau }=\sqrt{\frac{\alpha }{r}}=\sqrt{\frac{2GM}{r}}$$
(97)
$$\alpha =\frac{dV}{d\tau }=\frac{dV}{dr}\frac{dr}{d\tau }=\frac{\alpha }{2r^2}=\frac{GM}{r^2}$$
(98)
$$\tau =\sqrt{\frac{r}{\alpha }}𝑑r=\frac{2}{3\sqrt{\alpha }}(r^{\frac{3}{2}}r_0^{\frac{3}{2}})$$
(99)
They are just the results of the Newtonian gravitational theory. There are no singularities when $`r>0`$ . But when $`r<GM/c^2r`$ , the relative speed is over light’s speed. When $`r0`$ , the relative speed becomes infinite. So the reference system falling free down the gravitational field is not yet a good reference system for the discussion of gravitational problems. The reason will be discuss in the third chapter.
In brief, at least for the spherical symmetry gravitational field, it is more rational to study gravitational problems in flat space-time than in curved space-time. By transforming the solution of gravitational field equation to discuss in flat space-time, the problems can become more rational and simple. Some puzzling problems just as the singularity problems of black holes, the flat problem of the universal early stage and so on would be expounded. So it is necessary for us to re-examine the conclusions of the current general relativity by transforming them into flat space-time to study in order to get more rational results.
2. The transformations in the general situations
Now let’s generally prove that it is possible to transform the gravitational theories described in the curved space-time into in the flat space-time. In the following discussion, the indexes of Egyptian letters are used to represent the 4-diamention quantities and the indexes of Latin letters represent the 3-diamention quantities. Let $`x^\alpha `$ represent the 4-diamention coordinates in curved space-time and $`x_0^\alpha `$ represent the 4-diamention coordinates in flat space-time. The 4-diamention linear elements in both space-times are individually
$$ds^2=dx_0^\alpha dx_0^\alpha =(dx_0^0)^2dx_0^idx_0^i$$
(100)
$$ds^2=g_{\alpha \beta }dx^\alpha dx^\beta =g_{00}(dx^0)^22g_{0i}dx^0dx^ig_{ij}dx^idx^j$$
(101)
In the formulas, $`x^0`$ and $`x_0^0`$ are the time components. The equation of geodetic line of a particle moving in a gravitational field is
$$\frac{d^2x^\alpha }{ds^2}+\mathrm{\Gamma }_{\beta \sigma }^\alpha \frac{dx^\beta }{ds}\frac{dx^\sigma }{ds}=0$$
(102)
For a certain gravitational field, suppose the metric tensor $`g_{\alpha \beta }`$ has been obtained by solving the Einstein’s equation of gravitation field or other equations based on curved space-time, we can get from the integrals of Eq.(102)
$$x^i=x^i(s)$$
(103)
as well as
$$x^0=x^0(s)$$
(104)
From Eq.(104), we can obtain
$$s=s(x^0)$$
(105)
Put it into Eq.(103), we get
$$x^i(s(x^0))=x^i(x^0)=x^i(t)$$
(106)
The particle’s velocity and acceleration are
$$\frac{dx^i(t)}{dt}=V^i(t)$$
(107)
$$\frac{d^2x^i(t)}{dt^2}=\frac{dV^i(t)}{dt}=\alpha ^i(t)$$
(108)
In addition, two independent equations can be obtained by eliminating time t in the three equations of (106)
$$\varphi _1(x^1,x^2)=0\varphi _2(x^1,x^3)=0$$
(109)
If $`x^i`$ are the coordinates in the Euclidean space, $`\varphi _1`$ and $`\varphi _2`$ represent the equations of two columnar surfaces with the axial lines meeting at right angles. The intersecting line of the two columnar surfaces determined by Eq.(109) presents the orbit of a particle moving in the 3-dimention Euclidean space. If $`x^i`$ are the non-Euclidean space coordinates, $`\varphi _1`$ and $`\varphi _2`$ represent the two 2-dimention curved surface equations in the non-Euclidean space. Their intersecting line also represents the orbit of a particle moving in the 3-dimention non-Euclidean space
On the other hand, as we know, any point on the 2-diamention non-Euclidean curved surface can find a one-to-one point in the 3-diamention Euclidean space, that is to say, the 2-dimention non-Euclidean curved surface can be inlaid into the 3-dimention Euclidean space. Therefore, any point at the interesting line of the two 2-dimention non-Euclidean curved surfaces can also find a one-to-one point in the 3-diamention flat Euclidean space. So we have the transformation relation between the points of geodetic line in the non-Euclidean space and the points in the Euclidean space
$$x^i=x^i(x_0^j)orx_0^i=x_0^i(x^j)$$
(110)
In order to get transformation relation between time $`t`$ and $`t_0`$ , by the condition $`ds^2=constant`$ , we have
$$ds^2=(dx_0^0)^2dx_0^idx_0^i=g_{00}(dx^0)^22g_{oi}dx^0dx^ig_{ij}dx^idx^j$$
(111)
or
$$c^2dt_0^2=c^2g_{00}dt^22cg_{0i}dtdx^ig_{ij}dx^idx^j+dx_0^idx_0^i$$
(112)
By considering Eq.(106) and (110), each item on the right side of Eq.(112) can be expressed as the function of time t, so we get transformation relation of time
$$t_0=\sqrt{g_{00}2g_{0i}\frac{dx^i}{cdt}g_{ij}\frac{dx^i}{cdt}\frac{dx^j}{cdt}+\frac{x_0^i}{x^l}\frac{x_0^i}{x^k}\frac{dx^l}{cdt}\frac{dx^k}{cdt}}𝑑t$$
(113)
i.e.,
$$t_0=t_0(t)ort=t(t_0)$$
(114)
After the transformation relations of space-time coordinates are obtained, the equation (107) in the non-Euclidean space can be transformed into that in the Euclidean space. From Eq.(107),(110) and (114), we have
$$\frac{x^i}{x_0^j}\frac{dx_0^j}{dt_0}=\frac{dx^i}{dt}=V^i(t)=V^i(t(t_0))=V^i(t_0)$$
(115)
This is an equation set of three variables and one order about $`dx_0^j/dt_0`$ . We can obtain the velocity and acceleration in flat space by solving the equation set
$$\frac{dx_0^i}{dt_0}=V_0^i(t_0)$$
(116)
$$\frac{d^2x_0^i}{dt_0^2}=\frac{dV_0^i(t_0)}{dt_0}=a_0^i(t_0)$$
(117)
For the particle with static mass $`m`$ , the momentum in flat space is
$$\stackrel{}{P}_0=\frac{m\stackrel{}{V}_0}{\sqrt{1V_0^2/c^2}}$$
(118)
We get
$$\frac{d\stackrel{}{p}_0}{dt}=\frac{m}{\sqrt{1V_0^2/c^2}}\frac{d\stackrel{}{V}_0}{dt}+\frac{m(\stackrel{}{V}_0\stackrel{}{\alpha }_0)\stackrel{}{V}_0}{c^2(1V_0^2/c^2)^{3/2}}=\stackrel{}{F}$$
(119)
Put Eq. 116 and 117 into Eq.(119), we get the dynamic equation and force of the particle in the gravitational field in flat space. In general, they are different from the Newtonian theory.
As for photon, we can obtain the corresponding equations (116) and (117) from its geometric equation in curved space, then define photon’s momentum in the same form of Eq.(56). After that, the dynamic equation of photon can be established in flat space-time. But it is unnecessary for us to discuss it nay more here. In this way, we have achieved the transformation of gravitation’s descriptions from curved space-time to flat space-time.
In the same, the gravitation’s descriptions can also be transformed from flat space-time into curved space-time. From Eq.(119) in flat space-time, we can get by solving the equation
$$x_0^i=x_0^i(t_0)$$
(120)
By introducing arbitrary coordinate transformation
$$x_0^\alpha =x_0^\alpha (x^\beta )$$
(121)
we have
$$\frac{dx_0^\alpha }{ds^2}=\frac{x_0^\alpha }{x^\beta }\frac{dx^\beta }{ds}$$
(122)
$$\frac{d^2x_o^\alpha }{ds^2}=\frac{^2x_0^\alpha }{x^\beta x^\sigma }\frac{dx^\beta }{ds}\frac{dx^\sigma }{ds}+\frac{x_0^\alpha }{x^\beta }\frac{d^2x^\beta }{ds^2}$$
(123)
Using Eq.(116) and (121), we get from Eq.(101)
$$ds=(c^2\frac{dx_0^i}{dt_0}\frac{dx_0^i}{dt_0})^{\frac{1}{2}}dt_0=(c^2V_0^iV_0^i)^{\frac{1}{2}}dt_0=A(t_0(x^\beta ))dt_0=A(x^\beta )dt_0$$
(124)
Put it into Eq.(122), use Eq.(116) and (121), we have
$$\frac{dx_0^\alpha }{dt_0}=A^1\frac{x_0^\alpha }{x^\beta }\frac{dx^\beta }{ds}=V_0^\alpha (t_0(x^\beta ))=V_0^\alpha (x^\beta )$$
(125)
Here $`V_0^0=c`$. Eq.(125) is the equation set of four variables and one order about $`dx^\beta /ds`$ . We can obtain from Eq.(125)
$$\frac{dx^\alpha }{ds}=B_\alpha (x^\beta )$$
(126)
Put Eq.(124) into Eq.(123) and using Eq.(126), we get
$$A^{}2\frac{d^2x_0^\alpha }{dt_0^2}=B_\beta B_\sigma \frac{^2x_0^\alpha }{x^\beta x^\sigma }+\frac{x_0^\alpha }{x^\beta }\frac{d^2x^\beta }{ds^2}A^1\frac{dA^1}{dt_0}\frac{dx_0^\alpha }{dt_0}$$
(127)
By using Eq.(124)and(126) again, we have
$$\frac{dA^1}{dt_0}=\frac{A^1}{x^\beta }\frac{dx^\beta }{dt_0}=\frac{A^1}{x^\beta }\frac{dx^\beta }{ds}\frac{ds}{dt_0}=\frac{A^1}{x^\beta }B_\beta A$$
(128)
We can write $`\alpha _0^i(t_0)=\alpha _0^i(t_0(x^\beta ))`$ and have $`\alpha _0^0=0`$ in Eq.(117). By considering Eq.(125) and (128), Eq.(127) can be written as:
$$\frac{x_0^\alpha }{x^\beta }\frac{d^2x^\beta }{ds^2}=(A^2F_0^i+B_\beta V_0^i\frac{A^1}{x^\beta }B_\beta B_\sigma \frac{^2x_0^\alpha }{x^\beta x^\sigma })$$
(129)
It is an equation set of four variables and one order about $`d^2x^\beta /ds^2`$ , in which all coefficients are the function of $`x^\beta `$ , so we can obtain
$$\frac{d^2x^\alpha }{ds^2}=K_\alpha (x^\beta )$$
(130)
The formula can be re-written as
$$\frac{d^2x^\alpha }{ds^2}K_\alpha (\frac{dx^\beta }{ds}\frac{dx^\sigma }{ds})^1\frac{dx^\beta }{ds}\frac{dx^\sigma }{ds}=\frac{d^2x^\alpha }{ds^2}K_\alpha B_\beta ^1B_\sigma ^1\frac{dx^\beta }{ds}\frac{dx^\sigma }{ds}=0$$
(131)
Let
$$\mathrm{\Gamma }_{\beta \sigma }^\alpha =K_\alpha B_\beta ^1B_\sigma ^1$$
(132)
Eq.(131) becomes
$$\frac{d^2x^\alpha }{ds^2}+\mathrm{\Gamma }_{\beta \sigma }^\alpha \frac{dx^\beta }{ds}\frac{dx^\sigma }{ds}=0$$
(133)
Regarding $`\mathrm{\Gamma }_{\beta \sigma }^\alpha `$ as the Christoffel sign, Eq.(133) is just the geodetic line equation in the new reference system. From the definition
$$\mathrm{\Gamma }_{\beta \sigma }^\alpha =\frac{1}{2}g^{\alpha \beta }(\frac{g_{\rho \sigma }}{x^\beta }+\frac{g_{\rho \beta }}{x^\sigma }\frac{g_{\beta \sigma }}{x^\rho })$$
(134)
we known that the number of independent $`g_{\alpha \beta }/x^\sigma `$ is just the same as the number of independent $`\mathrm{\Gamma }_{\beta \sigma }^\alpha `$ . Meanwhile, $`g_{\alpha \beta }`$ can be got from the following formula
$$g_{\alpha \beta }=\frac{g_{\alpha \beta }}{x^\sigma }𝑑x^\sigma $$
(135)
So the metric tensors can be determined by Eq.(134) and (135). It is noted that if Eq.(121) is put into Eq.(101) directly, we have
$$ds^2=g_{\alpha \beta }^{}dx^\alpha dx^\beta $$
(136)
136 The metric tensor $`g_{\alpha \beta }`$ is different from that shown in Eq.(135) in general. The metric tensor $`g_{\alpha \beta }^{}`$ shown in Eq.(136) is the Euclidean metric in essence for it can return to the original form Eq.(101) by an inverse transformation. But the metric tensor $`g_{\alpha \beta }`$ determined by Eq.(134) and (135) can not in general, so they are the non-Euclidean metrics in general.
Up to now we have achieved the transformation of gravitation’s descriptions between the curved space-time and the flat space-time, and proved their equivalence. The difference is that in curved space-time particles move along the geodetic lines without forces acting on them, but in flat space-time particles move along the non-geodetic lines acted by gravitation. What kinds of descriptive forms are taken depends on convenience in principle, but as shown above, the practical results should be considered.
3. Discussions on some foundational concepts
1. Is the space-time curved or flat after all when gravitational field exists?
This is first a problem of measurement. Whether can we answer this problem by the direct measurement? The answer is negative. Even thought the space-time is curved when the gravitational field exists, we can not detect it by the direct measurement. This is owing to the fact that before the measurement we have to define standard ruler and clock. But only in flat space-time, we can do them. In curved space-time, we have no definitions of standard ruler and clock. When the ruler and clock defined in flat space-time are put into gravitational fields, they would change or become curved synchronously with the fields, so that the measurements can not show the changes of curved level of space-time. The ruler and clock in the gravitational fields can not free themselves from the effects of gravitational fields, so it is impossible to show whether space-time is curved or flat when gravitational field exists by the direct measurement.
What we can do is to use indirect methods, for example, to observe the orbits of test particles orbit or the red shifts of spectral lines in the gravitational fields to decide the curved level of space-time. However, as shown above, we can describe the orbit of test particle in gravitational field by either geometric equation in curved space-time, or the dynamic equation in flat space-time. We can also explain the red shifts of spectral lines in the gravitational fields as the results of time delay or the potential energy’s changes. Two methods are equal to each other. So it is obvious that space-time itself can not be designated as curved or flat actually. The reality is what kind of reference systems, curved or flat, we choose to describe it. If the curved reference system is chosen, the space-time is curved. If the flat reference system is chosen, the pace-time becomes flat. It is meaningless to talk about space-time itself curved or flat. So we should only use the concept of curved or flat reference system, in spite of the concept of curved or flat space-time.
2. The equivalent principle
According to the weak equivalent principle, gravitational mass and inertial mass are equivalent to each other. Let $`m`$ represent static mass, represent inertial motion mass, we have
$$m_1=\frac{m}{\sqrt{1V^2/c^2}}$$
(137)
Let $`m_G`$ represent gravitational mass, comparing Eq.(22) with the Newtonian formula of gravitation and considering the relation $`\stackrel{}{L}=\stackrel{}{V}\times \stackrel{}{r}`$ , we get
$$m_G=m\sqrt{1\frac{V^2}{c^2}}[1+\frac{3(\stackrel{}{V}\times \stackrel{}{n})^2}{c^2}]$$
(138)
Here $`\stackrel{}{n}`$ is the unit radius vector. It is obvious $`m_1m_G`$ in general situations. Only when $`V=0`$ , they are equal to each other. In fact, all completed Eotvos experiments only prove that gravitational mass and inertial mass are equal to each other when the testing bodies on the two ends of cantilever beam are at rest each $`other^{(4)}`$. It has not yet be proved that they are equivalent when there exists relative motion between them. It should be noted that the formula (138) is the result of Einstein’s field equation, showing that gravitational mass and inertial mass are not equivalent actually when the factor of speed is considered.
3. The principle of general relativity
The paper’s conclusions are completely based on the Einstein equation of gravitational field, no any new hypothesis is introduced besides transforming the theory to the flat reference system to discuss. According to the general theory of relativity, it is equivalent to discuss physical problems in any reference system in nature. No one is more superior. Since we can discuss gravitational problems in any reference system, we can also discuss them in flat reference system. However, the results show that flat reference system seems more superior for the discussion of gravitational problems. The results contradict the principle of general relativity. So we have to discuss this problem further.
Einstein established special relativity that denied the existence of absolutely static reference system. Later, he put forwards the principle of general relativity trying to cancel the superior position of the inertial reference systems. If we consider the principle of general relativity as that the motion equations are covariant, or the basic forms of the motion equations are the same in any reference system, the principle is all right. However, it does not mean that the concrete forms of motion equations and their solutions are the same. In special relativity, the 4-diamention coordinate transformation means that the relative speed is introduced. Let $`p_\mu `$ represents the 4-diamention momentum, $`F_\mu `$ represents the 4-diamention force. The basic form of motion equation $`dp_\mu /dt=F_\mu `$ is unchanged when an inertial reference system is transformed into another inertial reference system moving in a uniform speed.. But the concrete forms of the 4-diamention force $`F_\mu `$ and particle’s motion, as well as and some physical quantities would change. For example, length constrict, time delay, moving mass increasing and the form of force changing and so on, though according to special relativity, these changes only have relative meanings.
In general relativity, the coordinate transformations involve more problems. At present, it is considered that a solution of gravitational field equation can still represent the same field after the solution has been transformed into another new reference system. This conclusion is worthy of further discussion and consideration. If the coordinate transformation is carried out in the 3-diamention space, there is no any problem. But in the 4-diamention space-time, because time is involved, the situation is completely different. In the general theory of relativity, the 4-diamention coordinate transformation means that acceleration or non-inertial reference system is introduced. According to the principle of equivalence, non-inertial reference system is equal to gravitational field. The transformation from one non-inertial reference system to another means that a gravitational field is changed to another. So the coordinate transformations would change physical processes. Speaking clearly, for a gravitational field with a determinate form, suppose we have obtained the solution by solving the Einstein’s field equation, if the solution is transformed into another reference system, the form of the solution would change, though the basic form of field equation is unchanged. In the light of the principle of equivalence, it means that a new gravitational field is introduced and the original solution loses its meaning. Therefore, a determinate gravitational field can only corresponds to a determinate metric, arbitrary coordinate transformation is forbidden according to the principle of equivalence. Unless the same results can be reached in new reference system for all problems, but this is impossible in general. For example, we can not calculate the perihelion precession of Mercury and other experiments and get the same results in the Edington or Kruskal coordinate systems. This is just the reason why the energy of gravitational field can not be defined well in the current theory. For a definite gravitational field with a certain of symmetry, we can only define its energy in the definite metric with the same symmetry. But it is allowed to transform the solutions into the inertial or flat reference systems to discuss. In this case, what is done is to transform the geodesic lines into the dynamic equations of gravitation without any attached force being introduced. Or speaking more clearly, there exist no relativity and arbitrariness in the description of gravitation. A certain of absoluteness is needed for us to describe gravitation. We should establish a united standard for gravitation. Only based on the flat reference system, we can do it.
It can be seen from discussions above that though the Einstein’s theory of gravitation has obtained great succession, there still exist some problems in its theoretical and logical foundation which need to be cleared and reformed so that the theory can become more rational. It is obvious that we need renewing some ideas about the essence of space-time and gravitation. The author will discuss them in detail later.
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# Untitled Document
Supported by the DFG-project HI 412/5-2
Formal dimension for semisimple symmetric spaces
Bernhard Krötz
Abstract
If $`G`$ is a semisimple Lie group and $`(\pi ,)`$ an irreducible unitary representation of $`G`$ with square integrable matrix coefficients, then there exists a number $`d(\pi )`$ such that
$$(v,v^{},w,w^{})\frac{1}{d(\pi )}v,v^{}w^{},w=_G\pi (g).v,w\overline{\pi (g).v^{}.w^{}}d\mu _G(g).$$
The constant $`d(\pi )`$ is called the formal dimension of $`(\pi ,)`$ and was computed by Harish-Chandra in \[HC56, 66\].
If now $`H\backslash G`$ is a semisimple symmetric space and $`(\pi ,)`$ an irreducible $`H`$-spherical unitary $`(\pi ,)`$ belonging to the holomorphic discrete series of $`H\backslash G`$, then one can define a formal dimension $`d(\pi )`$ in an analogous manner. In this paper we compute $`d(\pi )`$ for these class of representations.
Keywords: Holomorphic discrete series, highest weight representation, formal dimension, formal degree, spherical representation, $`c`$-functions. AMS classification: 22E46, 43A85.
Introduction
Let $`H\backslash G`$ be a semisimple irreducible simply connected non-compact symmetric space admitting relative holomorphic discrete series, i.e., there exists a unitary highest weight representation $`(\pi _\lambda ,_\lambda )`$ of $`G`$ and a non-zero $`H`$-invariant hyperfunction vector $`\nu _\lambda ^\omega `$ such that
$$\frac{1}{d(\lambda )}:=\frac{1}{|\nu ,v_\lambda |^2}_{HZ\backslash G}|\nu ,\pi _\lambda (g).v_\lambda |^2d\mu _{HZ\backslash G}(HZg)$$
is finite. Here $`v_\lambda `$ denotes a highest weight vector, $`Z`$ the center of $`G`$ and $`\mu _{HZ\backslash G}`$ a $`G`$-invariant measure on the homogeneous space $`HZ\backslash G`$. Note that $`v_\lambda `$ and $`\nu `$ are unique up to scalar multiple as well as $`\nu ,v_\lambda 0`$. Therefore the number $`d(\lambda )`$ is well defined and we call it the formal dimension of the spherical highest weight representation $`(\pi _\lambda ,_\lambda )`$. We remark here that our definition of the formal dimension generalizes Harish-Chandra notion in the “group case”, i.e., where $`G=G_0\times G_0`$ and $`H=\mathrm{\Delta }(G)=\{(g,g):gG_0\}`$ for a simply connected hermitian Lie group $`G_0`$ (cf. \[HC56\] and Remark III.5 below).
Note that the constants $`d(\lambda )`$ determine the part of the Plancherel measure on $`H\backslash G`$ which corresponds to the relative holomorphic discrete series. Thus the explicit knowledge of the formal dimensions gives us a better understanding of the Plancherel Theorem on $`H\backslash G`$ which was recently obtained by van den Ban-Schlichtkrull and Delorme (cf. \[BS97,99\], \[De98\]).
Let $`(g,\tau )`$ be the symmetric Lie algebra attached to $`H\backslash G`$ and write $`g=hq`$ for the $`\tau `$-eigenspace decomposition. If $`g=kp`$ is a $`\tau `$-invariant Cartan decomposition of $`g`$, then the algebraic characterization of $`H\backslash G`$ admitting relative holomorphic discrete series is $`z(k)q0`$. Symmetric Lie algebras $`(g,\tau )`$ having this property are called compactly causal (cf. \[HiÓl96\]). In the group case, i.e., $`(g,\tau )=(g_0g_0,\sigma )`$ with $`\sigma (X,Y)=(Y,X)`$ the flip involution, this just means that $`g_0`$ is hermitian. We remark here that the formal dimension in the group case was computed by Harish-Chandra (cf. \[HC56\]).
In this paper we derive the formula for the formal dimension $`d(\lambda )`$ for compactly causal symmetric spaces. For the special class of Cayley type spaces this problem has been dealt with by Chadli with Jordan algebra methods (cf. \[Ch98\]). The approach presented here is general and purely representation theoretic.
Our key result is the Averaging Theorem (cf. Theorem II.16) which asserts that for large parameters $`\lambda `$ the $`H`$-integral over $`v_\lambda `$ converges. More precisely, for large parameters $`\lambda `$ we prove that
$$_H\pi _\lambda (h).v_\lambda d\mu _H(h)=\frac{v_\lambda ,v_\lambda }{\nu ,v_\lambda }c(\lambda +\rho )\nu ,$$
where the left hand side has to be understood as a vector valued integral with values in the Fréchet space of hyperfunction vectors and
$$c(\lambda )=_{\overline{N}HAN}a_H(\overline{n})^{(\lambda +\rho )}𝑑\mu _{\overline{N}}(\overline{n})$$
denotes the $`c`$-function of the non-compactly causal $`c`$-dual space $`H^c\backslash G^c`$ (cf. \[HiÓl96\]).
To obtain the formula for the formal degree $`d(\lambda )`$, we plug in the relation for $`\nu `$ obtained from the Averaging Theorem in the definition of $`d(\lambda )`$ and obtain for large parameters:
$$d(\lambda )=d(\lambda )^Gc(\lambda +\rho ),$$
where $`d(\lambda )^G`$ is the formal dimension of $`(\pi _\lambda ,_\lambda )`$ for the relative discrete series on $`G`$ (cf. Theorem III.6). Using some ideas of Ólafsson and Ørsted (cf. \[ÓØ91\]) we prove the analytic continuation of our formula for $`d(\lambda )`$ (cf. Theorem IV.15).
The $`c`$-function $`c(\lambda )`$ can be written as a product
$$c(\lambda )=c_0(\lambda )c_\mathrm{\Omega }(\lambda ),$$
where $`c_0(\lambda )`$ is the $`c`$-function of a certain Riemannian symmetric subspace of $`H\backslash G`$ and $`c_\mathrm{\Omega }(\lambda )`$ is the $`c`$-function of the real form $`\mathrm{\Omega }`$ of the bounded symmetric domain $`𝒟G/K`$. In particular we have
$$d(\lambda )=d(\lambda )^Gc_0(\lambda +\rho )c_\mathrm{\Omega }(\lambda +\rho ).$$
The ingredients in this formula of $`d(\lambda )`$ are known: Harish-Chandra computed $`d(\lambda )^G`$ in \[HC56\], Gindikin and Karpelevič $`c_0(\lambda )`$ (cf. \[GiKa62\]) and finally Ólafsson and the author computed $`c_\mathrm{\Omega }(\lambda )`$ in \[KrÓl99\] (see also \[Fa95\], \[Gr97\] for earlier results in important special cases).
In the final section we give applications of our results to spherical holomorphic representation theory. Recall that a unitary highest weight representation $`(\pi _\lambda ,_\lambda )`$ of $`G`$ extends naturally to a holomorphic representation of the maximal open complex Ol’shanskiĭ semigroup $`S_{\mathrm{max}}^0=GExp(iW_{\mathrm{max}}^0)`$ (cf. \[Ne99b, Sect. XI.2\]). If $`(\pi _\lambda ,_\lambda )`$ is an $`H`$-spherical unitary highest weight representation of $`G`$, then we define its spherical character by
$$\mathrm{\Theta }_\lambda :S_{\mathrm{max}}^0C,s\frac{v_\lambda ,v_\lambda }{|\nu ,v_\lambda |^2}\pi _\lambda (s).\nu ,\nu .$$
Note that $`\mathrm{\Theta }_\lambda `$ is an $`H`$-biinvariant holomorphic function on $`S_{\mathrm{max}}^0`$. On the other hand on $`S_{\mathrm{max}}^0HAN`$ one defines the spherical function with parameter $`\lambda a_C^{}`$ by
$$\phi _\lambda :S_{\mathrm{max}}^0HANC,s_Ha_H(sh)^{\lambda \rho }𝑑\mu _H(h),$$
whenever the right hand side makes sense(cf. \[FHÓ94\] or \[KNÓ98\]). For large parameters $`\lambda `$ we prove the long searched relation of Ólafsson (cf. \[Ól98, Open Problem 7(1)\])
$$(sS_{\mathrm{max}}^0HAN)\mathrm{\Theta }_\lambda (s)=\frac{1}{c(\lambda +\rho )}\phi _{\lambda +\rho }(s)$$
(cf. Theorem V.4). Finally we want to point out that the results of this paper are a major step towards a proof of the Plancherel Theorem of $`G`$-invariant Hilbert spaces of holomorphic functions on $`G`$-invariant subdomains of the Stein variety $`\mathrm{\Xi }_{\mathrm{max}}^0=G\times _HiW_{\mathrm{max}}^0`$ (cf. \[Ch98\], \[HiKr98, 99b\], \[HÓØ91\], \[Kr98, 99b\], \[KNÓ97\], \[Ne99a\].)
I am grateful to J. Faraut and J. Hilgert who both made my stay in the stimulating atmosphere of Paris VI possible such that I could work on this problem. Also I want to thank G. Ólafsson for many exciting discussions on the subject and proofreading the manuscript. Finally I want to thank the referee for his very careful work.
I. Causal symmetric Lie algebras
This subsection is a brief introduction to causal symmetric Lie algebras. Purely algebraic definitions of “causality” are given and the basic notation on the algebraic level is introduced.
Definition I.1. Let $`g`$ denote a finite dimensional Lie algebra over the real numbers.
(a) A symmetric Lie algebra is a pair $`(g,\tau )`$, where $`\tau `$ is an involutive automorphism of $`g`$. We set
$$h:=\{Xg:\tau (X)=X\}\text{and}q:=\{Xg:\tau (X)=X\}$$
and note that $`g=h+q`$. We call $`(g,\tau )`$ irreducible, if $`\{0\}`$ and $`g`$ are the only $`\tau `$-invariant ideals of $`g`$.
(b) We denote by $`g_C`$ the complexification of $`g`$. If $`\tau `$ is a involution on $`g`$, we also denote by $`\tau `$ the complex linear extension of $`\tau `$ to a endomorphism of $`g_C`$.
(c) The $`c`$-dual $`g^c`$ of $`(g,\tau )`$ is defined by $`g^c=h+iq`$.
(d) If $`g`$ is semisimple, then there exists a Cartan involution $`\theta `$ of $`g`$ which commutes with $`\tau `$ (cf. \[Be57\] or \[KrNe96, Prop. I.5(iii)\]). We write $`g=k+p`$ for the corresponding Cartan decomposition. By subscripts we indicate intersections, for example $`h_k=hk`$ etc. Since $`\tau `$ and $`\theta `$ commute, we have $`g=h_k+h_p+q_k+q_p`$. The prescription $`\theta ^c:=\theta \tau _{g^c}`$ defines a Cartan involution on $`g^c`$ and we denote by $`g^c=k^c+p^c`$ the corresponding Cartan decomposition of $`g^c`$.
If $`Vg`$ is a subspace, then we set $`z(V)=\{XV:(YV)[X,Y]=0\}`$.
Definition I.2. Let $`(g,\tau )`$ be an irreducible semisimple symmetric Lie algebra and $`\theta `$ a Cartan involution of $`g`$ commuting with $`\tau `$. Then we call $`(g,\tau )`$
(CC) compactly causal if $`z(q_k)\{0\}`$. (NCC) non-compactly causal, if $`(g^c,\tau _{g^c})`$ is (CC). (CT) of Cayley type, if it is both (CC) and (NCC).
Lemma I.3.Let $`(g,\tau )`$ be a symmetric Lie algebra. Then the following assertions hold:
(i) The symmetric Lie algebra $`(g,\tau )`$ is compactly causal if and only if it belongs to one of the following two types: (1) The Lie algebra $`g`$ is simple hermitian and $`z(k)q`$. (2) The subalgebra $`h`$ is simple hermitian and $`(g,\tau )(hh,\sigma )`$, where $`\sigma `$ denotes the flip involution $`\sigma (X,Y)=(Y,X)`$.
(ii) If $`(g,\tau )`$ is compactly causal, then (a) $`z(k)q`$ is one-dimensional, (b) every maximal abelian subspace $`bq_k`$ is maximal abelian in $`q`$ and $`h_p+q_k`$.
Proof. (i) This follows from \[HiÓl96, Lemma 1.3.5, Th. 1.3.8\] or \[KrNe96, Prop. V.6\].
(ii) This is a consequence of \[HiÓl96, Prop. 3.1.11\].
Remark I.4. (a) From the view point of convex geometry and complex analysis the compactly causal symmetric spaces are the natural generalization of hermitian groups in the symmetric space setting (cf. \[HiÓl96\], \[KrNe96\], \[KNÓ97, 98\] and \[Ne99b\]). The compactly and non-compactly causal symmetric Lie algebras have been classified; for a complete list see \[HiÓl96, Th. 3.2.8\].
(b) Suppose that $`H\backslash G`$ is a simply connected symmetric space associated to an irreducible semisimple symmetric Lie algebra $`(g,\tau )`$. If $`(g,\tau )`$ is compactly causal, then Lemma I.3(ii)(b) implies that the symmetric space $`H\backslash G`$ admits relative holomorphic discrete series (cf. \[FJ80\]). The converse is also true. This result seems to us to be well known. But since we do not know a proof in the literature, we added a proof in Appendix B (cf. Lemma B.1).
Let $`(g,\tau )`$ be compactly causal. Recall that this implies in particular that $`g`$ is hermitian (cf. Lemma I.3(i)).
We choose a maximal abelian subalgebra $`iaq_k`$ and extend $`ia`$ in $`k`$ to a compactly embedded Cartan subalgebra $`t`$ of $`g`$. Recall from Lemma I.3(ii)(b) that $`a`$ is maximal abelian in $`iq`$ and $`p^c`$. Then $`t=t_h+ia`$ and $`z(k)qia`$. By Lemma I.3 we know that $`z(k)q=RZ_0`$ is one-dimensional and by \[Hel78, Ch. VIII, §7\] we can normalize $`Z_0`$ in such a way that $`Spec(adZ_0)=\{i,0,i\}`$ holds. We denote by $`\widehat{\mathrm{\Delta }}=\widehat{\mathrm{\Delta }}(g_C,t_C)`$ the root system of $`g_C`$ with respect to $`t_C`$ and by $`\mathrm{\Delta }=\mathrm{\Delta }(g^c,a)`$ the restricted root system of $`g^c`$ with respect to $`a`$. Note that $`\widehat{\mathrm{\Delta }}_a\backslash \{0\}=\mathrm{\Delta }`$. The corresponding root space decompositions are denoted by
$$g_C=t_C\underset{\widehat{\alpha }\widehat{\mathrm{\Delta }}}{}g_C^{\widehat{\alpha }}\text{and}g^c=az_h(a)\underset{\alpha \mathrm{\Delta }}{}(g^c)^\alpha .$$
A root $`\widehat{\alpha }\widehat{\mathrm{\Delta }}`$ is called compact if $`\widehat{\alpha }(Z_0)=0`$ and non-compact otherwise. Analogously one defines compact and non-compact roots in $`\mathrm{\Delta }`$. Write $`\widehat{\mathrm{\Delta }}_k`$ and $`\widehat{\mathrm{\Delta }}_n`$ for the set of all compact, resp. non-compact, roots in $`\widehat{\mathrm{\Delta }}`$. Analogously one defines $`\mathrm{\Delta }_k`$ and $`\mathrm{\Delta }_n`$.
Once and for all we fix now a positive system $`\widehat{\mathrm{\Delta }}^+`$ of $`\widehat{\mathrm{\Delta }}`$ such that
$$\widehat{\mathrm{\Delta }}_n^+:=\widehat{\mathrm{\Delta }}^+\widehat{\mathrm{\Delta }}_n=\{\widehat{\alpha }\widehat{\mathrm{\Delta }}_n:\widehat{\alpha }(Z_0)=i\}$$
holds. A positive system $`\mathrm{\Delta }^+`$ of $`\mathrm{\Delta }`$ is defined by $`\mathrm{\Delta }^+:=\widehat{\mathrm{\Delta }}^+_a\backslash \{0\}`$.
II. Spherical highest weight representations
In this section we briefly recall the classification of analytic and hyperfunction vectors of a a unitary highest weight representation $`(\pi _\lambda ,_\lambda )`$ of a simply connected compactly causal group $`(G,\tau )`$. Further we collect the basic facts of $`H`$-spherical highest weight representations. Then, after giving the definitions of the various $`c`$-functions associated to the non-compactly causal $`c`$-dual space $`(G^c,\tau )`$ of $`(G,\tau )`$, we prove the key result of the whole paper: The Averaging Theorem, which asserts that for large parameters $`\lambda `$ the $`H`$-integral over the highest weight vector converges in the Fréchet space of hyperfunction vectors. One obtains the up to scalar multiple uniquely determined $`H`$-spherical vector with a normalization constant which is given by a certain $`c`$-function.
Unitary highest weight representations
Recall that if $`G`$ is a simply connected Lie group associated to a symmetric Lie algebra $`(g,\tau )`$, then $`\tau `$ integrates to an involution on $`G`$, also denoted by $`\tau `$, and that the fixed point group $`G^\tau `$ is connected (cf. \[Lo69, Th. 3.4\].)
To a compactly causal symmetric Lie algebra $`(g,\tau )`$ we associate the following analytic objects:
$$\begin{array}{cc}G& \text{simply connected Lie group with Lie algebra }g,\\ G^c& \text{simply connected Lie group with Lie algebra }g^c,\\ G_C& \text{simply connected Lie group with Lie algebra }g_C,\\ H& \tau \text{-fixed points in }G,\\ H^c& \tau \text{-fixed points in }G^c,\\ H_C& \tau \text{-fixed points in }G_C,\\ K& \text{analytic subgroup in }G\text{ corresponding to }k,\\ K^c& \text{analytic subgroup in }G^c\text{ corresponding to }k^c,\\ K_C& \text{analytic subgroup in }G_C\text{ corresponding to }k_C,\\ H^0& \text{centralizer of }a\text{ in }H,\\ H^{c,0}& \text{centralizer of }a\text{ in }H^c,\\ Z& \text{center of }G.\end{array}$$
Note that even though both $`H`$ and $`H^c`$ are connected and have the same Lie algebra, they are in general not isomorphic. Recall that $`ZK`$.
If $`X`$ is a topological space and $`YX`$ is a subspace, then we denote by $`Y^0`$ or $`intY`$ the interior of $`Y`$ in $`X`$.
For each $`Xg_C`$ we denote by $`\overline{X}`$ the complex conjugate of $`X`$ with respect to the real form $`g`$.
Definition II.1. (Complex Ol’shanskiĭ semigroups, cf. \[Ne99, Ch. XI\]) Let $`(g,\tau )`$ be a compactly causal symmetric Lie algebra and $`\widehat{\mathrm{\Delta }}^+=\widehat{\mathrm{\Delta }}^+(g_C,t_C)`$ be the positive system from Section I.
(a) Associated to $`\widehat{\mathrm{\Delta }}^+`$ we define the maximal cone in $`t`$ by
$$\widehat{C}_{\mathrm{max}}=\{Xt:(\alpha \widehat{\mathrm{\Delta }}_n^+)i\alpha (X)0\}.$$
We set $`\widehat{W}_{\mathrm{max}}:=\overline{Ad(G).\widehat{C}_{\mathrm{max}}}`$ and note that $`\widehat{W}_{\mathrm{max}}`$ is a closed convex $`Ad(G)`$-invariant convex cone in $`g`$ admitting no affine lines and which is maximal with respect to these properties (cf. \[Ne99b, Sect. VII.3\]).
(b) Let $`G_1:=\mathrm{exp}_{G_C}(g)`$. By Lawson’s Theorem $`S_{\mathrm{max},1}:=G_1\mathrm{exp}_{G_C}(i\widehat{W}_{\mathrm{max}})`$ is a closed subsemigroup of $`G_C`$, the maximal complex Ol’shanskiĭ semigroup, and the polar map
$$G_1\times \widehat{W}_{\mathrm{max}}S_{\mathrm{max},1},(g,X)g\mathrm{exp}(iX)$$
is a homeomorphism (cf. \[La94, Th. 3.4\]).
Denote by $`S_{\mathrm{max}}`$ the universal covering semigroup of $`S_{\mathrm{max},1}`$ and write $`Exp:g+i\widehat{W}_{\mathrm{max}}S_{\mathrm{max}}`$ for the lifting of $`\mathrm{exp}_{G_C}_{g+i\widehat{W}_{\mathrm{max}}}:g+i\widehat{W}_{\mathrm{max}}S_{\mathrm{max},1}`$. Then it is easy to see that $`S_{\mathrm{max}}=GExp(i\widehat{W}_{\mathrm{max}})`$ and that there is a polar map
$$G\times \widehat{W}_{\mathrm{max}}S,(g,X)gExp(iX)$$
which is homeomorphism. We define the interior of $`S_{\mathrm{max}}`$ by $`S_{\mathrm{max}}^0:=GExp(i\widehat{W}_{\mathrm{max}}^0)`$. Note that $`S_{\mathrm{max}}^0`$ is an open semigroup ideal in $`S_{\mathrm{max}}`$ which carries a natural complex structure for which the semigroup multiplication is holomorphic. Further the prescription $`s=gExp(iX)s^{}=Exp(iX)g^1`$ defines on $`S_{\mathrm{max}}`$ the structure of an involutive semigroup. Note that the involution is antiholomorphic on $`S_{\mathrm{max}}^0`$.
Remark II.2. Let $`(\pi _\lambda ,_\lambda )`$ be a unitary highest weight representation of $`G`$ with respect to $`\widehat{\mathrm{\Delta }}^+`$ and highest weight $`\lambda it^{}`$. Denote by $`B(_\lambda )`$ the $`C^{}`$-algebra of bounded operators on $`_\lambda `$. Recall from \[Ne99b, Th. XI.4.8\] that $`(\pi _\lambda ,_\lambda )`$ has a natural extension to a holomorphic representation $`\pi _\lambda :S_{\mathrm{max}}B(_\lambda )`$ of $`S_{\mathrm{max}}`$, i.e., $`\pi _\lambda `$ is strongly continuous, holomorphic when restricted to $`S_{\mathrm{max}}^0`$ and satisfies $`\pi _\lambda (s^{})=\pi _\lambda (s)^{}`$ for all $`sS_{\mathrm{max}}`$.
Note that for $`X\widehat{W}_{\mathrm{max}}`$ one has $`\pi _\lambda (Exp(iX))=e^{id\pi _\lambda (X)}`$.
Definition II.3. Let $`G`$ be a Lie group and $``$ a Hilbert space. If $`(\pi ,)`$ is a unitary representation of $`G`$, then we call $`v`$ an analytic vector if the orbit map $`G,g\pi (g).v`$ is analytic. We denote by $`^\omega `$ the vector space of all analytic vectors of $`(\pi ,)`$. There is a natural locally convex topology on $`^\omega `$ for which the representation $`(\pi ^\omega ,^\omega )`$ of $`G`$ on $`^\omega `$ is continuous (cf. \[KNÓ97, Appendix\]). The strong antidual of $`^\omega `$ is denoted by $`^\omega `$ and the elements of $`^\omega `$ are called hyperfunction vectors. Note that there is a natural chain of continuous inclusions
$$^\omega ^\omega .$$
The natural extension of $`(\pi ,)`$ to a representation on the space of hyperfunction vectors is denoted by $`(\pi ^\omega ,^\omega )`$ and given explicitly by
$$\pi ^\omega (g).\nu ,v:=\nu ,\pi ^\omega (g^1).v.$$
Proposition II.4.Let $`(\pi _\lambda ,_\lambda )`$ be a unitary highest weight representation of $`G`$ with respect to $`\widehat{\mathrm{\Delta }}^+`$ and highest weight $`\lambda `$. Let $`Xint\widehat{W}_{\mathrm{max}}`$ be an arbitrary element. Then the analytic vectors of $`(\pi _\lambda ,_\lambda )`$ are given by
$$_\lambda ^\omega =\underset{t>0}{}\pi _\lambda (Exp(tiX))._\lambda $$
and the topology on $`_\lambda ^\omega `$ is the finest locally convex topology making for all $`t>0`$ the maps $`_\lambda _\lambda ^\omega ,v\pi _\lambda (Exp(tiX)).v`$ continuous.
Proof. \[KNÓ98, Prop. A.5\].
If $`\lambda it^{}`$ is dominant integral for $`\widehat{\mathrm{\Delta }}_k^+`$, we denote by $`(\pi _\lambda ^K,F(\lambda ))`$ the irreducible highest weight representation of $`K`$ with highest weight $`\lambda `$. Note that $`(\pi _\lambda ^K,F(\lambda ))`$ extends naturally to a holomorphic representation of the universal covering group $`\stackrel{~}{K_C}`$ of $`K_C`$ and which we denote by the same symbol.
Remark II.5. Recall that $`(\pi _\lambda ,_\lambda )`$ can be realized in the Fréchet space $`Hol(𝒟,F(\lambda ))`$ of $`F(\lambda )`$-valued holomorphic functions on the Harish-Chandra realization $`𝒟`$ of the hermitian symmetric space $`G/K`$. So let us assume $`_\lambda Hol(𝒟,F(\lambda ))`$. Then for all $`z𝒟`$ and $`vF(\lambda )`$ the point evaluation
$$_\lambda C,ff(z),v$$
is continuous, hence can be written as $`f(z),v=f,K_{z,v}^\lambda `$ for some $`K_{z,v}^\lambda _\lambda `$. One can show that all vectors $`K_{z,v}^\lambda `$ are analytic. Then, if $`\overline{_\lambda }`$ denotes the closure of $`_\lambda `$ in the nuclear Fréchet space $`Hol(𝒟,F(\lambda ))`$, then the mapping
$$r:_\lambda ^\omega Hol(𝒟,F(\lambda )),\nu r(\nu );r(\nu )(z),v=\nu (K_{z,v}^\lambda )$$
is a $`G`$-equivariant topological isomorphism onto its closed image $`imr=\overline{_\lambda }`$. In particular $`_\lambda ^\omega `$ is a nuclear Fréchet space (cf. \[Kr99a, Sect. III\] for all that).
Spherical representations
Definition II.6. Let $`G`$ be a Lie group, $`HG`$ a closed subgroup and $`(\pi ,)`$ a unitary representation of $`G`$. Then we write $`(^\omega )^H`$ for the set of all those elements $`\nu ^\omega `$ satisfying $`\pi ^\omega (h).\nu =\nu `$ for all $`hH`$. The unitary representation $`(\pi ,)`$ is called $`H`$-spherical if there exists a cyclic vector $`\nu (^\omega )^H`$ for $`(\pi ^\omega ,^\omega )`$.
For $`\lambda it^{}`$ dominant integral with respect to $`\widehat{\mathrm{\Delta }}_k^+`$ recall the definition of the generalized Verma module
$$N(\lambda ):=𝒰(g_C)_{𝒰(k_Cp^+)}F(\lambda )$$
which is a highest weight module of $`g`$ with respect to $`\widehat{\mathrm{\Delta }}^+`$ and highest weight $`\lambda `$ (cf. \[EHW83\]). We denote by $`L(\lambda )`$ the unique irreducible quotient of $`N(\lambda )`$.
Proposition II.7.Let $`(\pi _\lambda ,_\lambda )`$ be a unitary highest weight representation of $`G`$ with respect to $`\widehat{\mathrm{\Delta }}^+`$.
(i) If $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical, then $`(\pi _\lambda ^K,F(\lambda ))`$ is $`HK`$-spherical. In particular $`\lambda a^{}`$ and the highest weight vector $`v_\lambda _\lambda `$ is fixed by $`H^0`$. (ii) The restriction mapping
$$\mathrm{Res}:(^\omega )^HF(\lambda )^{HK},\nu \nu _{F(\lambda )}$$
is injective. In particular, $`dim(^\omega )^H1`$ and $`\nu ,v_\lambda 0`$ for $`\nu 0`$. Moreover if $`L(\lambda )=N(\lambda )`$, then $`\mathrm{Res}`$ is a bijection, i.e., $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical if and only if $`(\pi _\lambda ^K,F(\lambda ))`$ is $`HK`$-spherical.
Proof. (i) is a special case of \[KNÓ97, Prop. VI.5\] and (ii) a special case of \[Kr99a, Th. III.14\].
Remark II.8. In general it is not true that $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical if the minimal $`K`$-type $`(\pi _\lambda ^K,F(\lambda ))`$ is $`HK`$-spherical. For a counter example see \[Kr99a, Ex. III.16\].
The $`c`$-functions on the $`c`$-dual space $`H^c\backslash G^c`$.
To the positve system $`\mathrm{\Delta }^+=\mathrm{\Delta }^+(g^c,a)`$ we associate several subalgebras of $`g^c`$
$$\begin{array}{cc}n=_{\alpha \mathrm{\Delta }^+}(g^c)^\alpha ,& \overline{n}=_{\alpha \mathrm{\Delta }^{}}(g^c)^\alpha ,\\ & \\ n_n^\pm =_{\alpha \mathrm{\Delta }_n^\pm }(g^c)^\alpha ,& n_k^\pm =_{\alpha \mathrm{\Delta }_k^\pm }(g^c)^\alpha .\end{array}$$
Further we set
$$p^\pm :=\underset{\widehat{\alpha }\widehat{\mathrm{\Delta }}_n^+}{}g_C^{\widehat{\alpha }}\text{and}g(0):=h_k+iq_kg^c.$$
Remark II.9. (a) The subalgebras $`p^\pm `$ and $`k_C`$ of $`g_C`$ are invariant under complex conjugation with respect to $`g^c`$ and we have $`p^\pm g^c=n_n^\pm `$ as well as $`k_Cg^c=g(0)`$. Thus the decomposition $`g_C=p^+k_Cp^{}`$ induces a splitting in subalgebras of $`g^c`$
$$g^c=n_n^+g(0)n_n^{}.$$
(b) Recall that $`g^c=han`$. The $`han`$\- decomposition restricted to $`g(0)`$ coincides with an Iwasawa decomposition of $`g(0)`$ given by $`g(0)=k(0)an_k^+`$, where $`k(0):=hg(0)=k^cg(0)`$.
We let $`H_CK_C`$ act on $`H_C\times K_C`$ from the left by $`x.(h,k):=(hx^1,xk)`$ and denote by
$$M:=H_C\times _{H_CK_C}K_C$$
the corresponding quotient space. The $`H_CK_C`$-coset of an element $`(h,k)K_C\times H_C`$ is denoted by $`[h,k]`$. If $`\stackrel{~}{H_C}`$ and $`\stackrel{~}{K_C}`$ denote the universal coverings of $`H_C`$ and $`K_C`$, respectively, then we realize the universal cover $`\stackrel{~}{M}`$ of $`M`$ by
$$\stackrel{~}{M}=\stackrel{~}{H_C}\times _{(\stackrel{~}{H_C}\stackrel{~}{K_C})_0}\stackrel{~}{K_C}.$$
Further let $`P^\pm :=\mathrm{exp}_{G_C}(p^\pm )`$. Recall that $`p^\pm `$ are abelian and that the exponential mapping $`\mathrm{exp}_{G_C}_{p^\pm }:p^\pm P^\pm `$ is an isomorphism. In particular $`P^\pm `$ is simply connected.
Proposition II.10. (The $`H_CK_CP^+`$-decomposition) The following assertions hold:
(i) The multiplication mapping
$$M\times P^+G_C,([h,k],p_+)hkp_+$$
is a biholomorphic map onto its open image $`H_CK_CP^+`$. Furthermore: (a) The open submanifold $`H_CK_CP^+`$ is dense in $`G_C`$ with complement of Haar measure zero. (b) We have $`S_{\mathrm{max},1}H_CK_CP^+`$. (ii) If $`j:S_{\mathrm{max},1}M\times P^+`$ denotes the injection obtained from the isomorphism in (i), then $`j`$ lifts to an inclusion mapping $`\stackrel{~}{j}:S_{\mathrm{max}}\stackrel{~}{M}\times P^+`$.
Proof. (i) \[KNÓ97, Prop. II.6, Lemma III.7\].
(ii) Since $`\pi _1(S_{\mathrm{max},1})=\pi _1(G_1)Z(G)Z(K)`$, it suffices to show that $`\stackrel{~}{j}_K`$ is injective. We may assume that $`K\stackrel{~}{K_C}`$, since both $`K`$ and $`\stackrel{~}{K_C}`$ are simply connected and $`k`$ is a maximal compact subalgebra of $`k_C`$. Further $`K`$ normalizes $`P^+`$, and so establishing the injectivity of $`\stackrel{~}{j}_K`$ boils down to proving injectivity of $`K\stackrel{~}{M},k[\mathrm{𝟏},k]`$, which is obvious.
We denote by $`G(0)`$, $`A`$, $`N`$, $`\overline{N}`$, $`N_k^\pm `$ and $`N_n^\pm `$ the analytic subgroups of $`G^c`$ corresponding to $`g(0)`$, $`a`$, $`n`$, $`\overline{n}`$, $`n_k^\pm `$ and $`n_n^\pm `$.
Remark II.11. (a) In view of the Bruhat decomposition of $`\stackrel{~}{K_C}`$, we may identify $`AN_k^+`$ as a subgroup of $`\stackrel{~}{K_C}`$. Note that $`N=N_k^+N_n^+`$ and so every $`nN`$ can be written uniquely as $`n=n_kn_n`$ with $`n_kN_k^+`$ and $`n_nN_n^+`$. Thus we conclude from Proposition II.10(ii) that the map
$$H\times A\times N\stackrel{~}{M}\times P^+,(h,a,n_kn_n)([h,an_k],n_n)$$
is an analytic diffeormphism onto its image which we denote by $`HAN`$. Accordingly every element $`sHAN`$ can be written uniquely as $`s=h_H(s)a_H(s)n_H(s)`$ with $`h_H(s)H`$, $`a_H(s)A`$ and $`n_H(s)N`$ all depending analytically on $`sHAN`$.
(b) If $`𝒟p^+`$ denotes the Harish-Chandra realization of the hermitian symmetric space $`G/K`$ and $`\overline{𝒟}`$ its conjugate in $`p^{}`$, then we set $`\mathrm{\Omega }:=\overline{𝒟}n_n^{}`$. In the sequel we realize $`\mathrm{\Omega }`$ as a subset of $`N_n^{}`$ via the exponential mapping. Recall from \[KNÓ98, Lemma I.18\] that
$$H^cAN=\mathrm{\Omega }G(0)N_n^+\text{and}\overline{N}H^cAN=\mathrm{\Omega }N_k^{}.$$
On the other hand $`\mathrm{\Omega }`$ can also naturally be realized in $`\stackrel{~}{M}\times P^+`$. In particular we obtain that the submanifold $`\mathrm{\Omega }N_k^{}`$ of $`\overline{N}`$ is naturally included in $`\stackrel{~}{M}\times P^+`$. Denote this realization by $`\overline{N}HAN`$. Further the $`HAN`$-decomposition and the $`H^cAN`$-decomposition (cf. \[KNÓ97, Prop. II.4\]) coincide on $`\overline{N}HAN`$. In the sequel we will use this fact frequently without mentioning it.
(c) Let $`p:XH^cAN`$ the universal covering of $`H^cAN`$. Since $`X`$ is simply connected, there exists a natural regular map $`\pi :X\stackrel{~}{M}\times P^+`$ with $`\pi (X)=HAN`$. In particular, the prescription
$$K^cHAN:=\pi (p^1(K^cH^cAN))$$
defines an open submanifold of $`HAN`$.
Note that the exponential mapping $`\mathrm{exp}_{\stackrel{~}{K_C}}_a:aA`$ is an isomorphism, hence has an inverse which we denote by $`\mathrm{log}:Aa`$. For each $`\lambda a_C^{}`$ and $`aA`$ we set $`a^\lambda =e^{\lambda (\mathrm{log}a)}`$.
Definition II.12. (The $`c`$-functions) We write $`\rho `$, $`\rho _k`$ and $`\rho _n`$ for the elements of $`a^{}`$ given by $`\frac{1}{2}trad_n`$, $`\frac{1}{2}trad_{n_k^+}`$ and $`\frac{1}{2}trad_{n_n^+}`$, respectively. To $`\lambda a_C^{}`$ we associate the following $`c`$-functions:
$$\begin{array}{cc}\hfill c(\lambda )& :=_{\overline{N}(HAN)}a_H(\overline{n})^{(\lambda +\rho )}𝑑\mu _{\overline{N}}(\overline{n}),\hfill \\ \hfill c_\mathrm{\Omega }(\lambda )& :=_\mathrm{\Omega }a_H(\overline{n})^{(\lambda +\rho )}𝑑\mu _{N_n^{}}(\overline{n}),\hfill \end{array}$$
and
$$c_0(\lambda ):=_{N_k^{}}a_H(\overline{n})^{(\lambda +\rho _k)}𝑑\mu _{N_k^{}}(\overline{n})$$
provided the integrals converge absolutely (cf. \[FHÓ94\] and \[KNÓ98\]). We write $``$ for the set of all $`\lambda a_C^{}`$ for which the defining integral for $`c`$ converges absolutely. Accordingly we define $`_\mathrm{\Omega }`$ and $`_0`$. Note that $`c_0`$ is the $`c`$-function of the non-compact Riemannian symmetric space $`K(0)\backslash G(0)`$, where $`K(0):=G(0)^\tau `$.
For each $`\alpha \mathrm{\Delta }^+`$ let $`\stackrel{ˇ}{\alpha }a`$ be the corresponding coroot, i.e., $`\stackrel{ˇ}{\alpha }[(g^c)^\alpha ,(g^c)^{\tau \alpha }]`$ such that $`\alpha (\stackrel{ˇ}{\alpha })=2`$. Associated to $`\mathrm{\Delta }^+`$ we define two minimal cones in $`a`$ by
$$C_{\mathrm{min}}:=cone(\{\stackrel{ˇ}{\alpha }:\alpha \mathrm{\Delta }_n^+\})\text{and}\stackrel{ˇ}{C}_k:=cone(\{\stackrel{ˇ}{\alpha }:\alpha \mathrm{\Delta }_k^+\}).$$
Definition II.13. Let $`V`$ be a finite dimensional vector space and $`V^{}`$ its dual.
(a) If $`CV`$ is a convex set, then its limit cone is defined by $`limC=\{xV:x+CC\}`$. Note that $`limC`$ is a convex cone and that $`limC`$ is closed if $`C`$ is open or closed.
(b) If $`EV`$ is a subset, then its dual cone is defined by $`E^{}:=\{\alpha V^{}:\alpha _V0\}`$. Note that $`E^{}`$ is a closed convex cone in $`V^{}`$.
Theorem II.14.The various $`c`$-functions are related by
$$c(\lambda )=c_0(\lambda )c_\mathrm{\Omega }(\lambda )$$
and $`=_\mathrm{\Omega }_0`$. Further: (i) The domain of convergence $`_\mathrm{\Omega }`$ of $`c_\mathrm{\Omega }`$ is a tube domain $`_\mathrm{\Omega }=ia^{}+_{\mathrm{\Omega },R}`$ with
$$_{\mathrm{\Omega },R}=\{\lambda a^{}:(\alpha \mathrm{\Delta }_n^+)\lambda (\stackrel{ˇ}{\alpha })<2m_\alpha \},$$
where $`m_\alpha :=dim(g^c)^\alpha `$. Further for all $`\lambda _\mathrm{\Omega }`$ we have
$$c_\mathrm{\Omega }(\lambda )=\underset{\alpha \mathrm{\Delta }_n^+}{}B(\frac{\lambda (\stackrel{ˇ}{\alpha })}{2}\frac{m_\alpha }{2}+1,\frac{m_\alpha }{2}),$$
where $`B`$ denotes the Euler Beta-function. In particular: (a) $`\rho C_{\mathrm{min}}^{}_{\mathrm{\Omega },R}`$ and $`lim_{\mathrm{\Omega },R}=C_{\mathrm{min}}^{}`$. (b) The function $`c_\mathrm{\Omega }`$ is holomorphic on $`_\mathrm{\Omega }`$ and $`c_\mathrm{\Omega }_{_\mathrm{\Omega }+\mu }`$ is bounded for all $`\mu \rho C_{\mathrm{min}}^{}`$. (ii) The domain of convergence of $`c_0`$ is given by
$$_0=ia^{}+int\stackrel{ˇ}{C}_k^{},$$
$`c_0`$ is holomorphic on $`_0`$ and $`c_0_{_0+\mu }`$ is bounded for all $`\mu \rho _k+\stackrel{ˇ}{C}_k^{}`$.
Proof. The product formula $`c(\lambda )=c_0(\lambda )c_\mathrm{\Omega }(\lambda )`$ and the relation $`=_\mathrm{\Omega }_0`$ are a special case of \[KNÓ98, Lemma IV.5\].
(i) \[KrÓl99, Th. III.5\].
(ii) All this follows from the Gindikin-Karpelevic product formula for $`c_0`$ (cf. \[Hel84, Ch. IV, Th. 6.13\]).
The Averaging Theorem
Lemma II.15.The group $`H^0`$ is compact and up to normalization of Haar measures for all $`fL^1(H/H^0)`$ the following integration formulas hold: (i)
$$_Hf(hH^0)𝑑\mu _H(h)=_{\overline{N}(HAN)}f(h_H(\overline{n})H^0)a_H(\overline{n})^{2\rho }𝑑\mu _{\overline{N}}(\overline{n}).$$
(ii)
$$_Hf(hH^0)𝑑\mu _H(h)=_{K^c(HAN)}f(h_H(k)H^0)a_H(k)^{2\rho }𝑑\mu _{K^c}(k).$$
Proof. In \[KNÓ98, Lemma III.15(i)\] it is proved that $`H^{c,0}`$ is compact and exactly the same argument also yields that $`H^0`$ is compact. In view of this fact and our identifications of the various domains in the big complex manifold $`\stackrel{~}{M}\times P^+`$ (cf. Remark II.11), (i) follows from \[KNÓ98, Prop. 1.19\] and (ii) from \[Ól87, Lemma 1.3\].
Theorem II.16. (The Averaging Theorem) Let $`(\pi _\lambda ,_\lambda )`$ be a unitary highest weight representation of $`G`$ for which $`(\pi _\lambda ^K,F(\lambda ))`$ is $`HK`$-spherical. If $`v_\lambda `$ is a highest weight vector, then the vector valued integral
$$_H\pi _\lambda (h).v_\lambda d\mu _H(h)$$
with values in the Fréchet space $`_\lambda ^\omega `$ (cf. Remark II.5) converges and defines a non-zero $`H`$-fixed hyperfunction vector if and only if $`\lambda +\rho _\mathrm{\Omega }`$. If this condition is satisfied and $`0\nu (_\lambda ^\omega )^H`$, then
$$_H\pi _\lambda (h).v_\lambda d\mu _H(h)=\frac{v_\lambda ,v_\lambda }{\nu ,v_\lambda }c(\lambda +\rho )\nu .$$
Proof. Step 1: The analytic function $`S_{\mathrm{max}}^0HAN_\lambda ,s\pi _\lambda (s).v_\lambda `$ extends to an analytic function $`F:HAN_\lambda `$ and is given ecplicitly by $`F(s)=a_H(s)^\lambda \pi _\lambda (h_H(s)).v_\lambda `$.
In fact since $`d\pi _\lambda (X).v_\lambda =0`$ for all $`Xn`$, the standard argument of differentiating yields
$$\pi _\lambda (s).v_\lambda =\pi _\lambda (h_H(s)a_H(s)n_H(s)).v_\lambda =\pi _\lambda (h_H(s)a_H(s)).v_\lambda =a_H(s)^\lambda \pi _\lambda (h_H(s)).v_\lambda ,$$
establishing Step 1.
Step 2: The integral exists if and only if $`\lambda +\rho _\mathrm{\Omega }`$.
Let $`Xint\widehat{W}_{\mathrm{max}}`$ be an arbitrary element and set $`a_t:=Exp(itX)`$ for all $`t>0`$. For each $`t>0`$ consider the possibly unbounded linear functional
$$f_t:_\lambda C,w_H\pi _\lambda (h).v_\lambda ,\pi _\lambda (a_t).wd\mu _H(h).$$
In view of Proposition II.4, we have to show that $`\lambda +\rho _\mathrm{\Omega }`$ is equivalent to $`f_t_\lambda ^{}`$ for all $`t>0`$.
Since $`v_\lambda `$ is fixed by $`H^0`$ (cf. Proposition II.7(i)), Step 1 and the integration formula of Lemma II.15(ii) yield
$$\begin{array}{cc}& _H\pi _\lambda (h).v_\lambda ,\pi _\lambda (a_t).wd\mu _H(h)\hfill \\ & =_{K^c(HAN)}\pi _\lambda (h_H(k)).v_\lambda ,\pi _\lambda (a_t).wa_H(k)^{2\rho }d\mu _{K^c}(k)\hfill \\ & =_{K^c(HAN)}\pi _\lambda (ka_H(k)^1).v_\lambda ,\pi _\lambda (a_t).wa_H(k)^{2\rho }d\mu _{K^c}(k)\hfill \\ & =_{K^c(HAN)}\pi _\lambda (a_tk).v_\lambda ,wa_H(k)^{(\lambda +2\rho )}d\mu _{K^c}(k).\hfill \end{array}$$
$`(2.1)`$
Recall from \[FHÓ94, Prop. 5.3\] that
$$_\mathrm{\Omega }=\{\lambda a_C^{}:_{K^c(HAN)}a_H(k)^{Re(\lambda +\rho )}𝑑\mu _{K^c}(k)<\mathrm{}\}.$$
$`(2.2)`$
In view of \[KNÓ98, Lemma III.15(ii)\], the set $`X_t:=\overline{a_t(K^cHAN)}`$ is a compact subset of $`HAN`$. In particular we find compact sets $`C_H^t`$, $`C_A^t`$, $`C_N^t`$ contained in $`H`$, $`A`$ and $`N`$, respectively, such that $`X_tC_H^tC_A^tC_K^t`$. Thus we conclude from Step 1 that
$$(w_\lambda )(xX_t)|\pi _\lambda (x).v_\lambda ,w|sup_{aC_A^t}a^\lambda v_\lambda w<\mathrm{}.$$
$`(2.3)`$
Hence, in view of (2.1), (2.2) and (2.3) the proof of Step 2 will be complete, provided we can show that for each element $`x`$ in the compact space $`X_t`$ we can find an open neighborhood $`UX_t`$ of $`x`$ and an element $`w_\lambda `$ such that $`inf_{yU}|\pi _\lambda (y).v_\lambda ,w|>0`$ holds. But this follows from $`\pi _\lambda (y).v_\lambda ,w=F(y),w`$ and the continuity of $`F`$.
Step 3: If the integral exists, then its value is $`\frac{v_\lambda ,v_\lambda }{\nu ,v_\lambda }c(\lambda +\rho )\nu `$.
By Step 1 we know that $`\lambda +\rho _\mathrm{\Omega }`$ in the case where the integral exists. Since $`\lambda `$ is a highest weight for an $`HK`$-spherical representation of $`K`$, it has to be dominant integral with respect to $`\mathrm{\Delta }_k^+`$, i.e., $`\lambda ,\alpha N_0`$ for all $`\alpha \mathrm{\Delta }_k^+`$. In particular $`c(\lambda +\rho )`$ exists (cf. Theorem II.14). Now by Step 2, we know that $`_H\pi _\lambda (h).v_\lambda d\mu _H(h)(_\lambda ^\omega )^H`$. Since $`dim(_\lambda ^\omega )^H1`$ (cf. Proposition II.7(ii)), it follows that $`_H\pi _\lambda (h).v_\lambda d\mu _H(h)=c\nu `$ for some constant $`cC`$. To determine $`c`$ we apply the integral to the element $`v_\lambda `$. With Step 1 and the integration formula of Lemma II.15(i) we compute
$$\begin{array}{cc}& _H\pi _\lambda (h).v_\lambda ,v_\lambda d\mu _H(h)\hfill \\ & =_{\overline{N}(HAN)}\pi _\lambda (h_H(\overline{n})).v_\lambda ,v_\lambda a_H(\overline{n})^{2\rho }d\mu _{\overline{N}}(\overline{n})\hfill \\ & =_{\overline{N}(HAN)}\pi _\lambda (\overline{n}a_H(\overline{n})^1).v_\lambda ,v_\lambda a_H(\overline{n})^{2\rho }d\mu _{\overline{N}}(\overline{n})\hfill \\ & =_{\overline{N}(HAN)}\pi _\lambda (\overline{n}).v_\lambda ,v_\lambda a_H(\overline{n})^{(\lambda +2\rho )}d\mu _{\overline{N}}(\overline{n})\hfill \\ & =v_\lambda ,v_\lambda _{\overline{N}(HAN)}a_H(\overline{n})^{(\lambda +2\rho )}𝑑\mu _{\overline{N}}(\overline{n})\hfill \\ & =v_\lambda ,v_\lambda c(\lambda +\rho ).\hfill \end{array}$$
This proves Step 3 and completes the proof of the theorem.
III. Representations of the relative discrete series
In this section we state and prove the Harish-Chandra - Godement Orthogonality relations for homogeneous spaces carrying an invariant measure. Then we give the definition of the formal dimension $`d(\lambda )`$ of a unitary highest weight representation $`(\pi _\lambda ,_\lambda )`$ which belongs to the relative discrete series of $`H\backslash G`$. Finally we derive the formula for $`d(\lambda )`$ for large values of $`\lambda `$.
Orthogonality Relations
Definition III.1. Let $`G`$ be a Lie group, $`Z`$ its center and $`\widehat{Z}`$ the group of unitary characters of $`Z`$.Let $`HG`$ be a closed subgroup. Suppose that $`HZ`$ is closed and that $`HZ\backslash G`$ carries an invariant positive measure $`\mu _{HZ\backslash G}`$. For a fixed $`\chi \widehat{Z}`$ we consider the Hilbert space of sections
$$\begin{array}{cc}\hfill \mathrm{\Gamma }_\chi ^2(H\backslash G)& =\{f:H\backslash GC:f\text{measurable},(zZ)(gG)f(Hzg)=\chi (z)f(Hg);\hfill \\ & f,f_\chi :=_{HZ\backslash G}|f(Hg)|^2d\mu _{HZ\backslash G}(HZg)<\mathrm{}\}.\hfill \end{array}$$
Let $`(\pi ,)`$ be an irreducible unitary $`H`$-spherical representation of $`G`$ with central character $`\chi `$. Then for all $`\nu (^\omega )^H`$ and $`v^\omega `$ we define a continuous section by
$$\pi _{v,\nu }:H\backslash GC,Hg\overline{\nu ,\pi (g).v}.$$
We say that $`(\pi ,)`$ belongs to the relative discrete series of $`H\backslash G`$, if there exists non-zero elements $`\nu (^\omega )^H`$ and $`v^\omega `$ such that $`\pi _{v,\nu }`$ belongs to $`\mathrm{\Gamma }_\chi ^2(H\backslash G)`$. We denote $`(^\omega )_2^H`$ the subspace of $`(^\omega )^H`$ which corresponds to the relative discrete series for $`H\backslash G`$.
In the proof of the following Proposition we adapt a nice idea of J. Faraut to our setting (cf. \[Gr96, Sect. III.3\]).
Proposition III.2. (Orthogonality Relations) Let $`G`$ be a Lie group with center $`Z`$. Then, if $`H`$ is a closed subgroup of $`G`$ such that $`HZ`$ is closed and $`HZ\backslash G`$ carries a positive $`G`$-invariant measure, then the following assertions hold:
(i) If $`(\pi ,)`$ belongs to the relative discrete series of $`H\backslash G`$ transforming under the central character $`\chi \widehat{Z}`$ and $`0\nu (^\omega )_2^H`$, then all matrix coefficients $`\pi _{v,\nu }`$, $`v^\omega `$, belong to $`\mathrm{\Gamma }_\chi ^2(H\backslash G)`$ and there exists a constant $`d(\pi ,\nu )`$ depending on the equivalence class of $`\pi `$ and on $`\nu `$ such that the mapping
$$T:^\omega \mathrm{\Gamma }_\chi ^2(H\backslash G),v\sqrt{d(\pi ,\nu )}\pi _{v,\nu }$$
extends to a $`G`$-equivariant isometry.
(ii) If $`(\pi ,)`$ and $`(\sigma ,𝒦)`$ are two inequivalent representations of the relative discrete series of $`HZ\backslash G`$ transforming under the same central character for $`Z`$, then for $`\nu (^\omega )_2^H`$ and $`\eta (𝒦^\omega )_2^H`$ one has
$$\pi _{v,\nu },\sigma _{w,\eta }=_{HZ\backslash G}\nu ,\pi (g).v\overline{\eta ,\sigma (g).w}d\mu _{HZ\backslash G}(HZg)=0$$
for all $`v^\omega `$ and $`w𝒦^\omega `$.
Proof. (i) (cf. \[Gr96, Sect. III.3\]) Let $`D:=\{v^\omega :\pi _{v,\nu }\mathrm{\Gamma }_\chi ^2(H\backslash G)\}`$ and consider the unbounded operator
$$S:D\mathrm{\Gamma }_\chi ^2(H\backslash G),v\pi _{v,\nu }.$$
Since $`\mu _{HZ\backslash G}`$ is $`G`$-invariant by assumption, the same holds for $`D`$ and therefore $`D`$ is dense in $``$ by the irreducibility of $`(\pi ,)`$. We define a positive hermitian form on $`D`$ by
$$(v|w):=v,w+S.v,S.w_\chi $$
$`(3.1)`$
for $`v,wD`$. Denote by $`\overline{D}`$ the Hilbert completion of $`D`$ with respect to $`(|)`$ and denote the extension of $`(|)`$ to its completion by the same symbol. Since $`\overline{D}`$ is continuously embedded into $``$, there exists a bounded selfadjoint injective operator $`AB()`$ such that $`imA=\overline{D}`$ and $`(A.v|w)=v,w`$ for all $`v`$, $`w\overline{D}`$. Since $`,_\chi `$ is $`G`$-invariant by the $`G`$-invariance of $`\mu _{HZ\backslash G}`$, it follows from (3.1) that $`A`$ commutes with $`\pi (G)`$. Now Schur’s Lemma applies and yields $`A=cid`$ for some constant $`c>0`$. Thus we deduce from (3.1) that
$$S.v,S.w_\chi =(\frac{1}{c}1)v,w$$
for all $`v,wD`$. In particular $`d(\pi ,\nu ):=\left(\frac{1}{c}1\right)>0`$. Moreover $`S`$ being weakly continuous, its extension to $`^\omega `$ coincides with $`\frac{1}{\sqrt{d(\pi ,\nu )}}T`$, concluding the proof of (i).
(ii) Let $`T_\pi :\mathrm{\Gamma }_\chi ^2(H\backslash G)`$ and $`T_\sigma :𝒦\mathrm{\Gamma }_\chi ^2(H\backslash G)`$ be the equivariant isometric embeddings from (i). If $`imT_\pi imT_\sigma \{0\}`$, then
$$T_\sigma ^{}T_\pi :𝒦$$
describes a non-trivial $`G`$-equivariant map. By Schur’s Lemma $`T_\sigma ^{}T_\pi `$ is a scalar multiple of an isometric isomorphism, contradicting the inequivalence of $`(\pi ,)`$ and $`(\sigma ,𝒦)`$.
Remark III.3. If $`H\backslash G`$ is a semisimple symmetric space, then the space $`(^\omega )_2^H=(^{\mathrm{}})_2`$ is finite dimensional (cf. \[Ba87, Th. 3.1\]). Then Proposition III.2(i) says that one can find an inner product on $`(^\omega )_2^H`$ such that
$$(^\omega )_2^H^\omega \mathrm{\Gamma }_\chi ^2(H\backslash G),\nu v\sqrt{d(\pi ,\nu )}\pi _{v,\nu }$$
extends to a $`G`$-equivariant isometry (with $`G`$ acting trivially on the first factor $`(^\omega )_2^H`$ of the tensor product).
The formal dimension
If $`G`$ denotes a unimodular locally compact group and $`LG`$ a closed unimodular subgroup, then we denote by $`\mu _{L\backslash G}`$ a positive right $`G`$-invariant measure on the homogeneous space $`L\backslash G`$.
Definition III.4. Let $`(\pi _\lambda ,_\lambda )`$ be an $`H`$-spherical unitary highest weight representation of $`G`$ and $`0\nu (_\lambda ^\omega )^H`$. If $`v_\lambda `$ is a highest weight vector for $`(\pi _\lambda ,_\lambda )`$, then the formal dimension $`d(\lambda )`$ of $`(\pi _\lambda ,_\lambda )`$ is defined by
$$\frac{1}{d(\lambda )}:=\frac{1}{|\nu ,v_\lambda |^2}_{HZ\backslash G}|\nu ,\pi _\lambda (g).v_\lambda |^2d\mu _{HZ\backslash G}(HZg).$$
Recall that $`\nu ,v_\lambda 0`$ and that the definition of $`d(\lambda )`$ is independent of the particular choice of $`v_\lambda `$ and $`0\nu (_\lambda ^\omega )^H`$ (cf. Proposition II.7(ii)).
The relation between the number $`d(\pi _\lambda ,\nu )`$ from Proposition III.2 and $`d(\lambda )`$ is given by $`d(\lambda )=\frac{|\nu ,v_\lambda |^2}{v_\lambda ,v_\lambda }d(\pi _\lambda ,\nu )`$. In particular, if $`\nu `$ is normalized by $`\frac{|\nu ,v_\lambda |^2}{v_\lambda ,v_\lambda }=1`$, then we have $`d(\lambda )=d(\pi _\lambda ,\nu )`$.
Remark III.5. The particular normalization of $`d(\lambda )`$ as in Definition III.4 is motivated from Harish-Chandra’s treatment of the “group case” (cf. \[HC56\]). The group case is defined by $`G=G_0\times G_0`$ and $`H=\mathrm{\Delta }(G)=\{(g,g):gG_0\}`$ for a simply connected hermitian Lie group $`G_0`$. Then we have a natural isomorphism
$$G_0H\backslash G,gH(g,1)$$
and the invariant measure $`\mu _{ZH\backslash G}`$ corresponds to a Haar measure $`\mu _{Z(G_0)\backslash G_0}`$ on $`Z(G_0)\backslash G_0`$.
The spherical unitary highest weight representations of $`G`$ are given by $`(\pi _\lambda \pi _\lambda ^{},_\lambda \widehat{}_\lambda ^{})`$ with $`(\pi _\lambda ,_\lambda )`$ a unitary highest weight representation of $`G_0`$ and $`(\pi _\lambda ^{},_\lambda ^{})`$ its dual representation. Recall that $`_\lambda \widehat{}_\lambda ^{}`$ is isomorphic to the space of Hilbert-Schmidt operators $`B_2(_\lambda )`$ on $`_\lambda `$ and that the corresponding analytic vectors are of trace class, i.e., $`B_2(_\lambda )^\omega B_1(_\lambda )`$ (cf. \[HiKr99a, App.\]). The up to scalar unique $`H`$-fixed hyperfunction vector is given by the conjugate trace:
$$\nu :B_2(_\lambda )^\omega C,A\overline{tr(A)}.$$
Further a highest weight vector for $`(\pi _\lambda \pi _\lambda ^{},_\lambda \widehat{}_\lambda ^{})`$ is given by $`v_\lambda v_\lambda ^{}`$. Then $`\nu ,v_\lambda v_\lambda ^{}=v_\lambda ,v_\lambda `$ and the expression for $`d(\lambda )`$ from Definition III.4 gives that
$$\frac{1}{d(\lambda )}=\frac{1}{|v_\lambda ,v_\lambda |^2}_{Z(G_0)\backslash G_0}|\pi _\lambda (g).v_\lambda ,v_\lambda |^2d\mu _{Z(G_0)\backslash G_0}(Zg).$$
Thus we see that our definition of the formal dimension coincides in the group case with the standard one introduced by Harish-Chandra (cf. \[HC56\]).
Theorem III.6.Let $`(\pi _\lambda ,_\lambda )`$ be an unitary highest weight representations of $`G`$ for which $`(\pi _\lambda ^K,F(\lambda ))`$ is $`HK`$-spherical. Assume that $`\lambda +\rho _\mathrm{\Omega }`$ and that $`(\pi _\lambda ,_\lambda )`$ belongs to the holomorphic discrete series of $`G`$. Then $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical, belongs to the relative discrete series of $`H\backslash G`$ and the formal degree $`d(\lambda )`$ is given by
$$d(\lambda )=d(\lambda )^Gc(\lambda +\rho ),$$
where $`d(\lambda )^G`$ is the formal dimension of $`(\pi _\lambda ,_\lambda )`$ relative to $`G`$.
Proof. Since $`\lambda +\rho _\mathrm{\Omega }`$ the asumptions of Theorem II.16 are satisfied and the theorem applies. Thus $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical and if $`0\nu (_\lambda ^\omega )^H`$ and $`v_\lambda `$ is a highest weight vector, then we have
$$\nu =\frac{\nu ,v_\lambda }{v_\lambda ,v_\lambda c(\lambda +\rho )}_H\pi _\lambda (h).v_\lambda d\mu _H(h).$$
$`(3.2)`$
If we insert (3.2) in the formula for $`\nu `$ in the definition of the formal dimension we obtain that
$$\begin{array}{cc}\hfill \frac{1}{d(\lambda )}& =\frac{1}{|\nu ,v_\lambda |^2}_{HZ\backslash G}|\nu ,\pi _\lambda (g).v_\lambda |^2d\mu _{HZ\backslash G}(HZg)\hfill \\ & =\frac{1}{v_\lambda ,v_\lambda ^2c(\lambda +\rho )^2}_{HZ\backslash G}_H_H\pi _\lambda (h_1).v_\lambda ,\pi _\lambda (g).v_\lambda \hfill \\ & \pi _\lambda (g).v_\lambda ,\pi _\lambda (h_2).v_\lambda d\mu _H(h_1)d\mu _H(h_2)d\mu _{HZ\backslash G}(HZg)\hfill \\ & =\frac{1}{v_\lambda ,v_\lambda ^2c(\lambda +\rho )^2}_{HZ\backslash G}_H_H\pi _\lambda (h_2h_1).v_\lambda ,\pi _\lambda (g).v_\lambda \hfill \\ & \pi _\lambda (h_2^1g).v_\lambda ,v_\lambda d\mu _H(h_1)d\mu _H(h_2)d\mu _{HZ\backslash G}(HZg)\hfill \\ & =\frac{1}{v_\lambda ,v_\lambda ^2c(\lambda +\rho )^2}_{HZ\backslash G}_H_H\pi _\lambda (h_1).v_\lambda ,\pi _\lambda (h_2g).v_\lambda \hfill \\ & \pi _\lambda (h_2g).v_\lambda ,v_\lambda d\mu _H(h_1)d\mu _H(h_2)d\mu _{HZ\backslash G}(HZg)\hfill \\ & =\frac{1}{v_\lambda ,v_\lambda ^2c(\lambda +\rho )^2}_H_{HZ\backslash G}_H\pi _\lambda (h_1).v_\lambda ,\pi _\lambda (h_2g).v_\lambda \hfill \\ & \pi _\lambda (h_2g).v_\lambda ,v_\lambda d\mu _H(h_2)d\mu _{HZ\backslash G}(HZg)d\mu _H(h_1)\hfill \\ & =\frac{1}{v_\lambda ,v_\lambda ^2c(\lambda +\rho )^2}_H_{Z\backslash G}\pi _\lambda (h_1).v_\lambda ,\pi _\lambda (g).v_\lambda \pi _\lambda (g).v_\lambda ,v_\lambda d\mu _{Z\backslash G}(Zg)d\mu _H(h_1).\hfill \end{array}$$
Thus if we apply the Harish-Chandra-Godement Orthogonality Relations for $`L^2(Z\backslash G)`$ and once more (3.2) we obtain that
$$\begin{array}{cc}\hfill \frac{1}{d(\lambda )}& =\frac{1}{d(\lambda )^G}\frac{1}{v_\lambda ,v_\lambda ^2c(\lambda +\rho )^2}v_\lambda ,v_\lambda _H\pi _\lambda (h).v_\lambda ,v_\lambda d\mu _H(h)\hfill \\ & =\frac{1}{d(\lambda )^G}\frac{1}{v_\lambda ,v_\lambda c(\lambda +\rho )^2}c(\lambda +\rho )v_\lambda ,v_\lambda =\frac{1}{d(\lambda )^Gc(\lambda +\rho )},\hfill \end{array}$$
as was to be shown.
IV. Analytic continuation in $`\lambda `$
In this section we prove the analytic continuation of the formula for the formal dimension $`d(\lambda )`$ from Theorem III.6. The proof is quite technical and we need some preparation on algebraic and analytic level.
Algebraic preliminaries
In this subsection we collect some facts concerning the fine structure theory of compactly causal symmetric Lie algebras. The results are mainly due to Ólafsson (cf. \[Ól91\]).
Lemma IV.1.Let $`(g,\tau )`$ be a compactly causal symmetric Lie algebra, then we can choose root vectors $`E_\alpha g_C^\alpha `$, $`\alpha \widehat{\mathrm{\Delta }}_n`$, such that the following conditions are satisfied: (1) $`\overline{E_\alpha }=E_\alpha `$. (2) $`\alpha (H_\alpha )=2`$ with $`H_\alpha =[E_\alpha ,E_\alpha ]`$. (3) $`\tau (E_\alpha )=E_{\tau \alpha }`$, where $`\tau \alpha =\tau \alpha `$.
Proof. Let $`\kappa `$ denote the Cartan-Killing form on $`g_C`$ and define a hermitian inner product on $`g_C`$ by $`X,Y:=\kappa (X,\theta (\overline{Y}))`$.
For each $`\alpha \widehat{\mathrm{\Delta }}_n^+`$ let $`0E_\alpha g_C^\alpha `$ be an arbitrary element of length $`1`$. Then define $`E_\alpha `$ by $`E_\alpha :=\overline{E_\alpha }`$. Thus (1) is satisfied. Now $`\tau (E_\alpha )CE_{\tau \alpha }`$ implies the existence of complex numbers $`c_\alpha `$ such that $`\tau (E_\alpha )=c_\alpha E_{\tau \alpha }`$. Now $`\tau `$ being an involutive implies $`c_\alpha c_{\tau \alpha }=1`$, further $`\tau `$ being an isometry implies that $`|c_\alpha |=1`$ and finally $`\tau `$ being complex linear implies that $`\overline{c_\alpha }=c_\alpha `$ for all $`\alpha \widehat{\mathrm{\Delta }}_n`$. Thus $`c_{\tau \alpha }=\overline{c_\alpha }=c_\alpha `$. For each complex number $`z=e^{i\phi }`$, $`\phi [0,2\pi [`$, of modulus $`1`$ we define $`z^{\frac{1}{2}}=e^{i\frac{\phi }{2}}`$. Thus redifining $`E_\alpha `$, $`\alpha \widehat{\mathrm{\Delta }}_n^+`$, by $`\overline{c_\alpha }^{\frac{1}{2}}E_\alpha `$, leaves (1) untouched and in addition satisfies (3).
Since $`g_C^\alpha p_C`$ for all $`\alpha \widehat{\mathrm{\Delta }}_n^+`$, we have $`\alpha ([E_\alpha ,E_\alpha ])>0`$, and so by rescaling $`E_\alpha `$ with an appropriate positive number we may in addition assume that (2) holds. This proves the lemma.
Let $`\widehat{\mathrm{\Gamma }}=\{\widehat{\gamma }_1,\mathrm{},\widehat{\gamma }_r\}`$ be a maximal system of strongly orthogonal, i.e., $`\widehat{\gamma }_j\pm \widehat{\gamma }_i`$ is never a root and $`\widehat{\mathrm{\Gamma }}\widehat{\mathrm{\Delta }}_n^+`$ has maximal many elements with respect to this property. In view of \[HiÓl96, Lemma 4.1.7\] or \[Ól91, Sect. 3\], we may choose $`\widehat{\mathrm{\Gamma }}`$ invariant under $`\tau `$.
For each $`1jr`$ we set $`\widehat{E}_j:=E_{\widehat{\gamma }_j}`$, $`\widehat{E}_j:=E_{\widehat{\gamma }_j}`$ and $`\widehat{X}_j:=i(\widehat{E}_j\widehat{E}_j)`$. According to \[HC56, Cor. to Lemma 8\], the space
$$e:=\underset{j=1}{\overset{r}{}}R\widehat{X}_j=\underset{j=1}{\overset{r}{}}Ri(\widehat{E}_j\widehat{E}_j)$$
is maximal abelian in $`p`$. Note that $`e`$ is $`\tau `$-invariant by the special choice of the non-compact root vectors (cf. Lemma IV.5(3)) and the $`\tau `$-invariance of $`\widehat{\mathrm{\Gamma }}`$.
We consider the Cayley transform
$$C=e^{i\frac{\pi }{4}ad(_{j=1}^r\widehat{E}_j+\widehat{E}_j)}$$
which is an automorphism of $`g_C`$. Finally we set $`\widehat{H}_j:=H_{\widehat{\gamma }_j}`$ for all $`1jr`$.
Lemma IV.2.The Cayley transform $`C`$ has the following properties:
(i) For all $`1jr`$ one has $`C(\widehat{X}_j)=\widehat{H}_j`$ and $`C(\widehat{H}_j)=\widehat{X}_j`$. (ii) We have $`i\frac{\pi }{4}(_{j=1}^r\widehat{E}_j+\widehat{E}_j)ih_p`$. In particular, one has (a) $`\tau C=C\tau `$, (b) $`\theta C=C^1\theta `$, (iii) The Cayley transform yields an isomorphism $`C:eC(e)`$ with $`C(e)it`$ a $`\tau `$-invariant subspace.
Proof. (i) This follows from $`sl(2,R)`$-reduction (cf. \[HC56, p. 584\], \[HiÓl96, Lemma A.3.2(3)\]).
(ii) It follows from $`g_C^{\widehat{\alpha }}p_C`$, for all $`\widehat{\alpha }\widehat{\mathrm{\Delta }}_n`$ and Lemma IV.1(1) that $`i\frac{\pi }{4}(_{j=1}^r\widehat{E}_j+\widehat{E}_j)ip`$. Further Lemma IV.1(3) and the $`\tau `$-invariance of $`\widehat{\mathrm{\Gamma }}`$ imply
$$\tau \left(\underset{j=1}{\overset{r}{}}\widehat{E}_j+\widehat{E}_j\right)=\underset{j=1}{\overset{r}{}}\tau (\widehat{E}_j)+\tau (\widehat{E}_j)=\underset{j=1}{\overset{r}{}}E_{\tau \widehat{\gamma }_j}+E_{\tau \widehat{\gamma }_j}=\underset{j=1}{\overset{r}{}}\widehat{E}_j+\widehat{E}_j.$$
Thus $`i\frac{\pi }{4}(_{j=1}^r\widehat{E}_j+\widehat{E}_j)ih_p`$. This proves (i).
(iii) This follows from (i) and (ii)(a).
Recall that $`e`$ is $`\tau `$-invariant and write $`b=eq`$ for the set of $`\tau `$-fixed points.
Lemma IV.3.Let $`c:=C(b)`$. Then $`ca`$ and the Cayley transform yields an isomorphism $`C:bc`$.
Proof. Since $`C(b)it`$ by Lemma IV.2(i), the fact that $`bq`$ and that $`C`$ commutes with $`\tau `$ (cf. Lemma IV.2(ii)) imply that $`C(b)i(tq)`$. But $`i(tq)=a`$ by the definition of $`a`$, proving the lemma.
Recall that $`b`$ is maximal abelian subspace of $`qp`$ (this follows for instance from the $`c`$-dual version of Lemma 4.1.9 in \[HiÓl96\]) and denote by $`\mathrm{\Sigma }=\mathrm{\Sigma }(g,b)`$ the set of roots of $`g`$ with respect to $`b`$. Recall that $`\mathrm{\Sigma }`$ is an abstract root system (cf. \[Sch84, Sect. 7.2\]). We write
$$g=z_g(b)\underset{\phi \mathrm{\Sigma }}{}g^\phi $$
for the corresponding root space decomposition. By Lemma IV.3, the Cayley transform induces a mapping $`C^t:a^{}b^{},\alpha \alpha C_b`$ and we set
$$\mathrm{\Sigma }_n=C^t(\mathrm{\Delta }_n)_b\text{and}\mathrm{\Sigma }_k=C^t(\mathrm{\Delta }_k)_b\backslash \{0\}.$$
Let $`\mathrm{\Gamma }=\{\frac{1}{2}(\widehat{\gamma }_j\tau \widehat{\gamma }_j):1jr\}`$ denote the restricted set of strongly orthogonal roots. Note that $`\mathrm{\Gamma }c^{}`$ by Lemma IV.2(i). Thus we can write $`\mathrm{\Gamma }=\{\gamma _1,\mathrm{},\gamma _s\}`$ for some $`1sr`$. For each $`1js`$ we define $`H_jc`$ by $`\gamma _j(H_j)=2`$ and $`\gamma _k(H_j)=0`$ for $`kj`$. We set $`X_j:=C(H_j)`$ for all $`1js`$. Then
$$b=\underset{j=1}{\overset{s}{}}RX_j.$$
As a final algebraic tool we need explicit information on the root system $`\mathrm{\Sigma }`$ which is provided by Ólafsson’s Theorem on double restricted root systems (cf. \[Ól91, Sect. 3\], \[HÓØ91, Prop. 3.1\]). For all $`1js`$ we set $`\psi _j:=C^t(\gamma _j)`$ and note that $`\psi _j(X_j)=2`$ since $`C(X_j)=H_j`$ (cf. Lemma IV.2(i), (ii)).
Finally we put $`\mathrm{\Sigma }^+:=C^t(\mathrm{\Delta }^+)_b\backslash \{0\}`$, $`\mathrm{\Sigma }_n^+:=\mathrm{\Sigma }_n\mathrm{\Sigma }^+`$ and $`\mathrm{\Sigma }_k^+:=\mathrm{\Sigma }_k\mathrm{\Sigma }^+`$.
Theorem IV.4. (Ólafsson) If $`(g,\tau )`$ is compactly causal, then the following assertions concerning the double restricted root system $`\mathrm{\Sigma }=\mathrm{\Sigma }(g,b)`$ hold: (i) The restricted root system has the following form
$$\mathrm{\Sigma }_k=\pm \{\frac{1}{2}(\psi _i\psi _j):i<j\}\pm \{\frac{1}{2}\psi _j:1js\}$$
and
$$\mathrm{\Sigma }_n^+=\{\frac{1}{2}(\psi _i+\psi _j):1i,js\}\{\frac{1}{2}\psi _j:1js\}.$$
The second sets in $`\mathrm{\Sigma }_k`$ and $`\mathrm{\Sigma }_n^+`$ are empty if and only if $`C^4=id`$. If further $`\psi _s`$ is chosen to be a simple root, then
$$\mathrm{\Sigma }_k^+\{\frac{1}{2}(\psi _i\psi _j):i<j\}\{\frac{1}{2}\psi _j:1js\}.$$
(ii) All $`\psi _j`$, $`1js`$, have the same length. (iii) The conjugacy classes of the restricted root system under the Weyl group associated to $`\mathrm{\Sigma }`$ are given by (1) $`\{\pm \frac{1}{2}(\psi _i\pm \psi _j):1i,js,ij\}`$ (2) $`\{\pm \psi _j:1js\}`$ (3) $`\{\pm \frac{1}{2}\psi _j:1js\}`$
Proof. (i) Let $`\widehat{\mathrm{\Sigma }}=\widehat{\mathrm{\Sigma }}(g,e)`$ be the restricted root system with respect to the maximal abelian subspace $`e`$ and $`\widehat{\mathrm{\Sigma }}_k`$, $`\widehat{\mathrm{\Sigma }}_n`$ defined as above. Write $`\widehat{\psi }_j:=C^t(\widehat{\gamma }_j)`$ for all $`1jr`$. Suppose first that $`g`$ is simple. Then for the analogous statement for $`\widehat{\mathrm{\Sigma }}`$ in stead of $`\mathrm{\Sigma }`$ and $`\widehat{\psi }_j`$ in stead of $`\psi _j`$, Harish-Chandra has proved in \[HC56, Lemma 13-16\] that $`\widehat{\mathrm{\Sigma }}_k,\widehat{\mathrm{\Sigma }}_n^+`$ are contained in the asserted subsets, Moore proved equality (cf. \[Mo64, Th. 2\]) and finally Korányi and Wolf have shown in \[KoWo65, Prop. 4.4 with Remark\] that the second set in $`\widehat{\mathrm{\Sigma }}_n^+`$ is empty if and only if $`C^4=id`$. Now taking restrictions to $`c`$ yields (i) for $`g`$ simple.
In the group case similar considerations lead to the same result.
(ii) This can be deduced from \[Mo64, Th. 2(2)\], but we propose here a much simpler proof. We use (i) and the fact that $`\mathrm{\Sigma }`$ is an abstract root system. As usual we write $`s_\psi `$, $`\psi \mathrm{\Sigma }`$, for the reflection associated to $`\psi `$. Then we obtain for all $`1ijs`$ that
$$\begin{array}{cc}\hfill s_{\frac{1}{2}(\psi _i+\psi _j)}(\psi _j)& =\psi _j\frac{2\psi _j,\frac{1}{2}(\psi _i+\psi _j)}{\frac{1}{2}(\psi _i+\psi _j),\frac{1}{2}(\psi _i+\psi _j)}\frac{1}{2}(\psi _i+\psi _j)\hfill \\ & =\psi _j\frac{2\psi _j,\psi _j}{\psi _i,\psi _i+\psi _j,\psi _j}(\psi _i+\psi _j).\hfill \end{array}$$
Thus it follows from (i) and $`s_{\frac{1}{2}(\psi _i+\psi _j)}(\psi _j)\mathrm{\Sigma }`$ that $`\frac{\psi _j,\psi _j}{\psi _i,\psi _i+\psi _j,\psi _j}\{\frac{1}{2},\frac{1}{4}\}`$. Interchanging $`i`$ and $`j`$ then yields $`\frac{\psi _j,\psi _j}{\psi _i,\psi _i+\psi _j,\psi _j}=\frac{1}{2}`$ or equivalently that $`\psi _j,\psi _j=\psi _i,\psi _i`$. This proves (ii).
(iii) In view of (i), we have for all $`1i,j,kr`$ that
$$\begin{array}{cc}& s_{\frac{1}{2}(\psi _i\pm \psi _j)}(\psi _j)=\psi _i,\hfill \\ & s_{\frac{1}{2}(\psi _i\pm \psi _j)}\left(\frac{1}{2}(\psi _j\pm \psi _k)\right)=\frac{1}{2}(\psi _i\pm \psi _k),\hfill \\ & s_{\psi _i}\left(\frac{1}{2}(\psi _i\pm \psi _j)\right)=\frac{1}{2}(\psi _i\pm \psi _j).\hfill \end{array}$$
$`(4.1)`$
This proves (iii).
From now on we assume that $`\psi _s`$ is a simple root. Then Theorem IV.4(i) says that
$$\mathrm{\Sigma }_n^+=\{\frac{1}{2}(\psi _i+\psi _j):1i,js\}\{\frac{1}{2}\psi _j:1js\}.$$
$`(4.2)`$
and
$$\mathrm{\Sigma }_k^+=\{\frac{1}{2}(\psi _i\psi _j):1i<js\}\{\frac{1}{2}\psi _j:1js\}.$$
$`(4.3)`$
Further it follows from Theorem IV.4(i) and the first formula in (4.1) that the Weyl group $`𝒲(\mathrm{\Sigma }_k)`$ of $`\mathrm{\Sigma }_k`$ acts on $`b`$ as the full permutation group of the $`X_j`$’s.
We write $`b^+=\{Xb:(\phi \mathrm{\Sigma }^+)\phi (X)0\}`$ for the Weyl chamber corresponding to $`\mathrm{\Sigma }^+`$. By (4.2) and (4.3) we then have
$$b^+=\{\underset{j=1}{\overset{s}{}}x_jX_j:0x_s\mathrm{}x_1\}.$$
Further let $`a^+:=\{Xa:(\alpha \mathrm{\Delta }^+)\alpha (X)0\}`$ and $`c^+:=a^+c`$. Note that $`C(b^+)=c^+`$ by the construction of $`\mathrm{\Sigma }^+`$.
Lemma IV.5.The following equality holds
$$C_{\mathrm{min}}^{}(\stackrel{ˇ}{C}_k^{})=(c^+)^{}(\stackrel{ˇ}{C}_k^{}),$$
where the stars $``$ are all taken in $`a^{}`$.
Proof. First recall some basic rules in dealing with convex cones (cf. \[Ne99b, Ch. V\]). If $`W`$ is a closed convex cone in an euclidean space $`V`$, then $`(W^{})^{}=W`$. Further for two closed convex cones $`W_1,W_2V`$ we have $`(W_1W_2)^{}=\overline{W_1^{}+W_2^{}}`$.
Let now the convex cone on the left hand side be denoted by $`W_1`$, the other one by $`W_2`$. Let $`p:ac`$ be the orthogonal projection with respect to the Cartan-Killing form. We claim that $`p(W_1^{})=p(W_2^{})`$. Assume first that no half roots in $`\mathrm{\Sigma }`$ occur. Then from the Cayley-transform analogs of (4.2) and (4.3) it follows that
$$p(W_1^{})=p(\overline{C_{\mathrm{min}}\stackrel{ˇ}{C}_k})=\underset{j=1}{\overset{s}{}}R^+H_j+\underset{j=1}{\overset{s1}{}}R^+(H_{j+1}H_j),$$
and
$$p(W_2^{})=p(\overline{c^+\stackrel{ˇ}{C}_k})=\left(\left(\underset{j=1}{\overset{s}{}}R^+H_j\right)\{\underset{j=1}{\overset{s}{}}x_jH_j:x_s\mathrm{}x_1\}\right)+\underset{j=1}{\overset{s1}{}}R^+(H_{j+1}H_j).$$
From these two equalities the claim follows in the case of no half roots in $`\mathrm{\Sigma }`$. The general case is easily deduced from this.
Let $`r:a^{}c^{},r(\lambda ):=\lambda _c`$ be the restriction map and note that $`r`$ is the dual map of the inclusion mapping $`ca`$. Since both $`W_1`$ and $`W_2`$ are closed, we have $`(W_{1,2}^{})^{}=W_{1,2}`$, and so
$$W_{1,2}_c=r(W_{1,2})=(p(W_{1,2}^{}))^{}.$$
Hence our claim implies that $`W_1_c=W_2_c`$. Thus $`W_1W_2`$ by the definition of $`W_1`$ and $`W_2`$.
For the converse inclusion we first note that an element $`\lambda \stackrel{ˇ}{C}_k^{}`$ belongs to $`W_1`$ if and only if $`\lambda (\stackrel{ˇ}{\beta })0`$, where $`\beta `$ is the highest root (this becomes clear from our construction of the positive systems). Recall that $`\widehat{\mathrm{\Gamma }}`$ can be constructed inductively starting with the highest root (cf. \[HC56, p. 108\]). Thus $`\beta =\gamma _1in\mathrm{\Gamma }`$. Hence $`W_1=(\gamma _1)^{}\stackrel{ˇ}{C}_k^{}`$, and so $`W_2W_1`$ since $`(c^+)^{}(\gamma _1)^{}`$.
Analytic preliminaries
Recall the definition of $`b^+`$ and set $`B^+:=\mathrm{exp}(b^+)`$.
Lemma IV.6. (Flensted-Jensen) Let $`L=Z_{HK}(b)`$. Then for the homegeneous space $`HZ\backslash G`$ the following assertions hold:
(i) The subgroups $`HZ`$ and $`LZ`$ of $`G`$ are closed and $`Z\backslash LZ`$ is compact. (ii) The mapping
$$\mathrm{\Phi }:B^+\times LZ\backslash KHZ\backslash G,(b,LZk)LZbk$$
is a diffeomorphism onto its open image. The image is dense with complement of Haar measure zero. (iii) Up to normalization of measures we have for all $`fL^1(HZ\backslash G)`$ the following integration formula
$$_{HZ\backslash G}f(HZg)𝑑\mu _{HZ\backslash G}(HZg)=_{Z\backslash K}_{b^+}f(HZ\mathrm{exp}(X)k)J(X)𝑑X𝑑\mu _{Z\backslash K}(Zk)$$
with
$$J(X)=\underset{\phi \mathrm{\Sigma }^+}{}\mathrm{cosh}(\phi (X))^{m_\phi ^+}\mathrm{sinh}(\phi (X))^{m_\phi ^{}},$$
where $`m_\phi ^\pm :=dim(\{Xg^\phi :\theta \tau (X)=\pm X\})`$.
Proof. (i) The closedness of $`HZ`$ and $`LZ`$ follows from the closedness of $`Ad(H)`$ and $`Z_{Ad(H)}(b)`$ in the adjoint group $`Ad(G)`$. Finally $`Z\backslash LZ`$ is a closed subgroup of the compact group $`Z\backslash Z(HK)`$ and hence compact.
(ii) \[Sch84, Prop. 7.1.3\].
(iii) It follows from \[FJ80, Th. 2.6\] or \[Sch84, Lemma 8.1.2\] that
$$J(X):=det\left(d\mathrm{\Phi }(X,LZk)\right)=\underset{\phi \mathrm{\Sigma }^+}{}\mathrm{cosh}(\phi (X))^{m_\phi ^+}\mathrm{sinh}(\phi (X))^{m_\phi ^{}}$$
for all $`Xb^+`$ and $`kK`$. Thus it follows from (ii) that
$$_{HZ\backslash G}f(HZg)𝑑\mu _{HZ\backslash G}(HZg)=_{LZ\backslash K}_{b^+}f(HZ\mathrm{exp}(X)k)J(X)𝑑X𝑑\mu _{LZ\backslash K}(Zk)$$
holds for all $`fL^1(HZ\backslash G)`$. In view of (i), we may replace the integration over $`LZ\backslash K`$ by an $`Z\backslash K`$-integral, proving (iii).
Lemma IV.7.Realize $`G`$ as a submanifold of $`\stackrel{~}{M}\times P^+`$ as in Proposition II.10(ii). Then for $`b=\mathrm{exp}_G(_{j=1}^sx_jX_j)B`$ and
$$\mu (b):=\mathrm{exp}_{\stackrel{~}{K_C}}\left(\underset{j=1}{\overset{s}{}}\frac{1}{2}\mathrm{log}\mathrm{cosh}(2x_j)H_j\right)A\stackrel{~}{K_C}$$
the following assertions hold: (i) We have $`b\{[h,\mu (b)]:h\stackrel{~}{H_C}\}\times P^+`$. (ii) If $`Xb^+`$, then $`\mathrm{log}\mu (\mathrm{exp}_G(X))c^+`$.
Proof. (i) This follows directly from \[HiÓl96, p. 210-211\].
(ii) Recall that $`X=_{j=1}^sx_jX_jb^+`$ if and only if $`0x_s\mathrm{}x_1`$. Now the assertion follows from (i) and the monotonicity of the mapping $`R^+R,x\mathrm{log}\mathrm{cosh}(x)`$.
Proof of the analytic continuation
Let $`(\pi _\lambda ,_\lambda )`$ be an $`H`$-spherical unitary highest weight representation of $`G`$. Further let $`\nu (_\lambda ^\omega )^H`$ an $`H`$-fixed hyperfunction vector and $`\nu _0=\nu _{F(\lambda )}F(\lambda )^{HK}`$. We normalize $`\nu `$ by setting $`\nu _0=1`$ and then $`v_\lambda `$ by $`|\nu ,v_\lambda |=1`$. Then we have
$$d(\lambda )=I(\lambda )^1\text{with}I(\lambda ):=_{HZ\backslash G}|\nu ,\pi _\lambda (g).v_\lambda |^2d\mu _{HZ\backslash G}(HZg).$$
Definition IV.8. On the non-compactly Riemannian symmetric space $`K(0)\backslash G(0)`$ we define the spherical function with parameter $`\lambda a_C^{}`$ by
$$\phi _\lambda ^0(g)=_{K(0)}a_H(gk)^{\lambda \rho _k}𝑑\mu _{K(0)}(k),$$
for all $`gG(0)`$.
Remark IV.9. Note that if $`\lambda a^{}`$ is the highest weight of an $`HK`$-spherical representation $`(\pi _\lambda ^K,F(\lambda ))`$ of $`\stackrel{~}{K_C}`$, then $`\phi _{\lambda +\rho _k}^0`$ extends to a holomorphic function on $`\stackrel{~}{K_C}`$ and we have
$$\phi _{\lambda +\rho _k}^0(k)=\pi _\lambda ^K(k).\nu _0,\nu _0$$
$`(4.4)`$
for all $`k\stackrel{~}{K_C}`$ (cf. \[Hel84, Ch. V, Th. 4.3\]).
Proposition IV.10.With the notation of Lemma IV.7 we have
$$I(\lambda )=\frac{1}{dimF(\lambda )}_{b^+}\phi _{\lambda +\rho _k}^0(\mu (\mathrm{exp}_G(X))^2)J(X)𝑑X$$
where $`J(X)`$ is given as in Proposition IV.6(iii).
Proof. (cf. \[HC56, p. 599\], \[Gr96, Prop. 10\]) In the sequel we identify $`b`$ with $`B`$ via the exponential mapping and for $`b=\mathrm{exp}_G(X)B^+`$ we set $`J(b):=J(X)`$. Then by Lemma IV.6(iii) we have
$$\begin{array}{cc}\hfill I(\lambda )& =_{HZ\backslash G}|\nu ,\pi _\lambda (g).v_\lambda |^2d\mu _{HZ\backslash G}(HZg)\hfill \\ & =_{Z\backslash K}_{B^+}|\nu ,\pi _\lambda (bk).v_\lambda |^2J(b)d\mu _B(b)d\mu _{Z\backslash K}(k).\hfill \end{array}$$
$`(4.5)`$
In view of Lemma II.10(ii), we can write each element in $`bB^+`$ as $`([h_C(b),\mu (b)],p_+(b))\stackrel{~}{M}\times P^+`$ with $`\mu (b)\stackrel{~}{K_C}`$. Now the same consideration as in the proof of Step 1 of Theorem II.16 yields for all $`bB^+`$ and $`kK`$ that
$$\begin{array}{cc}\hfill \nu ,\pi _\lambda (bk).v_\lambda & =\nu ,\pi _\lambda \left(([h_C(b),\mu (b)],p_+(b))k\right).v_\lambda \hfill \\ & =\nu ,\pi _\lambda ([\mathrm{𝟏},\mu (b)k],k^1p_+(b)k).v_\lambda =\nu ,\pi _\lambda (\mu (b)k).v_\lambda \hfill \\ & =\nu _0,\pi _\lambda ^K(\mu (b)k).v_\lambda .\hfill \end{array}$$
If we insert this expression for the matrix coefficient in (4.5), use Schur’s Orthogonality Relations for $`(\pi _\lambda ^K,F(\lambda ))`$ and the relation $`\pi _\lambda ^K(\mu (b))^{}=\pi _\lambda ^K(\mu (b))`$ (cf. Lemma IV.7), we arrive at
$$\begin{array}{cc}\hfill I(\lambda )& =_{B^+}_{Z\backslash K}|\nu _0,\pi _\lambda ^K(\mu (b)k).v_\lambda |^2J(b)d\mu _{Z\backslash K}(k)d\mu _B(b)\hfill \\ & =\frac{1}{dimF(\lambda )}_{B^+}\pi _\lambda ^K(\mu (b)).\nu _0,\pi _\lambda ^K(\mu (b)).\nu _0J(b)d\mu _B(b)\hfill \\ & =\frac{1}{dimF(\lambda )}_{B^+}\pi _\lambda ^K(\mu (b)^2).\nu _0,\nu _0J(b)d\mu _B(b).\hfill \end{array}$$
Now the assertion of the proposition follows from (4.4).
Lemma IV.11.Let $`V`$ be a finite dimensional real vector space, $`WV`$ be an open convex cone, $`\alpha _1,\mathrm{},\alpha _n,\beta _1,\mathrm{},\beta _mW^{}\backslash \{0\}`$ and $`p_1,\mathrm{},p_n,q_1\mathrm{},q_mN_0`$. For every $`\lambda V^{}`$ we define the integral
$$H(\lambda ):=_We^{\lambda (x)}\underset{j=1}{\overset{n}{}}(\mathrm{sinh}\alpha _j(x))^{p_j}\underset{j=1}{\overset{m}{}}\left(\mathrm{cosh}\beta _j(x)\right)^{q_j}d\mu _V(x).$$
Then $`H(\lambda )`$ converges if and only if $`\lambda +_{j=1}^np_j\alpha _j+_{j=1}^mq_j\beta _jintW^{}`$.
Proof. If $`q_1=\mathrm{}=q_m=0`$, then this is Lemma IV.6 in \[Kr98\]. The general case is easily obtained from this.
The following characterization of the relative discrete series by the parameter $`\lambda `$ is due to Hilgert, Ólafsson and Ørsted and was obtained in two steps (cf. \[ÓØ91, Th. 5.2\], \[HÓØ91, Th. 3.3\]). We present an essentially modified proof here, but we point out that it is not our objective to give new proofs of well-known facts. In the course of our arguments we obtain an important new estimate which is crucial for the analytic continuation of $`I(\lambda )`$.
Theorem IV.12. (Hilgert-Ólafsson-Ørsted) Let $`(\pi _\lambda ,_\lambda )`$ be an unitary highest weight representation of $`G`$ with $`(\pi _\lambda ^K,F(\lambda ))`$ being $`HK`$-spherical. Then $`(\pi _\lambda ,_\lambda )`$ belongs to the relative discrete series of $`H\backslash G`$ if and only if the condition
$$(\alpha \mathrm{\Delta }_n^+)\lambda +\rho ,\alpha <0$$
$`(\mathrm{RDS})`$
is satisfied.
Proof. Recall the definition of $`c^+`$, $`a^+`$ and the relation $`C(b^+)=c^+`$. Set $`A^+:=\mathrm{exp}_{G^c}(a^+)`$ and let $``$ denote an arbitrary norm on $`a`$. If we write $`(c^+)^{}`$, then the star $``$ is to be taken in $`a^{}`$.
Step 1: $`I(\lambda )<\mathrm{}`$, if $`\lambda +\rho int(c^+)^{}`$, the interior of $`(c^+)^{}`$.
Here we do not assume that $`\lambda a^{}`$ is dominant integral with respect to $`\mathrm{\Delta }_k^+`$, but only $`\lambda \stackrel{ˇ}{C}_k^{}`$. By Harish-Chandra’s estimates for spherical functions on non-compact Riemannian symmetric spaces, there exists constants $`c>0`$ and $`dN`$ such that
$$(\lambda \stackrel{ˇ}{C}_k^{})(aA^+)\phi _\lambda ^0(a)ca^{\lambda \rho _k}(1+\mathrm{log}a)^d$$
$`(4.6)`$
(cf. \[Wal88, 4.5.3\]). Note that $`J(X)e^{2\rho (C(X))}`$ for all $`Xb^+`$ by the formula for the Jacobian in Lemma IV.6(iii). Thus it follows for all $`\lambda \stackrel{ˇ}{C}_k^{}`$ and $`X=_{j=1}^sx_jX_jb^+`$ from (4.6) together with Lemma IV.7 that
$$\begin{array}{cc}\hfill \phi _{\lambda +\rho _k}^0(\mu (\mathrm{exp}_G(X))^2)J(X)& c\mu (\mathrm{exp}_G(X))^{2\lambda }(1+\mathrm{log}\mu (\mathrm{exp}_G(X))^2)^de^{2\rho (C(X))}\hfill \\ & ce^{2\lambda (C(X))}(1+2C(X))^de^{2\rho (C(X))}\hfill \\ & ce^{2(\lambda +\rho )(C(X))}(1+2C(X))^d.\hfill \end{array}$$
$`(4.7)`$
Now Proposition IV.11 shows that $`I(\lambda )<\mathrm{}`$ if $`\lambda +\rho int(c^+)^{}`$, proving our first step.
Step 2: $`\lambda +\rho int(c^+)^{}`$, if $`I(\lambda )<\mathrm{}`$.
Recall that $`\lambda `$ is supposed to be dominant integral with respect to $`\mathrm{\Delta }_k^+`$. Thus it follows from (4.4) and the fact that the $`HK`$-spherical vector $`\nu _0`$ has a non-zero $`v_\lambda `$-component (cf. \[Hel84, p. 537, (7)\]) that there is a constant $`c_\lambda >0`$ such that $`c_\lambda a^\lambda \phi _{\lambda +\rho _k}^0(a)`$ holds for all $`aA^+`$. Hence Lemma IV.7 implies that
$$(Xb^+)\frac{c_\lambda }{2}e^{2\lambda (C(X))}J(X)\phi _{\lambda +\rho _k}^0(\mu (\mathrm{exp}_G(X))^2)J(X).$$
In view of Proposition IV.10 and Lemma IV.11, we obtain $`\lambda +\rho int(c^+)^{}`$ if $`I(\lambda )<\mathrm{}`$. This proves our second step.
Step 3: If $`\lambda \stackrel{ˇ}{C}_k^{}`$, then $`\lambda `$ satisfies (RDS) if and only if $`\lambda +\rho int(c^+)^{}`$.
Note that $`\lambda `$ satisfies (RDS) means that $`\lambda +\rho intC_{\mathrm{min}}^{}`$. Now if $`\lambda \stackrel{ˇ}{C}_k^{}`$, then $`\lambda +\rho int\stackrel{ˇ}{C}_k^{}`$. Thus Step 3 follows from Lemma IV.5.
In view of Steps 1-3, it follows that $`I(\lambda )`$ is finite if and only if $`\lambda `$ satisfies the condition (RDS). The proof of the theorem will therefore be complete with
Step 4: If $`\lambda `$ satisfies (RDS), then $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical.
Let $`\kappa :G\stackrel{~}{K_C}/(\stackrel{~}{K_C}\stackrel{~}{H_C})_0`$ the canonical projection defined via the decomposition in Proposition II.10. Now the function
$$H\backslash GC,Hg\pi _\lambda ^K(\kappa (g)).v_\lambda ,\nu _0$$
generates an $`H`$-spherical module in the relative discrete series on $`H\backslash G`$ (cf. \[ÓØ91, Th. 5.2\]). This proves Step 4 and concludes the proof of the theorem.
The prescription
$$W:=intC_{\mathrm{min}}^{}\stackrel{ˇ}{C}_k^{}int(c^+)^{}$$
defines a convex cone in $`a^{}`$. We write $`T_W=ia^{}+W`$ for the associated tube domain in $`a_C^{}`$. Note that $`\rho _niz(k)^{}`$ by the construction of $`\mathrm{\Delta }_n^+`$ and so $`\rho _nW`$.
Lemma IV.13.The function $`I(\lambda )`$ extends naturally to a continuous function on the affine subtube $`T_W\rho `$, also denoted by $`I`$, and which is holomorphic when restricted to $`T_{W^0}\rho `$. If $`mN`$ is sufficiently large, then $`Wm\rho _nW\rho `$ and $`I_{T_Wm\rho _n}`$ is bounded.
Proof. First we show that $`Wm\rho _nW\rho `$ for large values of $`mN`$. Since $`\rho _nintC_{\mathrm{min}}^{}`$, we have $`\rho m\rho _nintC_{\mathrm{min}}^{}`$ provided $`mN`$ is sufficiently large. Further $`\rho _niz(k)^{}`$ shows that $`R.\rho _n\stackrel{ˇ}{C}_k^{}`$. Thus we have $`\rho m\rho _nW`$ if $`m`$ is chosen sufficiently large, proving our claim.
Recall the formula for $`I(\lambda )`$ from Proposition IV.10. Then (4.7) yields constants $`c>0`$, $`dN`$ such that
$$I(\lambda )\frac{c}{dimF(\lambda )}_{c^+}e^{2(\lambda +\rho )(X)}(1+2X)^d𝑑X$$
$`(4.8)`$
holds for some norm $``$ on $`a`$. Let $`\widehat{\rho }_k`$ denote the half sum of the roots in $`\widehat{\mathrm{\Delta }}_k^+`$ and recall Weyl’s Dimension Formula
$$dimF(\lambda )=\frac{\underset{\widehat{\alpha }\widehat{\mathrm{\Delta }}_k^+}{}\lambda +\widehat{\rho }_k,\widehat{\alpha }}{_{\widehat{\alpha }\widehat{\mathrm{\Delta }}_k^+}\widehat{\rho }_k,\widehat{\alpha }}.$$
In particular, we see that $`\lambda \frac{1}{dimF(\lambda )}`$ extends to a holomorphic map on $`T_W`$ and $`T_W\rho `$ which is bounded when restricted to $`T_Wm\rho _n`$ for all $`mN_0`$. Further for each fixed $`bB^+`$ the mapping
$$a_C^{}C,\lambda \phi _{\lambda +\rho _k}^0(\mu (b)^2)$$
is holomorphic. Now (4.8) together with Proposition IV.10 imply that $`I(\lambda )`$ extends to a continuous function on $`T_W\rho `$ which is holomorphic on $`T_{W^0}\rho `$ and bounded when restricted to $`T_Wm\rho _n`$ provided $`m`$ is chosen sufficiently large.
Lemma IV.14.If $`mN`$ is sufficiently large, then the function
$$T_{W^0}m\rho _nC,\lambda c(\lambda +\rho )$$
is holomorphic and bounded.
Proof. In view of $`\rho _niz(k)^{}`$, this is immediate from Theorem II.14.
Theorem IV.15. (The formal dimension for the relative holomorphic discrete series on a compactly causal symmetric space) Let $`H\backslash G`$ be a simply connected symmetric space associated to a compactly causal symmetric Lie algebra $`(g,\tau )`$ and $`(\pi _\lambda ,_\lambda )`$ be an unitary highest weight representations of $`G`$ for which $`F(\lambda )`$ is $`HK`$-spherical. Then the following assertions hold:
(i) The representation $`(\pi _\lambda ,_\lambda )`$ belongs to the relative discrete series for $`H\backslash G`$ if and only if the condition
$$(\alpha \mathrm{\Delta }_n^+)\lambda +\rho ,\alpha <0$$
$`(\mathrm{RDS})`$
is satisfied. (ii) If $`(\pi _\lambda ,_\lambda )`$ belongs to the relative discrete series of $`H\backslash G`$, then the formal dimension $`d(\lambda )`$ is given by
$$d(\lambda )=d(\lambda )^Gc(\lambda +\rho ),$$
where $`d(\lambda )^G`$ is the formal dimension of $`(\pi _\lambda ,_\lambda )`$ relative to $`G`$ and $`c`$ is the $`c`$-function of the non-compactly $`c`$-dual space $`H^c\backslash G^c`$ of $`H\backslash G`$ (cf. Theorem II.14). Here the right hand side has to be understood as an analytic continuation of a product of two meromorphic functions.
Proof. (i) Theorem IV.12.
(ii) Let $`\widehat{\rho }`$ denote the half sum of the elements in $`\widehat{\mathrm{\Delta }}^+`$ and recall Harish-Chandra’s condition for the relative discrete series on $`G`$
$$(\widehat{\alpha }\widehat{\mathrm{\Delta }}_n^+)\lambda +\widehat{\rho },\widehat{\alpha }<0$$
(cf. \[HC56, Lemma 29\]) as well as Harish-Chandra’s formula for the formal dimension $`d(\lambda )^G`$ of the relative discrete series on $`G`$
$$d(\lambda )^G=\frac{\underset{\widehat{\alpha }\widehat{\mathrm{\Delta }}^+}{}\lambda +\widehat{\rho },\widehat{\alpha }}{_{\widehat{\alpha }\widehat{\mathrm{\Delta }}^+}\widehat{\rho },\widehat{\alpha }}$$
(cf. \[HC56, Th. 4\]). In particular for $`mN`$ sufficiently large, the prescription $`\lambda \frac{1}{d(\lambda )^G}`$ defines a bounded holomorphic function on the affine tube $`T_{W^0}m\rho _n`$.
Now it follows from Lemma IV.13 and Lemma IV.14 that the function
$$f:T_{W^0}m\rho _nC,\lambda I(\lambda )c(\lambda +\rho )\frac{1}{d(\lambda )^G}$$
is holomorphic and bounded for $`m`$ sufficiently large. For such $`m`$ Theorem III.6 implies that $`f(\lambda )=0`$ for all $`\lambda W^0m\rho _n`$ which are dominant integral with respect to $`\mathrm{\Delta }_k^+`$. Thus the identity criterion of Proposition A.2 in Appendix A applies and yields $`f=0`$. We conclude in particular that $`I(\lambda )^1`$ defines a continuation of $`\lambda d(\lambda )^Gc(\lambda +\rho )`$ to a continuous function on $`T_W\rho `$ which is holomorphic when restricted to the interior $`T_{W^0}\rho `$. Since by definition $`d(\lambda )=I(\lambda )^1`$, the assertion in (ii) follows because $`\lambda `$ satisfies (RDS) if and only if $`\lambda T_W\rho `$.
The following result has already been obtained earlier by Faraut, Hilgert and Ólafsson in \[FHÓ94, Lemma 9.2\], but with a completely different type of arguments (see also Theorem II.14).
Corollary IV.16.Suppose that $`(g,\tau )=(hh,\sigma )`$ is of group type (cf. Lemma I.3(i)(2)). Then the domain of convergence $``$ for $`c`$ is given by
$$=ia^{}+(intC_{\mathrm{min}}^{})int\stackrel{ˇ}{C}_k^{}$$
and there exists a constant $`\gamma >0`$ only depending on the choice of the various Haar measures such that
$$c(\lambda )=\gamma \frac{1}{_{\alpha \mathrm{\Delta }^+}\lambda ,\alpha }$$
for $`\lambda `$.
Proof. In the following we use the notation of Remark III.5. Since $`(g,\tau )`$ is of group type we have $`d(\lambda )^G=d(\lambda )^{(G_0\times G_0)}=(d(\lambda )^{G_0})^2`$, and so it follows from Theorem IV.15(ii) that $`c(\lambda +\rho )=\frac{1}{d(\lambda )^{G_0}}`$ holds for the analytic continuations. In view of Harish-Chandra’s formula for $`d(\lambda )^{G_0}`$ (cf. \[HC56, Th. 4\]), this proves the corollary.
Problems . The discrete series on $`H\backslash G`$ are constructed by analytic methods, i.e., with generating functions (cf. \[FJ80\], \[MaOs84\], \[ÓØ91\]). But from the algebraic point of view there are still some interesting open problems.
(a) Using the classification sheme of unitary highest weight modules (cf. \[EHW83\]) together with the fine structure theory of compactly causal symmetric Lie algebras provided by Theorem IV.4 and \[Ól91\] one can check case by case that (RDS) implies that $`N(\lambda )=L(\lambda )`$. In view of Proposition II.7(ii), this gives a more algebraic proof of the fact that (RDS) implies that $`(\pi _\lambda ,_\lambda )`$ is $`H`$-spherical whenever $`(\pi _\lambda ^K,F(\lambda ))`$ is $`HK`$-spherical. The following questions are therefore natural: What is the algebraic impact of the condition (RDS)? Does there exists an analog of the Parthasarathy-condition (cf. \[EHW83. Prop. 3.9\]) for the symmetric space setting?
(b) Give a complete classification of $`H`$-spherical unitary highest weight representations. A first step in this direction might be Proposition II.7(ii) and Remark II.8.
V. Applications to holomorphic representation theory
In this final section we give a second application of the Averaging Theorem: We relate the spherical character of a spherical unitary highest weight representation of $`G`$ to the corresponding spherical functions on the $`c`$-dual space.
Spherical functions and character theory
Definition V.1. Let $`(\pi _\lambda ,_\lambda )`$ be an $`H`$-spherical unitary highest weight representation of $`G`$. If $`0\nu (_\lambda ^\omega )^H`$ and $`v_\lambda `$ is an highest weight vector, then we define the spherical character $`\mathrm{\Theta }_\lambda `$ of $`(\pi _\lambda ,_\lambda )`$ by
$$\mathrm{\Theta }_\lambda :S_{\mathrm{max}}^0C,s\frac{v_\lambda ,v_\lambda }{|\nu ,v_\lambda |^2}\pi _\lambda (s).\nu ,\nu .$$
Note that $`\mathrm{\Theta }_\lambda `$ is an $`H`$-biinvariant holomorphic function on $`S_{\mathrm{max}}^0`$ (cf. \[KNÓ97, Lemma V.6\]).
Remark V.2. The particular normalization of $`\mathrm{\Theta }_\lambda `$ has two reasons. First that it coincides in the group case (cf. Remark III.5) with the standard definition, and second because it has the best analytic properties for the assignments $`\lambda \mathrm{\Theta }_\lambda (s)`$, $`sS_{\mathrm{max}}^0`$ (as less poles as possible).
Definition V.3. (Spherical Functions) Recall the definition of the domain $`_\mathrm{\Omega }a_C^{}`$ (cf. Definition II.12). If $`\lambda _\mathrm{\Omega }`$, then the spherical function with parameter $`\lambda `$ is defined by
$$\phi _\lambda :S_{\mathrm{max}}^0HANC,s_Ha_H(sh)^{\lambda \rho }𝑑\mu _H(h)$$
(cf. \[FHÓ94\] or \[KNÓ98\]). Recall that the defining integrals converge absolutely if and only if $`\lambda _\mathrm{\Omega }`$ (cf. \[FHÓ94, Th. 6.3\]).
Theorem V.4.Let $`(\pi _\lambda ,_\lambda )`$ be an $`H`$-spherical unitary highest weight representation of $`G`$ such that $`\lambda +\rho _\mathrm{\Omega }`$ holds. Then the spherical character $`\mathrm{\Theta }_\lambda `$ of $`(\pi _\lambda ,_\lambda )`$ and the spherical function $`\phi _{\lambda +\rho }`$ are related by
$$(sS_{\mathrm{max}}^0HAN)\mathrm{\Theta }_\lambda (s)=\frac{1}{c(\lambda +\rho )}\phi _{\lambda +\rho }(s).$$
In particular, $`\phi _{\lambda +\rho }`$ extends naturally to a $`H`$-biinvariant holomorphic function on $`S_{\mathrm{max}}^0`$.
Proof. Since $`\lambda +\rho _\mathrm{\Omega }`$, the assumption of Theorem II.16 is satisfied and we can rewrite $`0\nu (_\lambda ^\omega )^H`$ as
$$\nu =\frac{\nu ,v_\lambda }{v_\lambda ,v_\lambda c(\lambda +\rho )}_H\pi _\lambda (h).v_\lambda d\mu _H(h).$$
Thus if we replace the first $`\nu `$ in the definition of $`\mathrm{\Theta }_\lambda `$ by this expression, we get for all $`sS_{\mathrm{max}}^0HAN`$ that
$$\begin{array}{cc}\hfill \mathrm{\Theta }_\lambda (s)& =\frac{v_\lambda ,v_\lambda }{|\nu ,v_\lambda |^2}\pi _\lambda (s).\nu ,\nu \hfill \\ & =\frac{1}{c(\lambda +\rho )}\frac{1}{v_\lambda ,\nu }_H\pi _\lambda (sh).v_\lambda ,\nu d\mu _H(h)\hfill \\ & =\frac{1}{c(\lambda +\rho )}\frac{1}{v_\lambda ,\nu }_H\pi _\lambda (h_H(sh)a_H(sh)n_H(sh)).v_\lambda ,\nu d\mu _H(h)\hfill \\ & =\frac{1}{c(\lambda +\rho )}\frac{1}{v_\lambda ,\nu }_H\pi _\lambda (a_H(sh)).v_\lambda ,\nu d\mu _H(h)\hfill \\ & =\frac{1}{c(\lambda +\rho )}_Ha_H(sh)^\lambda 𝑑\mu _H(h)\hfill \\ & =\frac{1}{c(\lambda +\rho )}\phi _{\lambda +\rho }(s),\hfill \end{array}$$
as was to be shown.
Remark V.5. (a) We remark here that the relation in Theorem V.4 was long time searched by G. Ólafsson (cf. \[Ól98, Open Problem 7(1)\]). For further interesting problems related to this subject we refer to \[Fa98\] and \[Ól98\].
(b) The analytic continuation of the relation in Theorem V.4 has been obtained in \[HiKr98\]. It has far reaching consequences for the theory of $`G`$-invariant Hilbert spaces of holomorphic functions on $`G`$-invariant subdomains of the Stein manifold $`\mathrm{\Xi }_{\mathrm{max}}^0=G\times _HiW_{\mathrm{max}}^0`$. In particular, it implies the Plancherel Theorem for these class of Hilbert spaces (cf. \[HiKr98\]). For further information related to this subject we refer to \[HiKr99b\], \[KNÓ97\], \[Kr98,99b\] and \[Ne99a\].
Appendix
A. An identity criterion for bounded analytic functions on tubes
Lemma A.1.Let $`\mathrm{\Pi }^+:=\{zC:Im(z)>0\}`$ be the upper half plane and $`H^{\mathrm{}}:=\{fHol(\mathrm{\Pi }^+):f_{\mathrm{}}<\mathrm{}\}`$ the Banach space of bounded holomorphic functions on it. Let $`\alpha >0`$ and $`N=\{n\alpha i:nN\}`$. Then the following identity assertion for elements $`f`$ of $`H^{\mathrm{}}(\mathrm{\Pi }^+)`$ holds: If $`f_N=0`$, then $`f=0`$.
Proof. Let $`D:=\{zC:|z|<1\}`$ and $`H^{\mathrm{}}(D)=\{fHol(D):f_{\mathrm{}}<\mathrm{}\}`$. Let $`fH^{\mathrm{}}(D)`$ and $`\{\beta _n:nN\}`$ be subset of zeros of $`f`$. Then it follows from \[Ru70, Th. 15.23\] that
$$f=0\text{if}\underset{n=1}{\overset{\mathrm{}}{}}(1|\beta _n|)=\mathrm{}.$$
$`(A.1)`$
We consider the Cayley transform
$$c:\mathrm{\Pi }^+D,z\frac{zi}{z+i},$$
which is an biholomorphic isomorphism, defining an isomorphism of Banach spaces
$$c_{}:H^{\mathrm{}}(D)H^{\mathrm{}}(\mathrm{\Pi }^+),f\stackrel{~}{f}=fc.$$
Let $`\alpha _n:=n\alpha i`$. Then we have
$$\beta _n:=c(\alpha _n)=\frac{n\alpha ii}{n\alpha i+i}=\frac{n\alpha 1}{n\alpha +1}.$$
Let $`N_0N`$ be such that $`n\alpha 1>0`$ for all $`nN_0`$. Then
$$\underset{n=1}{\overset{\mathrm{}}{}}(1|\beta _n|)\underset{n=N_0}{\overset{\mathrm{}}{}}(1\frac{n\alpha 1}{n\alpha +1})=\underset{n=N_0}{\overset{\mathrm{}}{}}\frac{2}{n\alpha +1}=\mathrm{}.$$
$`(A.2)`$
Thus if $`\stackrel{~}{f}H^{\mathrm{}}(\mathrm{\Pi }^+)`$ vanishes on all $`\alpha _n`$, $`nN`$, then $`f(\beta _n)=0`$ for all $`nN`$ and so $`f=0`$ by (A.1) and (A.2). Therefore $`\stackrel{~}{f}=c_{}(f)=0`$, proving the lemma.
Proposition A.2.Let $`ØWR^n`$ be an open convex cone, $`T_W:=R^n+iW`$ the associated tube domain in $`C^n`$ and $`H^{\mathrm{}}(T_W)=\{fHol(T_W):f_{\mathrm{}}<\mathrm{}\}`$ the space of bounded holomorphic functions on $`T_W`$. Let $`\mathrm{\Gamma }R^n`$ be a lattice. Then the following identity assertion holds:
$$(fH^{\mathrm{}}(T_W))f_{i(\mathrm{\Gamma }W)}=0f=0.$$
Proof. We prove the assertion by induction on the dimension $`nN`$.
If $`n=1`$, then $`\mathrm{\Gamma }=Z\alpha `$ for some $`\alpha >0`$ and $`W=R,R^+`$ or $`R^{}`$. If $`W=R`$, then the assertion follows from Liouville’s Theorem. In the two remaining cases the assertion follows from Lemma A.1.
Suppose now the assertion is true for all all dimensions less or equal to $`n1`$, $`n2`$. Let $`fH^{\mathrm{}}(R^n+iW)`$ be an element vanishing on $`i(\mathrm{\Gamma }W)`$. We have to show that $`f=0`$. Since $`W`$ is open, we find a basis $`e_1,\mathrm{},e_n`$ of $`R^n`$ which is contained in $`\mathrm{\Gamma }W`$. By the Identity Theorem for analytic functions, it suffices to prove the assertion for $`\mathrm{\Gamma }=Ze_1\mathrm{}Ze_n`$ and $`W=_{j=1}^nR^+e_j`$. Let $`\mathrm{\Gamma }_{n1}=Ze_1\mathrm{}Ze_{n1}`$ and $`W_{n1}=_{j=1}^{n1}R^+e_j`$. Write the variables $`zC^n`$ as tuples $`z=(z^{},z_n)`$ with $`z^{}=(z_1,\mathrm{},z_{n1})`$. By induction we obtain that $`f(z)=f(z^{},z_n)`$ does not depend on the $`z^{}`$-variable. Thus $`f(z)=F(z_n)`$ for some $`FH^{\mathrm{}}(\mathrm{\Pi }^+)`$ with $`F_{Ni}=0`$. Thus by the induction hypothesis $`F=0`$, and hence $`f=0`$ establishing the induction step.
B. A lemma on spherical highest weight modules
Throughout this subsection $`(g,\tau )`$ denotes a simple hermitian symmetric Lie algebra. Further we use the notation from Section I-II.
Lemma B.1.Suppose that $`(g,\tau )`$ is a simple hermitian symmetric Lie algebra and $`(G,\tau )`$ an associated simply connected Lie group. Set $`H=G^\tau `$ and assume that there exist a non-trivial $`H`$-spherical unitary highest weight representation $`(\pi _\lambda ,_\lambda )`$ of $`G`$. Then the symmetric Lie algebra $`(g,\tau )`$ has to be compactly causal.
Proof. Write $`g=kp`$ for a $`\tau `$-invariant Cartan decomposition of $`g`$ and let $`K`$ denote the analytic subgroup of $`G`$ corresponding to $`k`$.
By assumption we have $`(_\lambda ^\omega )^H\{0\}`$. In particular we can conclude that the module $`L(\lambda )`$ of $`K`$-finite vectors of $`(\pi _\lambda ,_\lambda )`$ admits nontrivial $`HK`$-fixed vectors. Recall that $`L(\lambda )`$ is the unique irreducible quotient of the generalized Verma module
$$N(\lambda )=𝒰(g_C)_{𝒰(k_Cp^+)}F(\lambda ).$$
In particular, there exists an element $`0v_0N(\lambda )^{HK}`$.
Recall that $`N(\lambda )`$ is $`k_C`$-isomorphic to $`𝒮(p^{})F(\lambda )`$, where the $`k_C`$-action on $`𝒮(p^{})F(\lambda )`$ is defined by
$$X.(pv):=[X,p]v+pX.v$$
$`(B.1)`$
for $`Xk_C`$, $`p𝒮(p^{})`$ and $`vF(\lambda )`$ (cf. \[EHW83\]).
In order to show that $`(g,\tau )`$ is compactly causal, we have to prove $`z(k)q`$. Assume the contrary, i.e., $`z(k)h`$. Recall the definition of the element $`Z_0z(k)`$ from Section I and set $`X_0:=iZ_0iz(k)`$. Then the spectrum of $`X_0`$, considered as an operator on the symmetric algebra $`𝒮(p^{})`$, is $`N_0`$, and we obtain a natural grading by homogeneous elements: $`𝒮(p^{})=_{n=0}^{\mathrm{}}𝒮(p^{})^n`$. Then $`N(\lambda )=_{n=0}^{\mathrm{}}𝒮(p^{})^nF(\lambda )`$ and we conclude from (B.1) that $`X_0`$ acts on $`𝒮(p^{})^nF(\lambda )`$ by $`n+\lambda (X_0)`$ times the identity. Write $`v=_{n=0}^{\mathrm{}}v_0^n`$ according to the decomposition $`N(\lambda )=_{n=0}^{\mathrm{}}𝒮(p^{})^nF(\lambda )`$. Since $`X_0i(hk)`$, the element $`v_0`$ is annihilated by $`X_0`$ and so we must have $`v_0=v_0^n`$ for some $`nN_0`$ with $`\lambda (X_0)=n0`$. But a necessary condition for $`L(\lambda )`$ to be unitarizable is $`\lambda (X_0)<0`$ (cf. \[Ne99b, Th. XI.2.37(ii)\]). This gives us a contradiction and proves the lemma.
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Bernhard Krötz Mathematisches Institut Technische Universität Clausthal Erzstraße 1 D-38678 Clausthal-Zellerfeld Germany mabk@math.tu-clausthal.de
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# Mass Enhancement in Narrow Band Systems
## I.Introduction
Several studies published over the last years have addressed the questions related to the existence of small polarons with itinerant properties in real systems . This issue is central for a possible description of high $`T_c`$ superconductivity in terms of (bi)polaronic models . In spite of being well defined quasiparticles, small polarons may loose their mobility either because of a dynamical dephasing between the charge carriers and their surrounding deformation field or because of the heaviness of the effective mass. These effects could however differ significantly according to the regime (adiabatic or antiadiabatic) and the strength of electron- phonon coupling characterizing the system . Theoretical investigations start generally from the Holstein Molecular Crystal Model , a fundamental tool which has revealed a rich variety of behaviors in the polaron landscape through the use of quantum Monte Carlo , density matrix renormalization group techniques , variational methods , cluster solutions and perturbative approaches . A numerical study of the polaron bandwidth in the first order of perturbative theory has proved that the phonon momentum dependence is a key feature of the Holstein Hamiltonian and that the lattice dimensionality strongly influences the bandwidth values . Unlike other properties such as ground state energy and effective mass, the bandwidth is not affected by second order corrections and therefore it provides a useful testing bench for alternative, non perturbative attacks on the polaron problem . Being aware of the importance that the intermolecular forces have in the narrowing of the polaron band, we report here on a perturbative numerical study of the mass enhancement in the strong coupling and antiadiabatic regime. The reasons why I choose to start from this regime are threefold: i) it is the easiest in the sense that the lattice deformation follows coherently the charge carriers and the abovementioned dephasing features can be ruled out, ii) the unit comprising electron and phonon dressing is a stable small polaron that is, the size of the quasiparticle is not significantly broadened in some portions of our parameter space, iii) this regime is relevant to several classes of narrow band materials whose charge carriers effective masses deserve accurate estimates. Although the present work assumes that the carriers are coupled to bosonic degrees of freedom having a vibrational origin, antiadiabatic conditions are likely to occur in excitonic systems where the characteristic frequency $`\overline{\omega }`$ could be easily of order of $`1eV`$ and the boson can therefore follow the electron eesentially without retardation. In these cases the carriers effective mass is expected to be only moderately enhanced with respect to the bare electron mass.
In Section II, the dispersive Holstein model is briefly reviewed while the numerical results are displayed in Section III. Some conclusions are drawn in Section IV.
## II. Holstein Model with Dispersive Phonons
My starting point is the real space - momentum space representation of the Holstein Hamiltonian which reads
$`H=`$ $`t{\displaystyle \underset{ij}{}}c_i^{}c_j+ϵ{\displaystyle \underset{j}{}}c_j^{}c_j+{\displaystyle \underset{𝐤}{}}\mathrm{}\omega _𝐤a_𝐤^{}a_𝐤`$ (2)
$`+{\displaystyle \frac{g}{\sqrt{N}}}{\displaystyle \underset{𝐤}{}}{\displaystyle \underset{j}{}}c_j^{}c_j(a_𝐤+a_𝐤^{})exp\left(i𝐤𝐫_j\right)`$
$`c_i^{}`$ ( $`c_i`$ ) creates (destroys) a tight binding electron at the $`ith`$ molecular site and $`t`$ is the hopping integral related to the bare electron half bandwidth $`D`$ by $`D=zt`$, $`z`$ being the coordination number. $`ϵ`$ is a reference electronic energy and $`𝐫_j`$ is the $`jth`$ lattice site vector. $`N`$ is the number of molecules in the lattice. It is understood that $`t`$ differs from zero only between first neighbors sites and this poses a constraint on the $`ij`$ sum in the first addendum. $`a_𝐤^{}`$ ( $`a_𝐤`$ ) creates (destroys) a k-phonon with frequency $`\omega _𝐤`$. The lattice dimensionality enters the problem through the phonon dispersion relations which have been obtained analytically by assuming first neighbors pairwise intermolecular forces both along the linear chain (1D), the square lattice (2D) and the simple cubic lattice (3D) :
$`\omega _{1D}^2(k)={\displaystyle \frac{\beta +\gamma }{M}}+{\displaystyle \frac{1}{M}}\sqrt{\beta ^2+A_x}`$ (3)
$`\omega _{2D}^2(𝐤)={\displaystyle \frac{\beta +2\gamma }{M}}+{\displaystyle \frac{1}{M}}\sqrt{\beta ^2+2B_{x,y}}`$ (4)
$`\omega _{3D}^2(𝐤)={\displaystyle \frac{\beta +3\gamma }{M}}+{\displaystyle \frac{1}{M}}\sqrt{\beta ^2+C_{x,y,z}}`$ (5)
$`A_x=\gamma ^2+2\gamma \beta c_x`$ (6)
$`B_{x,y}=\gamma ^2(1+c_xc_y+s_xs_y)+\beta \gamma (c_x+c_y)`$ (7)
$`C_{x,y,z}=\gamma ^2(3+2(c_xc_y+s_xs_y+c_xc_z+s_xs_z`$ (8)
$`+c_yc_z+s_ys_z))+2\beta \gamma (c_x+c_y+c_z)`$ (9)
where, $`c_x=cosk_x`$, $`c_y=cosk_y`$, $`c_z=cosk_z`$, $`s_x=sink_x`$ etc. $`\beta `$ is the intra-molecular force constant and $`\gamma `$ is the inter-molecular first neighbors force constant. Let’s define $`\omega _0^2=\mathrm{\hspace{0.17em}2}\beta /M`$ and $`\omega _1^2=\gamma /M`$ with $`M`$ being the reduced molecular mass. $`N`$ is the number of diatomic molecules in the lattice and $`g`$ is the local electron-phonon coupling constant. The adiabatic parameter is $`\mathrm{}\overline{\omega }/D`$, $`\overline{\omega }`$ being a characteristic phonon frequency which we take as the zone center frequency and whose expression is: $`\overline{\omega }^2=\omega _0^2+z\omega _1^2`$.
Throughout this paper we fix $`\mathrm{}\omega _0=\mathrm{\hspace{0.17em}100}meV`$ and $`t=15meV`$ so that the antiadiabatic condition $`\mathrm{}\overline{\omega }/D>1`$ is fulfilled in any dimensionality. Moreover, our perturbative approach requires the occurence of the condition $`D<g`$ . By applying the Lang-Firsov unitary transformation and the subsequent $`1/\lambda _0`$ expansion with $`\lambda _0g^2/(\mathrm{}\omega _0D)`$ being the ratio between polaron binding energy and electron half bandwidth , $`H`$ of eq.(1) transforms into $`\stackrel{~}{H}=\stackrel{~}{H}_0+\stackrel{~}{H}_P`$ with:
$`\stackrel{~}{H}_0`$ $`={\displaystyle \underset{𝐤}{}}\mathrm{}\omega _𝐤a_𝐤^{}a_𝐤+\left(ϵ{\displaystyle \frac{g^2}{N}}{\displaystyle \underset{𝐤}{}}(\mathrm{}\omega _𝐤)^1\right){\displaystyle \underset{j}{}}c_j^{}c_j`$ (11)
$`{\displaystyle \frac{g^2}{N}}{\displaystyle \underset{𝐤}{}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{exp\left(i𝐤(𝐫_i𝐫_j)\right)}{\mathrm{}\omega _𝐤}}c_j^{}c_jc_i^{}c_i`$
$`\stackrel{~}{H}_P`$ $`=t{\displaystyle \underset{ij}{}}exp[{\displaystyle \frac{2g^2}{N}}{\displaystyle \underset{𝐤}{}}(\mathrm{}\omega _𝐤)^2sin^2{\displaystyle \frac{(𝐤(𝐫_i𝐫_j))}{2}}]`$ (14)
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}\left[{\displaystyle \frac{g}{\sqrt{N}}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{a_𝐤^{}}{\mathrm{}\omega _𝐤}}(e^{i𝐤𝐫_i}e^{i𝐤𝐫_j})\right]^m`$
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left[{\displaystyle \frac{g}{\sqrt{N}}}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{a_𝐤}{\mathrm{}\omega _𝐤}}\left(e^{i𝐤𝐫_j}e^{i𝐤𝐫_i}\right)\right]^nc_i^{}c_j`$
$`\stackrel{~}{H}_0`$ is diagonal except for a second order term in the electron density operator which could cause an attractive electron-electron interaction . The perturbation $`\stackrel{~}{H}_P`$ displays the fundamental features of the polaronic quasiparticle as the hopping integral narrowing (first factor in eq.(3)) plus the peculiar mixing of fermionic and bosonic variables. At any electron-phonon interaction vertex $`m`$ ($`n`$) phonons can be emitted from (absorbed by) the cloud surrounding the electron provided the total crystal momentum is conserved. By choosing a transformed ground state with no phonons we see that the first order dispersive contribution $`E^{(1)}`$ to the ground state polaron band arises only from the $`n=m=0`$ term in eq.(3) hence from the zero phonon scattering process. In 3D and taking a lattice spacing $`a=|𝐫_i𝐫_j|=1`$, one finds
$`E^{(1)}(𝐩)`$ $`=2t(\mathrm{cos}p_x+\mathrm{cos}p_y+\mathrm{cos}p_z)`$ (16)
$`exp\left[{\displaystyle \frac{2g^2}{N}}{\displaystyle \underset{k_x}{}}\mathrm{sin}^2{\displaystyle \frac{k_x}{2}}{\displaystyle \underset{k_x,k_y}{}}(\mathrm{}\omega _𝐤)^2\right]`$
where the total crystal momentum p coincides with the electron momentum due to the absence of self-energy corrections. The second order perturbative contribution requires summation over all intermediate states having $`m`$ k-phonons more than the vacuum and one electron on a $`i`$ first neighbor of the $`j`$ initial site. Moreover, the final electronic position $`f`$ can either coincide with $`j`$ (this process does not introduce dispersive effects in the polaron band) or be a first neighbor of the $`i`$ site. The latter event is clearly dimension dependent: in 1D the final site is a second neighbor of $`j`$, in 2D $`f`$ can be either a second or a third neighbor of $`j`$ and in 3D, also the fourth neighbor site can be reached via hopping. While the detailed study of these dispersive effects (which can become relevant in adiabatic conditions) is postponed to a next paper we turn now to compute the polaron effective mass. It should be remarked that the second order corrections decrease the mass values calculated in first order perturbative theory which therefore should be meant to provide upper bounds for the polaron mass.
## III. Polaron Effective Masses
The polaron mass $`m^{}`$ can be obtained according to the definition
$$\frac{m^{}}{m_0}=\frac{zt}{^2E(𝐩)|_{𝐩=0}}$$
(17)
where $`m_0`$ is the bare band mass and the dispersive polaron band is given by eq.(4). The polaron binding energy has obviously no p-dispersion. Then, $`m^{}/m_0`$ is at first order independent of $`t`$ and coincident simply with the reciprocal of the band narrowing factor. This picture holds in the strong coupling regime here assumed. In Fig.1(a), the ratios $`m^{}/m_0`$ for the 1D, 2D and 3D are computed versus the first neighbor intermolecular force constant $`\omega _1`$. While the polaron masses strongly depend on the dimensionality $`d`$ and are very large at small $`\omega _1`$, the ratios become essentially $`d`$ independent in the upper portion of the parameter range and tend to converge to 2. The value of the polaron binding energy $`\lambda _0>1`$ signals that the energy gain associated with the lattice deformation is larger than the kinetic energy due to the tight binding hopping in the bare band. Therefore it is energetically convenient to the electron to be dressed by the phonon cloud and become a quasiparticle. Actually, in antiadiabatic regimes, the more restrictive condition concerning the lattice deformation $`\alpha _0g/(\mathrm{}\omega _0)>1`$ needs to be fulfilled to guarantee that our quasiparticle is a small polaron. While $`\lambda _0`$ and $`\alpha _0`$ refer to a system with dispersionless phonons it is clear that both polaron binding energy and lattice deformation parameter will change after switching on the intermolecular interactions. The role of the intermolecular couplings is not simply that of increasing the characteristic phonon frequency but rather that of establishing the correct Holstein model dependence on dimensionality. Ignoring the intermolecular couplings would yield the 1D polaron band $`\mathrm{\Delta }E_{1D}`$ larger than the 2D polaron band $`\mathrm{\Delta }E_{2D}`$ and $`\mathrm{\Delta }E_{2D}>\mathrm{\Delta }E_{3D}`$ which is clearly unphysical since the polaronic wave functions overlap is larger in higher dimensionality. This wrong trend would hold for any value of the intramolecular frequency $`\omega _0`$. Then, as observed by Holstein himself in his original paper , the phonon dispersion is a vital ingredient of the theory and this observation motivates our numerical investigation.
Because of our definitions $`\lambda _0`$ does not depend on $`d`$ whereas $`\alpha _0`$ is $`\sqrt{d}`$. In Fig.1(b) we see that the lattice deformation $`\alpha =N^1g_𝐤\mathrm{}\omega _{𝐤}^{}{}_{}{}^{1}`$ is in all dimensions a decreasing function of the intermolecular force constant and, in 1D, the system does not fulfill the small polaron condition at the largest $`\omega _1`$ values. This case has been presented to point out how the starting condition $`\alpha _0^{1D}=1.3`$ sets the 1D system rather in an intermediate coupling regime where a broadening of the polaron size can take place . Under these conditions the same perturbative method based on the Lang Firsov transformation becomes questionable . Below $`\overline{\omega }_1=\mathrm{\hspace{0.17em}48}meV`$ the polaron bandwidth inequalities $`\mathrm{\Delta }E_{3D}>\mathrm{\Delta }E_{2D}>\mathrm{\Delta }E_{1D}`$ are not satisfied as expected on general grounds hence, the dispersionless and the weakly dispersive Holstein Hamiltonian yield erroneous estimates of the effective masses. The straight line in Fig.1(b) marks therefore the lower bound for the intermolecular coupling which guarantees the validity of the model.
Increasing the electron-phonon coupling, see Fig.2(a), leads to a strong mass enhancement (particularly in 3D) at small $`\omega _1`$ while the mass ratios converge to 5 at large intermolecular coupling strenghts. Note (Fig.2(b)) that the polaron is small in all $`d`$ throughout the whole $`\omega _1`$ range hence the Lang Firsov method works well in this case. The threshold value for the validity of the Holstein model is set here at $`\overline{\omega }_1=\mathrm{\hspace{0.17em}59}meV`$. The inequalities $`m_{3D}^{}>m_{2D}^{}>m_{1D}^{}`$ keep on being satisfied for a portion of $`\omega _1`$ values above the threshold before convergence is achieved but second order perturbative terms (being larger in higher dimensionality) are expected to correct partly this trend. Figs.3 show that a stronger e-ph coupling, with $`\lambda _0=21.3`$, yields a mass ratio of $`25`$ and shifts the threshold $`\overline{\omega }_1`$ at $`65meV`$ pointing out the relationship between features of the phonon spectrum and strength of the $`g`$ coupling. Also in this case the polaron size remains small throughout the whole $`\omega _1`$ range (see Fig.3(b)) thus confirming the reliability of the Lang Firsov method in a strong e-ph coupling regime with antiadiabatic conditions.
Next, I have varied $`g`$ in the range 1 - 4 (in units of $`\mathrm{}\omega _0`$ and found, at any $`g`$, the minimum intermolecular coupling $`\overline{\omega }_1(g)`$ at which the bandwidth inequalities $`\mathrm{\Delta }E_{3D}>\mathrm{\Delta }E_{2D}>\mathrm{\Delta }E_{1D}`$ are satisfied. This criterion yields an empirical relation, $`\overline{\omega }_1(g)\overline{\omega }_1(1)(1+ln(g))`$, which allows one to obtain a reliable estimate of the polaron effective mass. In Fig.4 the 1D mass ratio is plotted versus the dimensionless $`g/(\mathrm{}\omega _0)`$ both in the first and second order of perturbative theory: it turns out that second order corrections are negligible in 1D systems with antiadiabatic conditions as those we have assumed. I want to point out that the mass values reported in Fig.4 correspond, at any $`g`$, to the minimum $`\omega _1`$ (the threshold) which ensure the smallness of the ground state polaron. Then, they are upper bounds for the 1D mass in the sense that the presence of larger intermolecular forces would yield lighter mass values. As expected on general grounds the small polaron solution is the ground state of the discrete Holstein model in the intermediate to strong e-ph-coupling regime here considered while, by decreasing the coupling, a continuous cross over to large polaron solutions can take place in 1D . We have however seen (Fig.1(b)) that the dispersive features of the phonon spectrum could affect the transition by inducing a spreading of the lattice deformation. Anyway, the smoothness of our $`m^{}`$ versus $`\alpha _0`$ curve (persisting also in the lower $`\alpha _0`$ range not displayed in Fig.4) confirms that no self trapping is found in 1D antiadiabatic regimes whereas recent variational and perturbative investigations signalling a rapid growth of $`m^{}`$ vs. e-ph coupling support the existence of the self trapping transition between polaron states of different structure in 1D adiabatic systems. In any case, phase transitions in Holstein models are ruled out being the ground state energy an analytic function of the e-ph coupling.
## IV. Concluding Remarks
I have presented the first results of a perturbative approach to the polaron problem which focusses on the lattice dimensionality effects. Having chosen the antiadiabatic regime of the Holstein Molecular Crystal Model we are confident of the accuracy of the first order perturbative theory for one dimensional systems with strong e-ph coupling whereas some significant second order corrections can occur in higher dimensionality . While a previous work had shown that a dispersionless Holstein model leads to i) erroneous estimates of the polaron bandwidth versus dimensionality and ii) unphysical divergences in the site jump probability , the present study reveals that the polaron effective mass is in all dimensions very sensitive to the strength of the forces which tie the molecules in the lattice. We obtain polaron masses between 2 and 25 times the bare band mass by varying the e-ph coupling in the range $`(12.5)`$ and these values become essentially dimension independent when the intermolecular forces are sufficiently strong. The molecular lattice structure has been described by means of a single parameter, the first neighbor intermolecular coupling, being understood that the range of the interactions should be extended in real systems if least squares fitting of the experimental phonon frequencies can provide effective values for the next neighbors and long range force constants . The antiadiabatic regime with strong e-ph coupling ensures the validity of the quasiparticle picture for the small polaron nonetheless we have seen that some broadening of the phonon cloud can arise at intermediate e-ph couplings for strong values of the intermolecular forces with consequent lowering of the lattice deformation parameter. This interesting feature suggests that the intermolecular forces influence the quasiparticle size and, incorporating the effects of the e-ph coupling, have a role in driving the continuous transition between large and small polaron.
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# THE GRAVITATIONAL LENSING IN REDSHIFT-SPACE CORRELATION FUNCTIONS OF GALAXIES AND QUASARS
## 1 INTRODUCTION
The correlation function is one of the fundamental quantities in searching into the physical origin of the universe. In large-scale redshift surveys, redshifts and spherical positions on the sky of luminous objects are used for estimating the spatial distribution of mass, but the former is distorted by inhomogeneity of the universe.
Two intrinsic distortion effect on the correlation function originate in the velocity field (Kaiser, 1987) and the cosmological warp (Ballinger, Peacock & Heavens, 1996; Matsubara & Suto, 1996; Matsubara, 2000). The former comes from the fact that an observed redshift corresponds to the recession velocity which is composed not only of the expansion of the universe but also of peculiar velocities. The latter distortion is brought about by the nonlinear mapping of the objects from the expanding curved space on a light cone onto a flat redshift space (Alcock & Paczyński, 1979).
In addition to those intrinsic distortions, there are secondary distortions which are due to perturbations of light rays. While the redshift is altered by the Sachs-Wolfe effect (Sachs & Wolfe, 1967), the spherical position is recast by the gravitational lensing (Schneider et al., 1992; Mellier, 1999; Bartelmann & Schneider, 1999). Although the Sachs-Wolfe effect is not so important, the gravitational lensing can affect the observable correlation function in forthcoming redshift surveys as this effect is efficient for high-redshift objects (Gunn, 1967).
The effect of the gravitational lensing on the angular functions $`w(\theta )`$ have been intensively investigated so far (Bartelmann & Schneider, 1999; Moessner et al., 1998; Kaiser, 1992; Villumsen, 1996). Recently, among others, cross correlations of galaxies at different redshifts (Moessner et al., 1998) are successfully applied to the commissioning data from the Sloan Digital Sky Survey (Jain et al., 2000), to single out the weak lensing effect.
As for the 3D correlation function in redshift space, a qualitative treatment of the lensing effect is given in Suto et al. (1999) to estimate the upper limit of the effect, using the phenomenological Dyer-Roeder distance (Dyer & Roeder, 1973). It is still not clear whether or not the lensing is actually efficient where the intrinsic correlation function is negligible on scales comparable to 100Mpc. The main purpose of this letter is to give a quantitative treatment of this issue, consistently including velocity and cosmological distortions, and consequently to show the weak lensing actually has detectable effects on 3D correlation function in redshift space when large-scale redshift surveys like Sloan Digital Sky Survey (SDSS) are available.
## 2 OBSERVABLE QUANTITIES
We take the homogeneous, isotropic FRW metric with scalar perturbations in longitudinal gauge:
$`ds^2=a^2(\tau )\left\{(1+2\varphi )d\tau ^2+(12\varphi )\left[d\chi ^2+S_{K}^{}{}_{}{}^{2}(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)\right]\right\},`$ (1)
where $`a`$ is the scale factor, and $`\tau `$ is the conformal time, $`d\tau =dt/a`$, and $`S_K`$ is the comoving angular distance of the spatial curvature $`K=\mathrm{\Omega }_0+\lambda _01`$. For example, $`S_K(\chi )=(K)^{1/2}\mathrm{sinh}\left[(K)^{1/2}\chi \right]`$ for open universe, $`K<0`$, and $`S_K(\chi )=\chi `$ for flat universe, $`K=0`$. We adopt the unit $`c=H_0=1`$ throughout this letter.
From the first-order Einstein equation of the metric in equation (1), the density contrast $`\delta (𝒙,\tau )`$ and the velocity field $`v^i(𝒙,\tau )`$ on scales much less than the curvature scale satisfy
$`\mathrm{}\varphi ={\displaystyle \frac{3\mathrm{\Omega }_0}{2}}{\displaystyle \frac{\delta }{a}},v^i={\displaystyle \frac{2}{3\mathrm{\Omega }_0}}aHf^i\varphi ,`$ (2)
where $`H=\dot{a}/a`$, $`f=d\mathrm{ln}D/d\mathrm{ln}a`$, and $`D`$ is the linear growth rate (Peebles, 1980). The Laplacian $`\mathrm{}`$ is taken with respect to comoving coordinates.
Let us consider a light ray emitted from an object at comoving coordinates $`(\chi ,\theta ,\phi ;\tau )`$, which an observer recieves at $`\tau _0`$. The conventional redshift $`z=a^11`$ is given by $`\chi =_0^z𝑑zH^1`$, but the actually observed redshift is changed by the line-of-sight peculiar velocity, $`V=n_iv^i`$, where $`n_i`$ is a line-of-sight unit vector, and also by the gravitational potential, $`\varphi `$. From the time-component of the geodesic equation of the light ray, the observed redshift $`z_\mathrm{s}`$ is given by (Sachs & Wolfe, 1967)
$`z_\mathrm{s}=z+(1+z)\left[V(\chi )V(0)\varphi (\chi )+\varphi (0)2{\displaystyle _{\tau _0\chi }^{\tau _0}}𝑑\tau {\displaystyle \frac{\varphi }{\tau }}\right],`$ (3)
where we abbreviate the function on a light cone as $`V(\chi )V(\chi ,\theta ,\phi ;\tau _0\chi )`$, and so as $`\varphi (\chi )`$. The integral is performed on the light cone for a fixed direction of line of sight.
Now, we consider the small angle approximation so that light rays are confined to a narrow cone around the polar axis, $`\theta 1`$, and introduce new coordinates $`\theta _1=\theta \mathrm{cos}\phi `$, $`\theta _2=\theta \mathrm{sin}\phi `$, following Kaiser (1998). Adopting the Born approximation, the angular components of the geodesic equation reduces to equation for the observed angular components $`\theta _{\mathrm{s}a}`$ ($`a=1,2`$):
$`\theta _{\mathrm{s}a}=\theta _a+{\displaystyle \frac{2}{S_K(\chi )}}{\displaystyle _0^\chi }𝑑\chi ^{}S_K(\chi \chi ^{})_a\varphi (\chi ^{}),`$ (4)
where $`_a=S_{K}^{}{}_{}{}^{1}/\theta _a`$.
The apparent luminosity of the light is magnified by a factor $`A(z)=|det(\theta _{\mathrm{s}b}/\theta _a)|=1+2\kappa `$, where $`\kappa `$ is a local convergence field of weak lensing (Schneider et al., 1992; Bernardeau et al., 1997; Kaiser, 1998) for a fixed redshift of source object:
$`\kappa (z,𝜽)={\displaystyle _0^\chi }𝑑\chi ^{}g(\chi ,\chi ^{})_a_a\varphi (\chi ^{});g(\chi ,\chi ^{}){\displaystyle \frac{S_K(\chi ^{})S_K(\chi \chi ^{})}{S_K(\chi )}}.`$ (5)
Due to the magnification, the observed apparent magnitude $`m_\mathrm{s}`$ is given by $`m_\mathrm{s}=m2.5\mathrm{log}_{10}A=m5\kappa /\mathrm{ln}10`$ (Broadhurst et al., 1995), where $`m`$ is the apparent magnitude in the absence of lensing.
The magnitude-limited number density in real space $`n_\mathrm{r}(z,𝜽;<m)`$ and that in observed redshift space $`n_\mathrm{s}(z_\mathrm{s},𝜽_\mathrm{s};<m_\mathrm{s})`$ are related by the number conservation equation, $`n_\mathrm{r}z^2dzd^2\theta =n_\mathrm{s}z_{\mathrm{s}}^{}{}_{}{}^{2}dz_\mathrm{s}d^2\theta _\mathrm{s}`$, while the number density is given by $`n(z,𝜽;<m)=[1+\delta (z,𝜽)]N(z;<m)/(4\pi z^2)`$, where $`\delta (z,𝜽)`$ is the density contrast and $`N(z;<m)`$ is the magnitude-limited number count per redshift. Evaluating the Jacobian, one obtains the relation between $`n`$ and $`n_\mathrm{s}`$, as well as the relation between density contrasts. The result contains the terms with $`\delta `$, $`V`$, $`\varphi `$ and their derivatives. For fluctuations on a scale $`k`$ in units of Hubble distance, such variables scales as $`Vk^1\delta `$, $`V\delta `$, $`\varphi k^2\delta `$, $`\varphi k^1\delta `$, and $`^2\varphi \delta `$, where $``$ schematically represents the spatial derivatives in comoving coordinates. Consistently to the small angle approximation, we neglect the fluctuations which scale as $`k^1`$ and $`k^2`$, because $`k`$ is large enough on scales below Hubble distance. Eventually, the distorted density contrast is given by
$`\delta _\mathrm{s}=\delta _\mathrm{r}{\displaystyle \frac{1+z}{H}}{\displaystyle \frac{V}{\chi }}+(5\alpha 2)\kappa `$ (6)
where $`\delta _\mathrm{r}(z,𝜽)`$ is the number density contrast of the objects in real space, and $`H(z)`$ is the Hubble parameter at $`z`$. The logarithmic slope of the number counts $`\alpha `$ at the limiting magnitude $`m`$ is given by $`\alpha (z,m)=\mathrm{log}_{10}N(z;<m)/m`$ (c.f., Moessner, Jain & Villumsen 1998). The first two terms of equation (6) depends strongly on radius $`z`$ and angles $`𝜽`$, while the last, the lens surface density, depends strongly on angles but very weakly on radius. The first term of equation (6) is the real density fluctuations, the second term is the velocity distortion (Kaiser, 1987; Matsubara & Suto, 1996), and the last term is the effect of weak lensing, which consists of the contribution from the modulation by magnification bias, $`5\alpha \kappa `$, and of the alternation of surface density by lensing, $`2\kappa `$. In the following, we denote each term in equation (6) as $`\delta _\mathrm{s}=\delta _\mathrm{r}+\delta _\mathrm{v}+\delta _\mathrm{l}`$.
## 3 CORRELATION FUNCTION
We consider the correlation function between two objects $`(z_1,𝜽_1)`$ and $`(z_2,𝜽_2)`$ in the small angle approximation, $`\theta |𝜽_1𝜽_2|1`$, and assume $`z_1z_2`$ without loss of generality. In the absence of the lensing term $`\delta _\mathrm{l}`$, the correlation function is given by Matsubara & Suto (1996), which generalize the work by Hamilton (1992) to high-redshift objects. In the coordinate system $`(z_1,z_2,\theta )`$, their result, with slight modification allowing difference of the kind of two objects, is expressed as follows:
$`\xi _{\mathrm{MS}}=\left[\delta _\mathrm{r}(1)+\delta _\mathrm{v}(1)\right]\left[\delta _\mathrm{r}(2)+\delta _\mathrm{v}(2)\right]`$
$`=\left[1+{\displaystyle \frac{1}{3}}(\beta _1+\beta _2)+{\displaystyle \frac{1}{5}}\beta _1\beta _2\right]\xi _0(x;\overline{z})P_0(\mu )`$
$`\left[{\displaystyle \frac{2}{3}}(\beta _1+\beta _2)+{\displaystyle \frac{4}{7}}\beta _1\beta _2\right]\xi _2(x;\overline{z})P_2(\mu )+{\displaystyle \frac{8}{35}}\beta _1\beta _2\xi _4(x;\overline{z})P_4(\mu ),`$ (7)
where $`\overline{z}=(z_1+z_2)/2`$, $`\beta _i=f(z_i)/b_i(z_i)\mathrm{\Omega }^{0.6}(z_i)/b_i(z_i)`$, and $`b_i(z_i)`$ is the bias parameter of object $`i`$ ($`i=1,2`$) at redshift $`z_i`$. Similarly, a bar for any variable means the evaluation at $`\overline{z}`$ and subscripts $`1,2`$ assume evaluation at objects 1 and 2, respectively. We denote the comoving separation $`x\sqrt{[S_K(\overline{\chi })]^2\theta ^2+(\chi _2\chi _1)^2}`$, the comoving cosine $`\mu (\chi _2\chi _1)/x`$, and $`P_n`$’s are Legendre polynomials, and
$`\xi _{2l}(x;\overline{z})=\overline{b}^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{k^2dk}{2\pi ^2}}j_{2l}(kx)P(k;\overline{z}),`$ (8)
where $`P(k;z)`$ is the power spectrum at redshift $`z`$. This formula is valid only for distant observer approximation, $`x\chi _1`$ and $`\theta 1`$. We have omitted the finger-of-God effect which is only important on scales less than $`10h^1\mathrm{Mpc}`$ or $`1000\mathrm{k}\mathrm{m}/\mathrm{s}`$. In the following, we are interested in the scales of 30Mpc or larger where lensing effect appears, so that we can safely ignore the nonlinear effect like finger-of-God effect.
Adopting the small angle approximation (e.g., Bernardeau et al. 1997; Kaiser 1998; Moessner et al. 1998), the correlations involving the lensing term $`\delta _\mathrm{l}`$ are obtained as
$`\xi _{\mathrm{rl}}=\delta _\mathrm{r}(1)\delta _\mathrm{l}(2)={\displaystyle \frac{3}{2}}\mathrm{\Omega }_0b_1(5\alpha _22)g(\chi _2,\chi _1)(1+z_1)\xi _\mathrm{p}[\theta S_K(\chi _1);z_1],`$ (9)
$`\xi _{\mathrm{ll}}=\delta _\mathrm{l}(1)\delta _\mathrm{l}(2)=\left({\displaystyle \frac{3}{2}}\mathrm{\Omega }_0\right)^2(5\alpha _12)(5\alpha _22)`$
$`\times {\displaystyle _0^{\chi _1}}d\chi (z)g(\chi _1,\chi )g(\chi _2,\chi )(1+z)^2\xi _\mathrm{p}[\theta S_K(\chi );z],`$ (10)
where $`d\chi (z)=dz/H(z)`$, $`\xi _\mathrm{p}`$ is the projected correlation function defined by
$`\xi _\mathrm{p}(y;z)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\xi (\sqrt{x^2+y^2};z)={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{kdk}{2\pi }}J_0(ky)P(k;z),`$ (11)
and $`\xi (r;z)=\xi _0(r;z)`$ is the correlation function in real space at $`z`$. The terms $`\delta _\mathrm{l}(1)\delta _\mathrm{r}(2)`$, $`\delta _\mathrm{l}(1)\delta _\mathrm{r}(2)`$ are zero for $`z_1z_2`$ in the small angle approximation. The velocity-lensing term, $`\xi _{\mathrm{vl}}=\delta _\mathrm{v}(1)\delta _\mathrm{l}(2)`$ is explicitly calculated to be zero, which is because the term $`\delta _\mathrm{v}`$ only depends on fluctuations along the line of sight which are smoothed out. If we convolve the above expression of $`\xi _{\mathrm{rl}}`$ and $`\xi _{\mathrm{ll}}`$ with a selection function along the line of sight, we obtain the form of angular correlation function with the effect of weak lensing (Bartelmann, 1995; Villumsen, 1996; Dolag & Bartelmann, 1997; Moessner et al., 1998; Moessner & Jain, 1998).
Thus the total correlation function in 3D redshift space is given by $`\xi _{\mathrm{tot}}(z_1,z_2,\theta )=\xi _{\mathrm{MS}}+\xi _{\mathrm{rl}}+\xi _{\mathrm{ll}}`$ for $`z_1z_2`$ and $`\theta 1`$. The first term dominates on scales much smaller than Hubble distance, while the last two terms dominate on scales comparable to the Hubble distance along the line of sight. Therefore, even though $`\xi _{\mathrm{MS}}`$ is valid only for $`z_1z_2z_1`$ (distant observer approximation), we can use the form $`\xi _{\mathrm{tot}}`$ even when $`z_1z_2z_1`$.
In Figure 1, we plot the total correlation function $`\xi _{\mathrm{tot}}`$ together with each component, $`\xi _{\mathrm{MS}}`$, $`\xi _{\mathrm{rl}}`$, and $`\xi _{\mathrm{ll}}`$. With the choice of the CDM-like initial power spectrum (Bardeen et al., 1986) with a shape parameter $`\mathrm{\Gamma }=0.25`$ and a linear amplitude $`\sigma _8=1`$, we use the fitting formula for the fully non-linear power spectrum of Peacock & Dodds (1996) for lensing correlations, $`\xi _{\mathrm{rl}}`$ and $`\xi _{\mathrm{ll}}`$. Linear predictions are also plotted in the lower panels. The nonlinearity of the intrinsic correlation $`\xi _{\mathrm{MS}}`$, which is only important for the region $`z_2z_1\stackrel{<}{}0.003`$, is ignored. We exemplify the low-density flat model with $`\mathrm{\Omega }_0=0.3`$, $`\lambda =0.7`$. In the upper panels, the slope of the number counts is assumed as $`\alpha =1`$, and the bias factor as $`b=1`$, regardless of the redshift. This example corresponds to $`z0.2`$ and $`m18`$ of galaxies as seen in Table 1, in which the slope $`\alpha `$ is calculated from the B-band luminosity function of APM galaxies (Loveday et al., 1992), and of quasars (Boyle et al., 1988). The limiting magnitudes assumed in Table 1 correspond to estimated SDSS redshift data of galaxies and quasars. In practice, the slope $`\alpha `$ can be observationally determined for individual catalogue of quasar or galaxy redshift surveys. Each component of correlations roughly scales as $`\xi _{\mathrm{MS}}\sigma _{8}^{}{}_{}{}^{2}b_1b_2`$, $`\xi _{\mathrm{rl}}\sigma _{8}^{}{}_{}{}^{2}\mathrm{\Omega }_0b_1(5\alpha _22)`$, $`\xi _{\mathrm{ll}}\sigma _{8}^{}{}_{}{}^{2}\mathrm{\Omega }_{0}^{}{}_{}{}^{2}(5\alpha _12)(5\alpha _22)`$ for other parameters and models.
In the lower panels of Figure 1, the galaxy-galaxy (G-G), galaxy-QSO (G-Q), and QSO-QSO (Q-Q) correlations are plotted, assuming the SDSS slope $`\alpha `$ of Table 1. The bias parameter is set $`b=1`$ and $`3`$ for galaxies and quasars, respectively. The line-of-sight separations are large enough in lower panels, so that the intrinsic clustering is negligible.
## 4 DISCUSSION
The absolute value of the intrinsic clustering component $`\xi _{\mathrm{MS}}`$ is a decreasing function of the line-of-sight separation, $`z_2z_1`$, except the vicinity of zero crossings. On the other hand, the lens-lens component $`\xi _{\mathrm{ll}}`$ is almost independent on the separation and the density-lens component $`\xi _{\mathrm{rl}}`$ is an increasing function. Those behaviors are understood by the fact that the weak lensing is efficient between the object and the observer. Thus, intrinsic clustering component dominates for small separations, while lens-lens and/or density-lens components dominates for large separations along the direction of line of sight.
The lower panels in Figure 1 show the region where lensing contribution dominates in the case of the SDSS magnitude limits to illustrate the typical magnitude of correlations. Are those lensing signals detectable? The statistical uncertainty in estimating the correlation function is given by $`(\delta \xi )^2=\mathrm{\Omega }/(\delta \mathrm{\Omega }N_1N_2)`$ (Peebles, 1980), where $`N_1`$ and $`N_2`$ are numbers of object in the bin used for redshifts $`z_1`$ and $`z_2`$, respectively, $`\mathrm{\Omega }`$ is the solid angle subtended by the survey area, and $`\delta \mathrm{\Omega }`$ is the fraction in the bin used for angle $`\theta `$. To increase the signal to noise ratio, it is desirable to use large bins for $`\theta `$. To be specific, we consider a bin $`[1^{},10^{}]`$, and theoretical curves are integrated accordingly, so that $`\delta \mathrm{\Omega }100\pi [\text{arcmin}^2]`$. The effective scale of this bin is given by $`\theta _{\mathrm{eff}}=10^{}/\sqrt{2}`$ for $`\xi \theta ^1`$. In the SDSS, $`\mathrm{\Omega }\pi \text{[str]}1.2\pi \times 10^7[\text{arcmin}^2]`$, and the estimated numbers of galaxies and quasars are $`N_\mathrm{G}=10^6`$ and $`N_\mathrm{Q}=1.7\times 10^5`$, respectively. Assuming we take sufficiently large bins of redshifts (this choice is similar to considering angular correlation functions), the consequent estimates of the statistical error are given by $`5.0\times 10^4`$ for G-G, $`8.5\times 10^4`$ for G-Q, and $`2.9\times 10^3`$ for Q-Q correlations, which are plotted in lower panels. The S/N ratios turn out to be about $`10`$, $`1.3`$, and $`0.22`$ for G-G, G-Q, and Q-Q correlations, respectively.
Therefore, the weak lensing in 3D correlation function of galaxies in the SDSS is definitely detectable, and the detection of galaxy-QSO cross-correlation is marginal, while the quasar correlation by lensing is below the noise level in the SDSS. In order to detect the QSO-QSO lensing effect, the sample should be at least 5 times larger than the SDSS, or parameters $`\sigma _8`$, $`\mathrm{\Omega }_0`$, $`b`$, $`\alpha `$ should be larger than assumed values.
In summary, we have obtained a theoretical prediction of correlation function in redshift space, taking into account the effect of weak lensing, together with velocity distortions and cosmological distortions on a light-cone. Each effect contributes differently to the correlation function, and is realistically detectable. Our result provides a fundamental link between theoretical models and the observed correlation function in the 3D redshift survey data. Besides the determination of the power spectrum itself, various cosmological parameters, especially the bias parameter, can be estimated by proper likelihood analyses, including KL transform of the correlation matrix (Vogeley & Szalay, 1996; Matsubara et al., 2000). One may also be tempted to assume the cosmological parameters before analysing data. In which case, the error originated in choosing the wrong cosmological model is roughly given by the order of the redshift $`z`$ times the error of cosmological parameters, since the Alcock-Paczyński effect is roughly proportional to $`z`$ up to $`z=1`$-$`2`$.
I am grateful to Bhuvnesh Jain for many helpful discussions. I would like to thank Yasushi Suto and Alex Szalay for stimulating discussions. I acknowledge support from JSPS Postdoctoral Fellowships for Research Abroad.
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# An Intuitive Hamiltonian for Quantum Search
## 1 Introduction
Quantum algorithms can, in theory at least, solve useful problems faster than classical algorithms. Two primary families of quantum algorithms in this regard are algorithms for factoring and discrete log , and Grover’s search algorithms with quadratic speed-up .
There are many variations on Grover’s original algorithm—counting, starting with partial data, multiple targets, et cetera. The algorithm is also surprisingly robust; although the original algorithm uses the Walsh-Hadamard transform, essentially any unitary operator will do just as well . Starting with a simple condition on what transform is used, we will show how Grover’s algorithm arises from a particularly simple—almost naive—intuition about quantum algorithms. Our ideas also generalize to variants of Grover’s algorithm.
There have been good explanations in the literature of how and why fast quantum search works: the initial state is slowly rotated (in two complex dimensions) into the target state by repeatedly applying a special operator known as Grover’s iterate. Using a more physics-based approach, Farhi & Gutmann describe an “analog” version of quantum search by means of a simple, time-independent Hamiltonian which transforms any initial state $`|\sigma `$ into some prespecified target state $`|w`$ in optimal time, provided that $`|\sigma `$ and $`|w`$ are not orthogonal. Their analog algorithm rotates the initial state into the target state in the same time that Grover’s “digital” algorithm does, yet their rotation strays far from the intermediate states reached in the original algorithm by applying Grover’s iterate.
We show here that a simple, time-independent Hamiltonian for a system of qubits results in time evolution matching Grover’s iterate exactly. This Hamiltonian also provides a nice, simple insight into the workings of the algorithm that is new, to the best of our knowledge.
Our Hamiltonian bears some resemblance to that of Farhi & Gutmann, although ours was conceived independently. Ours differs from theirs in important respects, however, and may not be as plausible physically, but it does closely coincide with the iterations of Grover’s algorithm, and thus gives a much closer simulation of a digital quantum circuit by an analog process and vice versa. While Farhi & Gutmann’s Hamiltonian is appealing from a physical point of view, ours is appealing from an algorithmic perspective. The heuristic it gives for Grover’s algorithm suggests that other digital quantum algorithms might be found by first looking at analog versions. (More work on analog algorithms has been done recently, see for example.)
### 1.1 Structure of the Paper
We give mathematical preliminaries in Section 2, including a brief description of Grover’s algorithm as described in . In Section 3 we show how Grover’s algorithm arises from a simple-minded approach to quantum search. The operator we describe there corresponds directly to our Hamiltonian, which in Section 4 we compare with that of Farhi and Gutmann , and show how it generates Grover’s iterate. Most of our work was done independently of before it came to our attention, so our approach to the problem is different. In Section 5, we suggest a Hamiltonian-based approach to quantum algorithms in general, and we present open problems.
## 2 Preliminaries
We work with linear operators over a fixed $`N`$-dimensional Hilbert space. A standard norm on operators is defined as
$$|A|=\underset{|v|=1}{sup}|Av|,$$
where $`A`$ is an operator and $`||`$ on the right hand side is the standard hermitian norm on the Hilbert space. This norm on operators satisfies $`|AB||A||B|`$. Clearly, all unitary operators have unit norm. The exponential map on operators is defined as
$$e^A=I+A+\frac{A^2}{2!}+\frac{A^3}{3!}+\mathrm{},$$
(1)
where $`I`$ is the identity operator.
If $`A`$ is skew hermitian ($`A^{}=A`$), then $`e^A`$ is unitary. Conversely, for any unitary $`U`$ there is a skew hermitian $`A`$ such that $`U=e^A`$. As in the case with the exponential function on scalars, we also have
$$e^A=\underset{k\mathrm{}}{lim}\left(I+\frac{A}{k}\right)^k.$$
(2)
For an $`n`$-qubit system we assume the standard basis of states $`|i`$ indexed by classical bit configurations $`i\{0,1\}^n`$. We use lower-case Roman letters to label basis states, and lower-case Greek letters to label other (arbitrary) states in the Hilbert space.
### 2.1 Grover’s Search Algorithm
Here we briefly review Grover’s search algorithm. Fix an integer $`n`$ and let $`N=2^n`$. Let $`f`$ be a Boolean-valued function on $`n`$-bit strings such that $`f(w)=1`$ for exactly one $`w`$ (the target), $`0w<N`$ (identifying strings with integers). A simple version of Grover’s algorithm is to find $`w`$ via a quantum algorithm where inputs to $`f`$ are stored in $`n`$ qubits, and $`f`$ is available as a black box function (oracle) that can be queried by the algorithm. Alternately, we may assume that $`f(y)`$ is efficiently computable given $`y`$, and embed the computation of $`f`$ into the quantum circuit.
In this setting, Grover’s algorithm (as described in or ) uses three $`n`$-qubit unitary transforms:
1. an arbitrary, easy-to-compute $`U`$ such that $`w|U|00`$,
2. the selected inverter $`I_0=_{0i<N}(1)^{i=0}|ii|=I2|00|`$, and
3. the selected inverter $`I_w=I2|ww|`$.
(Here, the formula $`i=0`$ in the exponent stands for its numerical truth value—$`1`$ for true, $`0`$ for false.) These combine to form Grover’s iterate
$$G=UI_0U^1I_w.$$
(3)
By adjusting $`U`$ by an appropriate phase factor, we can assume that $`w|U|0=x`$ for some real $`x>0`$. This adjustment leaves $`G`$ unchanged.
Suppose $`f`$ is as above with $`w`$ unique such that $`f(w)=1`$. The algorithm starts in the state $`|0`$ (all qubits cleared), then $`U`$ is applied to get the state
$$|\psi =U|0.$$
(4)
Next, $`G`$ is applied repeatedly to $`|\psi `$, approximately $`\frac{\pi }{4x}`$ times. At this point, the system will be very close to the state $`|w`$, so when we now measure the qubits we get $`w`$ with high probability. Note that
$$I_w=\underset{i}{}(1)^{f(i)}|ii|,$$
so $`I_w`$ can be simulated easily given access to $`f`$ alone and some extra work qubits.
In Grover’s original presentation, $`U=U^1`$ is the Walsh-Hadamard transform on $`n`$ qubits, and so
$$|\psi =2^{n/2}\underset{i}{}|i,$$
whence, $`x=w|\psi =2^{n/2}`$, which yields the quadratic speed-up in the search.
## 3 Quantum Search Revisited
The point of this section is to show how one might stumble upon Grover’s algorithm by taking a simplistic, almost naive, approach to quantum search. The intuition here is not geometric, as it is with Jozsa ; rather, it is purely algorithmic in flavor.
We start with the basic observation that if $`A`$ is a skew hermitian operator ($`A^{}=A`$) and $`0<ϵ<<1`$, then $`I+ϵA`$ approximates $`e^{ϵA}`$, which is unitary. Therefore, $`I+ϵA`$ approximates a plausible step in a quantum computation. By (2), we can approximate the action of $`e^A`$ on a state by repeatedly applying $`I+ϵA`$ to the state roughly $`1/ϵ`$ times. The smaller $`ϵ`$ is, the better the approximation. (In general, it is not certain that $`e^{ϵA}`$ is renderable by a small quantum circuit; it will be in the present case, though.)
A simple example is when $`A=|ij||ji|`$ for some $`i,j\{0,1\}^n`$, $`ij`$. Applying $`I+ϵA`$ to a state $`|\phi =_i\alpha _i|i`$ gives
$$(I+ϵA)|\phi =|\phi +ϵ\alpha _j|iϵ\alpha _i|j.$$
The operator alters $`|\phi `$ (viewed as a column vector) by adding an $`ϵ`$ fraction of its $`j`$th component into its $`i`$th component, and in exchange, subtracting an $`ϵ`$ fraction of its $`i`$th component from its $`j`$th component. In a sense, we are moving probability amplitude from state $`|j`$ to state $`|i`$. With arbitrary $`A`$, this swap may take place between many pairs of components of $`|\phi `$ at once.
Suppose we are given $`n`$, $`N`$, $`f`$, and $`w`$ as in Section 2.1. We start in the state $`|\psi =N^{1/2}_i|i`$, which we would like to transform to the target state $`|w`$. A promising way to do this, given our considerations above, is to pile positive probability amplitude onto $`|w`$ while taking it away from all the other states evenly. The real skew symmetric operator that does this is
$$A=\left[\begin{array}{ccccccc}0& \mathrm{}& 0& 1& 0& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0& 1& 0& \mathrm{}& 0\\ 1& \mathrm{}& 1& 0& 1& \mathrm{}& 1\\ 0& \mathrm{}& 0& 1& 0& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0& 1& 0& \mathrm{}& 0\end{array}\right]$$
expressed in the $`\{|i\}`$ basis, where the nonzero entries are all in the $`w`$th row and $`w`$th column. We see that
$$I+ϵA=\left[\begin{array}{ccccccc}1& \mathrm{}& 0& ϵ& 0& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 1& ϵ& 0& \mathrm{}& 0\\ ϵ& \mathrm{}& ϵ& 1& ϵ& \mathrm{}& ϵ\\ 0& \mathrm{}& 0& ϵ& 1& \mathrm{}& 0\\ \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& \mathrm{}& 0& ϵ& 0& \mathrm{}& 1\end{array}\right].$$
The $`ϵ`$’s on row $`w`$ have the effect of giving probability amplitude to $`|w`$ while removing it from all the other states evenly (the column of $`ϵ`$’s). The probability amplitude of $`|w`$ gains at the expense of an $`ϵ`$ fraction of all the other probability amplitudes. From this it is clear that if we start in state $`|\psi `$, where all the probability amplitudes are equal, and apply $`I+ϵA`$ (for some small $`ϵ`$) the right number of times, eventually the state $`|w`$ will dominate.
We note that, using bracket notation,
$$A=\sqrt{N}(|w\psi ||\psi w|).$$
The operator $`iϵA`$ acts as a Hamiltonian for the time evolution of the system from $`|\psi `$ to $`|w`$. As we’ll see in the next section, for the right value of $`ϵ`$, $`e^{ϵA}`$ is exactly two applications of Grover’s iterate.
## 4 Hamiltonians
In this section, we give a Hamiltonian for Grover’s algorithm, that is, an operator $`H`$ such that $`e^{iHt}`$ follows the course of the algorithm as $`t`$ increases. It is clear both by geometric considerations and by the last section that such an operator must exist. $`H`$ is analogous to a previous Hamiltonian $`H^{}`$ for quantum search found by Farhi & Gutmann which does not match Grover’s algorithm. We first briefly describe their results, then describe our Hamiltonian using their framework.
### 4.1 Farhi & Gutmann’s Hamiltonian
We are given $`n`$, $`N`$, $`f`$ and $`w`$ as above. Farhi & Gutmann describe a physical, analog way to do quantum search by first assuming that a Hamiltonian
$$H_w=E|ww|$$
is available that distinguishes the target state $`|w`$ from all others by giving it some positive energy $`E`$ (the other basis states have energy $`0`$). Let $`|\sigma `$ be some arbitrary unit vector in the Hilbert space (the “start” state). We assume $`|\sigma `$ is easy to prepare, so for example, $`|\sigma `$ may be $`|\psi `$ of equation (4). The goal is to evolve from $`|\sigma `$ into $`|w`$. To search for the state $`|w`$, we are allowed to add some “driver” Hamiltonian $`H_D`$ to $`H_w`$, provided that $`H_D`$ does not depend on the actual value of $`w`$ at all. They choose $`H_D=E|\sigma \sigma |`$, so their Hamiltonian is
$$H^{}=H_D+H_w=E\left(|\sigma \sigma |+|ww|\right),$$
where $`E`$ is some arbitrary positive value in units of energy. If $`|\sigma `$ and $`|w`$ are not orthogonal, then we can assume as before that $`\sigma |w=w|\sigma =x`$ for some $`x>0`$ by adjusting $`|\sigma `$ by an appropriate phase factor.
Applying $`e^{iH^{}t}`$ to the start state $`|\sigma `$ gives the time-evolution of the system,<sup>1</sup><sup>1</sup>1The evolution of a quantum system under a time-independent Hamiltonian $`H`$ is actually $`e^{iHt/\mathrm{}}`$. We choose units so that $`\mathrm{}=1`$, and so $`Et`$ is a unitless quantity. which stays in the two-dimensional subspace spanned by $`|\sigma `$ and $`|w`$. Restricting our attention to this subspace, it is easy to see that $`H^{}`$ has eigenvalues $`E(1\pm x)`$ with corresponding eigenvectors
$`|+^{}`$ $`=`$ $`(2+2x)^{1/2}(|\sigma +|w),`$
$`|^{}`$ $`=`$ $`(22x)^{1/2}(|\sigma |w).`$
A straightforward calculation yields
$$e^{iH^{}t}|\sigma =e^{iEt}\left[\mathrm{cos}(xEt)|\sigma i\mathrm{sin}(xEt)|w\right].$$
(5)
When $`t=\pi /(2Ex)`$, we have
$$e^{iH^{}t}|\sigma =ie^{i\pi /(2x)}|w$$
as desired.
Farhi & Gutmann observe that if the unit vector $`|\sigma `$ is chosen at random, then the expected value of $`x`$ is $`N^{1/2}`$, making $`t=O(N^{1/2}/E)`$. For constant $`E`$, this time is the same order of magnitude as Grover’s algorithm. They show that their time evolution is optimal up to an order of magnitude for any $`w`$-independent driver Hamiltonian $`H_D`$, even one that varies with time.
### 4.2 Another Hamiltonian
The time evolution of the system according to $`H^{}`$ strays far from the intermediate steps Grover’s algorithm. There surely is a Hamiltonian, however, whose time evolution matches the steps of Grover’s algorithm exactly, since each step of Grover’s algorithm essentially amounts to a rotation in a two-dimensional space. We show that this Hamiltonian can be described very simply: the operator $`iϵA`$ mentioned at the end of Section 3 is exactly the Hamiltonian in question, for an appropriate value of $`ϵ`$ which we will calculate.
The fact that Grover’s iterate can be rendered by a small quantum circuit then tells us that our intuition of Section 3 is justified: the incremental application of $`I+ϵA`$ indeed corresponds to a legitimate quantum algorithm.
Given $`n`$, $`N`$, $`f`$, $`|w`$, $`|\sigma `$, $`H_w`$ and $`H_D`$ as above, with $`w|\sigma =x>0`$, we define the Hamiltonian
$$H=\frac{2i}{E}[H_w,H_D]=2iEx(|w\sigma ||\sigma w|).$$
The rest of this section is devoted to proving the following
###### Theorem 1
Assume the special case where $`|\sigma =|\psi `$ and $`E=1`$. Restricted to the $`(|\sigma ,|w)`$-plane, $`e^{iH}`$ approximates Grover’s iterate $`G`$ to within $`O(N^{3/2})`$ in norm. In fact, $`e^{iHt_0}`$ exactly matches $`G`$ where
$$t_0=\frac{\pi 2\mathrm{arccos}x}{2x\sqrt{1x^2}}.$$
(6)
On the whole Hilbert space, $`e^{2iH}`$ approximates $`G^2`$ to within $`O(N^{3/2})`$, and $`e^{2iHt_0}=G^2`$.
For the moment, we allow $`E`$ to be any positive value and $`|\sigma `$ an arbitrary unit vector with $`0<w|\sigma =x<1`$. Restricting our attention to the subspace spanned by $`|\sigma `$ and $`|w`$, and letting $`\theta =\mathrm{arccos}x`$, the eigenvalues of $`H`$ are seen to be $`\pm \frac{1}{2}E\mathrm{sin}2\theta `$ with corresponding eigenvectors
$$|\pm =\frac{1}{\sqrt{2}\mathrm{sin}\theta }\left(e^{\pm i\theta }|\sigma |w\right).$$
Setting $`\eta =E\mathrm{sin}2\theta =2Ex\mathrm{sin}\theta `$, a routine calculation shows that
$`e^{iHt}|\sigma `$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}\theta }}\left[\mathrm{sin}(\theta \eta t)|\sigma +\mathrm{sin}(\eta t)|w\right],`$ (7)
$`e^{iHt}|w`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}\theta }}\left[\mathrm{sin}(\eta t)|\sigma +\mathrm{sin}(\theta +\eta t)|w\right].`$ (8)
If $`x`$ is small, $`\theta `$ will be close to $`\pi /2`$. For $`t=\theta /\eta =\theta /(E\mathrm{sin}2\theta )\pi /(2Ex)`$ we have $`e^{iHt}|\sigma =|w`$. That is, the system finds the target state in roughly the same time as with $`H^{}`$.
Comparing (5) and (7), we see that the quantum system evolves significantly differently under the two Hamiltonians $`H^{}`$ and $`H`$—by more than just a global phase factor. We now show how the latter evolution, run for a short time interval, matches a single step of Grover’s algorithm (one application of $`G`$). We now assume $`E=1`$ and $`|\sigma =|\psi =U|0`$ given by equation (4), with $`G`$ given by (3). We again set $`x=\sigma |w=\psi |w=\mathrm{cos}\theta >0`$, for some $`0<\theta <\pi /2`$.
We can express $`G`$ in the basis $`|\sigma ,|w`$:
$`G`$ $`=`$ $`U(I2|00|)U^1(I2|ww|)`$
$`=`$ $`(I2U|00|U^{})(I2|ww|)`$
$`=`$ $`(I2|\sigma \sigma |)(I2|ww|)`$
$`=`$ $`I+2|\sigma \sigma \left|+2\right|ww\left|4x\right|\sigma w|,`$
whence
$`G|\sigma `$ $`=`$ $`(14x^2)|\sigma +2x|w,`$
$`G|w`$ $`=`$ $`2x|\sigma +|w.`$
In view of (7), we solve the equation
$$\frac{\mathrm{sin}(\eta t)}{\mathrm{sin}\theta }=2x=2\mathrm{cos}\theta $$
for $`t`$ to get the solution
$$t_0=\frac{\pi 2\theta }{\eta }=\frac{\pi 2\theta }{\mathrm{sin}2\theta }=\frac{\pi 2\mathrm{arccos}x}{2x\sqrt{1x^2}},$$
(9)
as in (6). It is then easy to check that $`e^{iHt_0}|\sigma =G|\sigma `$ and that $`e^{iHt_0}|w=G|w`$.
Let $`S`$ be the subspace spanned by $`|\sigma `$ and $`|w`$ and let $`S^{}`$ be its orthogonal complement. Let $`P`$ be the orthogonal projection onto $`S^{}`$. We have just shown that $`e^{iHt_0}=G`$ restricted to $`S`$. For $`|\alpha S^{}`$, clearly $`G|\alpha =|\alpha `$, while $`e^{iHt_0}`$ leaves $`S^{}`$ pointwise fixed. Thus,
$$e^{iHt_0}=G+2P,$$
and since $`GP=PG=P=P^2`$, we have $`e^{2iHt_0}=G^2`$ as expected.
#### Remark.
By adding $`\frac{\pi }{t_0}P`$ to $`H`$, we get a slightly more complicated Hamiltonian $`\stackrel{~}{H}`$ such that $`e^{i\stackrel{~}{H}t_0}=G`$ on the whole Hilbert space.
Finally, we show how close $`t_0`$ is to $`1`$, assuming $`x<<1`$. Expanding $`t_0`$ as a power series in $`x`$, we get
$$t_0=1+\frac{2}{3}x^2+O(x^4),$$
and thus
$`\left|e^{iHt_0}e^{iH}\right|`$ $`=`$ $`\left|e^{iH(2x^2/3+O(x^4))}I\right|`$
$`=`$ $`\left|{\displaystyle \frac{2}{3}}iHx^2+O(x^4)H\right|`$
$`=`$ $`{\displaystyle \frac{2}{3}}x^3\sqrt{1x^2}+O(x^5)`$
$`=`$ $`{\displaystyle \frac{2}{3}}x^3+O(x^5).`$
The second equation comes from expanding the exponential as a power series. The third equation holds because $`|H|=|\frac{1}{2}\mathrm{sin}2\theta |=x\sqrt{1x^2}`$. When $`x=N^{1/2}`$, we see that $`e^{iH}`$ comes within $`O(N^{3/2})`$ of $`G+2P`$ in norm, and thus $`e^{2iH}`$ comes within $`O(N^{3/2})`$ of $`G^2`$.
###### Corollary 2
If $`t=\frac{\pi }{4}N^{1/2}`$ then $`e^{iHt}|\sigma =|w+O(1/N)`$.
### 4.3 Discussion
Farhi and Gutmann show that the Hamiltonian $`H^{}`$ finds the target in optimal time in the following sense: no other Hamiltonian of the same form—that is, $`H_D+E|ww|`$ where $`H_D`$ has no special dependence on $`w`$—can find $`|w`$ any faster, even if $`H_D`$ is allowed to depend on time. Our Hamiltonian is clearly not of this form, so their lower bounds aren’t directly applicable here. Indeed, it is only by the lower bounds shown for “digital” quantum search that we know that our Hamiltonian is optimal to simulate a small digital quantum circuit. It is an interesting question whether one can deduce the same lower bound by more direct means.
## 5 Further Research and Open Problems
We have seen how Grover’s algorithm can be described much more simply using a Hamiltonian than directly with unitary operators. We don’t present details them here, but variants of Grover’s original algorithm also admit simple Hamiltonian descriptions. There may be new, yet unknown quantum algorithms which are more easily described with Hamiltonians than with unitary operators, and which may indeed be first discovered by their Hamiltonians.
There are two principal challenges to fashioning a new quantum algorithm via a Hamiltonian:
1. finding an appropriate, and hopefully intuitive, Hamiltonian for the problem at hand, and
2. deciding how (or if) the time evolution governed by such a Hamiltonian approximates a true (digital) quantum algorithm given by a small quantum circuit.
In the case of Grover’s algorithm considered here, we were fortunate to achieve both goals. Grover’s algorithm came first, however, so we knew what to shoot for. Even so, the intuition provided in Section 3 may be useful for constructing new algorithms, or at least viewing other existing algorithms from a different angle.
Our results thus point to an important general question: when, given a Hamiltonian on a system of qubits, can the corresponding time evolution be simulated (even approximately) by a small quantum circuit? Is there an easy criterion, based on the structure of the Hamiltonian itself? Such a criterion would provide a new way to prototype new quantum algorithms via Hamiltonians.
Can Farhi & Gutmann’s original $`H^{}`$ be simulated efficiently by a quantum circuit?
Recently, Farhi, et al. show how to solve certain instances of SAT with slowly time-dependent Hamiltonians (adiabatic evolution). Their results provide good physical intuition. Is there a corresponding algorithmic intuition?
Can one find an intuitive Hamiltonian for a quantum factoring algorithm?
## 6 Acknowledgments
I would like to thank Frederic Green and Steven Homer for helpful discussions.
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# 1 Introduction
## 1 Introduction
It is well known that there are new sources of CP violation in supersymmetric theories which arise from the soft SUSY breaking sector of the theory. The normal size of such phases is O(1) and an order of magnitude estimate shows that such large phases would lead to a conflict with the current experimental limits on the electron and on the neutron electric dipole moment. The conventional ways suggested to avoid this conflict is either to assume that the phases are small or that the SUSY spectrum is heavy. However, recently it was demonstrated that this need not to be the case and indeed there could be consistency with experiment even with large CP violating phases and a light spectrum due to an internal cancellation mechanism among the various contributions to the EDMs. The above possibility has led to a considerable further activity and the effects of large CP phases under the cancellation mechanism have been investigated in dark matter with the EDM constraints, in $`g_\mu 2`$ and in other low energy physics phenomena.
In this paper we investigate the effects of large CP violating phases on nucleon stability in supersymmetric grand unification with baryon and lepton number violating dimension five operators. The main result of this analysis is that the dressing loop integrals that enter in the supersymmetric proton decay analysis are modified due to the effect of the large CP violating phases. The CP effects on the proton life time are most easily exhibited by considering $`R_\tau `$ defined in Eq.(22) which is the ratio of the p lifetime with phases and without phases. $`R_\tau `$ is largely independent of the GUT structure which cancels out in the ratio. Since the dressing loop integrals enter in the proton decay lifetime in both GUT and string models which contain the baryon and the lepton number violating dimension five operators, the phenomena of CP violating effects on the proton lifetime should hold for a wide range of models both of GUT and of string variety. However, for concreteness we will consider first the simplest SU(5) supersymmetric grand unified model, and then consider a non-minimal extension. As discussed above similar analyses should hold for a wider class of models and so what we do below should serve as an illustration of the general idea of the effect of large CP phases on the proton lifetime.
The outline of the paper is as follows: In Sec.2 we give a theoretical analysis of the effects of CP violating phases on proton decay in the minimal supersymmetric SU(5) model for specificity. In Sec.3 we discuss the numerical effects of the CP violating phases on $`R_\tau `$ under the EDM constraints. A non-minimal extension is also discussed and an analysis of the uncertainties in the predictions of the proton life time due to uncertainties in the quark masses, in $`\beta _p`$ and in the KM matrix elements is given. Conclusions are given in Sec.4.
## 2 Theoretical analysis of CP violating phases on proton decay in supersymmetric GUTs
In the minimal supergravity unified model (mSUGRA) the soft SUSY breaking can be parameterized by $`m_0`$, $`m_{\frac{1}{2}}`$, $`A_0`$, and $`\mathrm{tan}\beta `$, where $`m_0`$ is the universal scalar mass, $`m_{\frac{1}{2}}`$ is the universal gaugino mass, $`A_0`$ is the universal trilinear coupling all taken at the GUT scale, and $`\mathrm{tan}\beta =<H_2>/<H_1>`$ is the ratio of the Higgs VEVs where $`H_2`$ gives mass to the up quark and $`H_1`$ gives mass to the down quark and the lepton. In addition, the effective theory below the GUT scale contains the Higgs mixing parameter $`\mu `$ which enters in the superpotential in the term $`\mu H_1H_2`$. In the presence of CP violation one finds that the minimal model contains two independent CP violating phases which can be taken to be $`\theta _\mu `$, which is the phase of $`\mu `$ and $`\alpha _{A_0}`$ which is the phase of $`A_0`$.
For more general situations when one allows for non-universalities, the soft SUSY breaking sector of the theory brings in more CP violating phases. Thus unlike the case of mSUGRA here the $`U(1)\times SU(2)\times SU(3)`$ gaugino masses $`\stackrel{~}{m}_i`$ (i=1,2,3) can have arbitrary phases, i.e.,
$$\stackrel{~}{m}_i=|\stackrel{~}{m}_i|e^{i\xi _i}(i=1,2,3)$$
(1)
While in the universal case a field redefinition can eliminate the common phase of the gaugino masses, here one finds that the difference of the gaugino phases does persist in the low energy theory and in fact is found to be a useful tool in arranging for the cancellation mechansim to work for the satisfaction of the EDMs. In the following analysis we carry out an analysis of the proton decay with the most general allowed set of CP violating phases. The definition of the mass matrices for charginos, neutralinos and for squarks and sleptons have been explicitly exhibited in Ref. and we refer the reader to this paper for details. The focus of the present work is to analyze the effects of CP violating phases on p decay and to estimate its size. For the sake of concreteness we begin with a discussion of the simplest grand unification model, i.e., the minimal SU(5) model. However, the technique discussed here to include the CP effects on p decay can be used to anlayse the CP violating effects for any supersymmetric unified model with baryon and lepton number violating dimension five operators. This class includes string models.
As mentioned above we consider for concreteness and simplicity the minimal SU(5) model whose matter interactions are given by
$$W_Y=\frac{1}{8}f_{1ij}ϵ_{uvwxy}H_1^uM_i^{vw}M_j^{xy}+f_{2ij}\overline{H}_{2u}\overline{M}_{iv}M_j^{uv}$$
(2)
where $`M_u,M^{uv}`$ are the $`\overline{5},10`$ plet representations of SU(5), and $`H_1,H_2`$ are the $`\overline{5},5`$ of SU(5). After the breakdown of the GUT symmetry and integration over the Higgs triplet fields the effective dimension five interactions below the GUT scale which governs p decay is given by
$$_{5L}=\frac{1}{M}ϵ_{abc}(Pf_1^uV)_{ij}(f_2^d)_{kl}(\stackrel{~}{u}_{Lbi}\stackrel{~}{d}_{Lcj}(\overline{e}_{Lk}^c(Vu_L)_{al}\overline{\nu }_k^cd_{Lal})+\mathrm{})+H.c.$$
(3)
$$_{5R}=\frac{1}{M}ϵ_{abc}(V^{}f^u)_{ij}(PVf^d)_{kl}(\overline{e}_{Ri}^cu_{Raj}\stackrel{~}{u}_{Rck}\stackrel{~}{d}_{Rbl}+\mathrm{})+H.c.$$
(4)
where $`_{5L}`$ and $`_{5R}`$ are the LLLL and RRRR lepton and baryon number violating dimension 5 operators, V is the CKM matrix and $`f_i`$ are related to quark masses, and $`P_i`$ appearing in Eqs.(3) and (4) are the generational phases given by $`P_i=(e^{i\gamma _i})`$, $`_i\gamma _i=0`$ (i=1,2,3).
The baryon and the lepton number violating dimension five operators must be dressed by the chargino, the gluino and the neutralino exchanges to generate effective baryon and lepton number violating dimension six operators at low energy (Some examples of dressing loop diagrams are given in Fig.1). It is in this process of dressing of the dimension five operators that the CP violating phases of the soft SUSY breaking sector enter in the proton decay amplitude. The CP phases enter the dressings in two ways, via the mass matrices of the charginos, the neutralinos and the sfermions, and via the interaction vertices. Taking account of this additional complexity, the analysis for computing the proton decay amplitudes follows the usual procedure. Thus to dress the dimension five operators the squark and slepton fields must be eliminated in terms of their sources. As an example, the up squarks in the presence of CP violating phases can be eliminated using the relations
$`\stackrel{~}{u}_{iL}=2{\displaystyle [\mathrm{\Delta }_{ui}^LL_{ui}+\mathrm{\Delta }_i^{LR}R_{ui}]}`$
$`\stackrel{~}{u}_{iR}=2{\displaystyle [\mathrm{\Delta }_{ui}^RR_{ui}+\mathrm{\Delta }_i^{RL}L_{ui}]}`$ (5)
where
$$L_{ui}=\frac{\delta L_I}{\delta \stackrel{~}{u}_{iL}^{}},R_{ui}=\frac{\delta L_I}{\delta \stackrel{~}{u}_{iR}^{}}$$
(6)
Here $`L_I`$ is the sum of fermion-sfermion-gluino, fermion-sfermion-chargino and fermion-sfermion-neutralino interactions and
$`\mathrm{\Delta }_{ui}^L=[|D_{ui11}|^2\mathrm{\Delta }_{ui1}+|D_{ui12}|^2\mathrm{\Delta }u_{i2}]`$
$`\mathrm{\Delta }_{ui}^R=[|D_{ui21}|^2\mathrm{\Delta }_{ui1}+|D_{ui22}|^2\mathrm{\Delta }u_{i2}]`$ (7)
and
$`\mathrm{\Delta }_{ui}^{LR}=D_{ui11}D_{ui12}[\mathrm{\Delta }_{ui1}\mathrm{\Delta }u_{i2}]`$
$`\mathrm{\Delta }_{ui}^{RL}=D_{ui11}D_{ui12}^{}[\mathrm{\Delta }_{ui1}\mathrm{\Delta }u_{i2}].`$ (8)
Here $`\stackrel{~}{u}_{i1}`$ and $`\stackrel{~}{u}_{i2}`$ are the squark mass eigenstates for the squark flavors $`u_{i1}`$ and $`u_{i2}`$ and $`\mathrm{\Delta }_{u_i1}`$ and $`\mathrm{\Delta }_{u_i2}`$ are the corresponding propagators, and $`D_{ui}`$ is the diagonalizing matrix for the $`\stackrel{~}{u}_i`$ squarks, i.e.,
$$D_{u_i}^{}M_{\stackrel{~}{u}_i}^2D_{u_i}=diag(M_{\stackrel{~}{u}_{i1}}^2,M_{\stackrel{~}{u}_{i2}}^2)$$
(9)
We note the special arrangement of the complex quantities and their complex conjugates in Eqs.7 and 8. Specifically we note that while in the absence of CP phases $`\mathrm{\Delta }_{ui}^{LR}`$ = $`\mathrm{\Delta }_{ui}^{RL}`$ this is not the case in the presence of CP phases and in general one has $`\mathrm{\Delta }_{ui}^{LR}\mathrm{\Delta }_{ui}^{RL}`$ as is seen from Eqs.(8). $`L_{ui}`$ and $`R_{ui}`$ defined by Eq.(6) receive contributions from the chargino, the neutralino and the gluino exchanges.
Following the standard procedure one obtains the effective dimension six operators for the baryon and the lepton number violating interaction arising from dressing of the dimension five operators. From this effective interaction one obtains the proton lifetime decay widths for various modes using the effective Lagrangian methods. We limit ourselves here to the dominant decay mode $`p\overline{\nu }_iK^+`$. Including the CP violating effects the decay width for this process is given by
$`\mathrm{\Gamma }(p\overline{\nu }_iK^+)={\displaystyle \frac{\beta _p^2m_N}{M_{H_3}^232\pi f_\pi ^2}}(1{\displaystyle \frac{m_K^2}{m_N^2}})^2|𝒜_{\nu _iK}|^2A_L^2(A_S^L)^2`$
$`|(1+{\displaystyle \frac{m_N(D+3F)}{3m_B}})(1+𝒴_i^{tk}+(e^{i\xi _3}𝒴_{\stackrel{~}{g}}+𝒴_{\stackrel{~}{Z}})\delta _{i2}+{\displaystyle \frac{A_S^R}{A_S^L}}𝒴_1^R\delta _{i3})`$
$`+{\displaystyle \frac{2}{3}}{\displaystyle \frac{m_N}{m_B}}D(1+𝒴_3^{tk}(e^{i\xi _3}𝒴_{\stackrel{~}{g}}𝒴_{\stackrel{~}{Z}})\delta _{i2}+{\displaystyle \frac{A_S^R}{A_S^L}}𝒴_2^R\delta _{i3})|^2`$ (10)
where
$$𝒜_{\nu _iK}=(\mathrm{sin}2\beta M_W^2)^1\alpha _2^2P_2m_cm_i^dV_{i1}^{}V_{21}V_{22}[(\stackrel{~}{c};\stackrel{~}{d}_i;\stackrel{~}{W})+(\stackrel{~}{c};\stackrel{~}{e}_i;\stackrel{~}{W})]$$
(11)
In the above $`A_L(A_S)`$ are the long (short) suppression factors, D,F, $`f_\pi `$ are the effective Lagrangian parameters, and $`\beta _p`$ is defined by $`\beta _pU_L^\gamma =ϵ_{abc}ϵ_{\alpha \beta }<0|d_{aL}^\alpha u_{bL}^\beta u_{cL}^\gamma |p>`$ where $`U_L^\gamma `$ is the proton wavefunction. Theoretical determinations of $`\beta _p`$ lie in the range $`0.0030.03GeV^3`$. Perhaps the more reliable estimate is from lattice gauge calculations which gives $`\beta _p=(5.6\pm 0.5)\times 10^3GeV^3`$.
Aside from the explicit CP phases via the exponential factor $`e^{i\xi _3}`$ in Eq.(10), CP effects enter dominantly in $``$’s which are the dressing loop integrals. For the chargino exchange in the presence of CP violating phases one has
$`(\stackrel{~}{u}_i;\stackrel{~}{d}_j;\stackrel{~}{W})=32\pi ^2i{\displaystyle \underset{A=1,2}{}[\mathrm{\Delta }_{uai}^LS_{A1}^{}\mathrm{\Delta }_{uai}^{LR}ϵ_i^uS_{A2}^{}]}`$
$`\stackrel{~}{G}_A[\mathrm{\Delta }_{dj}^LU_{A1}^{}\mathrm{\Delta }_{dj}^{LR}ϵ_j^dU_{A2}^{}]`$ (12)
Here $`\stackrel{~}{G}_A`$ (A=1,2) are the propagators for the chargino mass eigenstates and the matrices U and S enter in the biunitary transformations to diagonlize the chargino mass matrix $`M_C`$ such that
$$U^{}M_CS^1=diag(\stackrel{~}{m}_{\chi _1^+},\stackrel{~}{m}_{\chi _2^+})$$
(13)
In Eq.(10) the quantities $`𝒴_i^{tk}`$ are the corrections due to the chargino exchanges involving third generation squarks, $`𝒴_{\stackrel{~}{g}}`$ is the contribution from the gluino exchange, $`𝒴_{\stackrel{~}{Z}}`$ is the contribution from the neutralino exchange, and $`𝒴_i^R`$ are the contributions from the dressing of the RRRR dimension 5 operators.
The gluino exchange contribution $`𝒴_{\stackrel{~}{g}}`$ is given by
$$𝒴_{\stackrel{~}{g}}=\frac{4}{3}\frac{P_1}{P_2}\frac{\alpha _3}{\alpha _2}\frac{m_uV_{11}}{m_cV_{21}V_{21}^{}V_{22}}\frac{(\stackrel{~}{u};\stackrel{~}{d};\stackrel{~}{g})(\stackrel{~}{d};\stackrel{~}{d};\stackrel{~}{g})}{(\stackrel{~}{c};\stackrel{~}{s};\stackrel{~}{W})+(\stackrel{~}{c};\stackrel{~}{\mu };\stackrel{~}{W})}$$
(14)
The function $``$ is defined by $`(\stackrel{~}{u};\stackrel{~}{d};\stackrel{~}{g})=f(\stackrel{~}{m}_u;\stackrel{~}{m}_d;\stackrel{~}{m}_g)`$ where f is defined by Eq.(19) below. The contributions $`𝒴_i^R`$ from the dresssing of the RRRR dimension five operators are given by
$$𝒴_1^R=\frac{P_1}{P_2}\frac{m_tm_dV_{11}V_{32}V_{33}^{}}{m_cm_bV_{21}V22V_{31}^{}}\frac{𝒬(\stackrel{~}{\tau };\stackrel{~}{t};\stackrel{~}{W})}{(\stackrel{~}{c};\stackrel{~}{b};\stackrel{~}{W})+(\stackrel{~}{c};\stackrel{~}{\tau };\stackrel{~}{W})}$$
(15)
and by
$$𝒴_2^R=\frac{P_1}{P_2}\frac{m_tm_sV_{31}V_{12}V_{33}^{}}{m_cm_bV_{21}V22V_{31}^{}}\frac{𝒬(\stackrel{~}{\tau };\stackrel{~}{t};\stackrel{~}{W})}{(\stackrel{~}{c};\stackrel{~}{b};\stackrel{~}{W})+(\stackrel{~}{c};\stackrel{~}{\tau };\stackrel{~}{W})}$$
(16)
where $`𝒬`$’s are defined as follows:
$$𝒬(\stackrel{~}{\tau };\stackrel{~}{t};\stackrel{~}{W})=32\pi ^2i(\frac{m_\tau }{\sqrt{2}M_W\mathrm{cos}\beta })\underset{A=1,2}{}\mathrm{\Delta }_\tau ^RU_{A2}^{}[ϵ_t\mathrm{\Delta }_{\stackrel{~}{t}}^RS_{A2}^{}\mathrm{\Delta }_{\stackrel{~}{t}}^{RL}S_{A1}^{}]\stackrel{~}{G}_A$$
(17)
The dressing loop integrals can be expressed in terms of the basic integral
$$f(\mu _i,\mu _j,\mu _k)=16\pi ^2i\mathrm{\Delta }_i\mathrm{\Delta }_j\stackrel{~}{G}_k$$
(18)
where
$$f(\mu _i,\mu _j,\mu _k)=\frac{\mu _k}{\mu _j^2\mu _k^2}[\frac{\mu _j^2}{\mu _i^2\mu _j^2}ln\frac{\mu _i^2}{\mu _j^2}\frac{\mu _k^2}{\mu _i^2\mu _k^2}ln\frac{\mu _i^2}{\mu _k^2}]$$
(19)
In the limit of no CP violation the analysis limits correctly to the previous results which do not include CP violating effects.
## 3 Numerical effects of CP violating phases on proton decay
We discuss now the numerical size of the effects of CP phases on the proton lifetime. As is obvious from our discussion above the proton lifetime is highly model dependent. Specifically there are two main factors that govern the lifetime of the proton. One of these depends on the nature of the GUT sector, i.e., if the GUT group is SU(5), SO(10), E<sub>6</sub>,..etc and on the nature of the GUT interaction, e.g., on the GUT Higgs structure, while the second factor that controls proton decay is the sparticle spectrum and the sparticle interactions that enter in the dressing loop integrals.
In the simplest SU(5) GUT model with two Higgs multiplets $`H_1`$ and $`H_2`$, GUT physics enters mainly via the Higgs triplet mass $`M_{H_3}`$. The proton decay lifetime is significantly affected if we were to change the GUT structure. Thus, for example, if one had many Higgs triplets, $`H_i`$ and $`\overline{H}_i`$ (i=1,..,n) where only the triplets $`H_1`$ and $`\overline{H}_1`$ couple with matter, i.e.,
$$W_3^{triplet}=\overline{H}_1J+\overline{J}H_1+\overline{H}_iM_{ij}H_j$$
(20)
then the effective interaction on eliminating the Higgs triplets is
$$W_4^{eff}=\overline{J}(M_{11}^1)J$$
(21)
Here one finds that the effective Higgs triplet mass is $`M_{H_3}^{eff}=((M^1)_{11})^1`$. The model with $`M_{H_3}`$ and $`M_{H_3}^{eff}`$ would have very similar CP violating effects for the same low energy sparticle spectrum. The reason is that the CP effects are largely governed by the nature of the low energy physics, e.g., the sparticle mass spectrum and the couplings of the sparticles with matter. Thus we expect similar size CP violating effects in models with different GUT structures but with similar size sparticle spectrum
To discuss the CP violating effects on the proton lifetime it is useful to consider the ratio $`R_\tau (p\overline{\nu }+K^+)`$ defined by
$$R_\tau (p\overline{\nu }+K^+)=\tau (p\overline{\nu }+K^+)/\tau _0(p\overline{\nu }+K^+)$$
(22)
Here $`\tau (p\overline{\nu }+K^+)`$ is the proton lifetime with CP violating phases and $`\tau _0(p\overline{\nu }+K^+)`$ is the lifetime without CP phases. This ratio is largly model independent. Thus most of the model dependent features such as the nature of the GUT or the string model would be contained mostly in the front factors such as the Higgs triplet mass, the quark masses, the $`A_S`$ and $`A_L`$ suppression factors all of which cancel out in the ratio. Similarly the quantity $`\beta _p`$ which is poorly known cancels out in the ratio as do the KM matrix elements.
We analyze $`R_\tau `$ under the constraints that CP violating phases obey the experimental limits on the electron and the neutron EDMs. For the electron and for the neutron the current experimental limits are
$$|d_e|<4.3\times 10^{27}ecm,|d_n|<6.3\times 10^{26}ecm$$
(23)
We are interested in the effects of large phases on the proton lifetime and for these to satisfy the EMD constraints we use the cancellation mechanism. In Fig.2 we present five cases where for different inputs the electron EDM is plotted as a function of $`\theta _\mu `$. An analysis of the neutron EDM for the same input is given in Fig.3. One finds the cancellation mechanism produces several regions where the EDM constraints are satisfied. In Fig.4 we give a plot of $`R_\tau `$ for the same set of inputs as in Figs. 2 and 3. The analysis shows that $`R_\tau `$ is a sensitive function of the CP phase $`\theta _\mu `$ and variations of a factor of around 2 can occur. We also note that both a suppression as well as an enhancement of the proton lifetime can occur as a consequence of the CP violating effects. Interestingly the largest CP effects on $`R_\tau `$ occur here at the points of maximum cancellation in the EDMs as may be seen by a comparison of Figs.2, 3 and 4. The variations in $`R_\tau `$ due to the phases arise because of constructive and destructive interference between the exchange contributions of chargino 1 ($`\chi _1^+`$) and chargino 2 ($`\chi _2^+`$) (see Fig.1). We give an illustration of this phenomenon in Table 1. The analysis of Table 1 exhibits the cancellation in the imaginary part of the amplitude for the decay process $`p\overline{\nu }+K^+`$ from chargino 1 and chargino 2, and this cancellation leads to an enhancement in the p lifetime ratio for this case.
It is possible to promote each of the cancellation points in Figs.2 and 3 into a trajectory in the $`m_0m_{\frac{1}{2}}`$ plane by scaling upwards by a common scale transformation
$$m_0\lambda m_0,m_{\frac{1}{2}}\lambda m_{\frac{1}{2}}$$
(24)
The size of the sparticle spectrum depends on the scale $`\lambda `$. In general, the larger the value of $`\lambda `$ the heavier is the sparticle spectrum and correspondingly smaller is the CP effect on the dressing loop as demonstrated in Fig.5. For the minimal SU(5) case one needs a relatively heavy spectrum with some of the sparticle masses $`1`$ TeV to stabilize the proton which has the current experimental limit for the $`p\overline{\nu }K`$ decay mode of $`\tau (p\overline{\nu }K)>5.5\times 10^{32}`$ yr. Because of the heaviness of the sparticle spectrum, the CP effects for the minimal SU(5) model are typically small, i.e., of order only a few percent. Larger CP effects can occur in non-minimal models where one has several Higgs triplets. Thus we consider an example where one has two pairs of heavy Higgs triplets with the Higgs triplet mass matrix given by
$$\left(\begin{array}{cc}0& a\mathrm{\Lambda }\\ a\mathrm{\Lambda }& M_2\end{array}\right)$$
(25)
Such a structure can arise, for example, in an SO(10) model with two $`\mathrm{𝟏𝟎}^{}𝐬`$ of Higgs and a 45 of Higgs with a superpotential of the type $`W_H=M_2\mathrm{𝟏𝟎}_{\mathrm{𝟐}𝐇}^\mathrm{𝟐}+\mathrm{\Lambda }\mathrm{𝟏𝟎}_{\mathrm{𝟏}𝐇}\mathrm{𝟒𝟓}_𝐇\mathrm{𝟏𝟎}_{\mathrm{𝟐}𝐇}`$. After the $`\mathrm{𝟒𝟓}`$ of Higgs develops a VEV $`<\mathrm{𝟒𝟓}_𝐇>`$=$`(a,a,a,0,0)\times i\sigma _2`$ one finds that only one pair of Higgs doublets remain massless while the Higgs triplets ($`\overline{H}_{t1},\overline{H}_{t2}`$) and ($`H_{t1},H_{t2}`$) have the mass matrix given by Eq.(25). In this case one has $`M_{eff}=a^2\mathrm{\Lambda }^2/M_2`$ and one can arrange for proton stability even with a light spectrum by an adjustment of the parameters $`a\mathrm{\Lambda }`$ and $`M_2`$.
Finally we discuss the current uncertainties in the proton lifetime predictions. Uncertainties arise from the errors in the quark masses, in $`\beta _p`$ and in the KM matrix elements. The largest source of uncertainties arises from the strange quark mass ($`m_s`$). There are several determinations of $`m_s`$: $`m_s=193\pm 59`$ MeV, $`m_s=200\pm 70`$ MeV, $`m_s=170\pm 50`$ MeV,$`m_s=155\pm 15`$ MeV, all evaluated at the scale 1GeV. We take for our average $`m_s=180\pm 50`$ MeV. For the charm quark mass ($`m_c`$) we use $`m_c=1.4\pm 0.2`$ GeV while for the bottom quark mass ($`m_b`$) we use $`m_b=4.74\pm 0.14`$ GeV. The contributions from the first generation quarks are small and are not the sources of any significant uncertainty in the p lifetime. The errors in the KM matrix elements are of a subleading order for the $`\overline{\nu }K`$ mode but are still significant enough to be included. We use the results of Ref. for the allowed ranges of the KM matrix elements. For $`\beta _p`$ we use the result of the lattice gauge analysis of Ref.. In Fig.6 we exhibit the error corridor for the proton lifetime for the case (1) of Fig.2 with $`M_2/a\mathrm{\Lambda }=0.01,M_2=M_G`$. One finds that given the current errors in the input data the predictions for the proton lifetime has an uncertainty of about a factor of 2 ($`1_{0.5}^{1.5}`$) on either side of the mean. A similar analysis holds for the cases(2-5) of Fig.2. We note that the uncertainties in the predictions of the proton lifetime is of the same order as the size of the CP violating effects. It is for this reason that we choose to exhibit the results of our analysis in Figs. 4 and 5 in terms of the ratio $`R_\tau `$ since the effects of the uncertainties cancel in the ratio. The analysis also shows that an improvement in the determination of the quark masses and of $`\beta _p`$ is essential for a more precise prediction of the proton lifetime in supersymmetric unification of the type discussed here. The reduction of the error in the prediction of p lifetime will also help to define the CP effects on proton decay when such a decay is experimentally observed.
In summary the CP violating effects on the proton lifetime are relatively large if the sparticle spectrum entering the dressing loop integrals is relatively light and the CP violating effects get progressively smaller as the scale of the sparticle spectrum entering the dressing loops gets progressively larger. The current experimental limits on the sparticle masses allow for a relatively light sparticle spectrum, i.e., significantly smaller than 1 TeV. This means that there exists the possibility of significant CP violating effects on the proton lifetime. However, the minimal SU(5) model does not support the scenario with a light spectrum and thus the CP violating effects for the case of the minimal model are small. However, for the non-minimal case proton stability can occur even for a relatively light spectrum due to suppression from a more complicated Higgs triplet sector. In these types of models CP violating effects can be significant.
## 4 Conclusion
In this paper we have investigated the effects of CP violating phases arising from the soft SUSY breaking sector of the theory on the proton decay amplitudes. It is found that the CP effects can increase or decrease the proton decay rates and that the size of their effect depends sensitively on the region of the parameter space one is in. Effects as large as a factor of 2 are seen to arise from CP violating phases in the part of the parameter space investigated and even larger effects in the full parameter space may occur. It is found that the CP violating effects in the minimal SU(5) model are typically small since a relatively heavy sparticle spectrum is needed to stabilize the proton in this case and a heavy spectrum suppresses the CP effects in the dressing loop integral. However, significantly larger CP effects on the p lifetime are possible in non-minimal models with more than one pair of Higgs triplets since in these models the proton can be stabilized with a relatively light sparticle spectrum. We also investigated the uncertainties in the p lifetime predictions due to uncertainties in the quark masses, in $`\beta _p`$ and in the KM matrix elements. We find that these uncertainties modify the proton lifetime by a factor of 2 around the mean value. The observations arrived at in this analysis would be applicable to a wide class of models, including GUT models and string models with dimension five baryon and lepton number violating operators
Note Added: After the paper was submitted for publication an improved limit on $`p\overline{\nu }_\mu K^+`$ mode of $`\tau (p\overline{\nu }_\mu K^+)>1.9\times 10^{33}`$ yr has been reported. The new limit does not affect the conclusions arrived at in this paper.
Acknowledgements
This research was supported in part by NSF grant PHY-9901057
| Table 1:CP effects on Chargino dressings. | | |
| --- | --- | --- |
| case | Chargino 1 | Chargino 2 |
| (Re$`A_{\nu _\mu }^{\stackrel{~}{W}_i}/|A_{\nu _{\mu 0}}|`$, Im$`A_{\nu _\mu }^{\stackrel{~}{W}_i}/|A_{\nu _{\mu 0}}|`$) | $`(.22,.89)`$ | $`(.026,.2)`$ |
| case | (Re$`A_{\nu _\mu }/|A_{\nu _{\mu 0}}|`$, Im$`A_{\nu _\mu }/|A_{\nu _{\mu 0}}|`$) | $`|A_{\nu _\mu }|/|A_{\nu _{\mu 0}}|`$ |
| sum 1&2 | $`(.25,.69)`$ | $`.74`$ |
Table caption: The table gives an analysis of the dressing loop integrals for dresssings with Charginos 1&2 (see Fig.1) and their sum for the case when $`m_0=71`$ GeV, $`m_{\frac{1}{2}}=148`$ GeV, $`\mathrm{tan}\beta =2`$, $`\theta _\mu =1.4`$, $`\xi _1=0.3`$, $`\xi _2=1.8`$, $`\xi _3=0`$, where all phases are in radians. The analysis shows cancellations in the dressings between Chargino 1 and Chargino 2 for the case with phases.
Figure Captions
Fig.1: Examples of the dressing of LLLL baryon-number-violating dimension five operators by chargino, gluino and neutralino exchanges that contribute to the proton decay. Cancellation among diagrams such as between $`\chi _1^+`$ and $`\chi _2^+`$ exchanges can lead to an enhancement of the proton lifetime. The dressings of the RRRR dimension five operators is also exhibited.
Fig.2: Plot of Log$`{}_{10}{}^{}|d_e|`$ vs $`\theta _\mu `$ exhibiting cancellations where the five curves correspond to the five sets of input for the parameters $`\mathrm{tan}\beta `$, $`m_0`$, $`m_{1/2}`$, $`\xi _1`$, $`\xi _2`$, $`\xi _3`$, $`\alpha _{A_0}`$, and $`A_0`$ given by (1)2,71,148, $`1.15`$,$`1.4`$,1.27, $`.4`$,4 (dotted), (2)2,71,148, $`.87`$,$`1.0`$,1.78, $`.4`$,4 (solid), (3)4,550,88, .5,$`1.55`$,1.5, .6,.8 (dashed), (4)4,750,88, 1.5,1.6,1.7, .6,.8 (long dashed), and (5)2,71,148, .55,1.,1.35, $`.4`$,4 (dot-dashed). All masses are in GeV and all phases are in radians.
Fig.3: Plot of Log$`{}_{10}{}^{}|d_n|`$ vs $`\theta _\mu `$ exhibiting cancellations where the five curves correspond to the five sets of input for the parameters $`\mathrm{tan}\beta `$, $`m_0`$, $`m_{1/2}`$, $`\xi _1`$, $`\xi _2`$, $`\xi _3`$, $`\alpha _{A_0}`$, and $`A_0`$ as given in Fig.2.
Fig.4: The ratio $`R_\tau `$ of the proton lifetime with phases and without phases as a function of $`\theta _\mu `$ for the five cases given in Fig.2.
Fig.5: The ratio $`R_\tau `$ as a function of the scaling factor $`\lambda `$ defined in the text. The four curves correspond to the four sets of input for the parameters $`\mathrm{tan}\beta `$, $`\xi _1`$, $`\xi _2`$, $`\xi _3`$, $`\theta _\mu `$, $`\alpha _{A_0}`$ and $`A_0`$ given by (1)2,$`1.15`$,$`1.4`$,1.27,$`1.7`$,$`.4`$,4 with $`m_0=71`$ and $`m_{1/2}=148`$ for the point of intersection with $`R_\tau `$ axis (dotted)., (2)2,$`.87`$,$`1.0`$,1.78,$`2.15`$,$`.4`$,4 with $`m_0=71`$ and $`m_{1/2}=148`$ for the first point (solid). (3)4,.5,$`1.55`$,1.5,1.56,.6,.8 with $`m_0=550`$ and $`m_{1/2}=88`$ for the first point (dashed). (4)4,1.5,1.6,1.7,$`1.56`$,.6,.8 with $`m_0=750`$ and $`m_{1/2}=88`$ for the first point (long dashed). All masses are in GeV and all phases are in radians. All trajectories satisfy edms constraints.
Fig.6: Exhibition of the uncertainties in the proton lifetime predictions due to uncertainties in the input data for case(1) of Fig.2 where we assumed $`M_2/a\mathrm{\Lambda }=0.01,M_2=M_G`$.
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# Higgs Boson Searches at LEP Up To √𝒔=𝟐𝟎𝟐 GeV
## 1 Standard Model Higgs Boson
The search for Higgs bosons at LEP shows no indication of a signal. This review includes the data from the outstanding performance of the LEP accelerator in 1999 with a collected luminosity of about 900 pb<sup>-1</sup> at $`\sqrt{s}=192`$ to 202 GeV $`^\mathrm{?}`$. The confidence levels $`CL_b`$ for a signal observation and $`CL_s`$ for setting mass limits are shown in Fig. 1. The resulting Standard Model (SM) Higgs boson mass limit is 107.9 $`\mathrm{GeV}/c^2`$ at 95% CL. The reconstructed mass distribution is shown in Fig. 2. In extensions of the SM the HZZ coupling might be weaker and thus the production cross section is reduced. Figure 2 shows limits on the reduction factor at 95% CL. Even if the SM cross section is reduced by a factor three, a Higgs boson mass up to 103 $`\mathrm{GeV}/c^2`$ is excluded.
## 2 MSSM Benchmark Results
The Minimal Supersymmetric extension of the Standard Model (MSSM) is the most attractive extension of the SM. The LEP experiments have searched for the reactions $`\mathrm{e}^+\mathrm{e}^{}\mathrm{hA}\mathrm{b}\overline{\mathrm{b}}\mathrm{b}\overline{\mathrm{b}}`$ and $`\mathrm{b}\overline{\mathrm{b}}`$$`\tau ^+\tau ^{}`$. No indication of a signal has been observed as shown in Fig. 3 for the example of $`m_\mathrm{h}m_\mathrm{A}`$, and the resulting confidence levels $`CL_b`$ and $`CL_s`$ are given. Figure 4 shows the so-called benchmark results in the MSSM for large mixing in the scalar top sector ($`m_\mathrm{h}`$-max) for $`CL_b`$ and $`CL_s`$, leading to mass limits of 88.3 and 88.4 $`\mathrm{GeV}/c^2`$ on the CP-even and CP-odd neutral Higgs bosons, and to an exclusion of the range $`0.7<\mathrm{tan}\beta <1.8`$ at 95% CL $`^\mathrm{?}`$.
## 3 A General MSSM Parameter Scan
Important reductions of the mass limits compared with benchmark results were reported for LEP1 and LEP2 data $`^\mathrm{?}`$. With increasing statistics the reduction was only 6 to 8 $`\mathrm{GeV}/c^2`$ including the 189 GeV data of one LEP experiment (DELPHI) $`^\mathrm{?}`$ and similar for OPAL $`^\mathrm{?}`$. Figure 4 shows new results from a MSSM parameter scan including the combined LEP data up to 202 GeV for $`CL_b`$ and $`CL_s`$, leading to mass limits of 86 and 87 $`\mathrm{GeV}/c^2`$ on the CP-even and CP-odd neutral Higgs bosons, and an exclusion of the range $`0.7<\mathrm{tan}\beta <1.8`$ at 95% CL. These limits are almost identical to the benchmark limits (Fig. 4), which is expected for the high luminosity of the combined data.
## 4 Charged Higgs Bosons
The search for charged Higgs bosons is performed in the framework of the general extension of the SM with two Higgs boson doublets. The combined results from the four LEP experiments for the reactions $`\mathrm{e}^+\mathrm{e}^{}`$$``$$`\mathrm{H}^+`$$`\mathrm{H}^{}`$$``$$`\mathrm{c}\overline{\mathrm{s}}`$$`\overline{\mathrm{c}}\mathrm{s}`$, cs$`\tau \nu `$, and $`\tau ^+\nu \tau ^{}\overline{\nu }`$ are presented in Fig. 6, resulting in the limit of 78.6 $`\mathrm{GeV}/c^2`$ at 95% CL, which is valid for any branching ratio Br($`\mathrm{H}^+`$ $``$ $`\tau ^+`$$`\nu `$).
## 5 Conclusions
The combination of the 1999 data from the four LEP experiments resulted in a large increase for the sensitivity of Higgs bosons; however, no indication of a signal has been found. Stringent mass limits are set for the SM Higgs boson, the neutral Higgs bosons of the MSSM, and charged Higgs bosons. Owing to the large luminosity of the combined data, a general scan of the MSSM parameters excludes neutral Higgs bosons below 86 GeV and the range $`0.7<\mathrm{tan}\beta <1.8`$.
## Acknowledgments
I would like to thank the organizers of the conference for their kind hospitality.
## References
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# Design, construction, and operation of SciFi tracking detector for K2K experiment
## 1 Introduction
### 1.1 K2K experiment
The Super-Kamiokande Collaboration has shown strong evidence for neutrino oscillations in an analysis of atmospheric neutrino data.
K2K is a long baseline neutrino experiment designed to allow more precise studies of neutrino oscillations. A high-purity muon neutrino beam, generated using the KEK 12 GeV proton synchrotron, is directed through a near detector system at KEK to the Super-Kamiokande site 250 km away. By comparing the $`\nu _\mu `$ beam flux and energy spectrum between the near and far detectors, we can investigate neutrino oscillations between $`\nu _\mu `$ and other neutrino flavors.
The near detector includes a 1kt water Cherenkov detector and a fine-grained detector (FGD), which in turn consists of the scintillating fiber (SciFi) detector, scintillating counters, a lead glass counter, and a muon range detector. Data taking for K2K began in March, 1999.
### 1.2 Scintillating fiber (SciFi) detector
The near neutrino detector is designed to measure the flux and the energy spectrum of the neutrino beam as it leaves KEK. The detector must provide good tracking capability, allowing discrimination between different types of interactions such as quasi-elastic or inelastic events. Its mass composition should be dominated by water, to allow cancellation of common systematic uncertainties with the Super-Kamiokande detector to the maximum extent possible. In order to study neutrino interactions in greater detail, we decided to use a SciFi detector constructed with fiber tracking layers interleaving water target tanks. SciFi detectors have been used in UA2 and CHORUS experiments . The K2K SciFi detector is the largest ever so far, and employs a simple design of optoelectronics readout and a calibration system using electro-luminescent (EL) plates.
## 2 Construction and operation of the K2K SciFi detector
### 2.1 Design and structure
A schematic overview of the SciFi detector is shown in Fig. 1.
It consists of 20 layers of 2.6 m $`\times `$ 2.6 m tracking modules, spaced 9 cm apart, each of which contains double layers of scintillating fiber sheets in both the horizontal and vertical directions (i.e., XXYY layers). The fiber sheets are coupled to image-intensifier tubes (IITs) which in turn are read out by CCD cameras (Fig. 2).
Between the fiber modules, there are 19 layers of water target contained in extruded aluminum tanks. Each target layer consists of 15 tank modules, and thus there are 285 in total. The width, height, and length of the tanks are 6 cm, 16 cm, and 240 cm, respectively. The thickness of the aluminum tank is 1.8 mm. The fiber is 3.7 m long and 0.692 mm in diameter.
### 2.2 Optimization of the scintillating fibers
#### 2.2.1 Measurement of the light yield of various fibers
We measured the light yield and attenuation length of various types of scintillating fibers using a photo-multiplier tube (PMT), irradiating the fiber with a <sup>90</sup>Sr $`\beta `$-ray source. The light yield, after propagation of the light over a distance $`x`$ in the fiber, is expressed by
$`Y(x)=Y_0\mathrm{exp}[x/\lambda ],`$ (1)
where $`Y_0`$ is an initial light yield at $`x=0`$ (IIT surface) and $`\lambda `$ is an attenuation length. We estimated $`Y_0`$ by extrapolation from measurements at $`x`$ = 1.5, 2.0, 2.5, and 3.0 m. The results are summarized in Table 1.
Kuraray SCSF-78M is a multi-cladding fiber, and the others have single-cladding. From these measurements, we also estimated the detection efficiency for minimum ionizing particles after light propagation over a distance $`x`$=4 m for a double-layered fiber sheet. Here, we took into account the quantum efficiencies of the IIT ($`22\%`$) and PMT ($`25\%`$), and the reflectivity at the far end of the fiber ($`70\%`$). When we required more than 99 % detection efficiency, there were only two candidates left. In order to minimize the readout area, we selected the smaller-diameter fiber, SCSF-78M (0.7 mm diameter).
#### 2.2.2 Aging
Scintillating fibers are known to have a finite lifetime, with fiber transparency diminishing with age. Aging is mainly caused by a chemical reaction of the core material (polystyrene) with oxygen in air. Fiber lifetime depends on temperature. We measured the decrease in light yield of SCSF-78M at several elevated temperatures and estimated the fiber lifetime at lower temperatures using an “Arrhenius” plot, lifetime as a function of temperature (Fig. 3).
Here we define a fiber lifetime as the number of days at which the light yield of the fiber drops to 90 % of its initial value. From the Arrhenius plot, the lifetimes of SCSF-78M at 15 C and 20 C are estimated to be 21,000 and 9,400 days, respectively. If we keep the temperature of the experimental hall less than 20 C, the decrease in the light yield is expected to be less than 10 % during the five years of the experiment. In fact, we attempt to maintain the SciFi detector temperature at 16 C.
### 2.3 Fiber sheet production
We had to fabricate double-layered fiber sheets before building the fiber tracking modules. Each fiber sheet is 370 cm long and 40 cm wide, consisting of a 260 cm long sensitive area, 80 cm long light guide, and 30 cm long IIT bundle segment (Fig. 4).
A fiber tracking module consists of 6 double-layered fiber sheets laid side by side to make a total width of 240 cm, in both horizontal and vertical directions. For 20 such modules, we needed to fabricate 240 sheets in total, plus 8 spares. This method of fiber sheet production is quite similar that used by the CHORUS Collaboration . To produce 3.7 m long fiber sheets, we prepared a 1.2 m diameter drum with a spiral groove whose pitch was 0.7 mm. We wound fiber stock along the groove until the sheet became 40 cm wide. Then, we coated the sheet with white paint. After it became dry, we wound the second layer fiber onto the first layer and painted it again. Finally, we cut the sheet, removed it from the drum, and painted its inner side.
A control system for fiber sheet production is shown in Fig. 5.
Fiber stock was wound onto the bobbin and threaded onto the winding machine. The drawing system, which consists of 2 rubber coated rollers, a rotary encoder, and a stepping motor, drew fiber from the bobbin. Because the two layers have slightly different circumferences on the drum, we applied a constant tension to each layer (150 gw to the first and 300 gw to the second layers, respectively) during winding. After being removed from the drum, the prestressed sheet then became flat due to the restoring force supplied by the outer fibers. At the downstream end of the tension controller, a laser sensor monitored the fiber diameter. When it detected a fiber section out of tolerance (692 $`\mu `$m $`\pm 18\mu `$m), the whole system stopped, and we cut away the segment. Guide pullies attached to an accurate linear guide, driven by a stepping motor, lead the fiber onto the groove of the drum. They were synchronized with rotation of the drum and moved 0.7mm for a turn of rotation. The drum was driven by an AC servo motor. The rotation speed was adjustable and the maximum speed was about 1 m/sec (0.265 rps). The whole fiber sheet production system was controlled by a conventional PC with 2 custom interface boards.
### 2.4 Fiber sheet modules
After all the fiber sheets were made, both ends of the sheets were polished by a milling machine and far ends of the sheets were sputtered with aluminum. We checked the quality of all the individual fibers by irradiating the sheets with <sup>90</sup>Sr $`\beta `$-source at the farthest point. The light yield of the fiber sheets was measured to be 8$`\pm `$1 photo-electrons.
Since the fiber sheets are very fragile, we needed a strong and light material to support the fiber sheet. We chose a honeycomb panel with a dimension of 2.6 m $`\times `$ 2.6 m $`\times `$ 1.6 cm, consisting of a 1.4 cm thick paper-honeycomb core and two 1 mm thick GFRP skins. The area density of the honeycomb board is 0.52 g/cm<sup>2</sup>, and their net mass is 30 kg per board.
We glued 6 fiber sheets side by side on one surface of the honeycomb panel in a horizontal direction, and 6 more sheets on the other side in a vertical direction (referring to the final orientation), using epoxy adhesive CY221/HY2967. Since we found that gas released by the glue actually accelerates fiber aging, we placed the glue in a vacuum chamber for 5 minutes before gluing. The positions and the straightness of the sheets were controlled within 0.5 mm by a ruler bar and alignment pins when the sheets are glued. The straightness of each sheet was measured by a micrometer every 10 cm with accuracy of 0.1 mm.
Since the fiber sheet modules were rather large (3.7 m $`\times `$ 3.7 m), we made a special aluminum jig (4 m $`\times `$ 4 m) to perform assembly, transportation, and installation of the modules safely and easily. Throughout the installation, only about 20 fibers were broken out of 274080.
### 2.5 Fiber bundles
In order to make effective use of the photosensitive area (10 cm diameter) of the IIT surface, fibers are bundled at the readout. All the fibers (1142 $`\times `$ 6 $`\times `$ 20 $`\times `$ 2 = 274080) are grouped at the readout end and glued into bundles, which are then attached to the 24 IIT’s. Ten consecutive fiber sheets along the beam direction are assigned to one fiber bundle. Thus, there are 6 fiber bundles in the upstream half of the detector, and another 6 fiber bundles in the downstream half, in both x (horizontal) and y (vertical) directions.
One fiber bundle contains 1142 $`\times `$ 10 (or 274080/24)= 11420 fibers. Fig. 6 shows the arrangement used to accumulate fibers and make a (half) bundle from (5) fiber sheets. The remaining half bundle is symmetrical with respect to the direction along the SciFi layer and the two halves are glued together.
In order to optically separate fiber sheets from different layers, acrylic films (200 $`\mu `$m thick) are inserted between individual fiber sheets in a bundle. Epoxy glue (CY221/HY2967) was used to fabricate the bundle. For good optical contact with the IIT surface, we trimmed and polished the surface of each bundle using a custom-made polishing machine. The detector geometry and construction logistics required us to perform bundle polishing operations in the experimental hall, after installation of all the SciFi modules, with very limited working space. A compact portable milling machine, with dual carbide and diamond bits and pre-programmed automatic feed control, was specially designed and built for this purpose. Following polishing, bundles were inspected for flaws using a microscope system.
### 2.6 Readout system
Our opto-electronics readout comprises two stages of image-intensifiers, an optical lens and a CCD camera(See Fig. 2). The system was designed to ensure simplicity and thus low cost. The first stage is an electrostatic image-intensifier (Hamamatsu V5502UX) with 100 mm diameter photocathode. The image is reduced in diameter by a factor of 23 % and the light yield is amplified by about 5. The bialkali photocathode has quantum efficiency 22% at 430nm, which is the peak wavelength in the emission spectrum of the scintillating fiber. The second stage is a micro-channel-plate(MCP) image-intensifier (Hamamatsu MCP-IIT:V1366GX). The gain in light yield is about 1000 at a typical HV value of 2600 volts. We operate the MCP-IIT with pulsed HV to reduce unwanted events. The gate width is set to 100 $`\mu `$s, which corresponds to the decay time of the fluorescent phosphor screen in the first stage IIT. The optical lens at the third stage reduces the image to 30 % and projects it onto the CCD camera (C3077). The CCD camera used here has 768 pixels and 493 pixels in horizontal $`(x)`$ and vertical $`(y)`$ directions, respectively, with pixel size 11 $`\mu `$m $`\times `$ 13 $`\mu `$m. One fiber is seen by about $`4\times 4`$ pixels.
The video signal from the CCD is fed to a flash-ADC module housed in a NIM crate, where $`(x,y)`$ positions of the CCD and the pulse height are digitized. The digital outputs are then sent to 12 VME FIFO modules. The readout time is about 30 msec.
### 2.7 EL calibration
To identify hit fibers from the CCD image, we have to know the correspondence between fiber and CCD coordinates. For this purpose, we illuminate selected fibers at regular intervals using an electro-luminescent plate (EL). We will refer to these position calibrations as “EL” calibrations hereafter.
We placed EL’s at the edge of the fiber sheets (See Fig. 6 in the previous subsection), where we coated the sheet with black paint and then scraped the coating off selected fiducial fibers (Fig. 7). Thus, light from the EL goes only into the fiducial fibers. We illuminate only one out of every 10 or 20 fibers and fibers at the edges of the sheets, to save labor and data size. Positions of the other fibers are interpolated from these fiducial measurements. We get the fiber positions within 100 $`\mu `$m by this method. We measured the x, y, and z positions of the edge fibers of each layer with accuracy of 1 mm during installation, along with other fiducial marks on the fiber modules, to relate fiber coordinates to the experiment coordinate system. Cosmic-ray muons were used to obtain the final alignment information.
## 3 Performance of the detector and analysis of the data
### 3.1 Hit definition and the pulse height
The energy loss of a minimum ionizing particle is estimated to be 0.19 MeV in the fiber sheet and generates about 8 photoelectrons (p.e.) on the IIT surface if the particle hits the fiber at the mid-point of SciFi detector (2.4 m from the readout end), and 6.5 p.e. if it hits at the end of the fiber (3.7 m from the readout).
Fig. 8 shows a typical IIT image corresponding to a particle track, read out by CCD: the dark spots represent hit CCD pixels. The fiber positions determined from EL calibration are also shown as small circles. To reconstruct a hit, the hit finding algorithm starts clustering CCD-pixels, finds all fibers which overlap the pixel cluster, and clusters the fibers. The hit is defined as a cluster of fibers which overlap the pixel clusters: the hit position is calculated for the center of gravity of the hit fibers (weighted by the number of overlapping pixels) in the fiber cluster.
The attenuation and reflection of photons inside the fiber was measured during the quality test as described in Section 2.4. The light yield is given by
$$Y(d)=Y_0(e^{\frac{d}{\lambda }}+Re^{\frac{2Ld}{\lambda }}),$$
(2)
where $`Y_0`$ is the normalization constant, $`d`$ is the distance from the IIT surface, $`\lambda `$ is the attenuation length (323 cm), $`R`$ is the reflectivity of aluminum coated side (0.74), and $`L`$ is the length of fiber (370 cm).
Fig. 9 shows the mean values of the number of pixels in a cluster ($`N_{pixel}`$) as a function of hit position. The IIT surface is located at $`x=d240`$ cm (the center of detector is located at $`x=0`$). The data shown in open circles are the measured values for the cosmic-ray muons and they agree well with the prediction by Eq. (2) (solid line). After the attenuation correction, the position dependence almost vanished, as shown in the figure (the solid circles). The $`N_{pixel}`$ value for all track hits are normalized to that at $`x=0`$.
Fig. 10 shows the typical $`N_{pixel}`$ distribution of data and a pixel simulation of cosmic-ray muons. $`N_{pixel}`$ is corrected for the attenuation length and the incident angle.
The response for a minimum ionizing particle (MIP) is calibrated using cosmic-ray muons. The $`N_{pixel}`$ distribution for cosmic-ray muons shows a peak at about 50 when it is normalized to the pulse height at $`x=0`$. This means 8 p.e. corresponds to $`N_{pixel}=50`$ at the peak position, which agrees well with LED test data.
### 3.2 SciFi Monte Carlo simulation - Pixel simulation
The Monte Carlo simulation for the SciFi detector proceeds in two steps. In the first, charged particles are tracked through the detector using GEANT. The energy loss of a particle (in GeV) is computed whenever the particle traverses a fiber and is then converted to the number of photoelectrons observed by the IIT/CCD imaging system, taking into account attenuation along the fiber, reflection at the aluminium coated end and the quantum efficiency of the IIT photocathodes.
The photoelectrons are then converted to an image on the CCD surface. Each photoelectron is offset from the center of the fiber according to a gaussian distribution measured using single photoelectrons during LED testing of the IITs. The size of the pixel cluster corresponding to a single photoelectron is then chosen from distributions parameterized from the LED data and the pixel cluster is built assuming circular symmetry. Finally, the tuning of the energy scale (i.e. how many photoelectrons correspond to the energy deposition from a minimum ionising particle) is carried out using cosmic-ray muons.
### 3.3 Track and vertex reconstruction
The track and vertex finding algorithm proceeds in 4 steps. First, the track finder searches for 2-dimensional (2D) track candidates (in the XZ or YZ planes) which hit at least 3 layers of SciFi modules. A 3D track is reconstructed by finding the best combination of two X-view and Y-view 2D tracks, comparing the starting and ending points and checking the overlap of the 2D tracks. A vertex is defined as an intersection of 3D tracks, and the vertex position is determined with full 3D $`\chi ^2`$ minimization fit. If there is only one track in the event, the mid-point of the nearest upstream water tube is taken as the event vertex. Finally, using the reconstructed vertex position, the track finder starts searching for extra tracks near the vertex, which are mainly short tracks. The last step is very important because most proton tracks from charge current (CC) interactions are short and may not be found in the first iteration of track finding.
The efficiency of track finding strongly depends on hit efficiency and the density of noise hits in the events. Hit efficiency ($`\epsilon _{hit}`$) was evaluated using cosmic-ray muons which penetrate whole SciFi modules as
$$\epsilon _{hit}=\frac{Numberofhitsonthetrack}{Numberofexpectedhitsonthetrack}.$$
(3)
In order to keep the hit efficiency reasonably high and to keep the noise rate as low as 0.5 %, we required a minimum value of $`N_{pixel}`$ for a hit. Fig. 11 shows the hit efficiency and the noise rate at various cut values on $`N_{pixel}`$. Based on this figure, we decided to require $`N_{pixel}2`$. As a result, the hit efficiency was found to be 92$`\pm `$2 % on average, with almost no noise contamination ($`<`$0.5 %)<sup>1</sup><sup>1</sup>1After upgrading the HV system of IIT in Summer, 1999, we found the hit efficiency to be 96 $`\pm `$ 2 % with the same noise level.. Fig. 12 shows the hit efficiency for each layer.
Fig. 13 shows the track finding efficiency, which was estimated using cosmic-ray muons and muon tracks from $`\nu _\mu `$ CC interactions generated in the Monte Carlo simulation. They agree with each other: 98$`\pm `$2 % for long tracks (traversing more than 7 SciFi layers) and 85$`\pm `$6 % for short tracks. Here, the horizontal axis is the number of layers which the muon passes through.
Vertex resolution estimated using Monte-Carlo events($`\nu _\mu `$ CC events) was 1.5 mm in the beam direction and 1.5 mm perpendicular to the beam.
The position resolution of SciFi detector was estimated from the residual distribution for hits. Since multiple Coulomb scattering effects are not negligible, we implemented a Kalman filtering/smoothing algorithm in order to take multiple scattering into account during the track fit. Fig. 14 shows a residual distribution of hits for cosmic-ray muons. The figure shows that hit position resolution is 0.8 mm.<sup>2</sup><sup>2</sup>2Using a simple straight-line fit instead of Kalman Filter method, the rms of the corresponding residual distribution was 1.5 mm. At present, performance is limited by multiple scattering effects in the water tank and imperfect alignment of the sheets. If we can correct for the detailed alignment parameters, such as the curvature and rotation of the sheets, we expect the hit position resolution to improve to about 0.5mm, which is dominated by multiple scattering effects.
### 3.4 Neutrino event reconstruction
Fig. 15 shows a typical $`\nu _\mu `$ CC quasi-elastic (qe) interaction ($`\nu _\mu +n\mu ^{}+p`$) candidate. The longer track is a muon, which hits the scintillating counters and lead glass counters and reaches the muon range detector, and the shorter track is a proton.
Fig. 16 shows the distribution of the number of the reconstructed tracks from the $`\nu _\mu `$ interaction vertex in $`\nu _\mu `$ event candidates. The total number of events in a Monte Carlo prediction is normalized to that in the data. It shows that charged particles in $`\nu _\mu `$ CC events are succesfully reconstructed and the data (solid circles) are well reproduced by Monte Carlo prediction (histogram). We estimated the fraction of tracks with misreconstructed trajectories to be less than 4 %.
### 3.5 Particle identification
Although tracking is the main purpose of the SciFi detector, we may use it for particle identification. Electrons can be identified from SciFi hit patterns, and are useful for identifying the $`\nu _e`$ events. The energy of an electron can be estimated from the number of hits in the shower, i.e. $`E_e`$ = $``$ 10 MeV/hit for SciFi detector, which is estimated from a Monte Carlo study and a beam test with a prototype detector. The energy resolution of electrons was measured to be about 28 % for 0.3 GeV in the beam test and it is estimated to be 15 % for 1.0 GeV electrons.
From LED runs and the beam test, the number of pixels ($`N_{pixel}`$) shows good linearity with the number of photoelectrons ($`N_{p.e.}`$). Fig. 17(a) shows the correlation of $`<N_{pixel}>`$ and $`N_{p.e.}`$ from the LED test. Fig. 17(b) shows the results of the beam test of the proton and pion beams. The horizontal axis is the $`p/M`$ of the particle and the vertical axis is the mean value of $`N_{pixel}`$. The simplified Bethe-Bloch formula reproduces the data points well.
$`N_{pixel}`$ of the tracks may be used to separate protons from muons or pions. A feasibility study was performed using $`\nu _\mu `$ CCqe event candidates, which contain one muon track and one proton track. Fig. 18 shows $`N_{pixel}`$ distribution for (a) muon and (b) proton candidates. The figure shows that the proton candidates have larger $`N_{pixel}`$ values because protons will lose more energy in the fiber sheet.
## 4 Summary
A very large scintillating fiber (SciFi) tracking detector for the K2K long baseline neutrino oscillation experiment has been in operation since March, 1999. Track finding efficiency is 98$`\pm `$2 % for long muon tracks (those which intersect more than 5 fiber planes), and 85$`\pm `$6 % for short tracks, which were estimated using cosmic-ray muons and a Monte Carlo simulation. The position resolution per layer is about 0.8 mm. The SciFi detector has demonstrated its capability for reconstruction of $`\nu _\mu `$ interactions. The pulse heights for cosmic-ray muons have been stable within 10 % after one year of operation. In addition, the SciFi detector offers the possibility of performing electron and proton identification.
## 5 Acknowledgements
We thank Mr. S. Ishikawa and Mr. T. Kawai of Nagoya University for designing and constructing the fiber spooling machine. We also thank Prof. K. Niwa and his group, especially Dr. T. Nakano, of Nagoya University for their cooperation. Thanks are also due to Mr. C. Lindenmeyer, Mr. P. Mulligan, and Mr. J. Roze, for the design, construction, and operation of the fiber bundle polishing machine. We also thank M. Powlowski, T. Raza, and T. Kato (SUNY) for helping the fiber sheet assembly.
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# Formation and evolution of dusty starburst galaxies I. A new method for deriving spectral energy distribution
## 1 Introduction
Recent observational studies have discovered possible dusty starburst candidates in various classes of galactic objects such as low z ultra-luminous infrared galaxies (ULIRGs) (e.g., Sanders et al. 1988; Sanders & Mirabel 1996), intermediate z ones (e.g., Tran et al. 1999), faint submillimeter sources detected by the Sub-millimeter Common-User Bolometer Array (SCUBA) (Holland et al. 1999) on the James Clerk Maxwell Telescope (Smail, Ivison, & Blain 1997; Hughes et al. 1998; Barger et al. 1998; Smail et al. 1998; Ivison et al. 1999; Lilly et al. 1999), extremely red objects (EROs) (Elston, Rieke, & Rieke 1988; Graham & Dey 1996; Cimatti et al. 1998; Dey et al. 1999; Smail et al. 1999), optically faint radio sources detected by the Very Large Array (VLA) (e.g., Richards et al. 1999), Lyman-break galaxies (Steidel et al. 1996; Lowenthal et al. 1997), and near-infrared emission line galaxies (e.g., Mannucci et al. 1998). These dusty starburst candidates are generally considered to provide valuable information on the formation and the evolution of galaxies, accordingly the origin and the nature of these galaxies have been extensively discussed in variously different contexts. These discussions include, for example, the importance of dissipative dynamics in the formation of elliptical galaxies at low and high redshift (Kormendy & Sanders 1992; Barnes & Hernquist 1992), an evolutionary link between starburst and active galactic nuclei (Norman & Scoville 1988), fueling of dusty gas to the central $``$ 100 pc for nuclear starburst in major galaxy mergers (e.g., Mihos & Hernquist 1996), physical correlations between morphological and photometric properties in ULIRGs (Bekki, Shioya, & Tanaka 1999; Bekki & Shioya 2000a), merging and clustering process of high $`z`$ dusty mergers within a hierarchical clustering scenario (e.g., Somerville & Primack 1998), and cosmic star formation history (e.g., Pei & Fall 1995; Madau et al. 1996; Meurer et al. 1997; Madau, Pozzetti, & Dickinson 1998; Pascarelle, Lanzetta, & Fernandez-Soto 1998; Blain et al. 1999; Steidel et al. 1999). One of important problems in extra-galactic astronomy is to settle the above discussions on high redshift dusty starburst galaxies, though observational analysis still has some difficulties in determining precisely redshifts (Sanders 1999), identifying optical counterparts (Richards 1999), and estimating the degree of dust extinction (e.g., Meurer, Heckman, & Calzetti 1999) for high $`z`$ dusty starburst galaxies.
Spectral energy distributions (SEDs) are generally considered to be one of essentially important factors that can determine the nature of dusty starburst galaxies. Accordingly several theoretical attempts have been made to derive the SEDs of galaxies with dusty starburst regions. Rowan-Robinson & Crawford (1989) discussed how the far-infrared spectra of a sample of galaxies observed by $`Infrared`$ $`Astronomical`$ $`Satellite`$ (IRAS) can be reproduced by a theoretical model with a given temperature of stars embedded by dust, the optical depth, and the density distribution (or geometry) of dust. Witt et al. (1991) investigated transfer processes of stellar light for models with variously different spherical geometries of dust with special emphasis on the effects of scattered light on photometric properties of galaxies and SEDs. Calzetti, Kinney, & Storchi-Bergmann (1994) analyzed ultraviolet (UV) and optical spectra of 39 starburst and blue compact galaxies and thereby provided an analytical formulation of the effects of dust extinction in galaxies and an effective extinction law for correcting the observed UV and optical spectral continua. Witt & Gordon (1996) investigated the radiative transfer processes of a central stellar source surrounded by a spherical, statistically homogeneous but clumpy two-phase scattering medium and found that the structure of dusty medium can greatly affect the conversion of UV, optical, and near-infrared radiation into thermal far-IR dust radiation in a dusty system. By comparing the observed SEDs of 30 starburst galaxies with theoretical radiative transfer models of dusty systems, Gordon, Calzetti, and Witt (1997) discussed the importance of geometry of stellar and gaseous components in determining the SED of a dusty starburst galaxy. Takagi, Arimoto, & Vansevi$`\stackrel{ˇ}{\mathrm{c}}`$ius (1999) investigated radiative transfer models with variously different ages of secondary starburst components and optical depths for dusty starburst galaxies and proposed a new method for estimating precisely the effect of ages of young starburst populations and that of dust attenuation on the shape of SED in UV and near$``$infrared bands. Efstathiou, Rowan-Robinson, & Siebenmorgen (2000) treated starburst galaxies as an ensemble of optically thick giant molecular clouds (GMCs) centrally illuminated by recently formed stars and thereby constructed a new radiative transfer model for calculating SEDs from UV to millimeter band of dusty starburst galaxies. By using this model, they discussed how the age and the star formation history of a dusty starburst galaxy can control the SED particularly for the prototypical starburst galaxy M82 and NGC 6090.
Although the above previous studies have succeeded in clarifying important dependences of SEDs on physical parameters such as spatial distribution of dust, geometries of stellar and gaseous components, and dust properties in dusty starburst galaxies, many of them did not discuss $`the`$ $`time`$ $`evolution`$ $`of`$ $`SEDs`$. Main reasons for $`some`$ previous theoretical studies’ not discussing the time evolution of SEDs are the following three. Firstly previous models did not follow the time evolution of stellar and gaseous distribution and accordingly could not derive the time evolution of SEDs in dusty starburst galaxies. Secondly, hydro-dynamical evolution of interstellar gas (e.g., time evolution of gaseous density) was not included in previous studies, and accordingly the time evolution of optical depth could not be derived. Thirdly, some previous studies did not consider chemical evolution of dusty interstellar medium and that of stellar components, they could not follow time evolution of metallicity and that of dust properties. Considering that low z infrared luminous galaxies with dusty starburst and high $`z`$ faint SCUBA sources with possible dusty starburst show very peculiar morphology (Sanders et al. 1988; Sanders & Mirabel 1996; Smail et al. 1998), spherical symmetric approximation or axisymmetric one adopted in previous theoretical and numerical studies for the stellar and gaseous distribution of a dusty galaxy should be also relaxed in order that the nature and the origin of low and high z dusty galaxies can be extensively investigated by theoretical studies.
The purpose of the present paper and our future papers is to investigate in detail the formation and the evolution of dusty starburst galaxies by using a new numerical code for deriving the time evolution of SEDs of these galaxies. We first perform numerical simulations that can follow dynamical evolution of stellar and gaseous components, star formation history, and chemical evolution for dusty starburst galaxies in an explicitly self-consistent manner. We then derive stellar and gaseous distribution and age and metallicity distribution of stellar populations and thereby calculate galactic SEDs. Furthermore we describe how the present numerical code is useful and helpful for investigating variously different physical properties of low and high z dusty galaxies. As an example, we here present the results of a major merger with dusty interstellar gas. We particularly demonstrate how dynamical evolution of stellar and gaseous components controls the time evolution of SEDs in dusty galaxy mergers and thus their photometric evolution.
The layout of this paper is as follows. In §2, we summarize numerical models used in the present study and describe in detail the methods for deriving the SEDs corrected by internal dust extinction. In this section, we also point out the limitations of the present numerical code. In §3, we present numerical results on the time evolution of morphology, SED, and photometric properties in a gas-rich major merger. In §4, we discuss the origin of high $`z`$ dusty starburst galaxies such as faint SCUBA sources and EROs. The conclusions of the preset study are given in §5.
## 2 Model
The most remarkable difference in deriving SEDs of dusty galaxies between the present model and previous models (e.g., Mazzei, Xu, & De Zotti 1992; Franceschini et al. 1994; Witt & Gordon 1996; Gordon st al. 1997; Guiderdoni et al. 1998) is that we derive SEDs of galaxies based on the result of numerical simulations that can follow both dynamical and chemical evolution of galaxies. The derivation of a SED for a galaxy with dusty starburst consists of the following three steps. In the first step we simultaneously derive age and metallicity distribution of stellar component in the galaxy, based on numerical simulations. Secondly, we use a stellar population synthesis code (e.g. Bruzal & Charlot 1993) and thereby calculate the SED of the galaxy, based on the derived age and metallicity distribution of stellar population in the galaxy. In this second step, the dust effects are not included. Thirdly, we consider extinction and reemission of stellar light by dusty interstellar gas and modify the SED of the galaxy, by using our new code. The details of our new method for modifying galactic SEDs are given in §2.3.
### 2.1 Merger model
#### 2.1.1 Initial conditions
We construct models of galaxy mergers between gas-rich disk galaxies with equal mass by using Fall-Efstathiou model (1980). The total mass and the size of a progenitor disk are $`M_\mathrm{d}`$ and $`R_\mathrm{d}`$, respectively. From now on, all the mass and length are measured in units of $`M_\mathrm{d}`$ and $`R_\mathrm{d}`$, respectively, unless specified. Velocity and time are measured in units of $`v`$ = $`(GM_\mathrm{d}/R_\mathrm{d})^{1/2}`$ and $`t_{\mathrm{dyn}}`$ = $`(R_\mathrm{d}^3/GM_\mathrm{d})^{1/2}`$, respectively, where $`G`$ is the gravitational constant and assumed to be 1.0 in the present study. If we adopt $`M_\mathrm{d}`$ = 6.0 $`\times `$ $`10^{10}`$ $`\mathrm{M}_{}`$ and $`R_\mathrm{d}`$ = 17.5 kpc as a fiducial value, then $`v`$ = 1.21 $`\times `$ $`10^2`$ km/s and $`t_{\mathrm{dyn}}`$ = 1.41 $`\times `$ $`10^8`$ yr, respectively. In the present model, the rotation curve becomes nearly flat at 0.35 radius with the maximum rotational velocity $`v_\mathrm{m}`$ = 1.8 in our units. The corresponding total mass $`M_\mathrm{t}`$ and halo mass $`M_\mathrm{h}`$ are 5.0 and 4.0 in our units, respectively, which means that the baryonic mass fraction of the initial disk galaxy is 0.2. The radial ($`R`$) and vertical ($`Z`$) density profile of a disk are assumed to be proportional to $`\mathrm{exp}(R/R_0)`$ with scale length $`R_0`$ = 0.2 and to $`\mathrm{sech}^2(Z/Z_0)`$ with scale length $`Z_0`$ = 0.04 in our units, respectively. The velocity dispersion of halo component at a given point is set to be isotropic and given according to the virial theorem. In addition to the rotational velocity made by the gravitational field of disk and halo component, the initial radial and azimuthal velocity dispersion are given to disk component according to the epicyclic theory with Toomre’s parameter (Binney & Tremaine (1987)) $`Q`$ = 1.2. The vertical velocity dispersion at given radius is set to be 0.5 times as large as the radial velocity dispersion at that point, as is consistent with the observed trend of the Milky Way (e.g., Wielen 1977). As is described above, the present initial disk model does not include any remarkable bulge components, and accordingly corresponds to ‘purely’ late-type spiral without galactic bulge. Although it is highly possible that galactic bulges greatly affect the chemical evolution of galaxy mergers, we however investigate this issue in our future papers.
The collisional and dissipative nature of the interstellar medium is modeled by the sticky particle method (Schwarz (1981)). It should be emphasized here that this discrete cloud model can at best represent the $`real`$ interstellar medium of galaxies in a schematic way. As is modeled by McKee & Ostriker (1977), the interstellar medium can be considered to be composed mainly of ‘hot’, ‘warm’, and ‘cool’ gas, each of which mutually interacts hydrodynamically in a rather complicated way. Actually, these considerably complicated nature of interstellar medium in disk galaxies would not be so simply modeled by the ‘sticky particle’ method in which gaseous dissipation is modeled by ad hoc cloud-cloud collision: Any existing numerical method probably could not model the $`real`$ interstellar medium in an admittedly proper way. In the present study, as a compromise, we only try to address some important aspects of hydrodynamical interaction between interstellar medium in disk galaxies and in dissipative mergers. More elaborated numerical modeling for real interstellar medium would be necessary for our further understanding of dynamical evolution in dissipative galaxy mergers. We assume that the fraction of gas mass ($`f_g`$) in a disk is set to be 0.5 initially. Actually, the gas mass fraction in precursor disks of a merger is different between galaxy mergers and depends on the epoch of the merging. For example, recent observational results on the total mass of molecular gas for faint SCUBA sources (e.g., Frayer et al. 1999) revealed that SMM J14011+0252 with some indications of major merging at z = 2.565 has $``$ 5.0 $`\times `$ $`10^{10}`$ $`M_{}`$. If the progenitor of this gas-rich high redshift merger has a mass of 6.0 $`\times `$ $`10^{10}`$ $`\mathrm{M}_{}`$ (i.e., the same as that of the Galaxy), the gas mass fraction is $``$ 0.42 for this galaxy. The gas mass fraction is observationally suggested to be larger for higher redshift mergers (e.g., Evans, Surace, & Mazzarella 1999). Guided by these observational results, we adopt the above value (0.5) as the initial gas mass fraction in order to discuss higher redshift galaxy mergers. Although the difference of gas mass fraction probably could yield a great variety of chemical and dynamical structures in mergers, we do not intend to consider this important difference for simplicity in the present paper and will address in our future paper. The radial and tangential restitution coefficient for cloud-cloud collisions are set to be 1.0 and 0.0, respectively. The number of particles for an above isolated galaxy is 10000 for dark halo, 10000 for stellar disk components, and 10000 for gaseous ones.
In all of the simulations of mergers, the orbit of the two disks is set to be initially in the $`xy`$ plane and the distance between the center of mass of the two disks, represented by $`r_{\mathrm{in}}`$, is assumed to be the free parameter which controls the epoch of galaxy merging. The pericenter distance, represented by $`r_\mathrm{p}`$, is also assumed to be the free parameter which controls the initial total orbital angular momentum of galaxy mergers. The eccentricity is set to be 1.0 for all models of mergers, meaning that the encounter of galaxy merging is parabolic. The spin of each galaxy in a pair merger is specified by two angle $`\theta _i`$ and $`\varphi _i`$, where suffix $`i`$ is used to identify each galaxy. $`\theta _i`$ is the angle between the $`z`$ axis and the vector of the angular momentum of a disk. $`\varphi _i`$ is the azimuthal angle measured from $`x`$ axis to the projection of the angular momentum vector of a disk onto $`x`$-$`y`$ plane. In the present study, we show the results of only one model with $`\theta _1`$ = 0.0, $`\theta _2`$ = $`150.0`$, $`\varphi _1`$ = 0.0, $`\varphi _2`$ = 0.0, $`r_\mathrm{p}`$ = 17.5 kpc, and $`r_{\mathrm{in}}`$ = 140 kpc: This model describes a nearly prograde-retrograde merger. The results of the models with variously different $`\theta _i`$, $`\varphi _i`$, $`r_\mathrm{p}`$ and $`r_{\mathrm{in}}`$ will be described in our future papers. The time when the progenitor disks merge completely and reach the dynamical equilibrium is less than 15.0 in our units for most of models and does not depend so strongly on the history of star formation in the present calculations.
#### 2.1.2 Global star formation
Star formation is modeled by converting the collisional gas particles into collisionless new stellar particles according to the algorithm of star formation described below. We adopt the Schmidt law (Schmidt 1959) with exponent $`\gamma `$ = 2.0 (1.0 $`<`$ $`\gamma `$ $`<`$ 2.0, Kennicutt (1989)) as the controlling parameter of the rate of star formation. The amount of gas consumed by star formation for each gas particle in each time step, $`\dot{M_\mathrm{g}}`$, is given as:
$$\dot{M_\mathrm{g}}C_{\mathrm{SF}}\times (\rho _\mathrm{g}/\rho _0)^{\gamma 1.0}$$
(1)
where $`\rho _\mathrm{g}`$ and $`\rho _0`$ are the gas density around each gas particle and the mean gas density at 0.48 radius of an initial disk, respectively. This density-dependent star formation model is similar to that of Mihos, Richstone, & Bothun (1992) and the probability approach of the gas consumption rate is similar to that described by Katz (1992) and O’Neil et al. (1998). In order to avoid a large number of new stellar particles with different mass, we convert one gas particle into one stellar one according to the following procedure. First we give each gas particle the probability, $`P_{\mathrm{sf}}`$, that the gas particle is converted into stellar one, by setting the $`P_{\mathrm{sf}}`$ to be proportional to the $`\dot{M_\mathrm{g}}`$ in equation (1) estimated for the gas particle. Then we draw the random number to determine whether or not the gas particle is totally converted into one new star. This method of star formation enables us to control the rapidity of star formation without increase of particle number in each simulation thus to maintain the numerical accuracy in each simulation. The $`C_{\mathrm{SF}}`$ in the equation (1) is the parameter that controls the rapidity of gas consumption by star formation. We determine the initial value of the $`C_{\mathrm{SF}}`$ so that mean star formation rate of an isolated late-type disk galaxy model with the typical gas mass fraction of $``$ 0.2 can become an order of 1 $`M_{}`$ $`\mathrm{yr}^1`$ for the first 1 Gyr evolution. The positions and velocity of the new stellar particles are set to be the same as those of original gas particles.
All the calculations related to the above dynamical evolution including the dissipative dynamics, star formation, and gravitational interaction between collisionless and collisional component have been carried out on the GRAPE board (Sugimoto et al. (1990)) at Astronomical Institute of Tohoku University. The parameter of gravitational softening is set to be fixed at 0.03 in all the simulations. The time integration of the equation of motion is performed by using 2-order leap-flog method. Energy and angular momentum are conserved within 1 percent accuracy in a test collisionless merger simulation. Most of the calculations are set to be stopped at T = 20.0 in our units unless specified.
### 2.2 Method for the SED of stellar populations
#### 2.2.1 Chemical enrichment
Chemical enrichment through star formation during galaxy merging is assumed to proceed both locally and instantaneously in the present study. The model for analyzing metal enrichment of each gas and stellar particle is as follows. First, as soon as a gas particle is converted into a new stellar one by star formation, we search neighbor gas particles locating within $`R_{\mathrm{chemi}}`$ from the position of the new stellar particle and then count the number of the neighbor gas particles, $`N_{\mathrm{gas}}`$. This $`R_{\mathrm{chemi}}`$ is referred to as chemical mixing length in the present paper, and represents the region within which the neighbor gas particles are polluted by metals ejected from the new stellar particle. The value of $`R_{\mathrm{chemi}}`$ relative to the typical size of a galaxy is set to be 0.01. Next we assign the metallicity of original gas particle to the new stellar particle and increase the metals of the each neighbor gas particle according to the following equation about the chemical enrichment:
$$\mathrm{\Delta }M_\mathrm{Z}=\{Z_iR_{\mathrm{met}}m_\mathrm{s}+(1.0R_{\mathrm{met}})(1.0Z_i)m_\mathrm{s}y_{\mathrm{met}}\}/N_{\mathrm{gas}}$$
(2)
where the $`\mathrm{\Delta }M_\mathrm{Z}`$ represents the increase of metal for each gas particle. $`Z_i`$, $`R_{\mathrm{met}}`$, $`m_\mathrm{s}`$, and $`y_{\mathrm{met}}`$ in the above equation represent the metallicity of the new stellar particle (or that of original gas particle), the fraction of gas returned to interstellar medium, the mass of the new star, and the chemical yield, respectively. The value of the $`R_{\mathrm{met}}`$ and that of $`y_{\mathrm{met}}`$ are set to be 0.2 and 0.03, respectively. Furthermore, the time, $`t_i`$, when the new stellar particle is created, is assigned to the new stellar particle in order to calculate the photometric evolution of merger remnants, as is described later. To verify the accuracy of the above treatment (including numerical code) for chemical enrichment process, we checked whether or not the following conservation law of chemical enrichment is satisfied for each time step in each test simulation:
$$\underset{\mathrm{star}}{}m_\mathrm{s}Z_i+\underset{\mathrm{gas}}{}m_\mathrm{g}Z_i=y_{\mathrm{met}}\underset{\mathrm{star}}{}m_\mathrm{s}$$
(3)
where $`m_\mathrm{g}`$, $`m_\mathrm{s}`$, and $`Z_i`$ are the mass of each gas particle, that of each stellar one, and the metallicity of each particle, respectively, and the summation ($``$) is done for all the gas particles or stellar ones. Strictly speaking, the above equation holds when the value of $`R_{\mathrm{met}}`$ is 0.0. Thus, in testing the validity of the present code of chemical enrichment, we set the value of $`R_{\mathrm{met}}`$ to be 0.0 and then perform a simulation for the test. We confirmed that the above equation is nearly exactly satisfied in our test simulations and furthermore that even if the $`R_{\mathrm{met}}`$ is not 0.0, the difference in the value of total metallicity between the left and right side in the above equation is negligibly small. Our numerical results do not depend so strongly on $`R_{\mathrm{chemi}}`$, $`R_{\mathrm{met}}`$, $`y_{\mathrm{met}}`$ within plausible and realistic ranges of these parameters.
#### 2.2.2 Population synthesis
It is assumed in the present study that the spectral energy distribution (SED) of a model galaxy is a sum of the SED of stellar particles. The SED of each stellar particle is assumed to be a simple stellar population (SSP) that is a coeval and chemically homogeneous assembly of stars. Thus the monochromatic flux of a galaxy with age $`T`$, $`F_\lambda (T)`$, is described as
$$F_\lambda (T)=\underset{\mathrm{star}}{}F_{\mathrm{SSP},\lambda }(Z_i,\tau _i)\times m_\mathrm{s},$$
(4)
where $`F_{\mathrm{SSP},\lambda }(Z_i,\tau _i)`$ and $`m_\mathrm{s}`$ are a monochromatic flux of SSP of age $`\tau _i`$ and metallicity $`Z_i`$, where suffix $`i`$ identifies each stellar particle, and mass of each stellar particle, respectively. The age of SSP, $`\tau _i`$, is defined as $`\tau _i=Tt_i`$, where $`t_i`$ is the time when a gas particle is converted into a stellar one. The metallicity of SSP is exactly the same as that of the stellar particle, $`Z_i`$, and the summation ($``$) in equation (4) is done for all stellar particles in a model galaxy.
A stellar particle is assumed to be composed of stars whose age and metallicity are exactly the same as those of the stellar particle and the total mass of the stars is set to be the same as that of the stellar particle. Thus the monochromatic flux of SSP at a given wavelength is defined as
$$F_{\mathrm{SSP},\lambda }(Z_i,\tau _i)=_{M_L}^{M_U}\varphi (M)f_\lambda (M,\tau _i,Z_i)𝑑M,$$
(5)
where $`M`$ is mass of a star, $`f_\lambda (M,\tau _i,Z_i)`$ is a monochromatic flux of a star with mass $`M`$, metallicity $`Z_i`$ and age $`\tau _i`$. $`\varphi (M)`$ is a initial mass function (IMF) of stars and $`M_U`$, $`M_L`$ are upper and lower mass limit of IMF, respectively. We here adopt the Salpter IMF with $`M_U`$ = 120$`M_{}`$ and $`M_L`$ = 0.1 $`M_{}`$. In this paper, we use the $`F_{\mathrm{SSP},\lambda }(Z_i,\tau _i)`$ of GISSEL96, which is the latest version of Bruzual & Charlot (1993).
### 2.3 Model of internal dust extinction
#### 2.3.1 Derivation of a SED modified by dust effects
Using the derived each stellar particle’s SED not corrected by dust ($`F_{\mathrm{SSP},\lambda }(Z_i\tau _i)\times m_\mathrm{s}`$ shown in equation (4) ) and stellar and gaseous distribution, we can obtain the SED of a galaxy corrected by dust. From now on the SED of a stellar particle, $`F_{\mathrm{SSP},\lambda }(Z_i,\tau _i)\times m_\mathrm{s}`$, is referred to as $`F_{sed,i}`$ for convenience. We first calculate dust extinction of star light for $`each`$ stellar particle and dust temperature for $`each`$ gaseous particle, based on the three dimensional spatial distribution of stellar and gaseous particles. We then sum each stellar particle’s SED corrected by dust extinction and the dust reemission of each gaseous particle. The method to derive the dust extinction and reemission for each particle consists of the following three steps. Firstly we assume that each stellar particle with the SED $`F_{sed,i}`$ emits rays with the SED for each of the rays equal with $`F_{sed,i}`$/$`N_{ray,i}`$, where $`N_{ray,i}`$ is the total number of the rays. We generate random numbers for each ray and thereby determine the direction of each ray from the stellar particle. Figure 1 represents a schematic explanation of this first step. $`N_{ray,i}`$ is set to be 10 for all stellar particles in the present study. Secondly, we investigate whether each of rays emitted from a stellar particle can penetrate a gas particle for all gas particles and thereby estimate dust extinction of stellar light for each of the rays. This investigation and estimation is done for all stellar particles. In this second step, we regard a ray as penetrating a gas particle and thus being affected by internal dust extinction if the ray passes through the region within the gas cloud radius ($`r_{cl}`$). This second step is also schematically shown in Figure 1. Here gas clouds with the masses of $`m_{cl}`$ are assumed to be spherical and the cloud radius $`r_{cl}`$ is changed following the relation (Kwan & Valdes 1987),
$$r_{cl}=k_{cl}\times 47(\frac{m_{cl}}{4\times 10^6\mathrm{M}_{}})^{1/3}pc$$
(6)
Although physical processes associated with the growth and the disruption of gas clouds have been investigated (e.g., Kwan 1979; Kwan & Valdes 1987; Olson& Kwan 1990a, b), it is not so clear which the most probable and physically reasonable value of $`k_{cl}`$ is for evolving galaxies. A probable value of $`k_{cl}`$ is $``$ 1.0 for disk galaxies (Kwan & Valdes 1987) whereas the typical cloud mass and size depend strongly on time in merging galaxies owing to coalescence and disruption of gas clouds (Olson& Kwan 1990a, b). Olson& Kwan (1990b) demonstrated that although most of the cloud-cloud collisions induced by galaxy merging is disruptive, several very large gas clouds with masses greater than $`10^7`$ $`M_{}`$ can form after merging owing to successive coalescence of colliding clouds. Furthermore the density of gas clouds in ULIRGs considered to be ongoing mergers is suggested to be very high from that in disk galaxies (Aalto et al. 1991a, b), which implies that the size and the mass might be also different from those of isolated disk galaxies. Accordingly it is reasonable for us to assume that $`k_{cl}`$ is a free parameter for the present study. In the present paper, however, we show only the results for the model with $`k_{cl}`$ = 4.68. The $`k_{cl}`$ dependence in major mergers will be described in our future papers.
Absorption of stellar light for each of rays from a stellar particle is modeled according to the following reddening formulation (Black 1987; Mazzei et al. 1992); $`E(BV)=N(\mathrm{H})/4.77\times 10^{21}\mathrm{cm}^2\times (Z_\mathrm{g}/0.02)`$, where $`N(\mathrm{H})`$ and $`Z_\mathrm{g}`$ are gaseous column density and gaseous metallicity, respectively. Since dust absorption is estimated for each of gas particles penetrated by a ray emitted from a stellar particle, the above $`N(\mathrm{H})`$ and $`Z_\mathrm{g}`$ correspond to column density of each gas cloud particle ($`N(\mathrm{H})_i`$) and metallicity of the particle ($`Z_i`$), respectively. The $`Z_i`$ is derived from our numerical simulations and the $`N(\mathrm{H})_i`$ is derived from the relation between cloud mass and size shown in equation (6). Using the derived redding $`E(BV)`$ and extinction law by Cardelli et al. (1989), we adopt the so-called screen model for rays from a stellar particle and calculate the dust absorption of the stellar particle. Thirdly, by assuming the modified black body radiation with the emissivity ($`ϵ`$) law $`ϵ\nu ^{\alpha _{\mathrm{em}}}`$, we determine the dust temperature of a gas particle such that total energy flux of dust absorption is equal to that of dust reemission. In the present study, we show the results of a model with $`\alpha _{\mathrm{em}}`$ = 2. We here do not include the albedo for dust grains, which means that only stars and new stars heat dust. Thus the SED of a galaxy consists of the stellar continuum modified by dust extinction and the reemission of dusty gas.
The method adopted in the present numerical code calculating dust absorption and reemission is essentially the same as the so-called ray tracing method adopted in previous theoretical calculations on radiative transfer problems (e.g. Witt 1977; Efstathiou & Rowan-Robinson 1990). This paper is the first step toward a sophisticated combination of N-body simulations with the already existing ray tracing method. Therefore there are still some problems of the present new code in calculating correctly the SED within feasible time scale for a simulation with the total particle number approximately equal to $`10^5`$. The most significant problem is on the time which we spent in calculating the SED. The time necessary for the SED calculation for each time step in a simulation can be roughly scaled to $`n_s\times n_g\times n_{ray}`$, where $`n_s`$, $`n_g`$, and $`n_{ray}`$ are total particle number of stars, that of gas, and that of rays emitted from each stellar particle. Therefore it is very time-consuming to derive a SED in each time step of a simulation, in particular, for large N-body simulations with both $`n_s`$ and $`n_g`$ larger than $`10^5`$: It is obviously our future work to develop a new method for calculating SEDs of such large N-body systems. Furthermore, our SED calculation becomes also time-consuming, unless we carefully choose the parameter value of $`n_{ray}`$. Although we can more correctly estimate SEDs of galaxy mergers with dusty starburst for the models with larger $`n_{ray}`$, we should adopt a feasible and reasonable value of $`n_{ray}`$, considering that the time for our SED calculation is roughly proportional to $`n_s\times n_g\times n_{ray}`$. In order to determine a plausible value of $`n_{ray}`$, we investigated the dependence of a SED on $`n_{ray}`$ (1 $``$ $`n_{ray}`$ $``$ 50) in a dusty starburst merger model and found that if $`n_{ray}`$ $``$ 10, the differences in SEDs between models with different $`n_{ray}`$ become negligibly small. We thus adopt $`n_{ray}`$ = 10 as a feasible and plausible value for our SED calculations in the present study. Although our new numerical method allows us to derive a SED for an arbitrary distribution of stellar and gaseous component for a dusty galaxy in numerical simulations, it still has some problems such as those described above. Both our method and previous ones (i.e., those by which numerical studies can precisely solve radiative transfer processes and estimate galactic SEDs within a certain spherical symmetric approximation or axisymmetric one) have advantages and disadvantages, and therefore we stress that the present numerical study is complementary to previous ones.
#### 2.3.2 Limitations of the present code
Our numerical model can derive a SED for a galaxy with an arbitrary distribution of stellar and gaseous components and accordingly enable us to examine a SED of a galaxy with very peculiar mass distribution, such as galaxy mergers and forming galaxies. However, our model has the following four limitations in investigating the time evolution of SEDs in galaxies. Firstly, we do not include physical processes related to dust production, grain destruction, and grain growth in interstellar medium in the present model and thus cannot follow the time evolution of dust composition. Dwek (1999) developed a new self-consistent model for the evolution of the compositions and abundances of elements and the dust in galaxies by including condensation, accretion, and destruction processes of grains in the model. By using this model, Dwek (1999) investigated the time evolution of dust composition (e.g., whether interstellar gas is silicate rich or not) and suggested that extinction law, which is an important factor for calculating a galactic SED, depends on the evolution. The method developed by Dwek (1999) is mainly for one-zone models of galactic chemical evolution, accordingly it can not be so simply applied to N-body numerical simulations. It is thus our future study to combine the method provided by Dwek (1999) with our numerical code and thereby to derive a galactic SED in a more self-consistent manner. Furthermore Chang, Schiano, & Wolfe (1987) investigated how the radiation from QSOs can affect the hydrodynamical evolution of dusty interstellar gas, considering ion-field emission, ion sputtering, coupling between gas ions and dust grains, photoionization, and photoelectrical processes. They found that owing to the high sputtering on the grain surface in dusty interstellar gas, most of the grains inside 7 kpc of a QSO host galaxy can be destroyed within $``$ 3 $`\times `$ $`10^7`$ yr. Although their results are not directly related to the dust destruction processes in dusty starburst galaxies, we consider that the essence of their results can be applied to the case of starburst galaxies. The time scale for intense starburst in ULIRGs formed by major mergers is suggested to be $``$ $`10^8`$ yr and the total luminosity of ULIRGs are similar to QSOs (e.g., Sanders et al. 1988), and accordingly the radiation from nuclear strong starburst can also destruct the dust grains in ULIRGs. If the dust grains are efficiently destructed by radiation of starburst, the total flux of far-infrared reemission from dusty starburst can be greatly decreased. Our present model does not include this effects of starburst radiation on dust grains, we stress that the total value of infrared and far-infrared luminosity can be overestimated in major mergers with dusty starburst.
Secondly, we do not consider the effects of small grains with the size between 1 and 10 nm, and consequently can not precisely estimate a SED around $`10^5`$ $`\mathrm{\AA }`$ in a dusty system. Devriendt, Guiderdoni, & Sadat (1999) included contributions from polycyclic aromatic hydrocarbons (PAHs), small grains, and big ones in their model and succeeded in reproducing reasonably well the SEDs of several ULIRGs. To construct multi-component dust model is accordingly the next step of our studies toward the more precise estimation of SEDs of dusty galaxies. Thirdly, the parameter $`k_{cl}`$ is fixed at a possibly reasonable value during a simulation, accordingly our model cannot follow the time evolution of cloud size and mass. Witt & Gordon (1996) demonstrated that the change of the structure of interstellar medium in a dusty system provides the change in effective optical depth and consequently controls the SED of the system. They furthermore suggested that a breakup of large interstellar cloud complexes into numerous smaller clouds and a intercloud medium of enhanced density in a galaxy merger can greatly affect the SED of the merger. Fourthly, we do not include the effects of albedo at all in the present study, though several authors have already investigated in detail the effects of dust albedo (e.g., multiple scattering of stellar light) on SEDs in dusty systems (e.g., Witt & Gordon 1996; Gordon et al. 1997). Although the total number of parameters increase if we include the effects of dust albedo in our numerical simulations, it is our future work to investigate the importance of dust albedo in determining galactic SEDs. Thus our numerical results should be carefully interpreted owing to the above four limitations of the code used in the present study.
### 2.4 Main points of analysis
Mainly important advantages of the present study are the following two. Firstly we can investigate simultaneously morphological, structural, kinematical, and photometric properties at a given time step in a galaxy merger with dusty starburst based on the derived SED of the merger. Secondly, by using the SED derived for a merger at a given z and considering the effects of k-correction, we can investigate how the morphological and photometric properties of the merger change with redshift. Accordingly we mainly investigate the time evolution of morphology, star formation, global colors, luminosity and $`A_\mathrm{V}`$ in a dusty starburst galaxy merger. Furthermore we investigate two-dimensional distributions of colors, luminosity, and $`A_\mathrm{V}`$, which have not been investigated at all in previous theoretical studies of dusty galaxies. In order to discuss apparent morphology of intermediate and high $`z`$ dusty galaxies, we also investigate how dust-enshrouded starburst galaxy mergers at z = 0.4, 1.0, and 1.5 can be seen in the Hubble Space Telescope ($`HST`$). The method to construct the synthesized $`HST`$ images of galactic morphology in the present study is basically the same as that described by Mihos (1995) and Bekki, Shioya, & Tanaka (1999). We adopt Madau’s model (Madau 1995) for intervening absorption of neutral hydrogen and consider k-correction of galactic SEDs in order to derive the synthesized images. For comparison, we also show how dust-enshrouded starburst galaxy mergers at z = 0.4, 1.0, and 1.5 can be seen in the SUBARU which is a Japanese large (8.2m) grand-based telescope. The results on the synthesized $`HST`$ and SUBARU images are presented in §3.3 and for discussing the origin of high $`z`$ faint SCUBA sources and EROs. In order to calculate the SEDs of the merger model at each redshift, we assume that mean ages of old stellar components initially in a merger progenitor disk at the redshift z=0.4, 1.0, and 1.5 are 7.14, 3.80, and 2.46 Gyr, respectively.
We here describe mainly the results of only one merger model with $`k_{cl}`$ = 4.68 (i.e., $`r_{cl}`$ is 0.011 in our units corresponding to 200 pc) because the main purpose of the present study is not to give dependences of SEDs on physical parameters of galaxy mergers but to demonstrate usefulness of the present code in studying theoretically dusty starburst galaxies. From now on this mode is referred to as the standard model. We do not intend to change the values of important free parameters such as $`M_\mathrm{d}`$, $`M_\mathrm{t}`$, $`r_{\mathrm{in}}`$, $`r_\mathrm{p}`$, $`\varphi _i`$, $`\theta _i`$, and $`f_g`$ in the present study. Dependences of the evolution of dusty starburst mergers on parameters will be described in our future papers (Bekki & Shioya 2000b). Although each of the results described in the following sections have much implications on the nature of ULIRGs and can be compared with up-to-date observational results such as near-infrared colors (e.g., Scoville et al. 1999), surface brightness distribution (e.g., Sanders et al. 1999; Scoville et al. 1999), and high resolution optical/near-infrared images of ULIRGs (e.g., Surace, Sanders, & Evans 1999), we do not intend to discuss so extensively the origin and the nature of ULIRGs in the present study: This is simply because the most important purpose of the paper is to demonstrate the importance and the usefulness of the present code in investigating dusty starburst galaxies. We will discuss the origin of ULIRGs in our future papers. In the followings, the cosmological parameters $`\mathrm{H}_0`$ and $`\mathrm{q}_0`$ are set to be 50 km s<sup>-1</sup> Mpc<sup>-1</sup> and 0.5 respectively.
## 3 Result
### 3.1 Evolution of morphological properties and SEDs
Figure 2 and 3 describe time evolution of morphology of the standard model for each of four components, dark halo, stars initially located within two disks, gas, and new stars formed during galaxy merging. From now on, for convenience, the time $`T_m`$ represents the time that has elapsed since the two disks begin to merge. As galaxy merging proceeds, the two disks are strongly disturbed to form a long tidal tail in the disk orbiting in a prograde sense at 0.6 $`<`$ $`T_m`$ $`<`$ 1.1 Gyr. A remarkable tidal tail is not developed in the disk orbiting in a retrograde sense. This one long tidal arm is characteristics of prograde-retrograde mergers. The two disks finally sink into the center of massive dark halos owing to dynamical friction during merging and consequently are completely destroyed by violent relaxation of galaxy merging (1.1 $`<`$ $`T_m`$ $`<`$ 1.7 Gyr). As a result of this, the two disks form an elliptical galaxies with the structure and kinematics similar to the observed ones ($`T_m`$ = 1.7 Gyr). During violent major merging, interstellar gas is very efficiently transferred to the central region of the two disks owing to gaseous dissipation of colliding gas clouds and gravitational torque. Gas accumulated in the central region of the merger is then consumed by massive starburst and consequently converted into new stars. As is shown in Figure 4, the star formation rate becomes maximum ($`378`$ $`M_{}`$ $`\mathrm{yr}^1`$) at $`T_m`$ = 1.3 Gyr when two disks of the merger become very close to suffering from violent relaxation. The maximum star formation rate is roughly two orders of magnitude larger than the mean star formation rate of an isolated disk in the present study and comparable to that required for explaining the strong infrared luminosity observed in ULIRGs. After the strong starburst, star formation rate rapidly declines ($``$ 1 $`M_{}\mathrm{yr}^1`$) within less than 1 Gyr after $`T_m`$ = 1.3, essentially because most of the gas is consumed up by the starburst. These results are qualitatively consistent with those in Mihos & Hernquist (1996).
We here stress that the above one-time starburst in the late phase of major merging is applied only to the present prograde-retrograde merger with the adopted inclination of two disks and internal structure of the disks. Mihos & Bothun (1998) demonstrated that physical details such as internal structure of merger progenitor disks and their initial gas mass fraction are important for the star formation history and the luminosity evolution in ULIRGs, by using imaging Fabry-P$`\stackrel{´}{\mathrm{e}}`$rot observations on four ULIRGs. They furthermore found that the spatial distribution of H$`\alpha `$ emission from starburst regions is very diverse in the four ULIRGs (i.e., some ULIRGs show the strong H$`\alpha `$ emission only in the central region and H$`\alpha `$ emission is quite extended in some ULIRGs) and accordingly suggested that several different factors play a role in triggering starburst in ULIRGs. Considering these important observational results, we suggest that the results on the star formation history and the resultant luminosity evolution described in the present study are only true for some ULIRGs.
Figure 5 and 6 give mass distribution projected onto $`x`$-$`y`$ plane (orbital plane) and $`x`$-$`z`$ plane at $`T_m`$ = 1.1 (pre-starburst phase) and 1.7 Gyr (post-starburst phase) for stars, gas, and new stars. New stars formed mainly by secondary massive starburst are more compactly distributed in the merger than old stars initially located within disks at the pre-starburst (weak starburst) epoch when the two cores in the merger have not yet merged with each other to form an elliptical galaxy. This is essentially because new stars experience much more gaseous dissipation when they were previously gaseous components. Consequently new stars with younger ages are more heavily obscured by dusty gas than old stars during galaxy merging. This result that young stellar components can be preferentially obscured by dust during evolution of galaxy mergers is suggested to be very important for understanding the nature of post-starburst galaxies detected by Smail et al. (1999) in SCUBA surveys of intermediate redshift clusters of galaxies (Shioya & Bekki 2000). Three very small compact stellar clusters composed mainly of new stars (and gas) can be seen above the two cores at $`T_m`$ = 1.1. This results suggests that efficient star formation during galaxy merging can occur not only in the central part of a merger but also in stellar clusters located in the outer part of the merger. As is shown in Figure 6, the morphology of the merger only $``$ 0.4 Gyr after the maximum starburst looks like an elliptical galaxy, which implies that the time scale within which a merger can be seen as an ULIRG with the very peculiar morphology is very short (less than 0.5 Gyr).
Figure 7 describes the SEDs of the merger which we derive based on the mass distribution of stellar and gaseous component shown in Figure 2 and 3 and the age and metallicity distribution of stellar populations of the merger at $`T_m`$ = 0.6, 1.1, 1.3 (the epoch of maximum starburst), 1.7, 2.3, and 2.8 Gyr. We can clearly observe how the dust extinction of interstellar gas can change the shape of the SED of the merger at each time by comparing the results of the model with dust extinction with those without dust extinction. The UV flux with the wavelength less than 3000 $`\mathrm{\AA }`$ rapidly increases during 1.1 $`<`$ $`T_m`$ $`<`$ 1.3 Gyr owing to a large number of young massive stars formed by the strong starburst whereas it decreases during 1.3 $`<`$ $`T_m`$ $`<`$ 2.8 Gyr because of very small star formation rate after the starburst in the model without dust extinction. The infrared and submillimeter fluxes become larger during 1.1 $`<`$ $`T_m`$ $`<`$ 1.3 Gyr in the model with dust extinction. This is firstly because star formation rate, which is closely associated with the total mount of stellar light absorbed by interstellar dust, becomes considerable higher owing to the efficient gas transfer to the central region of the merger and secondly because the density of dusty gas becomes also very high in the later phase of the merging (1.1 $`<`$ $`T_m`$ $`<`$ 1.3 Gyr) so that the gas can heavily obscure the strong starburst. After the maximum starburst at $`T_m`$ = 1.3 Gyr, both infrared and submillimeter fluxes rapidly decline (1.3 $`<`$ $`T_m`$ $`<`$ 1.7 Gyr). This is principally because most of interstellar gas indispensable for strong starburst and dust obscuration is rapidly consumed by star formation in the merger till $`T_m`$ $``$ 1.7 Gyr. Thus the time evolution of SED of a merger depends strongly on that of the star formation rate, which is basically controlled by dynamical evolution of the merger.
Figure 8 and 9 describe the mass distribution and the SED at the epoch of maximum starburst in the standard model, respectively. Figure 8 clearly shows that both new stars and gas are more centrally concentrated than old stellar components, which means that star light from new stars formed by the massive starburst are more heavily obscured by dusty interstellar medium than that of old stellar components. The SED at $`T_m`$ = 1.3 Gyr in Figure 9 is rather similar to the observed SED of typical ULIRGs, which implies that the present dusty starburst merger model can be observed as an ULIRG in the late of galaxy merging. These results in Figure 8 and 9 clearly demonstrate that starburst population in the merger is so heavily obscured by dusty interstellar medium (the mean $`A_V2.46`$ mag) that the dust reemission in far-infrared ranges ($`L_{IR}`$) becomes very strong at the maximum starburst ($`L_{IR}=1.59\times 10^{12}\mathrm{L}_{}`$). The reasons for this heavy dust extinction are firstly that column density of dusty interstellar medium becomes extremely high owing to the strong central accumulation of gas and secondly that the central gaseous metallicity, which is a measure of total amount of dust, also becomes rather large ($`0.04`$) because of efficient and rapid chemical evolution in galaxy merging with secondary starburst.
Based on the derived mass distribution and SEDs at the maximum starburst (at ULIRG phase), we can examine two-dimensional distribution of $`K`$ band surface brightness (Figure 10), rest-frame $`RK`$ color (11), and $`A_V`$ (12), all of which have been observationally investigated in a very extensive manner for the better understanding of the nature of ULIRGs. As is shown in Figure 10, near-infrared surface brightness is the highest (the smallest in number) in the central region within which most of young bright massive stars are located. This is essentially because most of young stars are formed in the central region at $`T_m`$ = 1.3 Gyr. The $`RK`$ color is also the largest in the central part of the merger (Figure 11), because the central starburst components are most heavily obscured by dust owing to the very high gaseous density there. Accordingly, this result implies that an ULIRG formed by major merging can show negative color gradients. As is shown in Figure 12, $`A_V`$ is larger in the central region of the merger than in the outer part, which indicates that dust absorption and its reemission is larger in the central region. This result is consistent with those described in Figure 11.
### 3.2 Photometric and color evolution
Figure 13 describes the time evolution of absolute magnitude from optical to near-infrared wavelength ($`M_V`$, $`M_R`$, and $`M_K`$) for the standard model with and without dust extinction. As galaxy merging proceeds, $`M_V`$, $`M_R`$, and $`M_K`$ gradually rise up owing to the increase of star formation rate (1.1 $`<`$ $`T_m`$ $`<`$ 1.3 Gyr). These $`M_V`$, $`M_R`$, and $`M_K`$ then become $`21.0`$, $`21.5`$, and $`24.2`$ mag, respectively, at the epoch of maximum star formation rate ($`T_m`$ = 1.3 Gyr) in the model with dust extinction. The difference in magnitude between the models with and without dust extinction at the epoch of maximum starburst is the largest for $`V`$ band (2.5 for $`V`$, 2.1 for $`R`$ and 1.3 for $`K`$ band), which reflects the fact that stellar light with shorter wavelength can be more greatly absorbed by dust. The mean value of $`M_K`$ in ULIRGs ($`25.2`$ mag; Sanders & Mirabel 1996) is about 1 mag brighter (smaller in number) than the $`M_K`$ of the present merger at $`T_m`$ = 1.3 Gyr, which probably means that initial total mass of the host disk of a typical ULIRG is about 2 times larger than that of the Galaxy adopted as a merger progenitor disk in the present study. Furthermore, irrespectively of wavelength, the difference in absolute magnitude between the two models in the preset study is the largest at the epoch of maximum starburst. This is essentially because young and very luminous stellar components formed by starburst are the most heavily obscured by dust at the maximum starburst epoch when a larger amount of dusty gas is transfered to form very high density region around the compact starburst. As the star formation rate rapidly declines, $`M_V`$, $`M_R`$, and $`M_K`$ also rapidly decline (1.3 $`<`$ $`T_m`$ $`<`$ 2.8 Gyr) because of the death of young and luminous stellar components (OB stars formed in starburst) and the aging of stellar populations. The present results thus suggest that time variation of absolute magnitude during galaxy merging becomes moderate owing to dust extinction.
Figure 14 shows the time evolution of $`VI`$, $`IK`$, and $`RK`$ colors in the standard model with and without dust extinction. Although the colors in the model without extinction become very blue at the epoch of maximum starburst, these colors are changed into redder ones for a starburst galaxy because of the heavy dust extinction in the merger. The color difference between the models with and without dust extinction is the most remarkable at the epoch of maximum starburst. Figure 15 describes the time evolution of infrared flux at 60 and 100 $`\mu `$m and submillimeter one at 450 and 850 $`\mu `$m for the model with dust extinction. Clearly these fluxes resulting mainly from dust reemission become maximum when the magnitude of starburst is the largest and the density of dusty gas in the central region of the merger is also considerably high. Furthermore these fluxes become an order of magnitude smaller than the maximum values within $``$ 1 Gyr after the maximum starburst. This result implies that the time scale during which a dusty starburst merger can be identified as a submillimeter source detected recently by the SCUBA is very short (an order of $`10^8`$ year). As is shown in Figure 16, the dust temperature of the model becomes maximum ($``$ 32 K) when the star formation rate becomes maximum and total amount of radiation from young massive stars formed by starburst becomes the largest. We here note that the dust temperature is the mean value for all gas particle for each of which the dust temperature is calculated. Figure 16 implies that the dust temperature in mergers with dusty starburst depends strongly on the strength of starburst.
### 3.3 Evolution with redshift
Figure 17, 18, and 19 describe how the difference in the epoch of galaxy merging affects the global colors ($`VI`$, $`IK`$, and $`RK`$, respectively) in observed frame for the standard model. For clarity, we give the dependence of colors on the merging epoch z only for the three different phases of the merger, the pre-starburst epoch ($`T_m`$ = 1.1 Gyr), the maximum starburst one (1.3), and the post-starburst one (1.7). As is shown in Figure 17, optical color $`VI`$ in observed frame is the reddest at the epoch of post-starburst ($`T_m`$ = 1.7 Gyr) at each redshift (i.e., at the epoch of galaxy merging). Furthermore, the $`VI`$ color difference between different redshifts is about 0.7 mag for the epoch of maximum starburst when infrared flux becomes the largest. As has been shown in Bekki et al (1999), the submillimeter flux at 850 $`\mu `$m at z $`>`$ 1.0 for the maximum starburst phase ($`T_m`$ = 1.3 Gyr) is a few mJy that is well above the current detection limit of the SCUBA ($``$ 2 mJy). These results thus imply that the observed color difference in $`VI`$ for the faint SCUBA sources (Smail et al. 1998) is due partly to the difference in redshifts between dusty starburst galaxies detected by the SCUBA. Irrespectively of redshifts, near-infrared colors $`IK`$ and $`RK`$ are the reddest for the post-starburst phase (Figure 18 and 19). Furthermore these colors are redder in higher z for the post-starburst phase ($`T_m`$ = 1.7 Gyr) because of the k-correction of stellar populations in this phase. It should be noted here that the $`RK`$ color for the post-starburst phase becomes larger than $``$ 6.0 around z = 1.5. These results imply that an ERO with the $`RK`$ color larger than 6.0 is more likely to be formed in the post-starburst phase of a gas-rich major merger with the strong dusty starburst and at higher redshifts (z $`>`$ 1 $``$ 2). These results on the redshift evolution of global colors are only for one merger model. Therefore, we stress that although the above results are useful and helpful for understanding the nature of the SCUBA sources and EROs, the results can be true only for some high z dusty mergers. It is our future work to investigate a much larger number of dusty merger models and thereby clarify the nature and the origin of high z dusty starburst galaxies such as the faint SCUBA sources and EROs.
## 4 Discussion
One of advantages of the present new numerical code is that we can investigate both dynamical and photometric evolution of dusty starburst galaxies in an explicitly self-consistent manner. Therefore we can address some important questions that classical one-zone models on galaxy evolution have not yet answered so clearly for the origin of dusty starburst galaxies. Here, we particularly enumerate three problems which the newly developed code is rather helpful and useful for clarifying and thus will be discussed in our future papers.
### 4.1 Origin of faint SCUBA sources
Recent observational studies with the SCUBA have revealed possible candidates of heavily dust-enshrouded starburst galaxies at intermediate and high redshift, which could be counterparts of low redshift ULIRGs (Smail, Ivison, & Blain 1997; Barger et al. 1998; Hughes et al. 1998; Smail et al. 1998; Ivison et al. 1999; Lilly et al. 1999). Although optical morphology of these submillimeter extra-galactic sources should be treated with caution owing to the absence of high resolution sub-mm imaging capability (Richards 1999), more than 50 % of those are suggested to show the indication of galactic interaction and merging (Smail et al. 1998). Furthermore observational studies of an extremely red object ERO J 164502-4626.4 (HR 10) with the redshift of 1.44 by the Hubble Space Telescope and the SCUBA have found that this high redshift galaxy is also a dust-enshrouded starburst galaxy with clear indication of galaxy merging/interaction (Graham & Dey 1996; Cimatti et al. 1998; Dey et al. 1998). Although low redshift ULIRGs are generally considered to be ongoing galaxy mergers with the triggered prominent nuclear activities (starburst or AGN) heavily obscured by dust (Sanders et al. 1988; Sanders & Mirabel 1996), the origin of the faint SCUBA sources is not so clearly understood. In particular, it is not clear whether the observed very large submillimeter luminosity, which could result either from active galactic nuclei (AGN) obscured by dust or from dusty starburst, is due essentially to physical processes associated closely with major galaxy merging.
Concerning this problem, one of tests to assess the validity of merger scenario of the faint SCUBA source formation is to compare the observed morphological, structural, and kinematical properties of the sources with those of a certain theoretical merger model $`at`$ $`the`$ $`considerably`$ $`strong`$ $`starburst`$ $`epoch`$ when submillimeter flux of the merger can exceed the detection limit of the SCUBA. Therefore, theoretical studies should investigate simultaneously the time evolution of submillimeter flux at 850 $`\mu `$m, morphology, structure, and kinematics of a dusty major galaxy merger in order to compare the results with the corresponding observational ones. Although classical one-zone models can investigate photometric evolution of dusty starburst galaxies at each wavelength and thus have contributed greatly to the understanding of dusty galaxies (e.g., Mazzei et al 1992), they cannot predict dynamical evolution of galaxies simultaneously. The present model, on the other hand, can predict not only photometric evolution from UV to submillimeter wavelength but also dynamical evolution in a dusty galaxy. For example, as is shown in the present study, we can predict optical and near-infrared morphology of a merger with the 850 $`\mu `$m flux larger than 2 mJy (the current detection limit of the SCUBA) at a given redshift. Thus we expect that future theoretical studies with our new code will provide valuable clues to the origin of the faint SCUBA sources.
### 4.2 Origin of EROs
Recent observational studies have discovered a significant number of high z EROs with $`RK`$ $``$ $`56`$ (Elston et al. 1988; Cimatti et al. 1998; Dey et al. 1999; Benítez et al. 1998; Cimatti et al. 1999; Smail et al. 1999; Soifer et al. 1999), and accordingly the nature and the origin of EROs have now been discussed very extensively both in observational studies and in theoretical ones. Thompson et al. (1999) argued that EROs are most likely to lie in the redshift range 1 $`<`$ $`z`$ $`<`$ 2 and they represent an important population in high redshift universe. Beckwith et al. (1998) argued that the surface density of EROs with $`RK6`$ mag and $`K19.75`$ mag is 0.14 $`\mathrm{arcmin}^2`$ for EROs with $`RK6`$ and $`K19.75`$. Furthermore Thompson et al. (1999) derived a surface density of EROs with $`RK^{}6`$ mag and $`K^{}19.0`$ of 0.039 $`\pm `$ 0.016 $`\mathrm{arcmin}^2`$ and estimated that the volume density of bright EROs to be as high as that of nearby Seyfert galaxies. There are mainly two possible interpretations for EROs (We here do not intend to discuss whether some EROs are actually low-mass Galactic stars such as main-sequence M stars brown dwarfs). One is that EROs are passively evolving red elliptical galaxies at high $`z`$ (e.g., those recently discovered by the VLT; Benítez et al. 1998), and the other is that EROs are dusty starburst galaxies with starburst components obscured heavily by dust (e.g., HR10; Dey et al. 1999). It still remains unclear which interpretation among the two is more plausible, essentially because spectroscopic studies of EROs have not been so accumulated yet which can reveal unambiguously the redshift of EROs and thus discriminate the effects of aging of stellar populations and those of dust extinction.
Recent observational results on EROs have begun to (or will soon begin to) provide valuable information concerning this problem. For example, Smail el al. (1999) discovered two EROs with $`IK6.0`$ and 6.8 among their SCUBA sources samples, which implies that the submillimeter telescopes with the improved detectability (the detection limit of the order of $`10^2`$ $`\mu `$Jy) can detect submillimeter flux indicative of the obscured dusty starburst in a significant number of EROs. Furthermore statistical spectroscopic studies of EROs by the already existing large grand-based telescopes (e.g., SUBARU with IRCS and FMOS) will reveal the strength of $`\mathrm{H}\alpha `$ emission that is not affected so strongly by dust extinction and thereby clarify star formation rate and history of EROs. It is also possible that the future improved $`HST`$ ACS and the 8m-class grand-based telescopes provide the detailed morphology of an ERO. Accordingly, one of important works for clarifying the origin of EROs is to compare emission line strengths, submillimeter flux, and morphological properties observed (or those that will be observed in near future) in EROs with those predicted by a certain theoretical dusty starburst model. In order to make this comparison possible, theoretical models of dusty starburst galaxies should investigate the time evolution of $`RK`$ color, $`K`$ band magnitude, emission lines, submillimeter flux, and morphology for a dusty starburst galaxy $`in`$ $`a`$ $`self`$-$`consistent`$ $`manner`$. Although the present code still cannot predict emission and absorption line strengths of a dusty starburst galaxy, our new code enables us to investigate most of the above important properties simultaneously in numerical simulations. Thus it is one of important issues for future numerical and theoretical studies with the present new code to clarify the origin of EROs.
## 5 Conclusion
We present a new code which is designed to derive a SED for an arbitrary spatial distribution of stellar and gaseous components for a dusty starburst galaxy. By using this code, we can calculate SEDs based on numerical simulations that can analyze simultaneously dynamical and chemical evolution, structural and kinematical properties, morphology, star formation history, and transfer of metals and dust in interstellar medium for a starburst galaxy. Accordingly, we can investigate variously different properties of starburst galaxies, such as effects of dynamical evolution on galactic SEDs, physical correlations between morphology and SEDs, photometric evolution from UV to submillimeter wavelength, two dimensional distribution of $`A_\mathrm{V}`$, and dependence of SEDs on line-of-site of an observer. Thanks to this code, we can furthermore try to clarify the origin of possible candidates of starburst galaxies, such as low $`z`$ ULIRGs (e.g., Sanders et al. 1988; Sanders & Mirabel 1996), intermediate z ones (e.g., Tran et al. 1999), faint SCUBA sources (Smail, Ivison, & Blain 1997; Barger et al. 1998; Hughes et al. 1998; Smail et al. 1998; Ivison et al. 1999; Lilly et al. 1999), EROs (Elston, Rieke, & Rieke 1988; Dey et al. 1999; Smail et al. 1999), optically faint radii sources detected by VLA (e.g., Richards et al. 1999), Lyman-break galaxies (Steidel et al. 1996; Lowenthal et al. 1997), and emission line galaxies (e.g., Mannucci et al. 1998). By using a new code developed in the present study, we try to answer the following seven questions in our forthcoming papers: (1) When and how a gas-rich major galaxy merger becomes an ULIRG during the dynamical evolution of the merger ? (2) How high z faint SCUBA sources with dusty starburst form and evolve? (3) What are physical conditions for high z dusty galaxies to become EROs ? (4) Are there any evolutionary links between high z possible dusty starburst galaxies, such as faint SCUBA sources, EROs, optically faint radio sources recently detected by VLA, emission line galaxies, and Lyman-break ones ? (5) When does a forming disk galaxy show the strongest submillimeter flux? (6) What physical processes can determine the shapes of SEDs observed in Lyman-break galaxies? (7) How dusty interstellar gas affects apparent morphology of intermediate and high $`z`$ dusty galaxies ? Although this code has some disadvantages in deriving very precise SEDs of dusty starburst galaxies, we believe that this code enables us to grasp some essential ingredients of physical processes related to galaxy formation with starburst at low and high z universe.
We are grateful for the referee G. D. Bothun for giving us valuable comments and suggestions that greatly improve our paper. Y.S. thank to the Japan Society for Promotion of Science (JSPS) Research Fellowships for Young Scientist.
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# Photon Propagation Around Compact Objects and the Inferred Properties of Thermally Emitting Neutron Stars
## 1. INTRODUCTION
Neutron stars in nature appear in various different flavors. They were first unambiguously discovered in radio wavelengths as rotation-powered pulsars and later on, in X-rays, as accretion-powered pulsars and bursters. Such systems, however, constitute only a small fraction of the expected number of neutron stars in the galaxy, as inferred from estimates of supernova rates (e.g., Kaspi 2000). As a result, ongoing searches exist for detecting neutron stars in different manifestations, e.g., as isolated stars accreting from the interstellar medium (e.g., Belloni et al. 1997) or simply cooling by thermal emission (e.g., Pavlov et al. 1996). Different emission mechanisms are thought to operate in different types of neutron stars and possibly even between different wavelength bands in a given system. For at least three distinct classes of neutron stars, which we discuss below, viable models exist in which radiation emerging directly from their surfaces is responsible for their high-energy spectra.
A number of slow X-ray pulsars, often called the anomalous X-ray pulsars (AXPs; Mereghetti & Stella 1995) form a distinct class of neutron stars with soft X-ray spectra but no radio or optical counterparts. Because of their high spin-down rates and association with supernova remnants (SNRs) they are thought to be young and related to the soft $`\gamma `$-ray repeaters, in their quiescent states (see Hurley 2000). It is still an open question whether such objects are neutron stars accreting from a fossil disk (e.g., van Paradijs et al. 1996; Chaterjee, Narayan, & Hernquist 2000), are powered by magnetospheric emission (e.g., Thompson & Duncan 1995, 1996), or even emitting thermally (Heyl & Hernquist 1998).
Another class of objects, potentially related to the AXPs are the compact, non-pulsing, soft X-ray sources that are being discovered within SNRs. The most recent of such compact objects is the central source in Cas A discovered with the Chandra X-ray Observatory (Tananbaum 1999; see also Pavlov et al. 2000; Chakrabarty et al. 2000). Their spectra and luminosities are typical of what is expected for cooling, young neutron stars (see also Pavlov et al. 1996; Zavlin et al. 1998, 1999). However, the absence of coherent pulsations in their X-ray brightness allows for the possibility that they are black holes accreting from fallback material (see, e.g., Pavlov et al. 2000; Chakrabarty et al. 2000).
Finally, weakly-magnetic ($`10^{10}`$ G) accreting neutron stars often show thermal emission from their surfaces during thermonuclear flashes, the so-called Type I X-ray bursts (see Lewin, van Paradijs, & Taam 1996). In six bursters, relatively coherent oscillations are detected during the bursts, at frequencies $`300`$ Hz (see, e.g., Strohmayer et al. 1996). The coherence and stability of the oscillations suggests that they occur at the neutron-star spin frequencies and are caused by the non-uniform pattern of burning in their surface layers.
Understanding the properties of the X-ray emission from such thermally emitting neutron stars, especially in connection to the presence or absence of detectable coherent pulsations, is crucial in assessing the viability of different models. Moreover, comparing model spectra to observations offers the possibility of measuring the radii of neutron stars and hence constraining the properties of neutron-star matter (see Lattimer & Prakash 2000 for a recent discussion). Recently, there has been significant progress in calculating spectra emerging from neutron-star atmospheres with various compositions (see, e.g., Pavlov et al. 1996; Rajagopal, Romani, & Miller 1997). Constraints on the amplitudes of oscillations from compact stars with non-uniform surface emission have also been studied in connection to the observed oscillation amplitudes in cooling radio pulsars (Page 1995), bursters (see, e.g., Miller, & Lamb 1998; Weinberg, Miller, & Lamb 2000), and AXPs (DeDeo, Psaltis, & Narayan 2000).
In this article, we address a number of issues related to the emission from a spinning neutron star with a non-uniform surface brightness. In particular, in §2 we study the effects of the three-dimensional geometry of the systems, phase averaging, and the general relativistic deflection of light. In §3, we demonstrate that, when the radiation emerges from a localized surface area on a rotating star, the radiation flux reaching an observer at infinity may be substantially different compared to the case of isotropic, spherically symmetric emission from a compact star with the same surface area (see also Zavlin, Shibanov, & Pavlov 1995; Leahy & Li 1995). As a result, when such effects are not taken explicitly into account, the hot-spot sizes, inferred from the observed fluxes and temperatures, may be significantly over- or under-estimated. We show that general-relativistic deflection typically reduces this discrepancy, depending on the compactness of the neutron star.
We also note that the surface area on the neutron star, from which the localized emission emerges, determines both the flux that reaches the observer at infinity and the amplitude of pulsations at the stellar spin frequency. In fact, for stars with the same local surface temperature and emerging spectrum, an anticorrelation is expected between the luminosity and pulsation amplitude. Therefore, simultaneous consideration of these two properties for a given system can offer strong constraints on the operating emission mechanisms. In §4, we investigate this property and its implications for the observations of compact X-ray sources with and without detectable pulsations.
## 2. THE BRIGHTNESS OF A SPINNING NEUTRON STAR
In this section, we calculate the phase-averaged radiation flux that reaches an observer at infinity from the surface of a spinning neutron star (or other compact object) with a non-uniform surface brightness, following the procedure outlined by Pechenick, Ftaclas, & Cohen (1983; see DeDeo et al. 2000 for the details of our implementation). Throughout this article, we neglect any interaction of photons with matter between the surface of the neutron star and the observer, as well as any polarization effects (see, e.g., Shaviv, Heyl, & Lithwick 1999). We also assume that the star is slowly rotating, so that its spacetime is described by the Schwarzschild metric. Finally, we set $`c=G=1`$, where $`c`$ is the speed of light and $`G`$ is the gravitational constant.
We specify, as a boundary condition, the specific intensity $`I(\theta ,\varphi ,\theta ^{})`$ of radiation emerging from the neutron-star surface, evaluated at the local rest frame and integrated over all photon energies. The polar coordinates $`(\theta ,\varphi )`$ determine the position on the stellar surface with respect to the rotation axis, while the angle $`\theta ^{}`$ is measured locally with respect to the radial direction. Hereafter, we assume that the emission is isotropic and hence that the specific intensity is independent of $`\theta ^{}`$. This choice leads to the strongest effects of general relativistic light deflection, even though it is not necessarily the most appropriate for thermal emission from a neutron star (see, e.g., Zavlin et al. 1998). Calculating in detail the beaming of radiation requires the knowledge of the temperature stratification of the neutron-star atmosphere and the solution of the resulting radiative-transfer problem, which is beyond the scope of this paper but will reported elsewhere.
The brightness distribution on the neutron star surface depends on the emission process under consideration (i.e., thermal cooling versus localized nuclear burning), the surface profile of its magnetic field (see, e.g., Heyl & Hernquist 1998), and the possible presence of lateral metallicity gradients (see, e.g., Pavlov et al. 2000). For the purposes of our analysis, considering a geometry of two antipodal hot spots of variable size with uniform surface brightness captures the relevant properties of these different processes. Following the notation used for radio pulsars (e.g., Lyne & Graham-Smith 1990), we denote the half-opening angle of each spot by $`\rho `$, its angular distance from the rotation pole by $`\alpha `$, and the emerging constant specific intensity by $`I_{\mathrm{NS}}`$.
The flux measured by an observer at distance $`d`$, whose polar coordinates with respect to the stellar rotation axis are denoted by $`(\beta ,\mathrm{\Phi })`$, is given by (Pechenick et al. 1983, eq. \[3.15\])
$`F_{\mathrm{}}(\beta ,\mathrm{\Phi })`$ $`=`$ $`I_{\mathrm{NS}}\left({\displaystyle \frac{R_{\mathrm{NS}}}{d}}\right)^2\left({\displaystyle \frac{M_{\mathrm{NS}}}{R_{\mathrm{NS}}}}\right)^2`$ (1)
$`\times \left(\sqrt{g_{00}}\right)^4{\displaystyle _0^{x_{\mathrm{max}}}}h[\theta (x,\beta );\rho ,\theta _0]x𝑑x.`$
Here, $`M_{\mathrm{NS}}`$ and $`R_{\mathrm{NS}}`$ are the neutron-star mass and radius, $`g_{00}(12M_{\mathrm{NS}}/R_{\mathrm{NS}})`$ for a Schwarzschild spacetime, $`x_{\mathrm{max}}(R_{\mathrm{NS}}/M_{\mathrm{NS}})/\sqrt{g_{00}}`$, and the function $`h(\theta ;\rho ,\theta _0)`$ is defined in Pechenick et al. (1983, eq. \[3.3\]). The angle $`\theta _0`$ measures the distance on the stellar surface between the center of one of the hot spots and the direction to the observer and is given by
$$\mathrm{cos}\theta _0=\mathrm{sin}\alpha \mathrm{sin}\beta \mathrm{cos}\mathrm{\Phi }+\mathrm{cos}\alpha \mathrm{cos}\beta .$$
(2)
Equation (1) is simply the integral of specific intensities, at the distance of the observer, over the impact parameters $`b=xM_{\mathrm{NS}}`$ of rays that are parallel at radial infinity. The angle $`\theta (x,\beta )`$ at which each parallel ray intersects the stellar surface is found by integrating the photon trajectory from radial infinity to the stellar surface, i.e.,
$$\theta =_0^{M_{\mathrm{NS}}/R_{\mathrm{NS}}}[x^2(12u)u^2]^{1/2}𝑑u$$
(3)
(Pechenick et al. 1983).
The average flux, measured by an observer at infinity, averaged over a time longer than the rotational period of the neutron star, simply is
$$F_{\mathrm{}}=\pi I_{\mathrm{NS}}\left(\sqrt{g_{00}}\right)^2A(\alpha ,\beta ,\rho ,p)\left(\frac{R_{\mathrm{NS}}}{d}\right)^2,$$
(4)
where $`pR_{\mathrm{NS}}/2M_{\mathrm{NS}}`$ and we have defined
$$A(\alpha ,\beta ,\rho ,p)\left(\frac{g_{00}}{8\pi ^2p^2}\right)_0^{2\pi }𝑑\mathrm{\Phi }_0^{x_{\mathrm{max}}}h[\theta (x,\beta );\alpha ,\theta _0]x𝑑x.$$
(5)
When the local radiation spectrum emerging from the stellar surface is that of a blackbody of temperature $`T_{\mathrm{NS}}`$, then $`\pi I_{\mathrm{NS}}=\sigma T_{\mathrm{NS}}^4`$, where $`\sigma `$ is the Stefan-Boltzmann constant.
If the emission were spherically symmetric, an observer at radial infinity would measure a specific intensity $`I_{\mathrm{}}=(g_{00})I_{\mathrm{NS}}`$ and a blackbody temperature $`T_{\mathrm{}}=(\sqrt{g_{00}})T_{\mathrm{NS}}`$. In the absence of any prior information about the opening angle of the hot spots and their orientation, an observer at infinity would therefore infer for the emitting region a surface area
$$S_{\mathrm{}}\frac{4d^2F_{\mathrm{}}}{I_{\mathrm{}}}=4\pi d^2\frac{F_{\mathrm{}}}{\sigma T_{\mathrm{}}^4}.$$
(6)
It is then customary to assume a given neutron-star mass and radius and correct for the effect of gravitational redshifts as (see, e.g., Lattimer & Prakash 2000)
$$S_{\mathrm{inf}}(g_{00})S_{\mathrm{}}=(g_{00})4\pi d^2\frac{F_{\mathrm{}}}{\sigma T_{\mathrm{}}^4}.$$
(7)
Given that the real surface area of two polar caps of half-opening angle $`\rho `$ is $`S_{\mathrm{pc}}=4\pi (1\mathrm{cos}\rho )R_{\mathrm{NS}}^2`$, we conclude that the error in the estimate of the emitting area is
$$\frac{S_{\mathrm{inf}}}{S_{\mathrm{pc}}}=\frac{1}{1\mathrm{cos}\rho }A(\alpha ,\beta ,\rho ,p).$$
(8)
As a consequence, estimating the radius of each of the two polar caps according to $`2\pi D_{\mathrm{inf}}^2S_{\mathrm{inf}}`$, leads to an error with respect to the real polar cap radius $`D_{\mathrm{pc}}=\rho R_{\mathrm{NS}}`$ equal to
$$\frac{D_{\mathrm{inf}}}{D_{\mathrm{pc}}}=\sqrt{\frac{2A(\alpha ,\beta ,\rho ,p)}{\rho ^2}}.$$
(9)
Finally, assuming that the whole neutron star is emitting uniformly and using equation (7) to infer its radius, results in an error equal to
$$\frac{R_{\mathrm{inf}}}{R_{\mathrm{NS}}}=A(\alpha ,\beta ,\rho ,p).$$
(10)
## 3. RESULTS
### 3.1. Uncertainties in the Inferred Surface Areas of the Emission Regions
We present in this section the uncertainties introduced to the estimates of the spectroscopically inferred polar-cap surface areas by the three-dimensional geometry of the problem, the effects of phase averaging, and the general relativistic light deflection.
Figures 1 and 2 show the ratio of the inferred to the intrinsic surface areas of the polar caps, for different opening
angles, neutron-star radii, and orientations of the rotation axis with respect to the polar caps and the observer. For neutron-stars that are not very relativistic ($`p=4`$), the inferred surface areas can be significantly over- or under-estimated, depending on the relative orientations of the polar caps and the observer. In order to understand this effect, we calculate explicitly the ratio $`S_{\mathrm{inf}}/S_{\mathrm{pc}}`$ for the limiting case of a Newtonian star, an infinitesimally small emitting area, and two specific orientations.
When $`\alpha =\beta =0^{}`$, one polar cap always appears at the geometric center of the stellar disk, and the flux measured at infinity is time independent. In this case, $`\mathrm{cos}\theta _0=1`$ and, by symmetry,
$$h[\theta (\chi ,\beta );\rho 0,\theta _0=0^{}]=2\pi .$$
(11)
Using these values in evaluating the integral (5), we obtain
$$A(\alpha =\beta =0^{},\rho 0,p\mathrm{})=2(1\mathrm{cos}\rho )$$
(12)
which gives $`S_{\mathrm{inf}}/S_{\mathrm{pc}}=2`$. When, on the other hand, $`\alpha =0^{}`$ and $`\beta =90^{}`$, $`\mathrm{cos}\theta _0=0`$, the polar caps are being viewed only at grazing angles. In this case,
$$h[\theta (\chi ,\beta );\rho 0,\theta _0=90^{}]\chi 0,$$
(13)
and
$$A(\alpha =0^{},\beta =90^{},\rho 0,p\mathrm{})0$$
(14)
and hence $`S_{\mathrm{inf}}/S_{\mathrm{pc}}0`$. As a result, for a Newtonian star and an infinitesimally small emitting area,
$$0\frac{S_{\mathrm{inf}}}{S_{\mathrm{pc}}}2.$$
(15)
For increasingly more compact neutron stars, i.e., for decreasing values of $`p`$, the error in the estimate of the emitting area decreases substantially, reaching $`5\%`$ when $`p=2`$. This is the result of the strong gravitational light bending near the neutron star surface, which efficiently redistributes the emitting photons to almost all directions of propagation, mimicking the spherically symmetric case.
Figure 3 shows the maximum and minimum of the ratio $`S_{\mathrm{inf}}/S_{\mathrm{pc}}`$ for all possible orientations of the hot spot and the observer, as a function of the fractional emitting surface area ($`S_{\mathrm{pc}}/4\pi R_{\mathrm{NS}}^2=2\mathrm{\Omega }_{\mathrm{pc}}/4\pi `$) of the neutron-star surface, for different neutron-star radii. As expected, in all cases, the maximum corresponds to $`\alpha =\beta =0^{}`$ and the minimum to $`\alpha =0^{}`$ and $`\beta =90^{}`$, and they both converge to unity as $`\mathrm{\Omega }_{\mathrm{pc}}2\pi `$. However, for small polar caps and typical neutron-star masses and radii, there is a factor of $`2`$ spread in the systematic uncertainty in the estimated area of the emitting region.
### 3.2. Trends Between the Pulse Fractions and Luminosities of Thermally Emitting Neutron Stars
The X-ray brightness of a spinning neutron-star with a non-uniform surface emission shows pulsations at the stellar spin frequency and its harmonics. Just as in the case of the flux observed at infinity, the amplitude of pulsations as well as the harmonic structure also depend strongly on the brightness distribution on the stellar surface and the degree of gravitational light bending. Therefore, inevitable
trends exist between the pulse fractions and the brightness of a source, as we discuss below.
For the case of two antipodal polar caps discussed here, the amplitude of pulsations from a neutron star of a given radius decreases with increasing polar-cap surface area and increasing mass. This is shown in Figure 4, where the maximum pulse fraction, defined as
$$PF\frac{F(\mathrm{\Phi })|_{\mathrm{max}}F(\mathrm{\Phi })|_{\mathrm{min}}}{F(\mathrm{\Phi })|_{\mathrm{max}}+F(\mathrm{\Phi })|_{\mathrm{min}}},$$
(16)
which occurs for an orthogonal rotator ($`\alpha =\beta =90^{}`$), is plotted against the fractional surface area of the emitting region, for different neutron-star radii.
As it has already been discussed extensively in the literature, for the isotropic beaming we use here, general relativistic light bending suppresses the pulse fraction below $`35\%`$, even for infinitesimally small polar-cap sizes. As a result, detection of a pulsation from a neutron star with a larger amplitude severely constrains any such models of emission from the stellar surface. However, these constraints may become even tighter if the observed brightness of the neutron star is also taken into account. For the same local radiation spectrum, small polar-cap sizes correspond to large pulse fractions but weak radiation fluxes and vice versa. As a result, in this case, the maximum pulse fraction for a bright source is significantly smaller compared to the maximum pulse fraction for a faint source.
Inferring the fractional surface area of the emitting region, and thus the source brightness, requires an a priori knowledge of the orientation of the polar caps and the observer with respect to the rotation axis, which are almost always unknown. For example, a given observed source brightness may be the result of a small polar cap viewed from a favorable orientation ($`\alpha \beta 0^{}`$) or of a larger polar cap viewed from an unfavorable orientation ($`\alpha 0^{}`$, $`\beta 90^{}`$). Therefore, since the observed brightness does not have an one-to-one correspondence with the fractional emitting surface area, the constraints plotted in Figure 4 cannot be compared directly with observations. However, although both configurations in the above example may produce the same pulse-average flux as measured by the observer, the first configuration typically produces a significantly smaller pulse fraction. As a result, the combination of these two properties, namely the source brightness and pulse fraction, provide us with a useful diagnostic tool, as we discuss below.
We can simultaneously account for all the above effects if we assume a particular model of the local radiation spectrum and search for the maximum pulse fraction as a function of the brightness of the source measured by an observer at infinity. As an example, we assume that the local radiation spectrum at each point on the polar caps is that of a blackbody of temperature $`T_{\mathrm{NS}}`$ with isotropic beaming. An observer at infinity, would observe a radiation flux $`F_{\mathrm{}}`$, and making the assumption that the emission is spherically symmetric, would infer a luminosity
$$L_{\mathrm{}}=4\pi d^2F_{\mathrm{}}.$$
(17)
According to equation , this luminosity is
$$L_{\mathrm{}}=4\pi R_{\mathrm{NS}}^2\sigma T_{\mathrm{NS}}^4\left(\sqrt{g_{00}}\right)^2A(\alpha ,\beta ,\gamma ,r),$$
(18)
where $`T_{\mathrm{NS}}=T_{\mathrm{}}/(\sqrt{g_{00}})`$. We define an apparent brightness as
$$l\left(\frac{4\pi d^2F_{\mathrm{}}}{10^{36}\text{erg s}^1}\right)\left(\frac{T_{\mathrm{}}}{0.5\text{keV}}\right)^4,$$
(19)
which is a function of only observed quantities and is independent of the neutron-star temperature in our model calculations. We use this quantity as a measure of the brightness of the star at infinity and, in Figure 5, we plot against it the maximum pulse fraction for a 10 km neutron star emitting a blackbody spectrum. This figure shows clearly that, even though emission from a neutron-star surface can in principle produce pulse fractions as large as $`30`$% (for very small hot spots), it can never produce, e.g., a pulse fraction of $`10`$% simultaneously with an inferred luminosity of $`10^{36}`$ erg s<sup>-1</sup> for a 0.5 keV blackbody temperature.
## 4. DISCUSSION
In this paper we have shown that the inferred properties of thermally emitting neutron stars, i.e., their emitting surface areas and pulse fractions, are significantly affected by the three-dimensional geometry of the systems, phase-averaging effects, and general relativistic light bending. The uncertainties introduced by these effects can be significant (with a spread of $`2`$), especially at the limit of small emitting surface areas. However, even when most of the neutron star surface is emitting, the uncertainties in estimating the neutron-star radius are of the same order as the $`5`$% level required (Lattimer & Prakash 2000) for constraining the equation of state of neutron-star matter.
We have also argued that for a given local radiation spectrum emerging from a bright spot on the stellar surface, faint sources can give rise to larger pulse fractions than brighter sources. The maximum pulse fraction as a function of the brightness of the neutron star, measured by $`l`$ (cf. eq. ), provides a diagnostic tool for distinguishing between different emission models. These derived constraints depend both on the local radiation spectrum (e.g., thermal versus non-thermal emission) and its beaming and are, therefore, different for different classes of models. Moreover, both quantities are measurable from spectral and timing observations and can be directly compared to the calculated constraints. The high observed pulse fractions of AXPs (see, e.g., Chakrabarty et al. 2000) and X-ray bursters (Strohmayer et al. 1998) have been shown to strongly constrain emission models and such constraints can become only tighter when the apparent luminosities of these systems are also taken into account. As an example, we consider the source RXS J1708$``$40, which has been identified as an AXP (Sugizaki et al. 1997). Fitting a blackbody spectrum to ASCA observations of this source gives a temperature of $`T_{\mathrm{}}=0.41\pm 0.03`$ keV and a blackbody flux of $`4\pi d^2F_{\mathrm{}}=3.16\times 10^{35}`$ erg s<sup>-1</sup> (Sugizaki et al. 1997; Chakrabarty et al. 2000). The inferred apparent brightness of this source is, therefore, $`l=0.8`$ and combined with the observed $`50`$% pulse fraction (Sugizaki et al. 1997) is inconsistent with isotropic thermal emission from the neutron-star surface (see Fig. 5).
Finally, our results also have a number of important implications for soft X-ray surveys, in supernova remnants, for young, cooling neutron stars that emit thermally. For a given spectral shape, the brightest sources, which would be more easily detectable, correspond to smaller pulse fractions, while strong pulsations are only expected from dimmer sources. This effect is shown in Figure 6, where the fraction of systems $`N_{\mathrm{obs}}/N_{\mathrm{tot}}`$ with a pulse fraction at infinity larger than a threshold $`PF_0`$ is plotted against the fractional emitting surface area. For this purpose, we assume a random orientation of the magnetic inclination and the inclination to the observer for a sample of systems and define
$$\frac{N_{\mathrm{obs}}}{N_{\mathrm{tot}}}(PF_0)=_{\alpha =0}^{\pi /2}_{\beta =0}^{\pi /2}X[PF(\alpha ,\beta );PF_0]\mathrm{sin}\alpha d\alpha \mathrm{sin}\beta d\beta ,$$
(20)
where the step function $`X(PF;PF_0)`$ is defined such that it is unity when $`PF>PF_0`$ and zero at all other values of the pulse fraction. For realistic neutron star masses and radii ($`p23`$), only a very small fraction of sources shows pulsations that are detectable at a significant level ($`1020`$%) and this fraction drops rapidly with increasing apparent luminosity.
If the central source in the remnant Cas A is a young, cooling neutron star, its surface brightness distribution cannot be uniform, as inferred from fitting thermal models to the observed countrate spectra (Pavlov et al. 2000; Chakrabarty et al. 2000). However, as Figure 6 shows, this property is not inconsistent with the $`30`$% upper limit on its pulse fraction (Chakrabarty et al. 2000). For the assumed isotropic beaming of the emerging radiation and the very small fractional surface area inferred for the central source in Cas A (Pavlov et al. 2000; Chakrabarty et al. 2000), less than half of the systems would have been detected with pulse fractions higher than the detection threshold, if $`p=4`$. For more realistic neutron-star properties ($`p=2,3`$), no system would show such a large pulse fraction. Therefore, the absence of detectable pulsations from this source is not a strong argument against its identification with a spinning neutron star.
We thank Deepto Chakrabarty, Lars Hernquist, and Ramesh Narayan for many useful discussions. D. P. acknowledges the support of a post-doctoral fellowship from the Smithsonian Insitution. F. Ö. acknowledges support of NSF Grant AST-9820686.
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# Quantum Games
## I Introduction
Game theory is the theory of decision making, which provides powerful tools for investigating situations in which several parties make decisions according to their personal interest . It gives an account of how the parties would decide in a situation which involves contest, rivalry, or struggle. Such have been found to be relevant to social sciences, biology, or economics. Of particular interest to the theory are games of incomplete information in which the parties may choose their plan of action with complete knowledge of the situation on rational grounds, but without knowing what decision the other parties have actually taken.
One important two player game is the so-called Prisoners’ Dilemma . In this game two players – in the following referred to as Alice and Bob – can independently decide whether they intend to ”cooperate” or ”defect”. Being well aware of the consequences of their decisions the players obtain a certain pay-off according to their respective choices. This pay-off provides a quantitative characterisation of their personal preferences. Both players are assumed to want to maximise their individual pay-off, yet they must pick their choice without knowing the other player’s decision. Fig. 1 indicates the pay-off of Alice and Bob. The players face a dilemma since rational reasoning in such a situation dictates the players to defect, although they would both benefit from mutual cooperation. As Alice is better off with defection regardless of Bob’s choice, she will defect. The game being symmetric, the same argument applies to Bob.
Formally, Alice has two basic choices, meaning that she can select from two possible strategies $`s_A=C`$ (cooperation) and $`s_A=D`$ (defection). Bob may also take $`s_B=C`$ or $`s_B=D`$. The game is defined by these possible strategies on the one hand, and on the other hand by a specification of how to evaluate the pay-off once the combination of chosen strategies $`(s_A,s_B)`$ is known, i.e., the utility functions mapping $`(s_A,s_B)`$ on a number . The expected pay-off quantifies the preference of the players.
In these lecture notes the idea of identifying strategic moves with quantum operations as introduced in Refs. is further developed. This approach appears to be fruitful in at least two ways . On the one hand several recently proposed applications of quantum information theory can already be conceived as competitive situations where several parties with more or less opposed motives interact. These parties may, for example, apply quantum operations on a bi-partite quantum system . In the same context, quantum cloning has been formulated as a game between two players . Similarly, eavesdropping in quantum cryptography can be regarded as a game between the eavesdropper and the sender, and there are similarities of the extended form of quantum versions of games and quantum algorithms . On the other hand a generalisation of the theory of decisions into the domain of quantum probabilities seems interesting, as the roots of game theory are partly in probability theory. In this context it is of interest to investigate what solutions are attainable if superpositions of strategies are allowed for.
Game theory does not explicitly concern itself with how the information is transmitted once a decision is taken. Yet, it should be clear that the practical implementation of any (classical) game inevitably makes use of the the exchange of voting papers, faxes, emails, ballots, and the like. In the Prisoners’ Dilemma, e.g., the two parties have to communicate with an advocate by talking to her or by writing a short letter on which the decision is indicated. Bearing in mind that a game is also about the transfer of information, it becomes legitimate to ask what happens if these carriers of information are taken to be quantum systems, quantum information being a fundamental notion of information.
By classical means a two player binary choice game may be played as follows: An arbiter takes two coins and forwards one coin each to the players. The players then receive their coin with head up and may keep it as it is (”cooperate”) or turn it upside down so that tails is up (”defection”). Both players then return the coins to the arbiter who calculates the players’ final pay-off corresponding to the combination of strategies he obtains from the players. Here, the coins serve as the physical carrier of information in the game. In a quantum version of such a game quantum systems would be used as such carriers of information. For a binary choice two player game an implementation making use of minimal resources involves two qubits as physical carriers.
## II Quantum Games and Quantum Strategies
Any quantum system which can be manipulated by two parties or more and where the utility of the moves can be reasonably quantified, may be conceived as a quantum game.
A two-player quantum game $`\mathrm{\Gamma }=(,\rho ,S_A,S_B,P_A,P_B)`$ is completely specified by the underlying Hilbert space $``$ of the physical system, the initial state $`\rho 𝒮()`$, where $`𝒮()`$ is the associated state space, the sets $`S_A`$ and $`S_B`$ of permissible quantum operations of the two players, and the utility functionals $`P_A`$ and $`P_B`$, which specify the utility for each player. A quantum strategy $`s_AS_A`$, $`s_BS_B`$ is a quantum operation, that is, a completely positive trace-preserving map mapping the state space on itself . The quantum game’s definition also includes certain implicit rules, such as the order of the implementation of the respective quantum strategies. Rules also exclude certain actions, as the alteration of the pay-off during the game.
The quantum games proposed in Refs. , , and can be cast into this form. Also, the quantum cloning device as described in can be said to be a quantum game in this sense. A quantum game is called zero-sum game, if the expected pay-offs sum up to zero for all pairs of strategies, that is, if $`P_A(s_A,s_B)=P_B(s_A,s_B)`$ for all $`s_AS_A`$, $`s_BS_B`$. Otherwise, it is called a non-zero sum game.
It is natural to call two quantum strategies of Alice $`s_A`$ and $`s_A^{}`$ equivalent, if $`P_A(s_A,s_B)=P_A(s_A^{},s_B)`$ and $`P_B(s_A,s_B)=P_A(s_A^{},s_B)`$ for all possible $`s_B`$. That is, if $`s_A`$ and $`s_A^{}`$ yield the same expected pay-off for both players for all allowed strategies of Bob. In the same way strategies $`s_B`$ and $`s_B^{}`$ of Bob will be identified.
A solution concept provides advice to the players with respect to the action they should take. The following solution concepts will be used in the remainder of this lecture. These definitions are fully analogous to the corresponding definitions in standard game theory .
A quantum strategy $`s_A`$ is called dominant strategy of Alice if
$$P_A(s_A,s_B^{})P_A(s_A^{},s_B^{})$$
(1)
for all $`s_A^{}S_A`$, $`s_B^{}S_B`$. Analogously we can define a dominant strategy for Bob. A pair $`(s_A,s_B)`$ is said to be an equilibrium in dominant strategies if $`s_A`$ and $`s_B`$ are the players’ respective dominant strategies. A combination of strategies $`(s_A,s_B)`$ is called a Nash equilibrium if
$`P_A(s_A,s_B)`$ $`P_A(s_A^{},s_B),`$ (3)
$`P_B(s_A,s_B)`$ $`P_B(s_A,s_B^{})`$ (4)
for all $`s_A^{}S_A`$, $`s_B^{}S_B`$. A pair of strategies $`(s_A,s_B)`$ is called Pareto optimal, if it is not possible to increase one player’s pay-off without lessening the pay-off of the other player.
A solution in dominant strategies is the strongest solution concept for a non-zero sum game. In the Prisoner’s Dilemma defection is the dominant strategy, as it is favourable regardless what strategy the other party picks. Typically, however, the optimal strategy depends on the strategy chosen by the other party. A Nash equilibrium implies that neither player has a motivation to unilaterally alter his or her strategy from this equilibrium solution, as this action will lessen his or her pay-off. Given that the other player will stick to the strategy corresponding to the equilibrium, the best result is achieved by also playing the equilibrium solution. The concept of Nash equilibria is of paramount importance to studies of non-zero-sum games. It is, however, only an acceptable solution concept if the Nash equilibrium is unique. For games with multiple equilibria the application of a hierarchy of natural refinement concepts may finally eliminate all but one of the Nash equilibria. Note that a Nash equilibrium is not necessarily efficient. In the Prisoners’ Dilemma for example there is a unique equilibrium, but it is not Pareto optimal, meaning that there is another outcome which would make both players better off.
## III Two-qubit quantum games
In the subsequent investigation we turn to specific games where the classical version of the game is faithfully entailed in the quantum game. In a quantum version of a binary choice game two qubits are prepared by a arbiter in a particular initial state, the qubits are sent to the two players who have physical instruments at hand to manipulate their qubits in an appropriate manner. Finally, the qubits are sent back to the arbiter who performs a measurement to evaluate the pay-off. For such a bi-partite quantum game the system of interest is a quantum system with underlying Hilbert space $`=_A_B`$, $`_A=_B=^2`$, and associated state space $`𝒮()`$. Quantum strategies $`s_A`$ and $`s_B`$ of Alice and Bob are local quantum operations acting in $`_A`$ and $`_B`$ respectively . That is, Alice and Bob are restricted to implement their respective quantum strategy $`s_A`$ and $`s_B`$ on their qubit only. In this step they may choose any quantum strategy that is included in the set of strategies $`S`$. They are both well aware of the set $`S`$, but they do not know which particular quantum strategy the other party would actually implement. As the application of both quantum strategies amounts to a map $`s_As_B:𝒮()𝒮()`$, after execution of the moves the system is in the state
$$\sigma =(s_As_B)(\rho ).$$
(5)
Particularly important will be unitary operations $`s_A`$ and $`s_B`$. They are associated with unitary operators $`U_A`$ and $`U_B`$, written as $`s_AU_A`$ and $`s_BU_B`$. In this case the final state $`\sigma `$ is given by
$$\sigma =(U_AU_B)\rho (U_AU_B)^{}.$$
(6)
From now on both the sets of strategies of Alice and Bob and the pay-off functionals are taken to be identical, that is,
$$S_A=S_B=S\text{ and }P_A=P_B=P,$$
(7)
such that both parties face the same situation.
The quantum game $`\mathrm{\Gamma }=(^2^2,\rho ,S,S,P,P)`$ can be played in the following way: The initial state $`\rho `$ is taken to be a maximally entangled state in the respective state space. In order to be consistent with Ref. let $`\rho =|\psi \psi |`$ with
$$|\psi =(|00+i|11)/\sqrt{2},$$
(8)
where the first entry refers to $`_A`$ and the second to $`_B`$. The two qubits are forwarded to the arbiter who performs a projective selective measurement on the final state $`\sigma `$ with Kraus operators $`\pi _{CC}`$, $`\pi _{CD}`$, $`\pi _{DC}`$, and $`\pi _{DD}`$, where
$`\pi _{CC}`$ $`=`$ $`|\psi _{CC}\psi _{CC}|,|\psi _{CC}=(|00+i|11)/\sqrt{2},`$ (10)
$`\pi _{CD}`$ $`=`$ $`|\psi _{CD}\psi _{CD}|,|\psi _{CD}=(|01i|10)/\sqrt{2},`$ (11)
$`\pi _{DC}`$ $`=`$ $`|\psi _{DC}\psi _{DC}|,|\psi _{DC}=(|10i|01)/\sqrt{2},`$ (12)
$`\pi _{DD}`$ $`=`$ $`|\psi _{DD}\psi _{DD}|,|\psi _{DD}=(|11+i|00)/\sqrt{2}.`$ (13)
According to the outcome of the measurement, a pay-off of $`A_{CC}`$, $`A_{CD}`$, $`A_{DC}`$, or $`A_{DD}`$ is given to Alice, Bob receives $`B_{CC}`$, $`B_{CD}`$, $`B_{DC}`$, or $`B_{DD}`$. The utility functionals, also referred to as expected pay-off of Alice and Bob, read
$`P_A(s_A,s_B)`$ $`=`$ $`A_{CC}\mathrm{tr}[\pi _{CC}\sigma ]+A_{CD}\mathrm{tr}[\pi _{CD}\sigma ]+A_{DC}\mathrm{tr}[\pi _{DC}\sigma ]+A_{DD}\mathrm{tr}[\pi _{DD}\sigma ],`$ (15)
$`P_B(s_A,s_B)`$ $`=`$ $`B_{CC}\mathrm{tr}[\pi _{CC}\sigma ]+B_{CD}\mathrm{tr}[\pi _{CD}\sigma ]+B_{DC}\mathrm{tr}[\pi _{DC}\sigma ]+B_{DD}\mathrm{tr}[\pi _{DD}\sigma ].`$ (16)
It is important to note that the Kraus operators are chosen in such a way that the classical game is fully entailed in the quantum game: The classical strategies $`C`$ and $`D`$ are associated with particular unitary operations,
$$C\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),D\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(17)
$`C`$ does not change the state at all, $`D`$ implements a ”spin-flip”. If both parties stick to these classical strategies, Eq. (15) and Eq. (16) guarantee that the expected pay-off is exactly the pay-off of the corresponding classical game defined by the numbers $`A_{CC}`$, $`A_{CD}`$, $`A_{DC}`$, $`A_{DD}`$, $`B_{CC}`$, $`B_{CD}`$, $`B_{DC}`$, and $`B_{DD}`$. E.g., if Alice plays $`C`$ and Bob chooses $`D`$, the state $`\sigma `$ after implementation of the strategies is given by
$$\sigma =(CD)(\rho )=|\psi _{CD}\psi _{CD}|,$$
(18)
such that Alice obtains $`A_{CD}`$ units and Bob $`B_{CD}`$ units pay-off (see Fig. 1). In this way the peculiarities of strategic moves in the quantum domain can be adequately studied. The players may make use of additional degrees of freedom which are not available with randomisation of the classical strategies, but they can also stick to mere classical strategies. This scheme can be applied to any two player binary choice game and is to a high extent canonical.
### A Prisoners’ Dilemma
We now investigate the solution concepts for the quantum analogue of the Prisoners’ Dilemma (see Fig. 1) ,
$`A_{CC}=B_{CC}=3,A_{DD}=B_{DD}=1,`$ (20)
$`A_{CD}=B_{DC}=0,A_{DC}=B_{CD}=5.`$ (21)
In all of the following sets of allowed strategies $`S`$ the classical options (to defect and to cooperate) are included. Several interesting sets of strategies and concomitant solution concepts will at this point be studied. The first three subsections involve local unitary operations only, while in the last subsection other quantum operations are considered as well.
#### 1 One-parameter set of strategies.
The first set of strategies $`S^{(CL)}`$ involves quantum operations $`s_A`$ and $`s_B`$ which are local rotations with one parameter. The matrix representation of the corresponding unitary operators is taken to be
$$U(\theta )=\left(\begin{array}{cc}\mathrm{cos}(\theta /2)& \mathrm{sin}(\theta /2)\\ \mathrm{sin}(\theta /2)& \mathrm{cos}(\theta /2)\end{array}\right)$$
(22)
with $`\theta [0,\pi ]`$. Hence, in this simple case, selecting strategies $`s_A`$ and $`s_B`$ amounts to choosing two angles $`\theta _A`$ and $`\theta _B`$. The classical strategies of defection and cooperation are also included in the set of possible strategies, $`CU(0)`$, $`DU(\pi )`$. An analysis of the expected pay-offs $`P_A`$ and $`P_B`$,
$`P_A(\theta _A,\theta _B)`$ $`=`$ $`3|\mathrm{cos}(\theta _A/2)\mathrm{cos}(\theta _B/2)|^2+5|\mathrm{cos}(\theta _B/2)\mathrm{sin}(\theta _A/2)|^2`$ (24)
$`+`$ $`|\mathrm{sin}(\theta _A/2)\mathrm{sin}(\theta _B/2)|^2,`$ (25)
$`P_B(\theta _A,\theta _B)`$ $`=`$ $`3|\mathrm{cos}(\theta _A/2)\mathrm{cos}(\theta _B/2)|^2+5|\mathrm{sin}(\theta _B/2)\mathrm{cos}(\theta _A/2)|^2`$ (26)
$`+`$ $`|\mathrm{sin}(\theta _A/2)\mathrm{sin}(\theta _B/2)|^2,`$ (27)
shows that this game is the classical Prisoners’ Dilemma game . The pay-off functions are actually identical to the analogous functions in the ordinary Prisoners’ Dilemma with mixed (randomised) strategies, where cooperation is chosen with the classical probability $`p=\mathrm{cos}^2(\theta /2)`$. The inequalities
$`P_A(D,s_B)`$ $``$ $`P_A(s_A,s_B),`$ (29)
$`P_B(s_A,D)`$ $``$ $`P_B(s_A,s_B)`$ (30)
hold for all $`s_A,s_BS^{(CL)}`$, therefore, $`(D,D)`$ is an equilibrium in dominant strategies and thus the unique Nash equilibrium. As explained in the introduction this equilibrium is far from being efficient, because $`P_A(D,D)=P_B(D,D)=1`$ instead of the Pareto optimal pay-off which would be 3.
#### 2 Two-parameter set of strategies
A more general set of strategies is the following two-parameter set $`S^{(TP)}`$. The matrix representation of operators corresponding to quantum strategies from this set is given by
$$U(\theta ,\varphi )=\left(\begin{array}{cc}e^{i\varphi }\mathrm{cos}(\theta /2)& \mathrm{sin}(\theta /2)\\ \mathrm{sin}(\theta /2)& e^{i\varphi }\mathrm{cos}(\theta /2)\end{array}\right)$$
(31)
with $`\theta [0,\pi ]`$ and $`\varphi [0,\pi /2]`$. Selecting a strategy $`s_A,s_B`$ then means choosing appropriate angles $`\theta _A,\varphi _A`$ and $`\theta _B,\varphi _B`$. The classical pure strategies can be realised as
$$CU(0,0)\text{ and }DU(\pi ,0).$$
(32)
This case has also been considered in Ref. . The expected pay-off for Alice, e.g., explicitly reads
$`P_A(\theta _A,\varphi _A,\theta _B,\theta _B)=3\left|\mathrm{cos}(\varphi _A+\varphi _B)\mathrm{cos}(\theta _A/2)\mathrm{cos}(\theta _B/2)\right|^2`$ (33)
$`+`$ $`5\left|\mathrm{sin}(\varphi _A)\mathrm{cos}(\theta _A/2)\mathrm{sin}(\theta _B/2)\mathrm{cos}(\varphi _B)\mathrm{cos}(\theta _B/2)\mathrm{sin}(\theta _A/2)\right|^2`$ (34)
$`+`$ $`\left|\mathrm{sin}(\varphi _A+\varphi _B)\mathrm{cos}(\theta _A/2)\mathrm{cos}(\theta _B/2)+\mathrm{sin}(\theta _A/2)\mathrm{sin}(\theta _B/2)\right|^2.`$ (35)
It turns out that the previous Nash equilibrium $`(D,D)`$ of $`S^{(CL)}`$ is no longer an equilibrium solution, as both players can benefit from deviating from $`D`$. However, concomitant with the disappearance of this solution another Nash equilibrium has emerged, given by $`(Q,Q)`$. The strategy $`Q`$ is associated with a matrix
$$QU(0,\pi /2)=\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right).$$
(36)
This Nash equilibrium is unique and serves as the only acceptable solution of the game. The astonishing fact is that $`P_A(Q,Q)=P_B(Q,Q)=3`$ (instead of $`1`$) so that the Pareto optimum is realised. No player could gain without lessening the other player’s expected pay-off. In this sense one can say that the dilemma of the original game has fully disappeared. In the classical game only mutual cooperation is Pareto optimal, but this pair of strategies does not correspond to a Nash equilibrium.
#### 3 General unitary operations
One can generalise the previous setting to the case where Alice and Bob can implement operations $`s_A`$ and $`s_B`$ taken from $`S^{(GU)}`$, where $`S^{(GU)}`$ is the set of general local unitary operations. Here, it could be suspected that the solution becomes more efficient the larger the sets of allowed operations are. But this is not the case. The previous Pareto optimal unique Nash equilibrium $`(Q,Q)`$ ceases to be an equilibrium solution if the set is enlarged: For any strategy $`s_BS^{(GU)}`$ there exists an optimal answer $`s_AS^{(GU)}`$ resulting in
$$(s_As_B)(\rho )=|\psi _{DC}\psi _{DC}|,$$
(37)
with $`\rho `$ given in Eq. (8). That is, for any strategy of Bob $`s_B`$ there is a strategy $`s_A`$ of Alice such that
$$P_A(s_A,s_B)=5\text{ and }P_B(s_A,s_B)=0:$$
(38)
Take
$`s_A\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),s_B\left(\begin{array}{cc}ib& a\\ d& ic\end{array}\right),`$ (43)
where $`a,b,c,d`$ are appropriate complex numbers. Given that Bob plays the strategy $`s_B`$ associated with a particular Nash equilibrium $`(s_A,s_B)`$, Alice can always apply the optimal answer $`s_A`$ to achieve the maximal possible pay-off. However, the resulting pair of quantum strategies can not be an equilibrium since again, the game being symmetric, Bob can improve his pay-off by changing his strategy to his optimal answer $`s_B^{}`$. Hence, there is no pair $`(s_A,s_B)`$ of pure strategies with the property that the players can only lose from unilaterally deviating from this pair of strategies.
Yet, there remain to be Nash equilibria in mixed strategies which are much more efficient than the classical outcome of the equilibrium in dominant strategies $`P_A(D,D)=P_B(D,D)=1`$. In a mixed strategy of Alice, say, she selects a particular quantum strategy $`s_A`$ (which is also called pure strategy) from the set of strategies $`S_A`$ with a certain classical probability. That is, mixed strategies of Alice and Bob are associated with maps of the form
$$\rho \sigma =\underset{i,j}{}p_A^{(i)}p_B^{(j)}(U_A^{(i)}U_B^{(j)})\rho (U_A^{(i)}U_B^{(j)})^{},$$
(44)
$`p_A^{(i)},p_B^{(i)}[0,1]`$, $`i,j=1,2,\mathrm{},N`$, with
$$\underset{i}{}p_A^{(i)}=\underset{j}{}p_B^{(j)}=1.$$
(45)
$`U_A^{(i)}`$ and $`U_B^{(j)}`$ are local unitary operators corresponding to pure strategies $`s_A^{(i)}`$ and $`s_B^{(j)}`$.
The map given by Eq. (61) acts in $`_A`$ and $`_A`$ as a doubly stochastic map, that is, as a completely positive unital map . As a result, the final reduced states $`\mathrm{tr}_B[\sigma ]`$ and $`\mathrm{tr}_A[\sigma ]`$ must be more mixed than the reduced initial states $`\mathrm{tr}_B[\rho ]`$ and $`\mathrm{tr}_A[\rho ]`$ in the sense of majorisation theory . As the initial state $`\rho `$ is a maximally entangled state, all accessible states after application of a mixed strategy of Alice and Bob are locally identical to the maximally mixed state $`1/\mathrm{dim}(_A)=1/\mathrm{dim}(_B)`$, which is a multiple of $`1`$.
The following construction, e.g., yields an equilibrium in mixed quantum strategies: Allow Alice to choose from two strategies $`s_A^{(1)}`$ and $`s_A^{(2)}`$ with probabilities $`p_A^{(1)}=1/2`$ and $`p_A^{(2)}=1/2`$, while Bob may take $`s_B^{(1)}`$ or $`s_B^{(2)}`$, with
$`s_A^{(1)}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ $`s_A^{(2)}\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right),`$ (51)
$`s_B^{(1)}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),`$ $`s_B^{(2)}\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right).`$ (56)
His probabilities are also given by $`p_B^{(1)}=1/2`$ and $`p_B^{(2)}=1/2`$. The quantum strategies of Eq. (51) and Eq. (56) are mutually optimal answers and have the property that
$`P_A(s_A^{(i)},s_B^{(i)})=0,`$ $`P_B(s_A^{(i)},s_B^{(i)})=5,`$ (58)
$`P_A(s_A^{(i)},s_B^{(3i)})=5,`$ $`P_B(s_A^{(i)},s_B^{(3i)})=0,`$ (59)
for $`i=1,2`$. Due to the particular constraints of Eq. (58) and Eq. (59) there exists no other mixed strategy for Bob yielding a better pay-off than the above mixed strategy, given that Alice sticks to the equilibrium strategy. This can be seen as follows. Let Alice use this particular mixed quantum strategy as above and let Bob use any mixed quantum strategy
$$s_B^{(1)},\mathrm{},s_B^{(N)}$$
(60)
together with $`p_A^{(1)},\mathrm{},p_A^{(N)}`$. The final state $`\sigma `$ after application of the strategies is given by the convex combination
$$\sigma =\underset{i=1,2}{}\underset{j}{}p_A^{(i)}p_B^{(j)}(s_A^{(i)}s_B^{(j)})(\rho ),$$
(61)
This convex combination can only lead to a smaller expected pay-off for Bob than the optimal pure strategy $`s_B^{(k)}`$ in Eq. (60), $`k\{1,\mathrm{},N\}`$. Such optimal pure strategies are given by $`s_B^{(1)}`$ and $`s_B^{(2)}`$ as in Eq. (56) leading to an expected pay-off for Bob of $`P_B(s_A,s_B)=2.5`$; there are no pure strategies which achieve a larger expected pay-off. While both pure strategies $`s_B^{(1)}`$ and $`s_B^{(2)}`$ do not correspond to an equilibrium, the mixed strategy where $`s_B^{(1)}`$ and $`s_B^{(2)}`$ are chosen with $`p_B^{(1)}=1/2`$ and $`p_B^{(2)}=1/2`$ actually does. Nash equilibria consist of pairs of mutually optimal answers, and only for this choice of Bob the original mixed quantum strategy of Alice is her optimal choice, as the same argument applies also to her, the game being symmetric.
This Nash equilibrium is however not the only one. There exist also other four-tuples of matrices than the ones presented in Eq. (51) and Eq. (56) that satisfy Eq. (58) and Eq. (59). Such matrices can be made out by appropriately rotating the matrices of Eq. (51) and Eq. (56). In the light of the fact that there is more than one equilibrium it is not obvious which Nash equilibrium the players will realise. It is at first not even evident whether a Nash equilibrium will be played at all. But the game theoretical concept of the focal point effect helps to resolve this issue.
To explore the general structure of any Nash equilibrium in mixed strategies we continue as follows: let
$$U_A^{(1)},\mathrm{},U_A^{(N)}$$
(62)
together with $`p_A^{(1)},\mathrm{},p_A^{(N)}`$ specify the mixed strategy pertinent to a Nash equilibrium of Alice. Then there exists a mixed strategy $`U_B^{(1)},\mathrm{},U_B^{(N)}`$, $`p_B^{(1)},\mathrm{},p_B^{(N)}`$ of Bob which rewards Bob with the best achievable pay-off, given that Alice plays this mixed strategy. Yet, the pair of mixed strategies associated with
$$QU_A^{(1)}Q^{},\mathrm{},QU_A^{(N)}Q^{},QU_B^{(1)}Q^{},\mathrm{},QU_B^{(N)}Q^{}$$
(63)
with $`p_A^{(1)},\mathrm{},p_A^{(N)}`$, $`p_B^{(1)},\mathrm{},p_B^{(N)}`$ is another Nash equilibrium. This equilibrium leads to the same expected pay-off for both players, and is fully symmetric to the previous one. Doubly applying $`Q`$ as $`QQU_A^{(1)}Q^{}Q^{},\mathrm{},QQU_A^{(N)}Q^{}Q^{}`$ results again into a situation with equivalent strategies as the original ones. For a given Nash equilibrium as above the one specified by Eq. (63) will be called dual equilibrium.
However, there is a single Nash equilibrium $`(R,R)`$ which is the only one which gives an expected pay-off of $`P_A(R,R)=P_B(R,R)=2.25`$ and which is identical to its dual equilibrium: it is the simple map
$$\rho \sigma =1/\mathrm{dim}().$$
(64)
Indeed, there exist probabilities $`p_A^{(1)},\mathrm{},p_A^{(N)}`$ and unitary operators $`U_A^{(1)},\mathrm{},U_A^{(N)}`$ such that $`_ip_A^{(i)}(U_A^{(i)}1)\rho (U_A^{(i)}1)^{}=1/\mathrm{dim}()`$ . If Alice has already selected $`s_A=R`$, the application of $`s_B=R`$ will not change the state of the quantum system any more.
Assume that Eq. (62) and Eq. (63) are associated with equivalent quantum strategies. This means that they have to produce the same expected pay-off for all quantum strategies $`s_B`$ of Bob. If Alice and Bob apply $`s_As_B`$ they get an expected pay-off according to Eq. (15) and Eq. (16); if Alice after implementation of $`s_A`$ manipulates the quantum system by applying the local unitary operator $`Q1`$, they obtain
$`P_A^{}(s_A,s_B)`$ $`=`$ $`A_{DD}\mathrm{tr}[\pi _{CC}\sigma ]+A_{DC}\mathrm{tr}[\pi _{CD}\sigma ]+A_{CD}\mathrm{tr}[\pi _{DC}\sigma ]+A_{CC}\mathrm{tr}[\pi _{DD}\sigma ],`$ (66)
$`P_B^{}(s_A,s_B)`$ $`=`$ $`B_{DD}\mathrm{tr}[\pi _{CC}\sigma ]+B_{DC}\mathrm{tr}[\pi _{CD}\sigma ]+B_{CD}\mathrm{tr}[\pi _{DC}\sigma ]+B_{CC}\mathrm{tr}[\pi _{DD}\sigma ].`$ (67)
The only $`s_A`$ with the property that $`P_A^{}(s_A,s_B)=P_A(s_A,s_B)`$ and $`P_B^{}(s_A,s_B)=P_B(s_A,s_B)`$ for all $`s_B`$ is the map given by Eq. (64).
In principle, any Nash equilibrium may become a self-fulfilling prophecy if the particular Nash equilibrium is expected by both players. It has been pointed out that in a game with more than one equilibrium, anything that attracts the players’ attention towards one of the equilibria may make them expect and therefore realise it . The corresponding focal equilibrium is the one which is conspicuously distinguished from the other Nash equilibria. In this particular situation there is indeed one Nash equilibrium different from all the others: it is the one which is equivalent to its dual equilibrium, the map which simply maps the initial state on the maximally mixed state. For all other expected pay-offs both players are ambivalent between (at least) two symmetric equilibria. The expected pay-off the players will receive in this focal equilibrium –
$$P_A(R,R)=P_B(R,R)=2.25$$
(68)
– is not fully Pareto optimal, but it is again much more efficient than the classically achievable outcome of 1 .
#### 4 Completely positive trace-preserving maps corresponding to local operations
In this scenario both Alice and Bob may perform any operation that is allowed by quantum mechanics. That is, the set of strategies $`S^{(CP)}`$ is made up of $`(s_A,s_B)`$, where both $`s_A`$ and $`s_B`$ correspond to a completely positive trace-preserving map
$$(s_As_B)(\rho )=\underset{i}{}\underset{j}{}(A_iB_j)\rho (A_iB_j)^{}$$
(69)
corresponding to a local operation, associated with Kraus operators $`A_i`$ and $`B_j`$ with $`i,j=1,2,\mathrm{}`$ . The trace-preserving property requires $`_iA_i^{}A_i=1`$ and $`_iB_i^{}B_i=1`$. This case has already been mentioned in Ref. . The quantum strategies $`s_A`$ and $`s_B`$ do no longer inevitably act as unital maps in the respective Hilbert spaces as before. In other words, the reduced states of Alice and Bob after application of the quantum strategy are not necessarily identical to the maximally mixed state $`1/\mathrm{dim}(_A)`$.
As already pointed out in Ref. , the pair of strategies $`(Q,Q)`$ of the two-parameter set of strategies $`S^{(TP)}`$ is again no equilibrium solution. It is straightforward to prove that the Nash equilibria of the type of Eq. (51) and Eq. (56) of mixed strategies with general local unitary operations are, however, still present, and each of these equilibria yields an expected pay-off of $`2.5`$.
In addition, as strategies do no longer have to be locally unital maps, it is not surprising that new Nash equilibria emerge: Alice and Bob may, e.g., perform a measurement associated with Kraus operators
$$A_1=|00|A_2=|11|,B_1=D|00|B_2=D|11|.$$
(70)
This operation yields a final state $`\sigma =(s_As_B)(\rho )=(|0101|+|1010|)/2`$. Clearly neither Alice nor Bob can gain from unilaterally deviating from their strategy.
One can nevertheless argue as in the previous case. Again, all Nash equilibria occur at least in pairs. First, there are again the dual equilibria from $`S^{(GU)}`$. Second, there are Nash equilibria $`(s_A,s_B)`$, $`s_As_B`$, with the property that $`(s_B,s_A)`$ is also a Nash equilibrium yielding the same expected pay-off. The only Nash equilibrium invariant under application of $`Q`$ and exchange of the strategies of the players is again $`(R,R)`$ defined in the previous subsection, which yields a pay-off $`P_A(R,R)=P_B(R,R)=2.25`$. This is the solution of the game is the most general case. While both players could in principle do better (as the solution lacks Pareto optimality), the efficiency of this focal equilibrium is much higher than the equilibrium in dominant strategies of the classical game. Hence, also in this most general case both players gain from using quantum strategies.
This study shows that the efficiency of the equilibrium the players can reach in this game depends on the actions the players may take. One feature, however is present in each of the considered sets: both players can increase their expected pay-offs drastically by resorting to quantum strategies.
### B Chicken
In the previous classical game – the Prisoners’ Dilemma – an unambiguous solution can be specified consisting of a unique Nash equilibrium. However, this solution was not efficient, thus giving rise to the dilemma. The situation of the players in the Chicken game ,
$`A_{CC}=B_{CC}=6,A_{CD}=B_{DC}=8,`$ (72)
$`A_{DC}=B_{CD}=2,A_{DD}=B_{DD}=0,`$ (73)
can be described by the matrix of Fig. 3.
This game has two Nash equilibria $`(C,D)`$ and $`(D,C)`$: it is not clear how to anticipate what the players’ decision would be. In addition to the two Nash equilibria in pure strategies there is an equilibrium in mixed strategies, yielding an expected pay-off $`4`$ .
In order to investigate the new features of the game if superpositions of classical strategies are allowed for, three set of strategies are briefly discussed:
#### 1 One-parameter set of strategies
Again, we consider the set of strategies $`S^{(CL)}`$ of one-dimensional rotations. The strategies $`s_A`$ and $`s_B`$ are associated with local unitary operators
$$U(\theta )=\left(\begin{array}{cc}\mathrm{cos}(\theta /2)& \mathrm{sin}(\theta /2)\\ \mathrm{sin}(\theta /2)& \mathrm{cos}(\theta /2)\end{array}\right)$$
(74)
with $`\theta [0,\pi ]`$,
$$CU(0)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),DU(\pi )=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(75)
Then as before, the quantum game yields the same expected pay-off as the classical game in randomised strategies. This means that still two Nash equilibria in pure strategies are present.
#### 2 Two-parameter set of strategies
The players can actually take advantage of an additional degree of freedom which is not accessible in the classical game. If they may apply unitary operations from $`S^{(TP)}`$ of the type
$$U(\theta ,\varphi )=\left(\begin{array}{cc}e^{i\varphi }\mathrm{cos}(\theta /2)& \mathrm{sin}(\theta /2)\\ \mathrm{sin}(\theta /2)& e^{i\varphi }\mathrm{cos}(\theta /2)\end{array}\right)$$
(76)
with $`\theta [0,\pi ]`$ and $`\varphi [0,\pi /2]`$ the situation is quite different than with $`S^{(CL)}`$. $`(C,D)`$ and $`(C,D)`$ with $`CU(0,0)`$ and $`DU(\pi ,0)`$ are no longer equilibrium solutions. E.g., given that $`s_A=D`$ the pair of strategies $`(D,Q)`$ with $`QU(0,\pi /2)`$ yields a better expected pay-off for Bob than $`(D,C)`$, that is to say $`P_B(D,Q)=8`$, $`P_B(D,C)=2`$. In fact $`(Q,Q)`$ is now the unique Nash equilibrium with $`P_A(Q,Q)=P_B(Q,Q)=6`$, which follows from an investigation of the actual expected pay-offs of Alice and Bob analogous to Eq. (33). This solution is not only the unique acceptable solution of the game, but it is also an equilibrium that is Pareto optimal. This contrasts very much with the situation in the classical game, where the two equilibria were not that efficient.
#### 3 Completely positive trace-preserving maps corresponding to local operations
As in the considerations concerning the Prisoner’s Dilemma game, more than one Nash equilibrium is present, if both players can take quantum strategies from the set $`S^{(CP)}`$, and all Nash equilibria emerge at least in pairs as above. The focal equilibrium is given by $`(R,R)`$, resulting in a pay-off of $`P_A(R,R)=P_B(R,R)=4`$, which is the same as the mixed strategy of the classical game.
## IV Summary and Conclusion
In these lecture notes the idea of implementing quantum operations as strategic moves in a game is explored . In detail, we investigated games which could be conceived as a generalisation into the quantum domain of a two player binary choice game. As a toy model for more complex scenarios we studied quantum games where the efficiency of the equilibria attainable when using quantum strategies could be contrasted with the efficiency of solutions in the corresponding classical game. We investigated a hierarchy of quantum strategies $`S^{(CL)}S^{(TP)}S^{(GU)}S^{(CP)}`$. Again , we found superior performance of quantum strategies as compared to classical strategies.
The nature of a game is determined by the rules of the game. In particular, the appropriate solution concept depends on the available strategic moves. Obviously, a player cannot make a meaningful choice without knowing the options at his or her disposal. So it comes to no surprise that also the actual achievable pay-off in such a game depends on the set of allowed strategies. Roughly speaking, one can say that the possibility of utilising strategies which are not feasible in the analogous classical game implicates a significant advantage. In the models studied in detail two kinds of ”dilemmas” were ”resolved”: (i) On the one hand there are quantum games with an efficient unambiguous solution, while in the classical analogue only an inefficient equilibrium can be identified. By taking advantage of appropriate quantum strategies much more efficient equilibria could be reached. In certain sets of strategies even a maximally efficient solution – the Pareto optimum – was attainable. (ii) On the other hand, there exist quantum games with a unique solution with a classical equivalent which offers two Nash equilibria of the same quality.
This paper deals with simple set-ups in which information is exchanged quantum-mechanically. The emphasis was to examine how situations where strategies are identified with quantum operations applied on quantum mechanical carriers of information are different from the classical equivalent. It is the hope that these investigations may enable us to better understand competitive structures in a game theoretical sense in applications of quantum information theory.
## V Acknowledgements
We would like to thank Maciej Lewenstein, Onay Urfalıo$`\overline{\mathrm{g}}`$lu, Joel Sobel, Tom Cover, Charles H. Bennett, Martin B. Plenio, and Uta Simon for helpful suggestions. We also acknowledge fruitful discussions with the participants of the A2 Consortial Meeting. This work was supported by the European Union and the DFG.
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# Ribbon-moves of 2-knots: the Farber-Levine pairing and the Atiyah-Patodi-Singer- Casson-Gordon-Ruberman 𝜂̃-invariants of 2-knots
## 1 Ribbon-moves of 2-knots
In this paper we discuss ribbon-moves of 2-knots. In this section we review the definition of ribbon-moves.
An (oriented) 2-(dimensional) knot is a smooth, oriented submanifold $`K`$ of $`S^4`$ which is diffeomorphic to the $`2`$-sphere. We say that 2-knots $`K_1`$ and $`K_2`$ are equivalent if there exists an orientation preserving diffeomorphism $`f:`$ $`S^4`$ $``$ $`S^4`$ such that $`f(K_1)`$=$`K_2`$ and that $`f|_{K_1}:`$ $`K_1`$ $``$ $`K_2`$ is an orientation preserving diffeomorphism. Let $`id:S^4`$ $``$ $`S^4`$ be the identity. We say that 2-knots $`K_1`$ and $`K_2`$ are identical if $`id(K_1)`$=$`K_2`$ and that $`id|_{K_1}:K_1K_2`$ is an orientation preserving diffeomorphism.
Definition 1.1. Let $`K_1`$ and $`K_2`$ be 2-knots in $`S^4`$. We say that $`K_2`$ is obtained from $`K_1`$ by one ribbon-move if there is a 4-ball $`B`$ of $`S^4`$ with the following properties.
(1) $`K_1(BK_1)=K_2(BK_2)`$.
This diffeomorphism map is orientation preserving.
(2) $`BK_1`$ is drawn as in Figure 1.1. $`BK_2`$ is drawn as in Figure 1.2.
We regard $`B`$ as (a close 2-disc)$`\times [0,1]\times \{t|1t1\}`$. We put $`B_t=`$(a close 2-disc)$`\times [0,1]\times \{t\}`$. Then $`B=B_t`$. In Figure 1.1 and 1.2, we draw $`B_{0.5},B_0,B_{0.5}`$ $`B`$.
We draw $`K_1`$ and $`K_2`$ by the bold line. The fine line denotes $`(B_t)`$.
$`BK_1`$ (resp. $`BK_2`$) is diffeomorphic to $`D^2(S^1\times [0,1])`$.
$`BK_1`$ has the following properties: $`B_tK_1`$ is empty for $`1t<0`$ and $`0.5<t1`$. $`B_0K_1`$ is diffeomorphic to $`D^2(S^1\times [0,0.3])(S^1\times [0.7,1])`$. $`B_{0.5}K_1`$ is diffeomorphic to $`(S^1\times [0.3,0.7])`$. $`B_tK_1`$ is diffeomorphic to $`S^1S^1`$ for $`0<t<0.5`$.
$`BK_2`$ has the following properties: $`B_tK_2`$ is empty for $`1t<0.5`$ and $`0<t1`$. $`B_0K_2`$ is diffeomorphic to $`D^2(S^1\times [0,0.3])(S^1\times [0.7,1])`$. $`B_{0.5}K_2`$ is diffeomorphic to $`(S^1\times [0.3,0.7])`$. $`B_tK_2`$ is diffeomorphic to $`S^1S^1`$ for $`0.5<t<0`$.
We do not assume which the orientation of $`BK_1`$ (resp. $`BK_2`$ ) is.
Figure 1.1.
Figure 1.2.
Suppose that $`K_2`$ is obtained from $`K_1`$ by one ribbon-move and that $`K_2^{}`$ is equivalent to $`K_2`$. Then we also say that $`K_2^{}`$ is obtained from $`K_1`$ by one ribbon-move.
If $`K_1`$ is obtained from $`K_2`$ by one ribbon-move, then we also say that $`K_2`$ is obtained from $`K_1`$ by one ribbon-move.
Definition 1.2. 2-knots $`K_1`$ and $`K_2`$ are said to be ribbon-move equivalent if there are 2-knots $`K_1=\overline{K}_1,\overline{K}_2,\mathrm{},\overline{K}_{p1},\overline{K}_p=K_2`$ ($`p𝐍,p2`$) such that $`\overline{K}_i`$ is obtained from $`\overline{K}_{i1}`$ $`(1<ip)`$ by one ribbon-move.
In this paper we discuss the following problems.
Problem 1.3. Let $`K_1`$ and $`K_2`$ be 2-knots. Consider a necessary (resp. sufficient, necessary and sufficient ) condition that $`K_1`$ and $`K_2`$ are ribbon-move equivalent.
In the author proved:
Theorem 1.4. ()
(1)If 2-knots $`K`$ and $`K^{}`$ are ribbon-move equivalent, then $`\mu (K)=\mu (K^{})`$.
(2)Let $`K_1`$ and $`K_2`$ be 2-knots in $`S^4`$. Suppose that $`K_1`$ are ribbon-move equivalent to $`K_2`$. Let $`W_i`$ be arbitrary Seifert hypersurfaces for $`K_i`$. Then the torsion part of $`\{H_1(W_1)H_1(W_2)\}`$ is congruent to $`GG`$ for a finite abelian group $`G`$.
(3)Not all 2-knots are ribbon-move equivalent to the trivial 2-knot.
(4)The inverse of (1) is not true. The inverse of (2) is not true.
## 2 Main results
Theorem 2.1. Let $`K`$ and $`K^{}`$ be 2-knots. Suppose that $`K`$ and $`K^{}`$ are ribbon-move equivalent. Then we have:
(1)There is an isomorphism $`\mathrm{Tor}H_1(\stackrel{~}{X}_K;𝐙)\mathrm{Tor}H_1(\stackrel{~}{X}_K^{};𝐙)`$ as $`𝐙[t,t^1]`$-modules.
(2)The Farber-Levine pairing on $`\mathrm{Tor}H_1(\stackrel{~}{X}_K;𝐙)`$ is equivalent to that on $`\mathrm{Tor}H_1(\stackrel{~}{X}_K^{};𝐙)`$.
Theorem 2.2. Let $`K`$ be a 2-knot. Suppose that $`K`$ is ribbon-move equivalent to the trivial knot. Then $`\stackrel{~}{\eta }(K,)`$ for $`𝐙_𝐝`$ is zero. ($`d𝐍`$ and $`d2`$).
Note. We review the Alexander module in §3. We review the Farber-Levine pairing in §4. We review the Atiyah-Patodi-Singer-Casson-Gordon-Ruberman $`𝐐/𝐙`$-valued $`\stackrel{~}{\eta }`$-invariants of 2-knots in §5.
## 3 The Alexander module
See for detail.
Let $`K`$ be a 2-knot $`S^4`$. Let $`N(K)`$ be the tubular neighborhood of $`K`$ in $`S^4`$. Let $`\alpha :\pi _1(\overline{S^4N(K)})H_1(\overline{S^4N(K)};𝐙)`$ be the abelianization. Note that any nonzero cycle $`xH_1(S^4N(K);𝐙)`$ is oriented naturally by using the orientation of $`K`$ and that of $`S^4`$. We define the canonical isomorphism $`\beta :H_1(\overline{S^4N(K)};𝐙)𝐙`$ by using these orientations of $`x`$. Let $`\stackrel{~}{X}_K^{\mathrm{}}`$ be the covering space associated with $`\beta \alpha :\pi _1(\overline{S^4N(K)})𝐙`$. We call $`\stackrel{~}{X}_K^{\mathrm{}}`$ the canonical infinite cyclic covering space of the complement $`\overline{S^4N(K)}`$ of $`K`$. Then $`H_i(\stackrel{~}{X}_K^{\mathrm{}};𝐙)`$ is regarded as a $`𝐙[t,t^1]`$-module by using the covering translations $`\stackrel{~}{X}_K^{\mathrm{}}\stackrel{~}{X}_K^{\mathrm{}}`$. This $`𝐙[t,t^1]`$-module $`H_i(\stackrel{~}{X}_K^{\mathrm{}};𝐙)`$ is called the Alexander module.
## 4 lk( , ) for 2-knots
See for detail.
Firstly we review of lk( , ) for closed oriented 3-manifolds. Let $`M`$ be a closed oriented 3-manifold. Let $`x,y`$Tor$`H_1(M;𝐙)`$. Let $`n`$ (resp. $`m`$) be a natural number. Let $`n`$ (resp. $`m`$) be the order of $`x`$ (resp. $`y`$). Let $`X`$ (resp. $`Y`$) be a circle embedded in $`M`$ such that $`[X]=x`$ (resp. $`[Y]=y`$). Let $`XY=\varphi `$. Then there is an immersion map $`f:FM`$ such that
(1) $`F`$ is an oriented compact surface and $`F`$ is one circle.
(2) $`f(\mathrm{I}ntF)`$ is an embedding.
(3) $`f(F)=X`$ and deg$`(f|_F)`$=$`n`$.
(4) $`f(\mathrm{I}ntF)`$ is transverse to $`Y`$.
Let $`f(F)Y=P_1\mathrm{}P_\alpha `$. (Note $`P_i`$ is a point.) We give $`P_i`$ a signature $`\epsilon _i`$ by using the orientation of $`f(F)`$, that of $`Y`$, and that of $`M`$.
Define lk$`(x,y)=\frac{1}{n}\mathrm{\Sigma }_{i=1}^\alpha \epsilon _i`$ $`𝐐/𝐙`$.
Proposition. lk(y,x)=lk(x,y).
There is an immersion map $`g:GM`$ such that
(1) $`G`$ is an oriented compact surface and $`G`$ is one circle.
(2) $`g(\mathrm{I}ntG)`$ is an embedding.
(3) $`f(G)=Y`$ and deg$`(g|_G)`$=$`m`$.
(4) $`f(\mathrm{I}ntG)`$ is transverse to $`X`$.
(5) Int$`F`$ is transverse to Int$`G`$.
Let $`g(G)X=Q_1\mathrm{}Q_\beta `$. (Note $`Q_j`$ is a point. ) We give $`Q_j`$ a signature $`\sigma _j`$ by using the orientation of $`g(G)`$, that of $`X`$, and that of $`M`$. Let $`f(\mathrm{I}ntF)g(\mathrm{I}ntG)=R_1\mathrm{}R_\gamma `$. (Note $`R_k`$ is a compact open 1-manifold. ) We give $`R_k`$ a signature $`\tau _k`$ by using the orientation of $`g(G)`$, that of $`f(F)`$, and that of $`M`$. Then we have:
lk$`(y,x)`$ $`=\frac{1}{m}\mathrm{\Sigma }_{j=1}^\beta \sigma _j`$ $`=\frac{1}{mn}\mathrm{\Sigma }_{k=1}^\gamma \tau _k`$ $`=\frac{1}{n}\mathrm{\Sigma }_{i=1}^\alpha \epsilon _i`$ $`=`$lk$`(x,y)`$.
Secondly we review lk( , ) for 2-knots. Let $`K`$ be a 2-knot. Let $`\stackrel{~}{X}_K^{\mathrm{}}`$ be the canonical infinite cyclic covering space of the complement $`\overline{S^4N(K)}`$ of $`K`$. Let $`x,y`$ Tor$`H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)`$. We define lk$`(x,y)`$.
Let $`p`$ be the natural projection map $`\stackrel{~}{X}_K^{\mathrm{}}X`$. Let $`V`$ be a Seifert hypersurface for $`K`$. Let $`V_\xi `$ be one connected component of $`p^1(V)=_{\mathrm{}}^{\mathrm{}}V_i`$. The natural inclusion map $`V_\xi \stackrel{~}{X}_K^{\mathrm{}}`$ induces the homomorphism $`\iota :H_1(V_\xi )H_1(\stackrel{~}{X}_K^{\mathrm{}})`$. Theorem 7.3 of and its proof essentially say that $`\iota :\mathrm{Tor}H_1(V_\xi )\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}})`$ is onto.
Let $`\widehat{V_\xi }`$ be the closed oriented 3-manifold which is obtained from $`V_\xi `$ by attaching a 3-dimensional 3-handle along $`V_\xi `$. The natural inclusion map $`V_\xi \widehat{V_\xi }`$ induces $`\gamma :H_1(V_\xi ;𝐙)\stackrel{}{}H_1(\widehat{V_\xi };𝐙)`$. Then the map $`(\iota \gamma ^1):\mathrm{Tor}H_1(\widehat{V_\xi };𝐙)\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)`$ is onto. Let $`(\iota \gamma ^1)(x^{})=x`$ and $`(\iota \gamma ^1)(y^{})=y`$. Define lk($`x,y`$) for the 2-knot $`K`$ to be lk$`(x^{},y^{})`$ for the 3-manifold $`\widehat{V_\xi }`$. Theorem 7.3 of and its proof essentially say that lk($`x,y`$) for the 2-knot $`K`$ is independent of the choice of $`V`$.
## 5 $`\stackrel{~}{\eta }(,)`$ of 2-knots
See $`\stackrel{~}{\eta }(,)`$ for for detail.
We firstly define $`\stackrel{~}{\eta }(,)`$ of closed oriented 3-manifolds. (See P.571 of .)
Let $`A`$ be an oriented compact 4-manifold. For $`d`$ an integer, set $`\omega =e^{\frac{2\pi i}{d}}`$, and suppose $`\varphi :AK(𝐙_𝐝,1)`$ is a continuous map. Then $`\varphi `$ corresponds to a class $`\varphi H^1(A;𝐙_𝐝)`$= Hom$`(H_1(A;𝐙);𝐙_𝐝)`$ and induces a cyclic cover $`\stackrel{~}{A}A`$ with a specific choice $`T:\stackrel{~}{A}\stackrel{~}{A}`$ of a generator of the covering translations. Let $`\overline{H}_2(A,\varphi )`$= $`\omega `$-eigenspace of $`T_{}`$ acting on $`H_2(\stackrel{~}{A};𝐂)`$. The intersection form on $`A`$ induces a Hermitian pairing on $`H_2(A;𝐂)`$; $`<x\alpha ,y\beta >`$=$`(xy)\alpha \overline{\beta }`$. Define $`\overline{\sigma }(A,\varphi )`$=the signature of $`<,>`$ restricted to $`\overline{H}_2(A,\varphi )`$.
Let $`M`$ be an oriented closed 3-manifold. Let $`\psi H^1(M;𝐙_𝐝)`$ be a homomorphism. By bordism theory, $`n(M,\psi )=(W,\varphi )`$ for a compact oriented 4-manifold $`W`$. (See e.g. ) Define $`\stackrel{~}{\eta }(M,\psi )=\frac{1}{n}(\overline{\sigma }(W,\varphi )\sigma (W))`$.
Secondly we define $`\stackrel{~}{\eta }(,)`$ of 2-knots. (See ) Let $`K`$ be a 2-knot. Take $`\stackrel{~}{X}_K^{\mathrm{}}`$, $`V`$, $`V_\xi `$, $`\iota :H_1(V_\xi ;𝐙)H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)`$, $`\widehat{V_\xi }`$, $`\gamma :H_1(V_\xi ;𝐙)H_1(\widehat{V_\xi };𝐙)`$ as in §4. Let $`\nu `$ be a homomorphism $`H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)𝐙_𝐤`$.
Then we have
$$H_1(\widehat{V_\xi };𝐙)\stackrel{\gamma ,}{}H_1(V_\xi ;𝐙)\stackrel{\iota }{}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)\stackrel{\nu }{}𝐙_𝐤$$
Then we have $`\nu \iota (\gamma ^1):H_1(\widehat{V_\xi };𝐙)𝐙_𝐤`$. Take $`\stackrel{~}{\eta }(\widehat{V_\xi },\nu \iota (\gamma ^1))𝐐`$. Let $`\pi `$ be the natural projection $`𝐐𝐐/𝐙`$. says that
$`\pi (\stackrel{~}{\eta }(\widehat{V_\xi },\nu \iota (\gamma ^1)))`$ is independent of the choice of $`V`$. We define $`\stackrel{~}{\eta }(K,\nu )𝐐/𝐙`$ to be $`\pi (\stackrel{~}{\eta }(\widehat{V_\xi },\nu \iota (\gamma ^1))).`$
Note. 3 etc. also apply the G-signature theorem and the $`\eta `$-invariants to $`n`$-knots ($`n=1,2,3,\mathrm{}`$).
## 6 Proof of main results
Let $`X`$ be $`\overline{S^4N(K)}`$. Let $`V`$ be a Seifert hypersurface for $`K`$. Suppose $`V`$ also denote $`\overline{VN(K)}`$. Let $`Y=XV`$. Let $`V\times [1,1]`$ be the tubular neighborhood of $`V`$ in $`X`$. Suppose $`Y`$ also denote $`\overline{X(V\times [1,1])}`$.
Let $`\stackrel{~}{X}_K^{\mathrm{}}`$ be the canonical infinite cyclic covering space. There is the natural projection map $`p:\stackrel{~}{X}_K^{\mathrm{}}X`$. Let $`p^1(Y)=_{\mathrm{}}^{\mathrm{}}Y_i`$. Let $`p^1(V)=_{\mathrm{}}^{\mathrm{}}V_i`$. Let $`p^1(V\times [1,1])=_{\mathrm{}}^{\mathrm{}}(V_i\times [1,1])`$. Let $`Y_i=(V_{i1}\times \{1\})(V_i\times \{1\})`$.
There is the Meyer-Vietoris exact sequences of Z-homology groups:
$$H_i(V\times \{1,1\};𝐙)$$
$$^{f_i}$$
$$H_i(V\times [1,1];𝐙)H_i(Y;𝐙)$$
$$^{g_i}$$
$$H_i(X;𝐙)$$
and
$$_{j=\mathrm{}}^{\mathrm{}}H_i(V_j\times \{1,1\};𝐙)$$
$$^{\stackrel{~}{f_i}}$$
$$(_{j=\mathrm{}}^{\mathrm{}}H_i(V_j\times [1,1];𝐙))(_{j=\mathrm{}}^{\mathrm{}}H_i(Y_j;𝐙))$$
$$^{\stackrel{~}{g_i}}$$
$$H_i(\stackrel{~}{X}_K^{\mathrm{}};𝐙)$$
Furthermore the second one is regarded as an exact sequence of $`𝐙[t,t^1]`$-modules by using the covering tranlations.
Claim 6.1. There is an exact sequence of $`𝐙[t,t^1]`$-modules:
$$_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times \{1,1\};𝐙)$$
$$^{\stackrel{~}{f^{\mathrm{tor}}}}$$
$$(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times [1,1];𝐙))(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(Y_j;𝐙))$$
$$^{\stackrel{~}{g^{\mathrm{tor}}}}$$
$$\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)$$
$$$$
$$0$$
where $`\stackrel{~}{f^{\mathrm{tor}}}=`$ $`\stackrel{~}{f_1}|_{_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times \{1,1\};𝐙)}`$
and $`\stackrel{~}{g^{\mathrm{tor}}}=`$ $`\stackrel{~}{g_1}|_{(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times [1,1];𝐙))(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(Y_j;𝐙))}`$.
Note. P.765 of says that proved $`|\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)|<\mathrm{}`$.
Proof. By using the covering translations, we regard $`_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times \{1,1\};𝐙)`$, $`(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times [1,1];𝐙))(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(Y_j;𝐙))`$, and $`\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)`$ as $`𝐙[t,t^1]`$-modules. Furthermore we regard $`\stackrel{~}{f^{\mathrm{tor}}}`$ and $`\stackrel{~}{g^{\mathrm{tor}}}`$ as homomorphisms of $`𝐙[t,t^1]`$-modules.
Take $`V_\xi `$, $`\iota `$ as in §4. Theorem 7.3 of and its proof essentially say that $`\iota :\mathrm{Tor}H_1(V_\xi )\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}})`$ is onto. Note $`\stackrel{~}{g^{\mathrm{tor}}}|_{\mathrm{Tor}H_1(V_\xi \times [1,1])}=\iota `$. Hence $`\stackrel{~}{g^{\mathrm{tor}}}`$ is onto. Therefore
$$(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times [1,1];𝐙))(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(Y_j;𝐙))$$
$$^{\stackrel{~}{g^{\mathrm{tor}}}}$$
$$\mathrm{Tor}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)$$
$$$$
$$0$$
is exact.
Let $`x(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times [1,1];𝐙))(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(Y_j;𝐙))`$ such that $`\stackrel{~}{g^{\mathrm{tor}}}(x)=0`$. Then there is $`y_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times \{1,1\};𝐙)`$ such that $`\stackrel{~}{f_1}(y)=x`$. Let $`n`$ be the order of $`x`$. Then $`\stackrel{~}{f_1}(ny)=nx=0`$. We prove that $`y_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times \{1,1\};𝐙)`$ in the following paragraphs.
There is the Meyer-Vietoris exact sequences of Q-homology groups:
(1)
$$H_i(V\times \{1,1\};𝐐)$$
$$^{f_i^Q}$$
$$H_i(V\times [1,1];𝐐)H_i(Y;𝐐)$$
$$^{g_i^Q}$$
$$H_i(X;𝐐)$$
and
(2)
$$_{j=\mathrm{}}^{\mathrm{}}H_i(V_j\times \{1,1\};𝐐)$$
$$^{\stackrel{~}{f_i^Q}}$$
$$(_{j=\mathrm{}}^{\mathrm{}}H_i(V_j\times [1,1];𝐐))(_{j=\mathrm{}}^{\mathrm{}}H_i(Y_j;𝐐))$$
$$^{\stackrel{~}{g_i^Q}}$$
$$H_i(\stackrel{~}{X}_K^{\mathrm{}};𝐐)$$
By using $`H_i(X;𝐐)H_i(S^1;𝐐)`$ and the sequence (1),
$`f_1^Q:H_1(V\times \{1,1\};𝐐)H_1(V\times [1,1];𝐐)H_1(Y;𝐐)`$ is isomorphism.
Let $`\pi _V:H_1(V\times [1,1];𝐐)H_1(Y;𝐐)H_1(V\times [1,1];𝐐)`$ be the natural projection. Let $`\pi _Y:H_1(V\times [1,1];𝐐)H_1(Y;𝐐)H_1(Y;𝐐)`$ be the natural projection. Suppose that the identity matrix $`E`$ represents
$`\pi _V\{f_1^Q|_{H_1(V\times \{1\};𝐐)}\}:`$ $`H_1(V\times \{1\};𝐐)H_1(V\times [1,1];𝐐)`$,
the identity matrix $`E`$ represents
$`\pi _V\{f_1^Q|_{H_1(V\times \{1\};𝐐)}\}:`$ $`H_1(V\times \{1\};𝐐)H_1(V\times [1,1];𝐐)`$,
a matrix $`A`$ represents
$`\pi _Y\{f_1^Q|_{H_1(V\times \{1\};𝐐)}\}:`$ $`H_1(V\times \{1\};𝐐)H_1(Y;𝐐)`$, and
a matrix $`B`$ represents
$`\pi _Y\{f_1^Q|_{H_1(V\times \{1\};𝐐)}\}:`$ $`H_1(V\times \{1\};𝐐)H_1(Y;𝐐)`$.
Then $`f_1^Q:H_i(V\times \{1,1\};𝐐)H_i(V\times [1,1];𝐐)H_i(Y;𝐐)`$ is represented by
$`P=\left(\begin{array}{cc}E& A\\ E& B\end{array}\right).`$
Since $`H_i(X;𝐐)H_i(S^1;𝐐)`$, det$`P0`$.
We regard $`H_i(V_j\times \{1,1\};𝐐)`$, $`_{j=\mathrm{}}^{\mathrm{}}H_i(V_j\times [1,1];𝐐)`$, and $`_{j=\mathrm{}}^{\mathrm{}}H_i(Y_j;𝐐)`$ as $`𝐐[t,t^1]`$-modules by using the covering translations. Then $`\stackrel{~}{f_1^Q}`$ is represented by
$`P(t)=\left(\begin{array}{cc}E& tA\\ E& B\end{array}\right).`$
Note $`P(1)=P`$. Then det$`P(1)=`$det$`P0`$. Hence det$`P(t)0`$. Hence $`\stackrel{~}{f_1^Q}`$ is injective. Hence we have:
if an element $`z`$ $`_{j=\mathrm{}}^{\mathrm{}}H_1(V_j\times \{1,1\};𝐙)`$
generates $`𝐙`$ $`_{j=\mathrm{}}^{\mathrm{}}H_1(V_j\times \{1,1\};𝐙)`$, $`\stackrel{~}{f_1}(z)0`$. Hence we have:
if $`y`$ $`_{j=\mathrm{}}^{\mathrm{}}H_1(V_j\times \{1,1\};𝐙)`$
generates $`𝐙`$ $`_{j=\mathrm{}}^{\mathrm{}}H_1(V_j\times \{1,1\};𝐙)`$, $`\stackrel{~}{f_1}(ny)0`$. Therefore
$`y_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_j\times \{1,1\};𝐙)`$. Therefore
$$_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}\{H_1(V_j\times \{1,1\};𝐙)\}$$
$$^{\stackrel{~}{f^{\mathrm{tor}}}}$$
$$(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}\{H_1(V_j\times [1,1];𝐙)\})(_{j=\mathrm{}}^{\mathrm{}}\mathrm{Tor}\{H_1(Y_j;𝐙)\})$$
$$^{\stackrel{~}{g^{\mathrm{tor}}}}$$
$$\mathrm{Tor}\{H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)\}$$
is exact. This completes the proof of Claim 6.1.
In order to prove our main theorems, we use the (1,2)-pass-moves for 2-knots. See for the (1,2)-pass-moves for 2-knots for detail.
Definition 6.2. Let $`K_1`$ and $`K_2`$ be 2-knots in $`S^4`$. We say that $`K_2`$ is obtained from $`K_1`$ by one (1,2)-pass-move if there is a 4-ball $`B`$ $`S^4`$ with the following properties. We draw $`B`$ as in Definition 1.1.
(1) $`K_1(BK_1)`$=$`K_2(BK_2)`$.
This diffeomorphism map is orientation preserving.
(2) $`BK_1`$ is drawn as in Figure 6.1. $`BK_2`$ is drawn as in Figure 6.2.
Figure 6.1.
Figure 6.2.
The orientation of the two discs in the Figure 6.1 (resp. Figure 6.2) is compatible with the orientation which is determined naturally by the $`(x,y)`$-arrows in the Figure. We do not assume which the orientations of the annuli in the Figures are.
Suppose that $`K_2`$ is obtained from $`K_1`$ by one (1,2)-pass-move and that $`K_2^{}`$ is equivalent to $`K_2`$. Then we also say that $`K_2^{}`$ is obtained from $`K_1`$ by one (1,2)-pass-move .
If $`K_1`$ is obtained from $`K_2`$ by one (1,2)-pass-move, then we also say that $`K_2`$ is obtained from $`K_1`$ by one (1,2)-pass-move .
2-knots $`K_1`$ and $`K_2`$ are said to be (1,2)-pass-move equivalent if there are 2-knots $`K_1=\overline{K}_1,\overline{K}_2,\mathrm{},\overline{K}_{p1},\overline{K}_p=K_2`$ $`(p𝐍,p2)`$ such that $`\overline{K}_i`$ is obtained from $`\overline{K}_{i1}`$ $`(1<ip)`$ by one (1,2)-pass-move.
In we proved:
Theorem 6.3 () Let $`K`$ and $`K^{}`$ be 2-knots. The following conditions (1) and (2) are equivalent.
(1) $`K`$ is (1,2)-pass-move equivalent to $`K^{}`$.
(2) $`K`$ is ribbon-move equivalent to $`K^{}`$.
Let $`K_<`$ and $`K_>`$ be 2-knots. Suppose $`K_<`$ is ribbon-move equivalent to $`K_>`$. By Theorem 6.3, $`K_<`$ is (1,2)-pass-move equivalent to $`K_>`$.
In order to prove Theorem 2.1 it suffices to prove Theorem 2.1 when $`K_<`$ is obtained from $`K_>`$ by one (1,2)-pass-move in a 4-ball $`BS^4`$.
Claim. There are Seifert hypersurfaces $`V_>`$ for $`K_>`$ and $`V_<`$ for $`K_<`$ such that:
(1) $`V_>(BV_>)=V_<(BV_<)`$.
This diffeomorphism map is orientation preserving.
(2) $`BV_>`$ is drawn as in Figure 6.3. $`BV_<`$ is drawn as in Figure 6.4.
Note. We draw $`B`$ as in Definition 1.1. We draw $`V_>`$ and $`V_<`$ by the bold line. The fine line means $`B`$.
$`BV_>`$ (resp. $`BV_<`$) is diffeomorphic to $`(D^2\times [2,3])(D^2\times [0,1])`$. We can regard $`(D^2\times [0,1])`$ as a 3-dimensional 1-handle which is attached to $`B`$. We can regard $`(D^2\times [2,3])`$ as a 3-dimensional 2-handle which is attached to $`B`$.
$`BV_>`$ has the following properties: $`B_tV_>`$ is empty for $`1t<0`$ and $`0.5<t1`$. $`B_0V_>`$ is diffeomorphic to $`(D^2\times [2,3])(D^2\times [0,0.3])(D^2\times [0.7,1])`$. $`B_{0.5}K_1`$ is diffeomorphic to $`(D^2\times [0.3,0.7])`$. $`B_tV_>`$ is diffeomorphic to $`D^2D^2`$ for $`0<t<0.5`$.
$`BV_<`$ has the following properties: . $`B_tV_<`$ is empty for $`1t<0.5`$ and $`0<t1`$. $`B_0V_<`$ is diffeomorphic to $`(D^2\times [2,3](D^2\times [0,0.3])(D^2\times [0.7,1])`$. $`B_{0.5}V_<`$ is diffeomorphic to $`(D^2\times [0.3,0.7])`$. $`B_tV_<`$ is diffeomorphic to $`D^2D^2`$ for $`0.5<t<0`$.
Figure 6.3.
Figure 6.4.
Proof of Claim. Put $`P=(`$the 3-manifolds in Figure 6.3$`)(B)`$. Note $`P=(`$the 3-manifolds in Figure 6.4$`)(B)`$. Put $`Q=K_>(S^4IntB^4)`$. Note $`Q=K_<(S^4IntB^4)`$. By applying the following Proposition to ($`PQ`$) and $`(S^4IntB^4)`$ the above Claim holds.
By using the obstruction theory, we have the following proposition. ( We can prove it by applying §III of . We can also prove it by generalizing Theorem 2,3 in P.49,50 of . ) The author gives a proof in the Appendix.
Proposition. Let $`X`$ be an oriented compact $`(m+2)`$-dimensional manifold. Let $`X\varphi `$. Let $`M`$ be an oriented closed $`m`$-dimensional manifold which is embedded in $`X`$. Let $`MX\varphi `$. Let $`[M]=0H_m(X;𝐙)`$. Then there is an oriented compact $`(m+1)`$-dimensional manifold $`P`$ such that $`P`$ is embedded in $`X`$ and that $`P=X`$.
Let $`V_{}`$ be a compact 3-manifold embedded in $`S^4`$ whose boundary is $`S^2S^2`$ with the following properties.
(1) $`V_{}(BV_{})=V_>(BV_>)=V_<(BV_<)`$.
(2) $`BV_{}B_0.`$ We draw $`BV_{}`$ as in Figure 6.5.
Figure 6.5.
Let $`N(V_{})`$ be the tubular neighborhood of $`V_{}`$ in $`S^4`$. Let $`X=\overline{S^4N(V_{})}`$. Let $`V_{}`$ also denote $`V_{}X`$. Let $`V_{}\times [1,1]`$ be the tubular neighborhood of $`V_{}`$ in $`X`$. Let $`Y_{}`$=$`\overline{X(V_{}\times [1,1])}`$.
Suppose the following maps are inclusion. ( $`\epsilon =1,1`$. ) The following two diagrams are commutative.
$`\begin{array}{ccc}V_<\times \{\epsilon \}& \stackrel{}{}& V_{}\times \{\epsilon \}\\ & & \\ Y_<& \stackrel{}{}& Y_{}\end{array}`$
$`\begin{array}{ccc}V_{}\times \{\epsilon \}& \stackrel{}{}& V_>\times \{\epsilon \}\\ & & ^{\alpha _>}\\ Y_{}& \stackrel{}{}& Y_>\end{array}`$
By the definition of $`V_{}`$, $`H_1(X;𝐙)=𝐙𝐙`$. Let $`f:\pi _1X𝐙`$ be an epimorphism. Then $`f`$ induces a smooth map $`\overline{f}:XS^1.`$ Let $`qS^1`$ be a regular value of $`\overline{f}`$. Suppose $`\overline{f}^1(q)=V_{}.`$ Take the covering space $`\stackrel{~}{X}`$ of $`X`$ associated with $`\overline{f}`$. Let $`\stackrel{~}{f}:\stackrel{~}{X}X`$ be the projection map. Put $`_{j=\mathrm{}}^{\mathrm{}}V_{,j}=\stackrel{~}{f}^1(V_{})`$. Put $`_{j=\mathrm{}}^{\mathrm{}}Y_{,j}=\stackrel{~}{f}^1(Y_{})`$. Let $`_{j=\mathrm{}}^{\mathrm{}}(V_{,j}\times [1,1])=\overline{f}^1(V_{}\times [1,1])`$. Suppose $`Y_{,j}=(V_{,j1}\times \{1\})(V_{,j}\times \{1\})`$.
Take $`V_{>,j}`$, $`V_{<,j}`$, $`\stackrel{~}{X}_{K_>}^{\mathrm{}}`$, $`\stackrel{~}{X}_{K_<}^{\mathrm{}}`$ as in §4. Suppose the following maps are inclusion. By the above commutative diagrams, the following diagrams are commutative.
$`\begin{array}{ccc}V_{<,j}\times \{\epsilon \}& \stackrel{}{}& V_{,j}\times \{\epsilon \}\\ & & \\ Y_{<,j}& \stackrel{}{}& Y_{,j}\end{array}`$
$`\begin{array}{ccc}V_{,j}\times \{\epsilon \}& \stackrel{}{}& V_{>,j}\times \{\epsilon \}\\ & & ^{\alpha _>}\\ Y_{,j}& \stackrel{}{}& Y_{>,j}\end{array}`$
By the definition of $`V_>`$, $`V_<`$, and $`V_{}`$, it holds that $`V_>`$ (resp. $`V_<`$) is obtained from $`V_{}`$ by attaching one 1-handle. By the definition of $`Y_>`$, $`Y_<`$, and $`Y_{}`$, it holds that $`Y_{}`$ is obtained from $`Y_>`$ (resp. $`Y_<`$) by attaching one 3-handle. Hence the above commutative diagrams induce the following commutative diagram ($`\epsilon =1,1`$).
$$\begin{array}{ccccc}\mathrm{Tor}H_1(V_{<,j}\times \{\epsilon \};Z)& \stackrel{\beta _{<,j},}{}& \mathrm{Tor}H_1(V_{,j}\times \{\epsilon \};Z)& \stackrel{\beta _{>,j},}{}& \mathrm{Tor}H_1(V_{>,j}\times \{\epsilon \};Z)\\ ^{\alpha _{<,j}}& & ^{\alpha _{,j}}& & ^{\alpha _{>,j}}\\ \mathrm{Tor}H_1(Y_{<,j};Z)& \stackrel{\gamma _{<,j},}{}& \mathrm{Tor}H_1(Y_{,j};Z)& \stackrel{\gamma _{>,j},}{}& \mathrm{Tor}H_1(Y_{>,j};Z)\end{array}$$
The above commutative diagram induces the following commutative diagram of Z-homology groups:
$`\begin{array}{ccccc}_{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(V_{<,j}\times \{1,1\}\right)& \stackrel{\beta _{<,j},}{}& _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(V_{,j}\times \{1,1\}\right)& \stackrel{\beta _{>,j},}{}& _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(V_{>,j}\times \{1,1\}\right)\\ ^{f_<}& & ^f_{}& & ^{f_>}\\ \left\{\begin{array}{c}_{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(V_{<,j}\times [1,1]\right)\\ \\ _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(Y_{<,j}\right)\end{array}\right\}& \stackrel{\gamma _{<,j},}{}& \left\{\begin{array}{c}_{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(V_{,j}\times [1,1]\right)\\ \\ _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(Y_{,j}\right)\end{array}\right\}& \stackrel{\gamma _{>,j},}{}& \left\{\begin{array}{c}_{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(V_{>,j}\times [1,1]\right)\\ \\ _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1\left(Y_{>,j}\right)\end{array}\right\}\end{array}`$
By using the covering tlanslations, the above commutative diagram is regarded as that of $`𝐙[t,t^1]`$-modules. By Claim 6.3, we have exact sequences of $`𝐙[t,t^1]`$-modules.
$`\begin{array}{ccc}_{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_{<,j}\times \{1,1\})& & _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_{>,j}\times \{1,1\})\\ ^{f_<}& & ^{f_>}\\ _{\mathrm{}}^{\mathrm{}}\left(\begin{array}{c}\mathrm{Tor}H_1(V_{<,j}\times [1,1])\\ \\ \mathrm{Tor}H_1(Y_{<,j})\end{array}\right)& \mathrm{and}& _{\mathrm{}}^{\mathrm{}}\left(\begin{array}{c}\mathrm{Tor}H_1(V_{>,j}\times [1,1])\\ \\ \mathrm{Tor}H_1(Y_{>,j})\end{array}\right)\\ ^{g_>}& & ^{g_>}\\ \mathrm{Tor}H_1(\stackrel{~}{X}_{K_<}^{\mathrm{}})& & \mathrm{Tor}H_1(\stackrel{~}{X}_{K_>}^{\mathrm{}})\\ & & \\ 0& & 0\end{array}`$
Therefore we have the following commutative diagram of $`𝐙[t,t^1]`$-modules.
$`\begin{array}{ccc}_{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_{<,j}\times \{1,1\})& \stackrel{a,}{}& _{\mathrm{}}^{\mathrm{}}\mathrm{Tor}H_1(V_{>,j}\times \{1,1\})\\ ^{f_<}& & ^{f_>}\\ _{\mathrm{}}^{\mathrm{}}\left(\begin{array}{c}\mathrm{Tor}H_1(V_{<,j}\times [1,1])\\ \\ \mathrm{Tor}H_1(Y_{<,j})\end{array}\right)& \stackrel{b,}{}& _{\mathrm{}}^{\mathrm{}}\left(\begin{array}{c}\mathrm{Tor}H_1(V_{>,j}\times [1,1])\\ \\ \mathrm{Tor}H_1(Y_{>,j})\end{array}\right)\\ ^{g_>}& & ^{g_>}\\ \mathrm{Tor}H_1(\stackrel{~}{X}_{K_<}^{\mathrm{}})& \stackrel{c,}{}& \mathrm{Tor}H_1(\stackrel{~}{X}_{K_>}^{\mathrm{}})\\ & & \\ 0& & 0\end{array}`$
In particular, $`\mathrm{Tor}H_1(\stackrel{~}{X}_{K_<}^{\mathrm{}})\mathrm{Tor}H_1(\stackrel{~}{X}_{K_>}^{\mathrm{}})`$ as $`𝐙[t,t^1]`$-module. This completes the proof of Theorem 2.1.(1).
Take $`V_{<,\xi }`$ $`V_{>,\xi }`$ as in §4. Then there is an orientation preserving diffeomorphism $`h:V_{<,\xi }V_{>,\xi }`$ with the following properties.
(1) The following diagram is commutative
$$\begin{array}{ccccc}V_{<,\xi }& & \stackrel{h}{}& & V_{>,\xi }\\ & & & & \\ & & V_{,\xi }& & \end{array}$$
where $``$ and $``$are inclusion.
(2) $`h_{}=b|_{\mathrm{Tor}H_1(V_{<,\xi };𝐙)}:\mathrm{Tor}H_1(V_{<,\xi };𝐙)\stackrel{}{}\mathrm{Tor}H_1(V_{>,\xi };𝐙).`$
Theorem 7.3 of \[Farber \] and its proof essentially say that
$`\mathrm{Tor}H_1(V_{<,\xi };𝐙)\stackrel{\iota _<}{}\mathrm{Tor}H_1(\stackrel{~}{X}_{K_<}^{\mathrm{}};𝐙)`$ and $`\mathrm{Tor}H_1(V_{>,\xi };𝐙)\stackrel{\iota _>}{}\mathrm{Tor}H_1(\stackrel{~}{X}_{K_>}^{\mathrm{}};𝐙)`$ are onto. Hence there are the following commutative diagram:
$$\begin{array}{ccc}\mathrm{Tor}H_1(V_{<,\xi };𝐙)& \stackrel{h_{},}{}& \mathrm{Tor}H_1(V_{>,\xi };𝐙)\\ ^{\iota _<,}& & ^{\iota _>,}\\ \mathrm{Tor}H_1(\stackrel{~}{X}_{K_<}^{\mathrm{}};𝐙)& \stackrel{c,}{}& \mathrm{Tor}H_1(\stackrel{~}{X}_{K_>}^{\mathrm{}};𝐙)\end{array}$$
We can define lk( , ) of $`K_{\mathrm{}}`$ by using $`V_{\mathrm{},\xi }`$ ($`\mathrm{}=<,>`$).
Let $`!=a,b.`$ Let $`x_!\mathrm{Tor}H_1(\stackrel{~}{X}_{K_<}^{\mathrm{}};𝐙)`$. Then there is $`x_!^{}\mathrm{Tor}H_1(V_{<,\xi };𝐙)`$ such that $`\iota _<(x_!^{})=x_!`$. Then $`\iota _>h(x_!^{})=c(x_!^{})`$. Therefore
lk$`(x_a,x_b)`$ for $`K_<`$
=lk$`(x_a^{},x_b^{})`$ for $`\widehat{V}_{<,\xi }`$ ($`\widehat{V}_{<,\xi }`$ is defined for $`V_{<,\xi }`$ as in §4.)
=lk$`(h(x_a^{}),h(x_b^{}))`$ for $`\widehat{V}_{>,\xi }`$ ($`\widehat{V}_{>,\xi }`$ is defined for $`V_{>,\xi }`$ as in §4.)
=lk$`(\iota _>h(x_a^{}),\iota _>h(x_b^{}))`$ for $`K_>`$
=lk$`(c(x_a^{}),c(x_b^{}))`$ for $`K_>`$
Therefore the Farber-Levine pairing for $`K_>`$ coincides with that for $`K_<`$. This completes the proof of Theorem 2.1.(2).
We next prove Theorem 2.2. Let $`\nu `$ be a homomorphism $`H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)𝐙_𝐝`$. We consider $`\stackrel{~}{\eta }(K,\nu )`$. Suppose $`K`$ is ribbon-move equivalent to the trivial 2-knot. Take $`\iota `$ and $`\gamma `$ as in §5. Then we have the following commutative diagram.
$$\begin{array}{ccc}H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)& \stackrel{\nu }{}& 𝐙_𝐝\\ _{\iota (\gamma ^1)}& & \\ H_1(\widehat{V_\xi };𝐙)& & \end{array}$$
By Theorem 2.1.(1), $`H_1(\stackrel{~}{X}_K^{\mathrm{}};𝐙)𝐙\mathrm{}𝐙`$. Hence
$$\begin{array}{ccc}𝐙\mathrm{}𝐙& \stackrel{\nu }{}& 𝐙_𝐝\\ _{\iota (\gamma ^1)}& & \\ H_1(\widehat{V_\xi };𝐙)& & \end{array}$$
By an elementary discussion on homomorphisms, we have:
$$\begin{array}{ccc}𝐙& \stackrel{\nu ^{}}{}& 𝐙_𝐝\\ _\zeta & & \\ H_1(\widehat{V_\xi };𝐙)& & \end{array}$$
The above homomorphism $`H_1(\widehat{V_\xi };𝐙)𝐙`$ is called $`\zeta `$. Then
$`\nu ^{}\zeta =`$ $`\nu \iota (\gamma ^1)`$ Then we can regard
$`\zeta `$Hom$`(H_1(\widehat{V_\xi };𝐙),𝐙)H^1(\widehat{V_\xi };𝐙)`$
$`\{\mathrm{homotopyclassesofmaps}\widehat{V_\xi }K(𝐙,1)\}`$
$`\nu ^{}\zeta `$Hom$`(H_1(\widehat{V_\xi };𝐙),𝐙_𝐝)H^1(\widehat{V_\xi };𝐙_𝐝)`$
$`\{\mathrm{homotopyclassesofmaps}\widehat{V_\xi }K(𝐙_𝐝,1)\}`$
The above diagram induces the following diagram.
$$\begin{array}{ccc}K(𝐙,1)& \stackrel{\nu ^{}}{}& K(𝐙_𝐝,1)\\ _\zeta & & \\ \widehat{V_\xi }& & \end{array}$$
The above diagram induces the following homomorphism.
$$\begin{array}{cccc}& \mathrm{\Omega }^3(K(𝐙,1))& \stackrel{\nu ^{}}{}& \mathrm{\Omega }^3(K(𝐙_𝐝,1))\\ & & & \\ & [(\widehat{V_\xi },\zeta )]& & [(\widehat{V_\xi },\nu ^{}\zeta )]\end{array}$$
By bordism theory $`\mathrm{\Omega }^3(K(𝐙,1))\mathrm{\Omega }^3(S^1)0`$. Therefore $`[(\widehat{V_\xi },\nu ^{}\zeta )]=0`$ Therefore $`(\widehat{V_\xi },\nu ^{}\zeta )=(W,\tau )`$ for a compact oriented 4-manifold $`W`$. Hence $`\stackrel{~}{\eta }(\widehat{V_\xi },\nu ^{}\zeta )=\frac{1}{1}(\overline{\sigma }(W,\tau )\sigma (W))`$. Hence $`\stackrel{~}{\eta }(\widehat{V_\xi },\nu ^{}\zeta )𝐙`$. Hence $`\stackrel{~}{\eta }(K,\nu )=\pi (\stackrel{~}{\eta }(\widehat{V_\xi },\nu ^{}\zeta ))=0`$. Therefore Theorem 2.2 holds.
Appendix
The author gives a proof of the following proposition.
Proposition. Let $`X`$ be an oriented compact $`(m+2)`$-dimensional manifold. Let $`X\varphi `$. Let $`M`$ be an oriented closed $`m`$-dimensional manifold which is embedded in $`X`$. Let $`MX\varphi `$. Let $`[M]=0H_m(X;𝐙)`$. Then there is an oriented compact $`(m+1)`$-dimensional manifold $`P`$ such that $`P`$ is embedded in $`X`$ and that $`P=X`$.
Proof. Let $`\nu `$ be the normal bundle of $`M`$ in $`X`$. By Theorem 2 in P.49 of $`\nu `$ is a product bundle. By using $`\nu `$ and the collar neighborhood of $`X`$ in $`X`$, we can take a compact oriented $`(m+2)`$-manifold $`NX`$ with the following properties.
(1) $`NM\times D^2`$. (Hence $`N=M\times S^1`$.)
(2) $`NX`$ $`=(N)(X)`$ $`=MX`$. (Hence (Int$`N)X=\varphi `$. )
Take $`X(`$Int$`N)`$. (Note $`X(`$Int$`N)X`$.) There is a cell decomposition:
$`X(\mathrm{Int}N)`$
$`=(N)(X)`$ (1-cells $`e^1)`$ (2-cells $`e^2)`$ (3-cells $`e^3)`$ (one 4-cell $`e^4)`$.
We can suppose that this decomposition has only one 0-cell $`e^0`$ which is in $`(N)(X)`$.
There is a continuous map $`s_0:(N)(X)S^1`$ with the following properties, where $`p`$ is a point in $`S^1`$.
(1) $`s_0(X)=p`$. (Hence $`s_0((N)(X))=p`$ and $`s_0(e^0)=p`$. )
(2) $`s_0|_N:M\times S^1S^1`$ is a projection map $`(x,y)y`$.
Let $`S_F^1`$ be a fiber of the $`S^1`$-fiber bundle $`N=M\times S^1`$. Since $`[M]=0H_m(X;𝐙)`$, $`[S_F^1]`$ generates $`𝐙`$ $`H_1(X`$Int $`N,X;𝐙)`$. (We can prove as in the proof of Theorem 3 in P.50 of )
Let $`f:H_1(X\mathrm{Int}N,X;𝐙)H_1(X\mathrm{Int}N,X;𝐙)/`$Tor be the natural projection map. Let $`\{f([S_F^1]),u_1,\mathrm{},u_k\}`$ be a set of basis of
$`H_1(X\mathrm{Int}N,X;𝐙)/`$Tor. Take a continuous map
$`s_1:(N)(X)`$(1-cells $`e^1)S^1`$
with the following properties.
(1) $`s_1|_{(N)(X)}=s_0`$
(2) $`s_1|_{e^0e^1}:e^0e^1S^1`$ satisfies the following condition: If $`f([e^0e^1])=`$
$`n_0f([S_f^1])+\mathrm{\Sigma }_{j=1}^kn_ju_j`$ $`H_1(X\mathrm{Int}N,X;𝐙)/`$Tor ($`n_{}𝐙`$),
then deg$`(s_1|_{e^0e^1})=n_0`$.
Note that, if a circle $`C`$ is nul-homologous in $`(N)(X)`$(1-cells $`e^1)`$, then deg$`(s_1|_C)=0`$.
Claim. There is a continuous map
$`s_2:(N)(X)`$(1-cells $`e^1)`$(2-cells $`e^2)S^1`$
such that $`s_2|_{(N)(X)(1\mathrm{cells}e^1)}`$ =$`s_1`$.
Proof. It is trivial that $`[e^2]=0`$ $`H_1((N)(X)`$(1-cells $`e^1);𝐙)`$. Hence deg($`s_1|_{e^2}`$) is zero. Hence $`s_1|_{e^2}`$ extends to $`e^2`$. Hence the above Claim holds.
The continuous map $`s_2`$ extends to a continuous map $`s:X(\mathrm{Int}N)S^1`$ since $`\pi _l(S^1)=0(l2)`$. We can suppose $`s`$ is a smooth map.
Let $`qp`$. Let $`q`$ be a regular value. Hence $`s^1(q)`$ be an oriented compact manifold. $`\{s^1(q)\}\{(N)X\}`$. Since $`qp`$, $`s^1(q)X=\varphi .`$ Hence $`\{s^1(q)\}N`$. Furthermore we have $`s^1(q)N=\{s^1(q)\}=M\times \{r\}`$, where $`r`$ is a point in $`S^1`$. By using $`N`$ and $`s^1(q)`$, Proposition holds.
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# Ribbon-moves of 2-links preserve the 𝜇-invariant of 2-links
Abstract. We introduce ribbon-moves of 2-knots, which are operations to make 2-knots into new 2-knots by local operations in $`B^4`$. (We do not assume the new knots is not equivalent to the old ones.)
Let $`L_1`$ and $`L_2`$ be 2-links. Then the following hold.
(1) If $`L_1`$ is ribbon-move equivalent to $`L_2`$, then we have
$$\mu (L_1)=\mu (L_2)$$
.
(2) Suppose that $`L_1`$ is ribbon-move equivalent to $`L_2`$. Let $`W_i`$ be arbitrary Seifert hypersurfaces for $`L_i`$. Then the torsion part of $`H_1(W_1)H_1(W_2)`$ is congruent to $`GG`$ for a finite abelian group $`G`$.
(3) Not all 2-knots are ribbon-move equivalent to the trivial 2-knot.
(4) The inverse of (1) is not true.
(5) The inverse of (2) is not true.
Let $`L=(L_1,L_2)`$ be a sublink of homology boundary link. Then we have: (i) $`L`$ is ribbon-move equivalent to a boundary link. (ii) $`\mu (L)=\mu (L_1)+\mu (L_2)`$.
We would point out the following facts by analogy of the discussions of finite type invariants of 1-knots although they are very easy observations. By the above result (1), we have: the $`\mu `$-invariant of 2-links is an order zero finite type invariant associated with ribbon-moves and there is a 2-knot whose $`\mu `$-invariant is not zero. The mod 2 alinking number of $`(S^2,T^2)`$-links is an order one finite type invariant associated with the ribbon-moves and there is an $`(S^2,T^2)`$-link whose mod 2 alinking number is not zero.
§1. Introduction
In this paper we discuss ribbon-moves.
An (oriented) (ordered) $`m`$-component 2-(dimensional) link is a smooth, oriented submanifold $`L=\{K_1,\mathrm{},K_m\}`$ of $`S^4`$, which is the ordered disjoint union of $`m`$ manifolds, each diffeomorphic to the $`2`$-sphere. If $`m=1`$, then $`L`$ is called a 2-knot. We say that 2-links $`L_1`$ and $`L_2`$ are equivalent if there exists an orientation preserving diffeomorphism $`f:`$ $`S^4`$ $``$ $`S^4`$ such that $`f(L_1)`$=$`L_2`$ and that $`f|_{L_1}:`$ $`L_1`$ $``$ $`L_2`$ is an order and orientation preserving diffeomorphism. Let $`id:S^4`$ $``$ $`S^4`$ be the identity. We say that 2-links $`L_1`$ and $`L_2`$ are identical if $`id(L_1)`$=$`L_2`$ and that $`id|_{L_1}:L_1`$ $`L_2`$ is an order and orientation preserving diffeomorphism.
We define ribbon-moves of 2-links.
Definition 1.1. Let $`L_1=(K_{1,1}\mathrm{}K_{1,m})`$ and $`L_2=(K_{2,1}\mathrm{}K_{2,m})`$ be 2-knots in $`S^4`$. We say that $`L_2`$ is obtained from $`L_1`$ by one ribbon-move if there is a 4-ball $`B`$ of $`S^4`$ with the following properties.
(1) $`L_1(BL_1)=L_2(BL_2)`$.
$`K_{1,j}(BK_{1,j})=K_{2,j}(BK_{2,j})`$
These diffeomorphism maps are orientation preserving.
(2) $`BL_1`$ is drawn as in Figure 1.1. $`BL_2`$ is drawn as in Figure 1.2.
We regard $`B`$ as (a close 2-disc)$`\times [0,1]\times \{t|1t1\}`$. We put $`B_t=`$(a close 2-disc)$`\times [0,1]\times \{t\}`$. Then $`B=B_t`$. In Figure 1.1 and 1.2, we draw $`B_{0.5},B_0,B_{0.5}`$ $`B`$. We draw $`L_1`$ and $`L_2`$ by the bold line. The fine line denotes $`B_t`$.
$`BL_1`$ (resp. $`BL_2`$) is diffeomorphic to $`D^2(S^1\times [0,1])`$.
$`BL_1`$ has the following properties: $`B_tL_1`$ is empty for $`1t<0`$ and $`0.5<t1`$. $`B_0L_1`$ is diffeomorphic to $`D^2(S^1\times [0,0.3])(S^1\times [0.7,1])`$. $`B_{0.5}L_1`$ is diffeomorphic to $`(S^1\times [0.3,0.7])`$. $`B_tL_1`$ is diffeomorphic to $`S^1S^1`$ for $`0<t<0.5`$.
$`BL_2`$ has the following properties: $`B_tL_2`$ is empty for $`1t<0.5`$ and $`0<t1`$. $`B_0L_2`$ is diffeomorphic to $`D^2(S^1\times [0,0.3])(S^1\times [0.7,1])`$. $`B_{0.5}L_2`$ is diffeomorphic to $`(S^1\times [0.3,0.7])`$. $`B_tL_2`$ is diffeomorphic to $`S^1S^1`$ for $`0.5<t<0`$.
We do not assume which the orientation of $`BL_1`$ (resp. $`BL_2`$ ) is.
Figure 1.1.
Figure 1.2.
Suppose that $`L_2`$ is obtained from $`L_1`$ by one ribbon-move and that $`L_2^{}`$ is equivalent to $`L_2`$. Then we also say that $`L_2^{}`$ is obtained from $`L_1`$ by one ribbon-move. If $`L_1`$ is obtained from $`L_2`$ by one ribbon-move, then we also say that $`L_2`$ is obtained from $`L_1`$ by one ribbon-move.
Definition 1.2. 2-knots $`L_1`$ and $`L_2`$ are said to be ribbon-move equivalent if there are 2-knots $`L_1=\overline{L}_1,\overline{L}_2,\mathrm{},\overline{L}_{p1},\overline{L}_p=L_2`$ ($`p𝐍,p2`$) such that $`\overline{L}_i`$ is obtained from $`\overline{L}_{i1}`$ $`(1<ip)`$ by one ribbon-move.
In this paper we discuss the following problems.
Problem 1.3. Let $`L_1`$ and $`L_2`$ be 2-links. Consider a necessary (resp. sufficient, necessary and sufficient ) condition that $`L_1`$ and $`L_2`$ are ribbon-move equivalent. In particular, is there a 2-knot which is not ribbon-move equivalent to the trivial 2-knot?
Note. (1) Of course all $`m`$-component ribbon 2-links are ribbon-move equivalent to the trivial $`m`$-component link. See Appendix for the definition of ribbon 2-links.
(2) By using §4 of , it is easy to prove that there is a nonribbon 2-knot which is ribbon-move equivalent to the trivial 2-knot.
Our motivation is as follows. We hope to investigate ‘link space’ $`E=\{f|f:S^2\mathrm{}S^2S^4`$ embeddings $`\}`$. In the case of 1-dimensional knots and links, we know that it is useful to investigate the space of immersions of circles in order to help investigate the space of embeddings. To discuss the space of immersions and that of embeddings is to discuss local moves (or knotting operations). In the case of 1-dimensional knots and links, we find many relations among ‘link space,’ local moves, invariants of links, and QFT. (See etc.) In 1-dimensional case, it is easy to find an unknotting operation. But high dimensional case, our first task is to define what kind of local moves we use. In this paper we discuss ribbon-moves as one of such moves.
This article is based on . After , the author discusses relations between ribbon-moves of 2-knots and the Levie-Farber pairing and the Atiyah-Patodi-Singer-Casson-Gordon-Ruberman $`\stackrel{~}{\eta }`$-invariants of 2-knots (see ). In the author discussed relations between local moves of $`n`$-knots and some invariants of $`n`$-knots.
§2. Main results
Theorem 2.1 Let $`L_1`$ and $`L_2`$ be 2-links in $`S^4`$. Suppose that $`L_1`$ is obtained from $`L_2`$ by one ribbon-move. Then there are Seifert hypersurfaces $`V_1`$ for $`L_1`$ and $`V_2`$ for $`L_2`$ such that ($`V_1`$, $`\sigma _1`$) is spin preserving diffeomorphic to ($`V_2`$, $`\sigma _2`$), where $`\sigma _i`$ is a spin structure induced from the unique one on $`S^4`$.
By using Theorem 2.1, we prove Theorem 2.2 and 2.3.
Theorem 2.2. If 2-links $`L`$ and $`L^{}`$ are ribbon-move equivalent, then $`\mu (L)`$=$`\mu (L^{})`$.
In §3 we define the $`\mu `$-invariant of 2-links.
Theorem 2.3. Let $`L_1`$ and $`L_2`$ be 2-links in $`S^4`$. Suppose that $`L_1`$ are ribbon-move equivalent to $`L_2`$. Let $`W_i`$ be arbitrary Seifert hypersurfaces for $`L_i`$. Then the torsion part of $`H_1(W_1)H_1(W_2)`$ is congruent to $`GG`$ for a finite abelian group $`G`$.
By using Theorem 2.2 we prove Corollary 2.4. By using Theorem 2.3 we also prove Corollary 2.4.
Corollary 2.4. Not all 2-knots are ribbon-move equivalent to the trivial 2-knot.
By using Theorem 2.3 we prove Corollary 2.5.
Corollary 2.5. There is a 2-knot $`K`$ such that $`\mu (K)=0`$ and that $`K`$ is not ribbon-move equivalent to the trivial 2-knot.
By using Theorem 2.2, we prove Corollary 2.6.
Corollary 2.6. The inverse of Theorem 2.3 is not true.
In §3-7 we prove the above results.
In §8 we prove that: Let $`L=(L_1,L_2)`$ be a sublink of homology boundary link. Then the following hold. (1) $`L`$ is ribbon-move equivalent to a boundary link. (2) $`\mu (L)=\mu (L_1)+\mu (L_2)`$.
In §9 we would point out the following facts by analogy of the discussions of finite type invariants of 1-knots although they are very easy observations. By Theorem 2.2, we have: the $`\mu `$-invariant of 2-links is an order zero finite type invariant associated with ribbon-moves and there is a 2-knot whose $`\mu `$-invariant is not zero. The mod 2 alinking number of $`(S^2,T^2)`$-links is an order one finite type invariant associated with the ribbon-moves and there is an $`(S^2,T^2)`$-link whose mod 2 alinking number is not zero.
§3. The $`\mu `$-invariant of 2-links
See §IV of for the spin structures and the $`\mu `$-invariant of closed spin 3-manifolds.
Definition. Let $`L=(K_1,\mathrm{},K_m)`$ be a 2-link. Let $`V`$ be a Seifert hypersurface for $`L`$. Note that $`V`$ is oriented so that the orientation is compatible with that on $`L`$. A spin structure $`\sigma `$ on $`V`$ is induced from the unique spin structure on $`S^4`$. Attach $`m`$ 3-dimensional 3-handles to $`V`$ along each component of the boundary. Then we obtain the closed oriented 3-manifold $`\widehat{V}`$. The spin structure $`\sigma `$ extends over $`\widehat{V}`$ uniquely. Call it $`\widehat{\sigma }`$. We define the $`\mu `$-invariant $`\mu (L)`$ of the 2-link $`L`$ to be the $`\mu `$-invariant $`\mu ((\widehat{V},\widehat{\sigma }))`$ $`𝐙_{\mathrm{𝟏𝟔}}`$ of the closed spin 3-manifold $`(\widehat{V},\widehat{\sigma })`$.
Claim. Under the above conditions $`\mu (L)`$ is independent of the choice of $`V`$.
Proof. P.580 of proved the above Claim when $`L`$ is a knot.
says:
Fact 3.1. () Let $`V`$ and $`V^{}`$ be Seifert hypersurfaces for $`L`$. Then we have: there are Seifert hypersurfaces $`V=V_1`$,$`V_2`$,…,$`V_{p1}`$,$`V_p`$ for $`L`$ with the following properties.
(1)The embedding map of $`V_p`$ is isotopic to that of $`V^{}`$, where we do not fix the boundary of the image. (Note. $`[VV^{}]`$ is not zero in general in $`H_3(S^4L;𝐙)`$. But we can set $`[V_1V_p]=0`$ $`H_3(S^4L;𝐙)`$. )
(2) For $`V_i`$ and $`V_{i+1}`$ ($`i=1,\mathrm{},p1`$), there is a compact oriented 4-manifold $`W_i`$ embedded in $`S^4`$ which has a handle decomposition
$`W_i=(V_i\times [0,1])`$ (one $`q`$-handle) $`(V_{i+1}\times [0,1])`$ ($`q\{1,2,3\}`$).
We give $`W_i`$ a spin structure induced from the unique one on $`S^4`$.
The following two spin structures on $`V_1`$ coincide one another. Call it $`\sigma _1`$.
(i) The spin structure induced from the unique one on $`S^4`$
(ii) The spin structure induced from the one on $`W_1`$.
The following two spin structures on $`V_p`$ coincide one another. Call it $`\sigma _p`$.
(i) The spin structure induced from the unique one on $`S^4`$
(ii) The spin structure induced from the one on $`W_{p1}`$.
The following three spin structures on $`V_i`$ coincide each other ($`i=2,\mathrm{},p1`$). Call it $`\sigma _i`$.
(i) The spin structure induced from the unique one on $`S^4`$.
(ii) The spin structure induced from the one on $`W_i`$.
(iii) The spin structure induced from the one on $`W_{i+1}`$.
The 3-dimensional closed oriented spin 3-manifolds $`(\widehat{V}_i,\widehat{\sigma }_i)`$ are defined from $`(V_i,\sigma _i)`$ as in the above Definition ($`i=1,\mathrm{},p`$). (See §IV of for the way to induce spin structures on manifolds from those on others. )
Let $`x,y`$ be arbitrary elements of $`H_2(W_i;𝐙)/\mathrm{Tor}`$. Let $`xy`$ be the intersection product.
We prove: $`xy=0.`$
There is an oriented closed surface $`F`$ embedded in $`W_i`$ which represents $`x`$. Since $`F`$ is embedded in $`S^4`$, $`[F][F]=0`$. Hence $`xx=0`$ for any element $`xH_2(W_i;𝐙)/\mathrm{Tor}`$. Hence $`xy=0`$ for arbitrary elements $`x,yH_2(W_i;𝐙)/\mathrm{Tor}`$. Hence the signature of the intersection form
$`H_2(W_i;𝐙)/\mathrm{Tor}\times H_2(W_i;𝐙)/\mathrm{Tor}𝐙`$$`(x,y)xy`$
is the zero map. Hence $`\sigma (W_i)=0`$.
Therefore $`\mu ((\widehat{V}_i,\widehat{\sigma }_i))\mu ((\widehat{V}_{i+1},\widehat{\sigma }_{i+1}))`$= $`\mu ((V_i,\sigma _i)((V_{i+1},\sigma _{i+1}))`$=
{mod 16 $`\sigma (W_i)\}=0`$. Hence $`\mu ((\widehat{V}_i,\widehat{\sigma }_i))=\mu (\widehat{V}_{i+1},\widehat{\sigma }_{i+1})`$. $`(i=1,\mathrm{}p1.)`$
Therefore $`\mu ((\widehat{V}_1,\widehat{\sigma }_1))`$= $`\mu ((\widehat{V}_2,\widehat{\sigma }_2))`$=…= $`\mu ((\widehat{V}_{p1},\widehat{\sigma }_{p1}))`$= $`\mu (\widehat{V}_p,\widehat{\sigma }_p)`$.
This completes the proof.
§4. Proof of Theorem 2.1
In order to prove Theorem 2.1, we introduce (1,2)-pass-moves of 2-links.
Definition 4.1. Let $`L_1=(K_{1,1}\mathrm{}K_{1,m})`$ and $`L_2=(K_{2,1}\mathrm{}K_{2,m})`$ be 2-knots in $`S^4`$. We say that $`L_2`$ is obtained from $`L_1`$ by one (1,2)-pass-move if there is a 4-ball $`B`$ $`S^4`$ with the following properties. We draw $`B`$ as in Definition 1.1.
(1) $`L_1(BL_1)`$=$`L_2(BL_2)`$.
$`K_{1,j}(BK_{1,j})=K_{2,j}(BK_{2,j})`$
These diffeomorphism maps are orientation preserving.
(2) $`BL_1`$ is drawn as in Figure 4.1. $`BL_2`$ is drawn as in Figure 4.2.
Figure 4.1.
Figure 4.2.
The orientation of the two discs in the Figure 4.1 (resp. Figure 4.2) is compatible with the orientation which is naturally determined by the $`(x,y)`$-arrows in the Figure. We do not assume which the orientations of the annuli in the Figures are.
Suppose that $`L_2`$ is obtained from $`L_1`$ by one (1,2)-pass-move and that $`L_2^{}`$ is equivalent to $`L_2`$. Then we also say that $`L_2^{}`$ is obtained from $`L_1`$ by one (1,2)-pass-move .
If $`L_1`$ is obtained from $`L_2`$ by one (1,2)-pass-move, then we also say that $`L_2`$ is obtained from $`L_1`$ by one (1,2)-pass-move .
2-knots $`L_1`$ and $`L_2`$ are said to be (1,2)-pass-move equivalent if there are 2-knots $`L_1=\overline{L}_1,\overline{L}_2,\mathrm{},\overline{L}_{p1},\overline{L}_p=L_2`$ $`(p𝐍,p2)`$ such that $`\overline{L}_i`$ is obtained from $`\overline{L}_{i1}`$ $`(1<ip)`$ by one (1,2)-pass-move.
Proposition 4.2. Let $`L`$ and $`L^{}`$ be 2-links. Then the following conditions (1) and (2) are equivalent.
(1) $`L`$ is (1,2)-pass-move equivalent to $`L^{}`$.
(2) $`L`$ is ribbon-move equivalent to $`L^{}`$.
It is obvious that Proposition 4.2 follows from Proposition 4.3.
Proposition 4.3. Let $`L`$ and $`L^{}`$ be 2-links. Then the following hold.
(1) If $`L`$ is obtained from $`L^{}`$ by one ribbon-move, then $`L^{}`$ is obtained from $`L`$ by one (1,2)-pass-move.
(2) If $`L`$ is obtained from $`L^{}`$ by one (1,2)-pass-move, then $`L^{}`$ is obtained from $`L`$ by two ribbon-move.
Proposition 4.3.(2) is obvious.
Proposition 4.3.(1) follows from Proposition 4.4 because: The pair of a manifold and a submanifold, $`(`$ the 4-ball, $`(`$the 2-link$`)(`$the 4-ball$`)`$$`)`$, in Figure 4.1 is included in the pair $`(`$ the 4-ball, $`(`$the 2-link$`)(`$the 4-ball$`)`$$`)`$ in Figure 4.4.
Proposition 4.4. Let $`L_1=(K_{1,1},\mathrm{},K_{1,m})`$ and $`L_2=(K_{2,1},\mathrm{},K_{2,m})`$ be 2-links in $`S^4`$. Then the following two conditions $`(I)`$ and $`(II)`$ are equivalent.
$`(I)`$ $`L_1`$ is equivalent to $`L_2`$.
$`(II)`$ There is a 4-ball $`B`$ $`S^4`$ with the following properties. We draw $`B`$ as in Definition 1.1.
(1) $`L_1(BL_1)=L_2(BL_2)`$.
$`K_{1,i}(BK_{1,i})=K_{2,i}(BK_{2,i})`$ for each $`i`$.
These diffeomorphism maps are orientation preserving.
(2) $`BL_1`$ is drawn as in Figure 4.3. $`BL_2`$ is drawn as in Figure 4.4.
Figure 4.3.
Figure 4.4.
The orientation of $`BL_2`$ is compatible with the orientation which is naturally determined by the $`(x,y)`$-arrows in the Figure 4.4.
It is obvious that Proposition 4.4 follows from Proposition 4.5.
Proposition 4.5. Let $`L_1=(K_{1,1},\mathrm{},K_{1,m})`$ and $`L_2=(K_{2,1},\mathrm{},K_{2,m})`$ be 2-links in $`S^4`$. Then the following two conditions $`(I)`$ and $`(II)`$ are equivalent.
$`(I)`$ $`L_1`$ is equivalent to $`L_2`$.
$`(II)`$ There is a 4-ball $`BS^4`$ with the following properties. We draw $`B`$ as in Definition 1.1.
(1) $`L_1(BL_1)=L_2(BL_2)`$.
$`K_{1,i}(BK_{1,i})`$=$`K_{2,i}(BK_{2,i})`$ for each $`i`$.
These diffeomorphism maps are orientation preserving.
(2) $`BL_1`$ is drawn as in Figure 4.5. $`BL_2`$ is drawn as in Figure 4.6.
Figure 4.5.
Figure 4.6.
We do not assume which the orientation of $`BL_1`$ (resp. $`BL_2`$ ) is.
Proposition 4.5 follows from Proposition 4.6 because: The pair of a manifold and a submanifold, $`(`$ the 4-ball, $`(`$the 2-link$`)(`$the 4-ball$`)`$$`)`$, in Figure 4.5 (resp. Figure 4.6) is made from the pair $`(`$ the 4-ball, $`(`$the 2-link$`)(`$the 4-ball$`)`$$`)`$ in Figure 4.7 (resp. Figure 4.8) by a rotation through $`90^{}`$ around an appropriate plane in the 4-ball in Figure 4.5 (resp. Figure 4.6) and by isotopy.
Proposition 4.6. Let $`L_1=(K_{1,1},\mathrm{},K_{1,m})`$ and $`L_2=(K_{2,1},\mathrm{},K_{2,m})`$ be 2-links in $`S^4`$. Then the following two conditions $`(I)`$ and $`(II)`$ are equivalent.
$`(I)`$ $`L_1`$ is equivalent to $`L_2`$.
$`(II)`$ There is a 4-ball $`BS^4`$ with the following properties. We draw $`B`$ as in Definition 1.1.
(1) $`L_1(BL_1)=L_2(BL_2)`$.
$`K_{1,i}(BK_{1,i})`$=$`K_{2,i}(BK_{2,i})`$ for each $`i`$.
These diffeomorphism maps are orientation preserving.
(2) $`BL_1`$ is drawn as in Figure 4.7. $`BL_2`$ is drawn as in Figure 4.8.
Figure 4.7.
Figure 4.8.
We do not assume which the orientation of $`BL_1`$ (resp. $`BL_2`$ ) is.
Proof of Proposition 4.6. We obtain $`L_2`$ from $`L_1`$ by an explicit isotopy, Figure 4.7 $``$ Figure 4.9 $``$ Figure 4.8. Note that the following Proposition 4.7 holds by an explicit isotopy. This completes the proof of Proposition 4.2-4.6.
Figure 4.9.
Proposition 4.7. Let $`L_1=(K_{1,1},\mathrm{},K_{1,m})`$ and $`L_2=(K_{2,1},\mathrm{},K_{2,m})`$ be 2-links in $`S^4`$. Then the following two conditions $`(I)`$ and $`(II)`$ are equivalent.
$`(I)`$ $`L_1`$ is equivalent to $`L_2`$.
$`(II)`$ There is a 4-ball $`BS^4`$ with the following properties. We draw $`B`$ as in Definition 1.1.
(1) $`L_1(BL_1)=L_2(BL_2)`$.
$`K_{1,i}(BK_{1,i})`$=$`K_{2,i}(BK_{2,i})`$ for each $`i`$.
These diffeomorphism maps are orientation preserving.
(2) $`BL_1`$ is drawn as in Figure 4.10. $`BL_2`$ is drawn as in Figure 4.11.
Figure 4.10.
Figure 4.11.
We do not assume which the orientation of $`BL_1`$ (resp. $`BL_2`$) is.
Note. Regard the operation,
‘t=0 of Figure 4.7 $``$ t=0 of Figure 4.8 $``$ t=0 of Figure 4.9,’
as an isotopy of (a part of) 1-knot. Then this operation is essentially same as the operation in the figure in the proof of Lemma 5.5 of .
Proof of Theorem 2.1. By Proposition 4.3.(1), $`L_1`$ is obtained from $`L_2`$ by one (1,2)-pass-move in a 4-ball $`B`$.
Claim 4.8. There are Seifert hypersurfaces $`V_1`$ for $`K_1`$ and $`V_2`$ for $`K_2`$ such that:
(1) $`V_1(BV_1)=V_2(BV_2)`$.
These diffeomorphism maps are orientation preserving.
(2) $`BV_1`$ is drawn as in Figure 4.12. $`BV_2`$ is drawn as in Figure 4.13.
Note. We draw $`B`$ as in Definition 1.1. We draw $`V_1`$ and $`V_2`$ by the bold line. The fine line means $`B`$.
$`BV_1`$ (resp. $`BV_2`$) is diffeomorphic to $`(D^2\times [2,3])(D^2\times [0,1])`$. We can regard $`(D^2\times [0,1])`$ as a 3-dimensional 1-handle which is attached to $`B`$. We can regard $`(D^2\times [2,3])`$ as a 3-dimensional 2-handle which is attached to $`B`$.
$`BV_1`$ has the following properties: $`B_tV_1`$ is empty for $`1t<0`$ and $`0.5<t1`$. $`B_0V_1`$ is diffeomorphic to $`(D^2\times [2,3])(D^2\times [0,0.3])(D^2\times [0.7,1])`$. $`B_{0.5}K_1`$ is diffeomorphic to $`(D^2\times [0.3,0.7])`$. $`B_tV_1`$ is diffeomorphic to $`D^2D^2`$ for $`0<t<0.5`$.
$`BV_2`$ has the following properties: . $`B_tV_2`$ is empty for $`1t<0.5`$ and $`0<t1`$. $`B_0V_2`$ is diffeomorphic to $`(D^2\times [2,3])(D^2\times [0,0.3])(D^2\times [0.7,1])`$. $`B_{0.5}V_2`$ is diffeomorphic to $`(D^2\times [0.3,0.7])`$. $`B_tV_2`$ is diffeomorphic to $`D^2D^2`$ for $`0.5<t<0`$.
Figure 4.12
Figure 4.13.
Proof of Claim. Put $`P=(`$the 3-manifolds in Figure 4.12$`)(B)`$. Note $`P=(`$the 3-manifolds in Figure 4.13$`)(B)`$. Put $`Q=L_1(S^4IntB^4)`$. Note $`Q=L_2(S^4IntB^4)`$. By applying the following Proposition to ($`PQ`$) and $`(S^4IntB^4)`$, Claim 4.8 holds.
The following Proposition is proved by using the obstruction theory. We give a proof although it is folklore.
Proposition. Let $`X`$ be an oriented compact $`(m+2)`$-dimensional manifold. Let $`X\varphi `$. Let $`M`$ be an oriented closed $`m`$-dimensional manifold which is embedded in $`X`$. Let $`MX\varphi `$. Let $`[M]=0H_m(X;𝐙)`$. Then there is an oriented compact $`(m+1)`$-dimensional manifold $`P`$ such that $`P`$ is embedded in $`X`$ and that $`P=X`$.
Proof. Let $`\nu `$ be the normal bundle of $`M`$ in $`X`$. By Theorem 2 in P.49 of $`\nu `$ is a product bundle. By using $`\nu `$ and the collar neighborhood of $`X`$ in $`X`$, we can take a compact oriented $`(m+2)`$-manifold $`NX`$ with the following properties.
(1) $`NM\times D^2`$. (Hence $`N=M\times S^1`$.)
(2) $`NX`$ $`=(N)(X)`$ $`=MX`$. (Hence (Int$`N)X=\varphi `$. )
Take $`X(`$Int$`N)`$. (Note $`X(`$Int$`N)X`$.) There is a cell decomposition:
$`X(\mathrm{Int}N)`$
$`=(N)(X)`$ (1-cells $`e^1)`$ (2-cells $`e^2)`$ (3-cells $`e^3)`$ (one 4-cell $`e^4)`$.
We can suppose that this decomposition has only one 0-cell $`e^0`$ which is in $`(N)(X)`$.
There is a continuous map $`s_0:(N)(X)S^1`$ with the following properties, where $`p`$ is a point in $`S^1`$.
(1) $`s_0(X)=p`$. (Hence $`s_0((N)(X))=p`$ and $`s_0(e^0)=p`$. )
(2) $`s_0|_N:M\times S^1S^1`$ is a projection map $`(x,y)y`$.
Let $`S_F^1`$ be a fiber of the $`S^1`$-fiber bundle $`N=M\times S^1`$. Since $`[M]=0H_m(X;𝐙)`$, $`[S_F^1]`$ generates $`𝐙`$ $`H_1(X`$Int $`N,X;𝐙)`$. (We can prove as in the proof of Theorem 3 in P.50 of )
Let $`f:H_1(X\mathrm{Int}N,X;𝐙)H_1(X\mathrm{Int}N,X;𝐙)/`$Tor be the natural projection map. Let $`\{f([S_F^1]),u_1,\mathrm{},u_k\}`$ be a set of basis of
$`H_1(X\mathrm{Int}N,X;𝐙)/`$Tor. Take a continuous map
$`s_1:(N)(X)`$(1-cells $`e^1)S^1`$
with the following properties.
(1) $`s_1|_{(N)(X)}=s_0`$
(2) $`s_1|_{e^0e^1}:e^0e^1S^1`$ satisfies the following condition: If $`f([e^0e^1])=`$
$`n_0f([S_f^1])+\mathrm{\Sigma }_{j=1}^kn_ju_j`$ $`H_1(X\mathrm{Int}N,X;𝐙)/`$Tor ($`n_{}𝐙`$),
then deg$`(s_1|_{e^0e^1})=n_0`$.
Note that, if a circle $`C`$ is nul-homologous in $`(N)(X)`$(1-cells $`e^1)`$, then deg$`(s_1|_C)=0`$.
Claim. There is a continuous map
$`s_2:(N)(X)`$(1-cells $`e^1)`$(2-cells $`e^2)S^1`$
such that $`s_2|_{(N)(X)(1\mathrm{cells}e^1)}`$ =$`s_1`$.
Proof. It is trivial that $`[e^2]=0`$ $`H_1((N)(X)`$(1-cells $`e^1);𝐙)`$. Hence deg$`(s_1|_{e^2})=0.`$ Hence $`s_1|_{e^2}`$ extends to $`e^2`$. Hence the above Claim holds.
The map $`s_2`$ extends to a continuous map $`s:X(\mathrm{Int}N)S^1`$ since $`\pi _l(S^1)=0(l2)`$. We can suppose $`s`$ is a smooth map.
Let $`qp`$. Let $`q`$ be a regular value. Hence $`s^1(q)`$ be an oriented compact manifold. $`\{s^1(q)\}\{(N)X\}`$. Since $`qp`$, $`s^1(q)X=\varphi .`$ Hence $`\{s^1(q)\}N`$. Furthermore we have $`s^1(q)N=\{s^1(q)\}=M\times \{r\}`$, where $`r`$ is a point in $`S^1`$. By using $`N`$ and $`s^1(q)`$, Proposition holds.
By Claim 4.8, there is a smooth transverse immersion $`F:V\times [1,2]S^4`$ such that $`F|_{V\times \{1\}}(V\times \{1\})=V_1`$ and $`F|_{V\times \{2\}}(V\times \{2\})=V_2`$. Give a spin structure $`\alpha `$ on $`V\times [1,2]`$ by using $`F`$. Then the following two spin structures on $`V_i`$ coincide one another. Call it $`\tau _i`$.
(i) the spin structure induced from the unique spin structure $`S^4`$
(ii) the spin structure induced from $`\alpha `$ on $`V\times [1,2]`$.
By using $`F`$, it holds that $`V_1`$ and $`V_2`$ are spin preserving diffeomorphism. This completes the proof of Theorem 2.1.
§5. Proof of Theorem 2.2
By Proposition 4.2, $`L`$ and $`L^{}`$ are (1,2)-pass-move equivalent. Take 2-links $`L=\overline{L}_1,\overline{L}_2,\mathrm{},\overline{L}_{p1},\overline{L}_p=L^{}`$ as in Definition 4.1. Obviously, it suffices to prove that $`\mu (\overline{L}_i)`$=$`\mu (\overline{L}_{i+1})`$ for each $`i`$ ($`1i<p`$). By Theorem 2.1 we have: There are Seifert hypersurfaces, $`V_{i,i+1}`$ for $`\overline{L}_i`$ and $`V_{i+1,i}`$ for $`\overline{L}_{i+1}`$, such that $`V_{i,i+1}`$ and $`V_{i+1,i}`$ are spin preserving diffeomorphism. Hence $`\mu (\overline{L}_i)`$=$`\mu (\overline{L}_{i+1})`$.
§6. The proof of Theorem 2.3
The following Fact 6.1 is an elementary fact.
Fact6.1. (Known) Let $`A,B,C,X`$ and $`Y`$ be a finite abelian group. Suppose $`ABXX`$ and $`BCYY`$. Then $`ACPP`$ for a finite abelian group $`P`$.
It is obvious that Theorem 2.3 follows from Theorem 2.1, Fact 6.1, and Proposition 6.2.
Proposition 6.2. Let $`V`$ and $`V^{}`$ be Seifert hypersurfaces for a 2-link $`L`$. Then the torsion part of $`H_1(V;𝐙)H_1(V^{};𝐙)`$ is congruent to $`GG`$ for a finite group $`G`$.
Proof. Take $`V_1,\mathrm{},V_p`$ and $`W_1,\mathrm{},W_p`$ as in Fact 3.1 and its proof. By using the Meyer-Vietoris sequence, we have Tor $`H_1(W_i;𝐙)`$ Tor $`\{H_1(V_i;𝐙)H_1(V_{i+1};𝐙)\}`$ ($`i=1,\mathrm{},p1`$). The manifold $`W_i`$ is a closed oriented 3-manifold embedded in $`S^4`$. Hence
Tor $`H_1(W_i;𝐙)G_iG_i`$ ——$`()`$
for a finite abelian group $`G_i`$. ( See e.g. . We give a proof in the following paragraph.) Hence Tor $`\{H_1(V_i;𝐙)H_1(V_{i+1};𝐙)\}`$ $`G_iG_i`$.
We give a proof for the above congruence $`()`$: By using the Meyer-Vietoris sequence $`H_i(W_i;𝐙)H_i(W_i;𝐙)H_i(\overline{S^4W_i};𝐙)H_i(S^4;𝐙)`$, Tor $`H_1(W_i;𝐙)`$ Tor $`\{H_1(W_i;𝐙)H_1(\overline{S^4W_i};𝐙)\}`$. By using the Meyer-Vietoris sequence $`H_i(W_i;𝐙)H_i(S^4;𝐙)H_i(S^4,W_i;𝐙)`$, $`H_1(W_i;𝐙)H_2(S^4,W_i;𝐙).`$ By the excision, $`H_2(S^4,W_i;𝐙)H_2(\overline{S^4W_i},W_i;𝐙).`$ By the Poincaré duality, $`H_2(\overline{S^4W_i},W_i;𝐙)H^2(\overline{S^4W_i};𝐙).`$ By the universal coefficient theorem, Tor$`H_1(\overline{S^4W_i};𝐙)`$ Tor$`H^2(\overline{S^4W_i};𝐙).`$ Hence Tor$`H_1(\overline{S^4W_i};𝐙)`$ Tor$`H_1(W_i;𝐙)`$. Hence Tor$`H_1(W_i;𝐙)`$ Tor$`H_1(\overline{S^4W_i};𝐙)`$ Tor$`H_1(\overline{S^4W_i};𝐙)`$. Hence the congruence $`()`$ holds.
By Fact 6.1, Tor $`\{H_1(V_1;𝐙)H_1(V_p;𝐙)\}`$ =$`GG`$ for a finite abelian group $`G`$. Hence Tor $`\{H_1(V;𝐙)H_1(V^{};𝐙)\}`$ =$`GG`$.
§7. The proof of Corollary 2.4, 2.5 and 2.6
Let $`K`$ be the 2-twist spun knot of a 1-knot $`A`$. Let $`M`$ be the 2-fold branched cyclic covering space of $`S^3`$ along $`A`$. By , $`\overline{MB^3}`$ is a Seifert hypersurface for $`K`$. Let $`S`$ be a Seifert matrix of $`K`$. By Lemma 12.1, Theorem 12.2, and Theorem 12.6 in Chapter XII of , there is a compact oriented 4-manifold $`X`$ with the following properties. (i) $`M=X`$. (ii) $`H_1(X;𝐙)0`$ (iii) $`H_3(X;𝐙)0`$ (iv) The intersection form $`H_2(X;𝐙)\times H_2(X;𝐙)𝐙`$ is represented by $`S+^tS`$. (Note: By using the Poincaré duality, the universal coefficient theorem, and the above conditions (ii) (iii), it holds that $`H_2(X;𝐙)`$ is torsion free.)
By the above fact (iv), the intersection form is even. By this fact, the above (ii), and P.27 of , it holds that $`X`$ is a spin manifold. Hence, for a spin structure $`\alpha `$ on $`M`$, $`\mu (M,\alpha )=`$ mod 16 $`\sigma (S+^tS)`$. (Note that there is a spin 3-manifold whose spin structure is more than one.)
(1) Let $`A`$ be the trefoil knot. Let $`S`$ be $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$. Then the intersection form of $`X`$ is represented by $`\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right)`$.
Hence we have:
(1.1) $`H_1(\overline{MB^3};𝐙)𝐙_\mathrm{𝟑}`$.
Hence $`H_1(\overline{MB^3};𝐙_\mathrm{𝟐})0`$. Hence $`\overline{MB^3}`$ has only one spin structure. Hence $`\mu (K)=\mu (M)=`$ mod 16 $`(\sigma \left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right))`$. Hence we have:
(1.2)$`\mu (K)=2`$.
(2) Let $`A`$ be the figure eight knot. Let $`S`$ be $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$. Then the intersection form of $`X`$ is represented by $`\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right)`$.
Then we have:
(2.1)$`H_1(\overline{MB^3};𝐙)𝐙_\mathrm{𝟓}`$.
Hence $`H_1(\overline{MB^3};𝐙_\mathrm{𝟐})0`$. Hence $`M`$ has only one spin structure. Hence $`\mu (K)=\mu (M)=`$ mod 16 $`(\sigma \left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right))`$. Hence we have:
(2.2) $`\mu (K)=0`$.
(3) Let $`K`$ be the 5-twist spun knot of the trefoil knot. Let $`M`$ be the Poincaré homology sphere. Then we have:
(3.1) There is a Seifert hypersurface for $`K`$ which is diffeomorphic to $`\overline{MB^3}`$. (See §65 of . )
(3.2) $`\mu (K)=\mu (M)=8`$. (See e.g. P.15 and P.67 of . )
The above (1.2) and Theorem 2.2 imply Corollary 2.4.
The above (1.1) (or (2.1)) and Theorem 2.3 imply Corollary 2.4.
The above (2.1), (2.2) and Theorem 2.3 imply Corollary 2.5.
The above (3.1), (3.2) and Theorem 2.2 imply Corollary 2.6.
§8. Any SHB link is ribbon-move equivalent to a boundary link
See P.640 of and P.536 of etc. for sublinks of homology boundary links ( i.e. SHB links ), homology boundary links and boundary links.
Theorem 8.1. Let $`L=(K_1,K_2)`$ be a 2-link. Let $`L`$ be a sublink of a homology boundary link. Then $`L`$ is ribbon-move equivalent to a boundary link.
To prove Theorem 8.1, we need lemmas. By the definition of SHB links (in P.536 of ) the following holds.
Lemma 8.1.1. Let $`L=(K_1,K_2)`$ be the 2-link in Theorem 8.1. There is a connected Seifert hypersurface $`V_i`$ for $`K_i`$ ( $`i=1,2`$ ) such that $`V_1V_2`$ is diffeomorphic to a disjoint union of 2-spheres $`S_1^2`$,…,$`S_\nu ^2`$.
We prove: Lemma 8.1.2. Let $`L=(K_1,K_2)`$ be the 2-link in Theorem 8.1. Then there is a 2-link $`L^{}=(K_1^{},K_2^{})`$ which is equivalent to $`L`$ satisfying the following condition: there is a Seifert hypersurface $`V_i^{}`$ for $`K_i^{}`$ ($`i=1,2`$) such that $`V_1^{}V_2^{}`$ is one 2-sphere $`S_0^2`$.
Proof of Lemma 8.1.2. Take $`V_1`$ and $`V_2`$ in Lemma 8.1.1. If $`\nu =0`$, then Theorem 8.1 holds. If $`\nu =1`$, then Lemma 8.1.2 holds. Suppose $`\nu 2`$.
We can suppose that $`S_1^2`$ and $`S_2^2`$ satisfy the following: There is a point $`p_1S_1^2`$, a point $`p_2S_2^2`$, and a path $`lV_1`$ such that (1) $`l=p_1p_2`$ (2) $`l(S_1^2\mathrm{}S_\nu ^2)=p_1p_2`$ (3) $`lK_1=\varphi .`$
Take a 4-dimensional 1-handle $`h^1S^4`$ whose core is $`l`$ such that $`h^1`$ is attached to $`V_2`$ along $`p_1p_2`$. Then $`h^1V_1`$ is a 3-dimensional 1-handle which is attached to $`S_1^2S_2^2`$ along $`p_1p_2`$. We carry out surgery on $`V_2`$ by using $`h^1`$. The new manifold is called $`V_2^{\mathrm{}}`$. Then $`V_2^{\mathrm{}}`$ is a connected Seifert hypersurface for $`K_2`$. When we carry out the surgery on $`V_2`$, we carry out surgery on $`S_1^2S_2^2`$ by using the 3-dimensional 1-handle $`h^1V_1`$. Then the result is a 2-sphere. Then $`V_1V_2^{\mathrm{}}`$ is $`(\nu 1)`$ 2-spheres. By the induction on $`\nu `$, Lemma 8.1.2 holds.
Lemma 8.1.3. Let $`L=(K_1,K_2)`$ be the 2-link in Theorem 8.1. Then there is a 2-link $`L^{\prime \prime }=(K_1^{\prime \prime },K_2^{\prime \prime })`$ which is equivalent to $`L`$ satisfying the following condition: there is a connected Seifert hypersurface $`V_i^{\prime \prime }`$ for $`K_i^{\prime \prime }`$ ($`i=1,2`$) such that $`V_1^{\prime \prime }V_2^{\prime \prime }`$ is one 2-disc $`D_0^2`$.
Proof of Lemma 8.1.3. Take $`V_1^{}`$ and $`V_2^{}`$ in Lemma 8.1.2. Take a point $`pK_1^{}=V_1^{}`$. Take a point $`qS_0^2=V_1^{}V_2^{}`$. Take a path $`lV_1^{}`$ such that (1) $`l=pq`$. (2) $`lS_0^2=q`$ (3) $`lK_1^{}=p.`$
Let $`N`$ be a tubular neighborhood of $`l`$ in $`V_1^{}`$, Then $`N`$ is a 3-ball. Note that $`NK_1`$ is a 2-disc, which is a tubular neighborhood of $`p`$ in $`K_1`$. Note that $`q`$Int $`N`$. Note that Int ($`NS_0^2`$) is in Int $`N`$. Then the following hold. (1) $`\overline{V_1^{}N}`$ $`V_2^{}`$ is a 2-disc. (2) $`(\overline{V_1^{}N})`$ is equivalent to $`K_1`$. (3) $`((\overline{V_1^{}N}),K_2)`$ is equivalent to $`L^{}`$ and hence to $`L`$. $`((\overline{V_1^{}N}),K_2)`$ is called $`L^{\prime \prime }=(K_1^{\prime \prime },K_2^{\prime \prime })`$. This completes the proof of Lemma 8.1.3.
Proof of Theorem 8.1. Take $`V_1^{\prime \prime }`$ and $`V_2^{\prime \prime }`$ in Lemma 8.1.3. We can suppose that $`D_0^2K_1^{\prime \prime }`$ and that $`D^2`$ Int$`V_2^{\prime \prime }`$. Take a 3-ball $`PV_2^{\prime \prime }`$ such that $`PK_2^{\prime \prime }`$ is a 2-disc and that $`D^2`$ Int $`P`$. Then $`V_1^{\prime \prime }(\overline{V_2^{\prime \prime }P})=\varphi `$. Let $`L^!`$ be a 2-link ($`K_1^{\prime \prime },(\overline{V_2^{\prime \prime }P})`$). Then $`L^!`$ is a boundary link. Furthermore $`L^!`$ is obtained from $`L^{\prime \prime }`$ by an operation that we fix $`K_1^{\prime \prime }`$ and that we move $`K_2^{\prime \prime }`$ to $`(\overline{V_2^{\prime \prime }P})`$ so that we fix $`K_2^{\prime \prime }(\overline{V_2^{\prime \prime }P})`$. This operation on $`L^!`$ is essentially same as a ribbon move. This completes the proof of Theorem 8.1. 3
Theorem 8.2 Let $`L=(K_1,K_2)`$ be an SHB 2-link. Then $`\mu (L)=\mu (K_1)+\mu (K_2)`$.
Proof. By Theorem 8.1, $`L`$ is ribbon-move equivalent to a boundary 2-link $`\overline{L}=(\overline{K_1},\overline{K_2})`$. Let $`\overline{V_i}`$ be a Seifert hypersurface for $`\overline{K_i}`$ such that $`\overline{V_1}\overline{V_2}=\varphi `$. Then $`\mu (\overline{L})`$=$`\mu (\overline{V_1}h^3)+\mu (\overline{V_2}h^3)`$, where $`h^3`$ is a 3-dimensional 3-handle which is attached to $`\overline{V_i}`$ along the 2-sphere $`\overline{V_i}`$. Hence $`\mu (\overline{L})`$=$`\mu (\overline{K_1})+\mu (\overline{K_2})`$. By Theorem 2.2, $`\mu (L)=\mu (\overline{L})`$ and $`\mu (K_i)=\mu (\overline{K_i})`$. Hence $`\mu (L)=\mu (K_1)+\mu (K_2)`$.
Problem 8.3.
(1) Let $`L=(K_1,K_2)`$ be a 2-link. Then does $`\mu (L)=\mu (K_1)+\mu (K_2)`$ hold?
(2) Is there an $`n`$-link which is not an SHB link ($`n2`$)?
§9. Discussions
We would point out the following facts by analogy of the discussions of finite type invariants of 1-knots (e.g. ) although they are very easy observations.
By using Theorem 2.2 we have: The $`\mu `$-invariant of 2-links is an order zero finite type invariant if we define ‘order of invariants’ by using ribbon-moves ( e.g. as follows ), and there is a 2-knot whose $`\mu `$-invariant is not zero.
We define order, for example, as follows. Let $`I_n`$ be the set of immersed $`m`$ 2-spheres with the conditions: (1) The set of singular points consists of double points. (2) Each component of the set of singular points is as in Figure 9.3. (3) The components of the set of singular points are $`n`$. Then $`I_0`$ is the set of $`m`$-component 2-links. Let $`v_i()G`$ be an invariant of elements of $`I_i`$, where $`G`$ is a group. Let $`X_0`$ be an element of $`I_{i+1}`$. Let $`X_+`$ and $`X_{}`$ be elements of $`I_i`$. Suppose that $`X_0`$, $`X_+`$ and $`X_{}`$ coincide in $`S^4B^4`$. Suppose that $`X_0B`$ is drawn as in Figure 9.3, $`X_+B`$ is drawn as in Figure 9.1, and $`X_{}B`$ is drawn as in Figure 9.2. In Figure 9.1, 9.2, 9.3, we do not assume the orientation of $`X_{}B`$ and that of $`B`$. If we have $`\{v_{i+1}(X)\}^2=\{v_i(X_+)v_i(X_{})\}^2`$ and $`v_i`$ is zero for $`i>p`$, then we call $`v_{}()`$ is an order $`p`$ invariant of 2-links.
We define a link-type invariant $`v()`$ of $`(S^2,T^2)`$-links. ( See for detail. See $`(S^2,T^2)`$-links for .) We call it the alinking number of $`(S^2,T^2)`$-links. Let $`L=(L_S,L_T)`$ be a $`(S^2,T^2)`$-link. Let $`\iota `$ be the map $`H^1(S^4L_S;𝐙)H^1(L_T;𝐙)`$ induced by the inclusion.
Define
$$v(L)=\{\begin{array}{cc}n\hfill & \text{if }H^1(L_T;𝐙)\text{/Im}\iota 𝐙(𝐙/(n𝐙))\text{ }(n2,n𝐍)\hfill \\ 1\hfill & \text{if }H^1(L_T;𝐙)\text{/Im}\iota 𝐙\hfill \\ 0\hfill & \text{if }H^1(L_T;𝐙)\text{/Im}\iota 𝐙𝐙\text{.}\hfill \end{array}$$
Then the mod 2 alinking number of $`(S^2,T^2)`$-links is an order one finite type invariant if we define ‘order of invariants’ by using ribbon-moves (e.g. as above), and there is an $`(S^2,T^2)`$-link whose mod 2 alinking number is not zero. (The proof is similar to the proof that the linking number of 2-component 1-links is an order one finite type invariant. See .)
Note. and etc. try to make a high-dimensional version of works on 1-links by Jones, Witten, Kontsevich, Vassiliev, etc. (in etc. )
Appendix
A ribbon 2-link is a 2-link $`L=(K_1,\mathrm{},K_m)`$ with the following properties. There is a self-transverse immersion $`f:D_1^3\mathrm{}D_m^3S^4`$ such that: (a)$`f(D_i^3)`$ coincides with $`K_i`$. (b)The singular point set $`X`$ consists of double points. (c)For each connected component $`X_i`$ of $`X`$, $`f^1(X_i)`$ is diffeomorphic to the two 2-discs. (d)Put $`\{f^1(X_i)\}`$=$`PQ`$. One of $`PQ`$ is included in the boundary of $`D_i^3`$ and another of $`PQ`$ is included in the interior of $`D_j^3`$ for integers $`i,j`$. ( We do not assume $`ij`$ nor $`i=j`$. )
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# Elliptic Curves of Twin-Primes Over Gauss Field and Diophantine Equations MR(1991) Subject Classification: 11R58; 11R27; 14H05; 11A55 Project Supported by the NNSFC (No. 19771052)
Let $`p,q`$ be twin prime numbers with $`qp=2`$ . Consider the elliptic curves
$$E=E_\sigma :y^2=x(x+\sigma p)(x+\sigma q).(\sigma =\pm 1)(1)$$
$`E=E_\sigma `$ is also denoted as $`E_+`$ or $`E_{}`$ when $`\sigma =+1`$ or $`1`$. Here the Mordell-Weil group and the rank of the elliptic curve $`E`$ over the Gauss field $`K=𝐐(\sqrt{}1)`$ (and over the rational field $`𝐐`$) will be determined in several cases; and results on solutions of related Diophantine equations and simultaneous Pellian equations will be given. The arithmetic constructs over $`𝐐`$ of the elliptic curve $`E`$ have been studied in , the Selmer groups are determined, results on Mordell-Weil group, rank, Shafarevich-Tate group, and torsion subgroups are also obtained. Some results on torsion subgroups in will be used here to determine $`E(K)`$.
Similarly to (1), some other special types of elliptic curves were studied by A. Bremner, J. Cassels, R. Strocker, J. Top, B. Buhler, B. Gross and D. Zagier (see \[3-5\]), $`e.g.,y^2=x(x^2+p),`$ $`y^2=(x+p)(x^2+p^2),`$ and $`y^2=4x^328x+25.`$ The last elliptic curve has rank 3 and is famous in solving the Gauss conjecture.
For any number field $`K`$, the Mordell-Weil theorem asserts that the $`K`$rational points of $`E`$ form a finitely generated abelian group (the Mordell-Weil group): $`E(K)E(K)_{tors}𝐙^r,`$ where $`E(K)_{tors}`$ is the torsion subgroup, $`r=\mathrm{rank}E(K)`$ is the rank (see ).
Theorem 1. Let $`E=E_\sigma `$ be the elliptic curve in (1), $`K=𝐐(\sqrt{}1)`$ the Gauss field.
(a) If $`p5(\text{mod}8),`$ then
$$\text{rank}E(K)=0,E(K)𝐙/2𝐙𝐙/2𝐙.$$
(b) If $`p3(\text{mod}8),`$ $`q=a^2+b^2,`$ $`(a+ϵ)^2+(b+\delta )^2=c^2`$ ($`ϵ,\delta =\pm 1,a,b,c𝐙),`$ then
$$\text{rank}E(K)=1,E(K)𝐙/2𝐙𝐙/2𝐙𝐙.$$
(c) $`\text{rank}E(K)3.`$
(d) If $`p,q`$ are as in (b), then
$$\text{rank}E_+(𝐐)=1,E_+(𝐐)𝐙/2𝐙𝐙/2𝐙𝐙.$$
In , Euler considered the Diophantine equation
$$\{\begin{array}{c}X^2+MY^2=S^2\\ X^2+NY^2=T^2\end{array}$$
for integers $`MN`$, and studied the problem of classification of the couples $`(M,N)`$ such that the above equation has non-trivial solutions, which is the famous“Euler Concordant Form problem”. E. Bell, T. Ono, K. Ono studied this problem further in \[8-10\]. For twin prime numbers $`p,q`$ we consider the following two Diophantine equations similar to the Euler’s :
$$\{\begin{array}{c}X^2pY^2=S^2\hfill \\ X^2qY^2=T^2\hfill \end{array}(I)$$
$$\{\begin{array}{c}X^2pY^2=2S^2\hfill \\ X^2qY^2=2T^2\hfill \end{array}(II)$$
If $`(X,Y,S,T)`$ is a primary solution of (I) (or (II)) (i.e., an integer solution with $`(X,Y)=1`$ and $`XY0`$), then by we know that $`(X^2/Y^2,XST/Y^3)`$ (or $`(X^2/Y^2,2XST/Y^3)`$) is a $`𝐐`$rational point of elliptic curve $`E_+`$ in (1). This fact leads us to the following theorem.
Theorem 2. If $`p5(\text{mod}8),`$ then equations (I) and (II) have no primary solution. (Furthermore, the only solution of them is $`(0,0,0,0)`$ )
For elliptic curve $`E=E_+`$ in (1) and the Gauss field $`K=𝐐(\sqrt{}1)`$, define
$$E(𝐐)^{}=\{(x,y)E(𝐐)|x<0\},$$
$$n_{2K}^{}=\frac{1}{2}\mathrm{\#}(2E(K)E(𝐐)^{}),$$
and let $`n(I)`$ denote the cardinal of the set of positive primary solutions of equation (I).
Theorem 3 (3.1) If $`n_{2K}^{}0,`$ then equation (I) has a primary solution.
(3.2) We have $`n_{2K}^{}n(I).`$ In particular, if $`n_{2K}^{}=\mathrm{}`$ then equation (I) has infinitely many primary solutions.
Finally, consider the simultaneous Pellian equations
$$\{\begin{array}{c}x^2py^2=\sigma \hfill \\ z^2qy^2=\sigma (\sigma =\pm 1)\hfill \end{array}(III)$$
This kind of equations has been studied in \[11-12\] (but $`\sigma =1`$ and $`\{p,q\}=\{m,n\}`$ are general integers). In particular, Rickert in obtained: if $`m=2,n=3`$ then the equation has no nontrivial integer solution (Nontrivial means $`y0`$). We find that equation (III) has an interesting relation with elliptic curve $`E`$ in (1): a nontrivial integer solution $`(a,b,c)`$ of (III) gives a $`𝐐`$rational point $`(1/b^2,ac/b^3)`$ of $`E`$. We obtain a general result on equation (III).
Theorem 4. (a) If $`p5(\text{mod}8),`$ then equation (III) has no nontrivial integer solution.
(b) If $`p3,5(\text{mod}8),`$ then equation (III) with $`\sigma =1`$ has no nontrivial integer solution.
References
\] QIU Derong and ZHANG Xianke, Mordell-Weil Groups and Selmer Groups of Two Types of Elliptic Curves, to appear
\] QIU Derong and ZHANG Xianke, Explicit determination of torsion subgroups of elliptic curves over multi-quadratic fields, Advances in Math., 28(1999), 475-478.
\] A. Bremner and J. W. S. Cassels, On the equation $`y^2=x(x^2+p),`$ Math. Comp. 42(1984), 257-264.
\] R. J. Stroeker and J. Top, On the equation $`y^2=(x+p)(x^2+p^2),`$ , Rocky Mountain J. of Math. 24(1994), 1135-1161.
\] J. P. Buhler, B. H. Gross and D. B. Zagier, On the conjecture of Birch and Swinnerton-Dyer for an elliptic curve of rank 3, Math. Comp. 44(1985), 473-481.
\] J. H. Silverman, The Arithmetic of Elliptic Curves, GTM 106, Springer-Verlag, 1986.
\] L. Euler, De binis formulis speciei $`xx+myy`$ et $`xx+nyy`$ inter se concordibus et disconcordibus, Opera Omnia Series I, 5(1780), 48-60, Leipzig-Berlin-Z$`\ddot{u}`$rich, 1944.
\] E. T. Bell, The problems of congruent numbers and concordant forms, Proc. Amer. Acad. Sci., 33(1947), 326-328.
\] T. Ono, Variations on a theme of Euler, Plenum, New York, 1994.
\] K. Ono, Euler’s condordant forms, Acta Arithmetica, LXXVIII. 2(1996), 101-123.
\] H. P. Schlickewei, The number of subspaces occurring in the p-adic subspace theorem in Diophantine approximation, J, Reine Angew. Math., 406(1990), 44-108.
\] J. Rickert, Simultaneous rational approximations and related Diophantine equations, Math. Proc. Cambridge Philos. Soc., 113(1993), 461-472.
TSINGHUA UNIVERITY
DEPARTMENT OF MATHEMATICAL SCIENCES
BEIJING 100084, P. R. CHINA
E-mail: xianke@tsinghua.edu.cn
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# Self-Quenched Dynamics
## 1 Introduction
The motion of random walkers in a random environment is one of the basic problems in the physics of disordered systems Alexander ; Sinai ; Fisher . It is known that the effect of the environment on the walker results in anomalous diffusion in some cases and logarithmically slow diffusion in others Alexander ; Sinai . Apart from their intrinsic interest, these simple models also find applications in several physical processes, such as the diffusion of electrons in a disordered medium Bernasconi or glassy activated dynamics Fisher . In particular, the Sinai model Sinai has been extensively studied from this point of view. It is known that for this model, the walker becomes logarithmically slow, moving as $`R(t)\mathrm{log}^2(t)`$ where $`R(t)`$ is the average distance at time $`t`$ of a walker from the launching point. Apart from this, two-time aging dynamics studied in this model also provide close analogies to glassy phenomenology Fisher .
In another framework, a “trapping” model (monthus and references therein) was introduced to provide a simple example of the glass transition. A single walker explores a random landscape of energy traps with e.g. an exponential distribution, from which it can escape through activated hops. At low temperatures, it was shown that the system cannot reach a steady state due to the unfavorable competition between the depth of the visited traps and the time needed to escape from them. The corresponding slow dynamics, and disappearance of an equilibrium state was interpreted as a glass transition. In this mean-field model, introducing an interaction between the walker and the random energy landscape was shown to not change the results monthus . However, the fact that the trapping time distribution is a power-law provides a natural (statistical) history dependence. (The dynamics is controlled by the deepest energy well visited so far). In the model we introduce here, the walker modifies its environment; aging and slow dynamics result purely because of this interaction.
Active walker models where the walker and the environment mutually affect each other have been studied earlier in different contexts. One such model — the Eulerian Walkers Model (EWM) — has been studied Eulwalk within the framework of self-organized criticality (SOC) Bak . The “landscape” here is defined by an arrow at each site. The walker follows the direction of the arrow at its site after which the direction of the arrow is changed according to some fixed rules. Besides many correspondences between the EWM and the Abelian Sandpile Model of SOC ddhar , it was also shown that the motion of the walker was sub-diffusive in two dimensions, i.e. $`R(t)t^{1/3}`$. In one dimension, $`Rt^{1/2}`$ due to a very simple organization of the landscape under the rules. For $`d>2`$ it was argued that the walker is diffusive. Models of mutually interacting walkers and landscape have also been studied extensively in the context of pattern formation and biophysical applications, where the emphasis is on the patterning of the medium under the influence of multiple walkers activ\_w ; Czirok .
The Self-Avoiding Walk (SAW) deGennes , where the walker is obliged to avoid its former path, and its modifications, can also be regarded as active walks in the sense that the path of the walker in these models is influenced by the trace it has made. For later purposes it is useful to mention here the so called “True” Self-Avoiding Walk (TSAW) where the walker’s probability to go to an already visited site is a strongly decaying function of the number of visits to that particular site TSAWorig .
The aim of this paper is to study a simple model of an active walker which exhibits logarithmically slow dynamics entirely due to the local interaction of the walker with a self-created random environment – a kind of self-organized trapping or self-burying effect. Besides the above mentioned connections with the physics of glasses, the model in one dimension may also be regarded as a very simplified version of a recently introduced model for the slow dynamics of sheared granular media, where the wandering of shear bands and the related restructuring of the material was shown to lead to extremely slow relaxation processes and inhomogeneous aging TKKR .
## 2 Definition of the SQW model
Our model is defined as follows. A walker can move on a hyper-cubic lattice in a $`d`$ dimensional space with linear size $`L`$ and $`N=L^d`$ number of sites. Periodic boundary conditions are imposed. A variable $`s_i`$ chosen from a uniform random distribution $`[0,1]`$ is initially assigned to every site $`i`$ on the lattice. We call the site with the walker the active site.
The only permitted elementary moves for the walker are to the nearest neighbours. At every time step a new random number $`s_i(t)`$ (uniformly sampled between $`0`$ and $`1`$) is assigned to the active site $`i`$. If this random number is larger than the value of the $`s`$ variable of all the nearest neighbours, then the same site remains active. Otherwise, the activity moves to the neighbouring site with the largest value of $`s`$ and the same procedure is repeated. We have chosen here a uniform distribution for $`s`$, however, this specific form can be shown to play no role in the time evolution of the active site. In the following we will refer to this model as a Self-Quenched Walk (SQW).
As time goes on, the average value of $`s`$ decreases and as a result, the probability that the activity moves to one of the neighbours decreases. The definition of the model is thus quite inconvenient in terms of numerical simulation since intervals when nothing happens grow longer and longer with time. However, it is easy to circumvent this difficulty. Let $`\sigma `$ be $`s_j`$, the largest $`s`$ value amongst the neighboring sites of the active one, $`i`$. For the activity to move to $`j`$, the $`s_i`$ value has to be smaller than $`\sigma `$ which is an event with probability $`\sigma `$. Thus the waiting time before a move is a Poisson process with a characteristic time $`\tau =1/\sigma `$. Once $`s_i<\sigma `$, the evolution of the activity is deterministic. Moreover, the distribution of $`s_i`$ when the activity moves to $`j`$ is uniform between 0 and $`\sigma `$. Therefore, we can directly reproduce the evolution of the model in terms of the number of moves $`n`$ rather than in time $`t`$. It turns out that the variable $`n`$ is also more convenient for the analytical treatment.
Similarly, since the $`s`$ values quickly evolve towards 0, as $`t`$ increases, it is more convenient to use an equivalent parameterization introducing $`r\mathrm{log}(s)`$. The uniform distribution of $`s`$ between 0 and $`\sigma `$ implies that $`r`$ is distributed with a density $`e^{\rho r}`$ for $`r>\rho =\mathrm{log}(\sigma )`$. Alternatively, we note that $`r\rho `$ is a random variable exponentially distributed from $`0`$ to $`\mathrm{}`$. This reformulation allows simulations to be carried out over practically unlimited times without loosing any accuracy. Moreover, as we will see below, $`r`$ is the appropriate scale for providing an accurate description of the long time regime.
Let us thus consider the motion of the activity (the walker) as a function of the number of moves $`n`$, in a potential $`V(i,n)`$ \[where the value of the potential $`V(i,n)=r_i(n)`$\]. The walker moves to the neighbouring site with the smaller value of $`r`$ \[since $`r\mathrm{log}(s)`$\] after having changed the value of $`V`$ on the site it was on. The above statement can be made more quantitative in the following two coupled equations of evolution for the walker and the $`r`$-landscape:
$`{\displaystyle \frac{d𝐗}{dn}}`$ $`=`$ $`V(𝐗(n),n)+\eta (n)`$ (1)
$`{\displaystyle \frac{V(𝐱,n)}{n}}`$ $`=`$ $`\lambda \delta ^d(𝐱𝐗(n))`$ (2)
where $`\eta (n)=0`$ and $`\eta (n)\eta (n^{})=C\delta (nn^{})`$. Here $`𝐗(n)`$ is the position of the walker at ’time’ $`n`$.
Eq. 1 quantifies the rules of the SQW in any dimension. When the walker is on a slope, it moves down towards the valley. When it is in a flat region, it moves to any of the nearest neighbours with equal probability (this is the reason for the uncorrelated noise term). Eq. 2 accounts for the increase of the $`r`$-landscape at the position of the walker. In this simplified continuum description, we have neglected the randomness in the distribution of the local increments in $`r`$, and only retained the average value $`\lambda `$. However, as we shall see below, all the essential features of the problem are contained in the equations.
## 3 Correspondence with the TSAW model
It turns out that the Langevin equations (Eqs. 1 and 2) of the SQW in terms of $`n`$ and $`r`$ are the same as that of the TSAW as a function of time and position TSAWorig . The definition of the TSAW model is the following: The walk takes place on a $`d`$-dimensional hypercubic lattice. At any step the traveller may move to any of the $`2d`$ nearest neighbours of the lattice site he is at. The probability of stepping to site $`i`$ depends on the number of times $`n_i`$ this site has already been visited and is given by
$$p_j=\mathrm{exp}(gn_j)\left[\underset{j=1}{\overset{2d}{}}\mathrm{exp}(gn_j)\right]^1,$$
(3)
where the sum runs over all $`2d`$ nearest neighbours of the current position of the walker, and $`g`$ is a positive parameter which measures the intensity with which the walk avoids itself. Note that the sum over $`i`$ of $`p_i`$ is equal to $`1`$, meaning that the traveller never stays at the same site. This is similar to the SQW when time is incremented in units of $`n`$.
In the SQW the value of $`s`$ decreases exponentially on average since each time the walker visits the site we multiply the $`s`$ value by a random number taken from a uniform distribution between $`0`$ and $`1`$. The value of $`r=\mathrm{log}(s)`$ thus increases linearly with the number of times the traveller has visited this site.
Further, in the SQW, whenever the walker can move, it goes deterministically to the neighbouring site with the higher value of the potential. This is realized by the $`g\mathrm{}`$ limit of the TSAW model.
The above mapping thus ensures that in the continuum limit the two models are governed by the same Langevin equations. The role of the preassigned probability (Eqn. 3) in the TSAW is played by the self-organized evolution of the walker and the potential.
The TSAW has been studied exhaustively by means of numerical simulations JCAdA ; BerPie , Flory theory FamDaoud , scaling analysis Pietronero and later by exact calculations BToth .
The critical dimension of the TSAW problem is $`d_c=2`$ above which the mean field solution applies and the traveller’s asymptotic behavior is not influenced by the interaction with its former path and performs basically a Brownian motion. Below two dimensions the trace of the walker is a fractional Brownian motion. The root mean square distance from the origin increases as
$$|𝐗(n)𝐗(n^{})|^2^{1/2}|nn^{}|^\nu $$
(4)
with
$$\nu =\{\begin{array}{cc}2/(d+2)\hfill & \text{for }d2\hfill \\ 1/2\hfill & \text{for }d>2\hfill \end{array}$$
(5)
Thus in one dimension the walker is super-diffusive with a Hurst exponent of $`\nu _{d=1}=2/3`$ which is due to the repulsive interaction with its former path. This asymptotic behavior is numerically verified for the SQW in Fig. 1.
We can define other related exponents through the scaling relations: $`XbX`$, $`nb^zn`$ and $`Vb^\chi V`$, where $`\chi `$ is the so called roughness exponent of the $`r`$ landscape and $`z`$ is the dynamic exponent. Equations 1 and 2 then predict the following values for these exponents:
$$\begin{array}{cccc}\chi =1/2\hfill & z=3/2\hfill & \text{ in }\hfill & d=1\hfill \\ \chi =0\hfill & z=d\hfill & \text{ for }\hfill & d2.\hfill \end{array}$$
(6)
We have confirmed the value of $`\chi `$ and $`z`$ by measuring the width of the $`r`$ landscape in our model.
## 4 Roughness of the potential
The width of a self-affine interface is defined as the root mean square fluctuation of the interface from its mean value. In a number of growth models, this width obeys the Family-Vicsek scaling FS with a dynamic exponent $`z`$ such that the overall roughness follows
$$w(n)L^\chi \phi \left(\frac{n}{L^z}\right)$$
(7)
Fig. 2 shows our numerical determination of the width of the $`r`$ landscape for five different system sizes in one and three dimensions. The collapse with the above mentioned value of the exponents indicates that the $`r`$ landscape is self-affine in one dimension. However, the growth exponent $`\beta `$ describing the roughening of the landscape fluctuations for early times ($`n<n^{}L^z`$) as $`wn^\beta `$ does not obey the Family-Vicsek scaling $`\beta =\chi /z`$. It is instead given by
$$\beta =(\chi +1/2)/(1+\chi ),$$
(8)
a formula typical for growth models with extremal dynamics supri . Thus our model, though it does not contain a global extremum criterion, belongs to the class of extremal growth models.
The reason why the interface growth (or the $`r`$ landscape evolution) is similar to extremal dynamics is the following. The activity is a ’random’ walk trapped by the maxima of the $`r`$-landscape. Before the activity can escape, the landscape has to be filled. In order to escape from a region of extent $`\stackrel{~}{l}`$, the number of moves to be made is of the order of the size of the valley, $`\stackrel{~}{l}^d`$, times its typical depth $`\stackrel{~}{l}^\chi `$. The growth is pinned everywhere except in the immediate vicinity of the walker. The walker itself is, however, in a hierarchically ordered valley structure with maxima of increasing heights. Thus the interface progresses jerkily just as in other extremal growth models. The relation $`z=d+\chi `$ predicted by Eqn. 2 is also known to occur in various extremal growth models Maslov ; Tang .
The correspondence with extremal dynamics does not hold above $`d_c=2`$. This is because above two dimensions, the walker is no longer trapped by surrounding maxima. It can also find its way around them instead of over them. As a result the interface is no longer rough and $`\chi =0`$. However, though the width saturates to a system size-independent value, the exponent $`\beta `$ is non-zero for the following trivial reason. In the random initial state, in $`n`$ steps, a number $`n`$ sites grow by, say, a uniform amount $`h`$. Therefore the width $`w`$ scales trivially as $`whn^{1/2}L^{1/2}`$. The exponent $`\beta `$ is hence equal to $`1/2`$ for all $`d2`$.
## 5 Real time behaviour
We now turn to the behaviour of the walker in real time. In order to understand this, we first note that the mean value of $`r`$ in the steady state increases linearly with the number of moves. Hence the mean increase of $`r`$ per site in the steady state is $`1/N`$, where $`N=L^d`$ is the number of sites in the system.
We now compute $`r^{}(n)`$ (the $`r`$ value of the active site as a function of $`n`$) and observe two regimes that we describe by the scaling assumption:
$$r^{}(n)n^\alpha \psi \left(\frac{n}{L^z}\right)$$
(9)
with $`\psi (a)a^0`$ for $`a1`$ and $`\psi (a)a^{1\alpha }`$ for $`a1`$. Therefore, for long times $`r^{}(n)nL^{z(1\alpha )}`$. However, in the long time regime, $`r^{}`$ has to increase at the same rate as the mean velocity of the front $`r`$, and hence, for $`nL^z`$, $`r^{}(n)=n/N=nL^d`$. This imposes
$$z(1\alpha )=d\mathrm{or}\alpha =1\frac{d}{z}=\frac{\chi }{d+\chi }$$
(10)
i.e. $`\alpha =1/3`$ for $`d=1`$ and $`0`$ for $`d2`$ consistently with our numerics as shown in Fig. 3.
Since, in the late stage regime, $`r^{}`$ increases as $`n/N`$, $`s^{}`$ decreases as $`e^{n/N}`$. Therefore the expectation value of the real time lapse between two consecutive moves $`t(n)t(n1)`$ is $`1/s^{}`$ or
$$t(n)t(n1)e^{n/N}$$
(11)
$$t(n)=\frac{e^{(n+1)/N}1}{e^{1/N}1}Ne^{n/N}$$
(12)
This law refers in fact only to the mean value of $`t`$. However, the distribution of each increment being exponentially distributed (a Poisson process), the central-limit theorem applies, and the relative standard deviation of $`t`$ with respect to its mean value vanishes. Figure 4 shows the numerically determined time as a function of the number of moves $`n`$.
The average value of $`s`$ as a function of time is thus
$$s\mathrm{exp}(n/N)\mathrm{exp}(\mathrm{log}(t/N))N/t$$
(13)
for long times.
Therefore in real time the walker is logarithmically slow with $`R(t)\mathrm{log}(t)^{2/3}`$ in one dimension and $`R(t)\mathrm{log}(t)^{1/2}`$ in higher dimensions. The logarithmic dependence of the RMS distance of the walker is just the consequence of Eq. 13 as a result of which, the probability of making a jump to a neighbouring site decreases as $`1/t`$. It is thus valid in any dimension. The value of the exponent of the log in one dimension, is however a non-trivial consequence of the coupling between the walker and the medium which induces long-range memory effects.
## 6 Conclusion
In summary, we have introduced and studied a simple model of a walker interacting with its environment. By choosing the correct measures for describing the time and the potential, we could map the SQW problem to the TSAW model and thus use the known exact results in the latter case, to describe the motion of the SQW walker. In addition, we have also studied the emerging landscape. Though the rules for the SQW are entirely local, a relationship with the so called extremal models could be established. The critical dimension is $`d_c=2`$ below which the potential landscape gets self-affine and the walker super-diffusive in terms of moves. In real time the walker is logarithmically slow in any dimension.
Acknowledgment: This work was partially supported by OTKA T029985.
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# Untitled Document
EFFICIENCY OF NEUTRON DETECTION OF
SUPERHEATED DROPS OF FREON-22
Mala Das, B. Roy, B. K. Chatterjee and S. C. Roy
Department of Physics, Bose Institute
93 / 1 A. P. C. Road, Calcutta 700009, India
## Abstract
Neutron detection efficiency of superheated drops of Freon-22 for neutrons obtained from a 3 curie Am-Be neutron source has been reported in this paper. Although Freon-22 having lower boiling point than many other similar liquids (e.g.Freon-12, Freon-114, Isobutane ) is expected to be more sensitive to neutrons ,it has not been reported so far and therefore this paper constitutes the first report on the subject. Neutron detection efficiency of both Freon-22 and Freon-12 have been determined from the measured nucleation rate using the volumetric method developed in our laboatory. The result shows that neutron detection efficiency of Freon-22 for the neutron energy spectrum obtained from an Am-Be source, is almost double, while the life time is 58.6% smaller than that of Freon-12, for a particular neutron flux of that source.
1. INTRODUCTION
A liquid maintained at the same state above its boiling point is said to be superheated. It is a metastable state of the liquid and can be nucleated to form vapour by the deposition of small energy by ions, charged particles or by any heterogeneous nucleation sites such as gas pockets, impurities etc. The superheated drops,suspended in gel can be used to detect neutrons through the nucleation induced by the recoil nuclei in the medium . The recoil nuclei are produced by collision of neutrons with the nuclei constituting the superheated drops. The application of superheated drop detector (SDD) in neutron dosimetry has already been established (Apfel et al., 1984, 1989; Ing, 1986) and several other potential applications of SDD in neutron research has been discussed (Apfel, 1979a, 1979b, 1981; Chakraborty et al., 1990 ). Apfel (1992) has developed and characterised a passive superheated drop dosemeter using a volumetric technique for neutron monitoring of personnel and in accelerartor applications. These type of dosimeters are also commercially available from Apfel Enterprises Inc.,USA. Practical application of the SDD demands that the sample must be reasonably stable against spontaneous nucleation due to background radiation and other environmental effects, and at the same time as much as possible sensitive to neutrons with energy spectrum of interest. Among the different important features of the detector, the systematic quantitative evaluation of the sensitivity of SDD has been studied for some liquids (e.g. Freon-114, Freon-12, Isobutane, Freon 142B) (Roy et al., 1987) and the response function was reported (Lo et al., 1988). This paper is an attempt to investigate more sensitive liquid for neutron detection. Of the liquids investigated so far, Freon-12 is considered to be the most sensitive liquid for neutron detection. Since the boiling point of Freon-22 is much lower than Freon-12, we expect Freon-22 to be more sensitive to neutrons than Freon-12. To the best knowledge of the authors, no such investigation has been reported with Freon-22. An accurate method of determining the nucleation rate of SDD has been developed by Roy, et al., (1997b) using a relative manometer. Some other studies on the neutron detection sensitivity and detector response were made by Ing and Birnboim (1984), Ing (1986), Ipe et al., (1988) and Biro et al., (1990). Nath et al.,(1993) measured the neutron dose equivalent to patients undergoing high energy x-ray and electron radiotherapy beams using a SDD device. They employed a passive method for measuring the total volume of neutron induced bubbles by displacing an eqivalent volume of gel into a graduated pipette. The method for the determination of efficiency of detection of neutrons by vapour nucleation of superheated drops using volumetric method has been developed in our laboratory by Roy et al., (1997a). We in this work, present the efficiency of neutron detection, maximum nucleation rate and life-time of SDD of Freon-22 irradiated by neutrons from an Am-Be source for a particular neutron flux and compare it with those of Freon-12. Some of the physical parameters of SDD of Freon-12 and Freon-22 are presented in table-1. Comparison is made with Freon-12, since Freon-12 is the most well studied liquid in neutron detection. The paper has been organised to present a brief outline of the method of bubble formation, description of volumetric method and measurement, results and discussion.
2. THEORY OF BUBBLE FORMATION
The free energy required to form a spherical vapour bubble of radius r in a liquid is given by (Roy et al., 1987)
$$G=4\pi r^2\gamma \left(T\right)\frac{4}{3}\pi r^3\left(p_vp_o\right)$$
(1)
where $`\gamma (T)`$ is the liquid-vapour interfacial tension,$`P_v`$ is vapour pressure of the superheated liquid and $`P_o`$ is the ambient pressure. The difference $`p_v`$ -$`p_o`$ is called the degree of superheat of a given liquid. One can see from equation (1) that G is maximum at
$$r=2\gamma \left(T\right)/\left(p_vp_o\right)=r_c$$
(2)
where $`r_c`$ is called the critical radius. When a bubble grows to the size of the critical radius it becomes thermodynamically unstable and grows very fast till the entire liquid droplet vaporises. The minimum amount of energy (W) needed to form a vapour bubble of critical size $`r_c`$ as given by Gibbs (1875) from reversible thermodynamics is
$$W=16\pi \gamma {}_{}{}^{3}\left(T\right)/3(p_vp_o)^2$$
(3)
which is supplied by the energy deposition dE/dx being the energy deposited per unit distance travelled by the nuclei in the liquid by the recoil nucleus in a path length of 2$`r_c`$ inside the droplet.
3. PRINCIPLE OF THE VOLUMETRIC METHOD
The present method utilizes the superheated drops suspended in a dust free viscous elastic gel. The excess pressure required to form a vapour bubble of diameter 1mm inside this gel matrix is usually less than 1mm of mercury (as observed in a separate experiment). Hence the volume of the bubble trapped inside the gel and that at atmospheric pressure are almost equal. This is why the volume change upon nucleation would be the same whether the bubble formed is trapped inside the gel or liberated from it. Upon nucleation the increase in volume of the droplet would displace the trapped air inside a vial containing sample. In the present volumetric method described in detail in the next section, the measurement of rate of change of volume has been made using such an air displacement system. The superheated drops suspended in the gel do not touch each other physically and they nucleate in a random manner, independent of each other. Under such conditions it can be shown that the number of the drops and hence the volume of superheated liquid would decay exponentially with time. The rate of change of nucleated volume varies exponentially with a time constant $`\tau `$ (lifetime). The life time of these droplets in presence of a neutron flux and the efficiency of neutron detection has been studied by measuring the volume of vapour formed upon nucleation.
If neutrons of flux $`\psi `$ are incident on superheated drops of volume V, liquid density $`\rho _L`$ and molecular weight M the vaporization rate is given by
$$\frac{dV}{dt}=V\psi \frac{N_A\rho _L}{M}\eta dn_i\sigma _i$$
(4)
| where | N<sub>A</sub> | = | Avogadro Number |
| --- | --- | --- | --- |
| | d | = | average droplet volume |
| | n<sub>i</sub> | = | number of nuclei of the ith element of the molecule whose |
| | | | neutron nucleus elastic scattering cross section is $`\sigma _i`$ |
| | $`\eta `$ | = | efficiency of neutron detection. |
Due to nucleation by neutrons, the change in position (h) of the water column along a horizontal glass tube which is connected to the sample vial (arrangement is given in next section), could be measured with respect to time (t). The present set up is made completely free from any leakage. So the rate of increase of the volume of vapour during nucleation should be equivalent to the rate of decrease of the volume of the superheated liquid. So, we have this equation
$$\rho _VA\frac{dh}{dt}=\rho _L\frac{dV}{dt}=\rho _LV\psi \frac{N_A\rho _L}{M}\eta dn_i\sigma _i$$
(5)
| where | A | = cross section of the horizontal tube |
| --- | --- | --- |
| | $`\rho _V`$ | = density of Freon vapour |
| | V | = volume of superheated liquid at any instant of time t. |
Integrating and solving equations (3) one can obtain
$$\frac{hA}{m}=a\left[1\mathrm{exp}b\left(tt_o\right)\right]$$
(6)
where a = $`\frac{\rho _LV_o}{\rho _Vm}`$ and
$$b=\frac{1}{\tau }=\psi \frac{N_A\rho _L}{M}\eta dn_i\sigma _i$$
(7)
$`\frac{hA}{m}`$ is the volume of accumulated vapour in time t per unit mass of the sample containing Superheated drops and gel, m being the mass of total sample (gel + superheated drops). Therefore the efficiency of neutron detection $`\eta `$ is given by
$$\eta =\frac{bM}{\psi N_A\rho _Ldn_i\sigma _i}$$
(8)
V<sub>o</sub> = initial volume of the Freon drops. t<sub>o</sub> = initial time at which the experiment has been started. $`\tau `$ = life-time of SDD in presence of source. and the product ab gives the maximum nucleation rate. The equation (6) has been scaled to per unit mass of the sample for standardization in making comparison of two samples. Fitting equation (6) for different values of $`\frac{hA}{m}`$ and t, constants a, b are obtained. Knowing b, the life time $`\tau `$ (=1/b) in presence of neutron flux $`\psi `$ is obtained. The neutron detection efficiency $`\eta `$ can be found out from equation (8) for a known flux of neutrons
4. EXPERIMENT
The experiment was performed with 3Ci Am-Be neutron source, which has an energy distribution of neutrons with peak near 3.5 MeV. The source was placed inside a chamber through which neutrons coming in a fixed range of direction were used to irradiate the SDDs. Two separate sets of experiments were done with Freon-12 and Freon-22 respectively at about 30<sup>o</sup>C and at atmospheric pressure. The superheated drops are usually suspended in an immiscible gel. The gel that we used here, was a homogeneous mixture of some ultrasonic gel and glycerol in suitable proportion. The detail of the preparation of the sample was given elsewhere (Roy et al., 1997b). The experimental apparatus consists of long glass tube of cross section 0.1573 sq.cm, placed horizontally on a graduated platform. The tube contained a coloured water column as an indicator of the volume of the vapour formed on nucleation. One end of the glass tube was connected to the glass vial containing the sample, by means of rubber tube. The glass vial and the horizontal tube were at the same height. So the pressure inside and outside the tube were equal, e.g. both were at the same atmospheric pressure. Therefore the displacement of the water column of 1 cm length due to nucleation would be directly related to the volume of the vapour formed. These displacements of the water column along the glass tube were measured as a function of time. The gamma sensitivity of sample of Freon-22 was tested with <sup>241</sup>Am (59.54keV), <sup>60</sup>Co (1170keV,1332kev), <sup>137</sup>Cs (662keV) gamma sources. The sample was placed either just in contact or at a close distance to the source. The sensitivity of the sample for 4.43MeV gamma ray present in <sup>241</sup>Am-Be source has also been tested by substantially reducing the neutron dose.
5. RESULTS
The chemical formula and other physical parameters including critical radius (r<sub>c</sub> ) and minimum energy required for nucleation (W) as calculated using equations (2) and (3) respectively, for Freon-12 and Freon-22 have been listed in Table-1. The experimentally observed variation of volume of vapour formed, $`\frac{hA}{m}`$ (per unit mass of the sample) during nucleation as a function of time for Freon-12 and Freon-22 are shown in Fig.1 and Fig.2 respectively. The measured data on rate of nucleation and life-time of superheated drops of Freon-12 and Freon-22 in presence of neutrons from a 3 Ci Am-Be source and the efficiency of neutron detection are given in table-2. The samples irradiated by gamma sources did not show any noticeble nucleation at the experimental temperature.
6. DISCUSSION
From the tabulated results it is seen that the degree of superheat (defined by the difference of vapour pressure of the superheated liquid $`p_v`$ and the ambient pressure $`p_o`$ ) attained for Freon-22 at room temperature (30<sup>o</sup>C) and at atmospheric pressure is about 38% larger than that of Freon-12. This indicates that in presence of neutrons the life time of Freon-22 should be smaller than that of Freon-12. Our experimental results show that the maximum nucleation rate of Freon-22 is about 38% larger and the life time of Freon-22 is about 58.6% smaller than that of Freon-12 for a fixed neutron flux from an Am-Be source. As can be seen from the chemical formula that both of Freon-12 and Freon-22 contain carbon, chlorine and fluorine while Freon-22 contains one hydrogen replacing one chlorine in Freon-12. The neutron-nucleus elastic scattering cross-section for Freon-12 is 1.89 barns while that of Freon-22 is 1.69 barns at neutron energy 3.5 MeV. Although the probability of interaction of a neutron with nuclei of Freon-12 is larger than Freon-22, our experimental results (table-2) show that the present prepared sample of Freon-22 is about twice as efficient to detect neutrons than Freon-12. This is due to the fact that as seen from equation (3), the energy required for nucleation (W) decreases with degree of superheat and as is evident fron table-1, the saturation vapour pressure of Freon-22 at 30<sup>o</sup>C is higher than that of Freon-12, therefore at this temperature Freon-22 will attain a higher degree of superheat which means a smaller amount of energy deposition is required for nucleation than that in Freon-12. So, although less number of recoil nuclei are available in Freon-22 from neutron nuclei elastic scattering, the percentage of nuclei capable of deposition of energy greater than W for Freon-22 must be larger than that in Freon-12. As a result the efficiency of neutron detection of SDD of Freon-22 is larger than that of Freon-12. From the experiment with gamma sources it is clear that the SDD based on Freon-22 is insensitive to gamma at the experimental temperature. The high efficiency of neutron detection and the insensitivity towards gamma makes Freon-22 a suitable superheated drop detector for neutron detection.
REFERENCES
Apfel R.E. (1979a) Detector and dosimeters for neutrons and other radiations. US Patent 4,143,274.
Apfel R.E. (1979b) The superheated drop detector. Nucl. Inst. Meth. 162, 603.
Apfel R.E. (1981) Photon-insensitive,thermal to fast neutron detector. Nucl. Inst. Meth. 179, 615.
Apfel R.E. (1992) Characterisation of new passive superheated drop (bubble) dosemeters. Rad. Prot. Dos. 44, 343.
Apfel R.E. and Lo Y.C. (1989) Practical neutron dosimetry with superheated drops. Health Phys. 56, 79.
Apfel R.E. and Roy S.C.(1984) Investigation on the applicability of superheated drop detector in neutron dosimetry. Nucl. Inst. Meth. 219, 582.
Biro T.,Kelemen A and Pavlicsek I. (1990) Acoustic Detection of neutrons by bubble detectors. Nucl. Tracks Radiat. Meas. 17, 587.
Chakraborty K.,Roy P.,Vaijapurkar S.G. and Roy S.C. (1990) Study on neutron spectrometer using superheated drop detector. Proc. of 7th National Conference on Particles and Tracks, Jodhpur pp 133.
Gibbs J.W. (1875) Translations of the Connecticut Academy III, p.108.
Ing H. (1986) The status of the bubble-damage polymer detector. Nuclear Tracks. 12, 49.
Ing H. and Birnboim H.C. (1984) A bubble damage polymer detector for neutrons. Nucl. Tracks and Radiat. Meas. 8, 285.
Ipe N.E., Busick D.D. and Pollock R.W. (1988) Factors affecting the response of the bubble detector BD-100 and a comparison of its response to CR-39. Rad. Prot. Dos. 23, 135.
Lo Y.C. and Apfel R.E. (1988) Prediction and experimental confirmation of the response function for neutron detection using superheated drops. Phys. Rev. A 38, 5260.
Nath R, Meigooni A.S., King C.R., Smolen S. and d’Errico F. (1993) Superheated drop detector for determination of neutron dose equivalent to patients undergoing high energy x-ray and electron radiotherapy. Medical Phys. 20, 781.
Roy B, Chatterjee B.K., Das Mala and Roy S.C. (1997a) Study on nucleating efficiency of superheated droplets by neutrons. Radiation Physics and Chemistry (accepted for publication).
Roy B., Chatterjee B.K. and Roy S.C.(1997b) An accurate method of measuring life time of superheated drops using differential manometer. Radiation Measurements (accepted for publication).
Roy S.C., Apfel R.E. and Lo Y.C. (1987) Superheated drop detector: A potential tool in neutron research. Nucl. Inst. Meth. A255, 199-206.
Table-1 : A comparison of the physical parameters of Freon-12 and Freon-22
| | Freon-12 | Freon-22 |
| --- | --- | --- |
| 1. ChemicalFormula | CCl<sub>2</sub>F<sub>2</sub> | CHClF<sub>2</sub> |
| 2. Molecular weight | 120.91 | 80.47 |
| 3. Boiling Point | -29.79 <sup>o</sup>C | -40.75 <sup>o</sup>C |
| 4. Surface tension ($`\gamma `$) | | |
| (at 30<sup>o</sup>C) dyn/cm | 9 | 8 |
| 5. Vapour Pressure (p<sub>v</sub>) | | |
| dyn / sq. cm. | 7.4556 $`\times `$ 10<sup>6</sup> | 1.14777 $`\times `$ 10<sup>7</sup> |
| 6. Density ($`\rho _L`$) gm / cc | 1.293 | 1.175 |
| 7. Degree of superheat | | |
| (p<sub>v</sub> \- p<sub>o</sub>) dyn / sq. cm | 6.441639$`\times `$ 10<sup>6</sup> | 1.0463739$`\times `$ 10<sup>7</sup> |
| 8. Critical radius | | |
| $`r_c=\frac{2\gamma (T)}{\left(p_vp_o\right)}`$ cm | 2.79$`\times `$ 10<sup>-6</sup> | 1.53$`\times `$ 10<sup>-6</sup> |
| 9. Minimum energy required (W) keV | | |
| to form a vapour bubble of size r<sub>c</sub> | 0.184 | 0.049 |
| (W = $`\frac{16\pi \gamma ^3(T)}{3(p_vp_o)^2}`$) | | |
Table - 2: Observed results on nucleation
(Neutron flux = 2.5374$`\times `$10<sup>7</sup> /cm<sup>2</sup>/s; Peak neutron energy = 3.5 MeV).
| | Freon-12 | Freon-22 |
| --- | --- | --- |
| 1.Initial Nucleation rate<sup>1</sup> (cm<sup>3</sup>/gm/s) | 3.0268 $`\times `$ 10<sup>-4</sup> | 4.9205 $`\times `$ 10<sup>-4</sup> |
| 2. Lifetime $`\tau `$ (s) | 3089.59 | 1279.13 |
| 3. Neutron detection efficiency<sup>2</sup> | 0.2825$`\%`$ | 0.5588$`\%`$ |
1 Initial nucleation rate is $`\left(\frac{1}{m}\right)`$$`\frac{dV}{dt}`$ at the initial time, where m is the mass of the sample and $`\frac{dV}{dt}`$ is the rate of volume change upon nucleation.
2 neutron detection efficiency is the percentage of neutrons causing nucleation.
FIGURE CAPTIONS
Fig. 1: Observed variation of volume of vapour formed as a function of time for Freon-12; $`\frac{hA}{m}`$ = vapour volume per unit mass of the sample.
Fig. 2: Observed variation of volume of vapour formed as a function of time for Freon-22; $`\frac{hA}{m}`$ = vapour volume per unit mass of the sample.
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# Enhanced Quantum Estimation via Purification
## I Introduction
An essential step of information processing, either classical or quantum, is the read-out of data. Reading a classical bit does not pose fundamental problems. A single measurement suffices to reveal the content $`0`$ or $`1`$ of a classical memory cell. However, already in the classical case this information will be lost eventually due to dissipative processes. In fact each classical dynamic computer memory requires permanent cycles of reading and refreshing in order to prevent this dissipation of data.
This situation changes considerably in the quantum domain. As soon as a classical apparatus registers a signature of a quantum system its original quantum state changes. This measurement process is effectively formulated by the von-Neumann rules . In the following we will accept this description of a quantum measurement rather than questioning its philosophical implications. Neither will we try to derive the von-Neumann process using quantum dynamics. For the present article we take this specific quantum feature of the measurement process for granted and consider its impact on quantum information processing.
As a consequence of the von-Neumann description it is impossible to unseal the complete content of a single quantum bit in form of the corresponding quantum state. This has lead to investigations whether the state of a quantum bit can be at least optimally estimated from a finite number of identical copies. The intriguing result of this research has been that optimal estimations cannot be realized by single measurements on each qubit separately. Rather one has to see the finite ensemble as one composite system which allows us to perform a generalized measurement .
However, the realization of such measurements is a difficult task. First, we might get the identical copies of our qubit sequentially in certain situations. Just imagine the situation of “quantum debugging”, where a quantum algorithm is run several times on a quantum computer. In this case, the sequential estimation of a specific qubit could be used as a test for the perfect performance of the algorithm. Second, the physical implementation of the generalized measurement remains an open problem so far.
Hence one has to ask how close single measurements can come the optimal ones , in particular, if we combine them with a strategy. That is, if we learn from each measurement step and adapt the next one appropriately. Such schemes have been discussed for pure qubits in Ref. . But pure qubits are rather ideal. As soon as they are connected to an environment they will start to decohere and gradually lose their original information content. Hence the original pure state will turn into a mixed quantum state depending on the decoherence model .
Is it possible to refresh such decohered qubits as in the classical case? One answer to this is quantum error correction . But it is also possible to purify a finite ensemble of qubits if we sacrifice some of them. This has been shown in Ref. : The remaining qubits after the purification step have indeed a higher single-qubit fidelity than the unpurified ones.
In the present paper we shall investigate if such a purification step also allows us to retrieve more information about the originally pure state of the qubits. Does it make sense to spend some of the decohered qubits for the purification step and to estimate the underlying pure state now from a smaller ensemble using single adaptive measurements? What will be the influence of the entanglement between the purified qubits which is an unavoidable consequence of the purification protocol presented in Ref. . These questions lie at the center of the present paper.
In Section II we recall the basic elements of single qubit purification needed for our purpose. We then combine this with an adaptive estimation strategy in Section III. Our results are discussed in Section IV and we conclude with Section V.
## II Purification of qubits
In this section we briefly review the single qubit purification proposed in Ref. and introduce our notation. We assume that Alice has a source of single qubits in the pure state
$$|1_\stackrel{}{n}\mathrm{cos}\frac{\mathrm{\Theta }}{2}|1+\mathrm{sin}\frac{\mathrm{\Theta }}{2}e^{i\mathrm{\Phi }}|0$$
(1)
defined by the basis states $`\{|0,|1\}`$ and the point $`\stackrel{}{n}=(\mathrm{sin}\mathrm{\Theta }\mathrm{cos}\mathrm{\Phi },\mathrm{sin}\mathrm{\Theta }\mathrm{sin}\mathrm{\Phi },\mathrm{cos}\mathrm{\Theta })`$ on the Bloch sphere. These qubits are sent to Bob via a depolarizing channel and hence Bob receives qubits described by a mixed state
$$\rho =c_1|1_\stackrel{}{n}1|+c_0|0_\stackrel{}{n}0|,$$
(2)
in which $`|1_\stackrel{}{n}`$ occurs with probability $`c_11/2`$ and the corresponding orthogonal state
$$|0_\stackrel{}{n}\mathrm{cos}\frac{\mathrm{\Theta }}{2}|0\mathrm{sin}\frac{\mathrm{\Theta }}{2}e^{i\mathrm{\Phi }}|1$$
(3)
with probability $`c_0`$, so that $`c_0+c_1=1`$.
If Alice sends an even number of $`N`$ qubits to Bob, the corresponding product state reads
$$\rho ^N\rho \rho \mathrm{}\rho .$$
(4)
Note that the qubits are assumed to be distinguishable.
The optimal purification protocol proposed in Ref. shows how one can distill $`MN`$ purer qubits from $`\rho ^N`$ by a measurement performed on the $`N`$ qubits. This measurement projects the state $`\rho ^N`$ on the output state
$$\rho _{out}=\rho _M\left(|\mathrm{\Psi }^{}\mathrm{\Psi }^{}|\right)^{(NM)/2}$$
(5)
with probability
$`p_M`$ $`=`$ $`\left[\left({\displaystyle \genfrac{}{}{0pt}{}{N}{\frac{(NM)}{2}}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{N}{\frac{(NM)}{2}1}}\right)\right]`$ (7)
$`\times \left[c_0c_1\right]^{(NM)/2}{\displaystyle \frac{c_1^{M+1}c_0^{M+1}}{c_1c_0}}`$
for values $`M=0,2,\mathrm{},N`$. Hence one obtains $`M`$ qubits in the entangled state
$`\rho _M`$ $`=`$ $`{\displaystyle \frac{c_1c_0}{c_1^{M+1}c_0^{M+1}}}(M+1)`$ (9)
$`\times {\displaystyle }{\displaystyle \frac{d\mathrm{\Omega }}{4\pi }}n(\theta )^M(|\mathrm{\Psi }(\theta ,\varphi )\mathrm{\Psi }(\theta ,\varphi )|)^M`$
with
$$n(\theta )=c_1\mathrm{cos}^2(\theta /2)+c_0\mathrm{sin}^2(\theta /2)$$
(10)
and the normalized state
$$|\mathrm{\Psi }(\theta ,\varphi )=\sqrt{c_1}\frac{\mathrm{cos}(\theta /2)}{\sqrt{n(\theta )}}|1_\stackrel{}{n}+\sqrt{c_0}\frac{\mathrm{sin}(\theta /2)}{\sqrt{n(\theta )}}e^{i\varphi }|0_\stackrel{}{n}.$$
(11)
The remaining $`NM`$ qubits are maximally entangled in pairs of Bell states $`|\mathrm{\Psi }^{}=\left(|1|0|0|1\right)/\sqrt{2}`$ and therefore carry no information about the original qubit.
The $`M`$ qubits in state $`\rho _M`$ are indeed purer than the original qubits described by $`\rho `$, Eq. (2). That is, the reduced density operator $`\mathrm{Tr}_{M1}(\rho _M)`$ obtained by tracing over $`M1`$ qubits lies closer to the pure state $`|1_\stackrel{}{n}`$ than $`\rho `$, which means that the single qubit fidelity
$`f_M`$ $`=`$ $`_\stackrel{}{n}1|\mathrm{Tr}_{M1}(\rho _M)|1_\stackrel{}{n}`$ (12)
$`=`$ $`{\displaystyle \frac{1}{M}}\left[{\displaystyle \frac{(M+1)c_1^{M+1}}{c_1^{M+1}c_0^{M+1}}}{\displaystyle \frac{c_1}{c_1c_0}}\right]`$ (13)
for $`M>0`$ fulfills $`f_Mc_1`$. Note that $`c_1=_\stackrel{}{n}1|\rho |1_\stackrel{}{n}`$ represents the fidelity of the original mixed state $`\rho `$, Eq. (2). For the case $`M=0`$ the output state Eq. (5) consists only of $`|\mathrm{\Psi }^{}`$ states which carry no information about the original qubit, i.e., $`f_0=1/2`$. However, for the mean fidelity we still find
$$\underset{\genfrac{}{}{0pt}{}{M=0}{\text{(even)}}}{\overset{N}{}}p_Mf_Mc_1$$
(14)
which demonstrates purification.
We emphasize, however, that the single qubit fidelity $`f_M`$ does not yet quantify the effects of a specific measurement sequence performed on the purified qubits. In fact they are entangled and a measurement on one qubit will change the state of all others. This point of view will be important for the next section.
## III Purification and estimation
We will now look at this purification protocol from a different perspective. So far it was shown in Ref. that a purified qubit lies closer to the original qubit prepared in the state $`|1_\stackrel{}{n}`$. However, the question arises whether the $`M`$ qubits in the state $`\rho _M`$ are also purer in the sense that they allow us to extract more information about the state $`|1_\stackrel{}{n}`$ than the $`N`$ unpurified qubits in state $`\rho ^N`$.
The scenario we have in mind is the following. Alice prepares $`N`$ pure qubits and sends them to Bob via a depolarizing channel. Hence Bob obtains $`N`$ mixed qubits as described by Eq. (2). He now has two possibilities to extract the information originally encoded by Alice. First, he can simply use the $`N`$ mixed qubits which are not correlated and estimate the pure qubit prepared by Alice. Second, he sacrifices some of the mixed qubits for a purification procedure as described in Section II. Hence after purification Bob obtains $`MN`$ entangled qubits with probability $`p_M`$, Eq. (7). These qubits have an improved single-qubit fidelity $`f_M`$, Eq. (13).
Our question is whether the purified qubits will also allow us to improve the estimation of Alice’s original qubits. It is important to note that our estimation only involves single qubit measurements. We will perform $`M`$ single measurements on the entangled state $`\rho _M`$, Eq. (9). This means that the measurement of one qubit affects the state of the remaining qubits via the von-Neumann projection postulate. Despite this fact we shall show that it is possible to improve the estimation of Alice’s qubit if we combine the purification protocol with an adaptive measuring scheme for single qubits.
### A Measurements on single qubits
A single measurement will give us only little information about the corresponding qubit. But if we not only have one qubit we can proceed by measuring further qubits thus receiving more and more information about their quantum state. If we do this using an adaptive strategy for choosing the measurement directions we can give a good estimate of the underlying quantum state .
This adaptive estimation strategy was originally formulated for uncorrelated qubits. However, the method does work for purified qubits as well. The only difference in this case is that the purification output or estimation input state is a non-separable $`M`$-qubit state. Therefore, the state of the qubits depends on the history of previous measurements on the system. One also has to take into account that the purification procedure is a probabilistic process. Hence purification results in several possible entangled states which must be weighted by the correct probabilities.
Let us start with the simpler case. Bob decides not to purify the qubits received but rather performs measurements on the finite ensemble described by $`\rho ^N`$, Eq. (4). Hence the qubits are uncorrelated and the $`n`$th polarization or spin measurement in direction $`(\theta _n,\varphi _n)`$ is defined by the projector $`|\theta _n,\varphi _n\theta _n,\varphi _n|`$ with
$$|\theta _n,\varphi _n\mathrm{cos}\frac{\theta _n}{2}|1+\mathrm{sin}\frac{\theta _n}{2}e^{i\varphi _n}|0.$$
(15)
For each of the $`N`$ qubits he finds the result $`1`$ with probability
$$p_1(\theta _n,\varphi _n)=\theta _n,\varphi _n|\rho |\theta _n,\varphi _n$$
(16)
and the result $`0`$ with probability
$$p_0(\theta _n,\varphi _n)=\overline{\theta _n,\varphi _n}|\rho |\overline{\theta _n,\varphi _n}=1p_1(\theta _n,\varphi _n)$$
(17)
using the state
$$|\overline{\theta _n,\varphi _n}\mathrm{cos}\frac{\theta _n}{2}|0\mathrm{sin}\frac{\theta _n}{2}e^{i\varphi _n}|1$$
(18)
orthogonal to $`|\theta _n,\varphi _n`$, Eq. (15). After $`n`$ measurements the density operator of the remaining qubits reads $`\rho ^{(Nn)}`$.
This simple situation changes considerably if Bob purifies the $`N`$ qubits before he starts measuring them. The purification procedure projects $`\rho ^N`$ on the output state $`\rho _{out}`$, Eq. (5). With probability $`p_M`$, Eq. (7), he has lost $`NM`$ qubits since they are now prepared in $`|\mathrm{\Psi }^{}`$ Bell states which carry no information about the original qubit $`|1_\stackrel{}{n}`$. This information is now concentrated in $`M`$ entangled qubits whose quantum state $`\rho _M`$, Eq. (9), has a rather complicated structure. In particular, due to the entanglement a measurement performed on one qubit now influences the quantum state of the rest.
In the $`n`$th measurement Bob finds the result $`1`$ with probability
$$\stackrel{~}{p}_1(\theta _n,\varphi _n)=\mathrm{Tr}[|\theta _n,\varphi _n\theta _n,\varphi _n|\stackrel{~}{\rho }_{Mn+1}]$$
(19)
where the trace has to be taken over all $`Mn+1`$ qubits described by the conditioned density operator $`\stackrel{~}{\rho }_{Mn+1}`$. Conditioned here means that finding the result $`1`$ leads to the normalized density operator
$$\stackrel{~}{\rho }_{Mn}=\frac{\theta _n,\varphi _n|\stackrel{~}{\rho }_{Mn+1}|\theta _n,\varphi _n}{\stackrel{~}{p}_1(\theta _n,\varphi _n)}$$
(20)
for the remaining $`Mn`$ qubits. Analogously, the result $`0`$ occurs with probability
$$\stackrel{~}{p}_0(\theta _n,\varphi _n)=1\stackrel{~}{p}_1(\theta _n,\varphi _n)$$
(21)
and the corresponding reduced density operator reads
$$\stackrel{~}{\rho }_{Mn}=\frac{\overline{\theta _n,\varphi _n}|\stackrel{~}{\rho }_{Mn+1}|\overline{\theta _n,\varphi _n}}{\stackrel{~}{p}_0(\theta _n,\varphi _n)}.$$
(22)
Hence we clearly see that the quantum state $`\stackrel{~}{\rho }_{Mn}`$ after $`n`$ measurements strongly depends on the results ($`1`$ or $`0`$) obtained for the chosen measurement directions $`(\theta _1,\varphi _1)`$, …, $`(\theta _n,\varphi _n)`$. In the notation for the conditioned density operator $`\stackrel{~}{\rho }_{Mn}`$ we should therefore include these measurement directions and the corresponding results. But we leave this out for clarity of notation. However, both situations — (i) the separable case for qubits not purified as well as (ii) the non-separable case for purified qubits — can be combined with an adaptive measurement scheme.
### B Adaptive measurements
In the last paragraph we have described how the measurements work and how they change the quantum state of all $`N`$ qubits. However, we have not yet discussed how to choose the measurement direction $`(\theta _n,\varphi _n)`$. In order to get an optimized qubit estimation we will apply an adaptive procedure . That is, we will keep track of the information obtained from measurements $`1`$, $`2`$, …, $`n1`$ in order to design the $`n`$th measurement. Various realizations of such adaptive measurements have been discussed in Ref. . Here we will shortly review the essential features.
Let us assume that we have already performed $`n1`$ measurements. Our present knowledge about the underlying state of the qubit is then represented by the estimated density operator
$$\rho _{n1}^{(est)}=𝑑\mathrm{\Omega }w_{n1}(\theta ,\varphi )|\theta ,\varphi \theta ,\varphi |$$
(23)
with a normalized probability density $`w_{n1}`$ on the Bloch sphere, i.e., $`𝑑\mathrm{\Omega }w_{n1}(\theta ,\varphi )=1`$. Hence an update of the estimated density operator from one measurement to the next really means an update of its probability density. This can be done in the following way. If we perform a polarization measurement in direction $`(\theta _n,\varphi _n)`$ we will find the result $`1`$ with probability
$$P_1(\theta ,\varphi |\theta _n,\varphi _n)=\left|\theta _n,\varphi _n|\theta ,\varphi \right|^2$$
(24)
and the result $`0`$ with probability
$$P_0(\theta ,\varphi |\theta _n,\varphi _n)=1P_1(\theta ,\varphi |\theta _n,\varphi _n)$$
(25)
conditioned on the general qubit state
$$|\theta ,\varphi \mathrm{cos}\frac{\theta }{2}|1+\mathrm{sin}\frac{\theta }{2}e^{i\varphi }|0.$$
(26)
Therefore, if we measure $`i=0`$ or $`1`$ we can update the probability density $`w_{n1}`$ according to Bayes rule resulting in the probability density
$$w_n(\theta ,\varphi )=𝒩P_i(\theta ,\varphi |\theta _n,\varphi _n)w_{n1}(\theta ,\varphi )$$
(27)
after $`n`$ measurements. The normalization constant reads
$$𝒩^1=𝑑\mathrm{\Omega }P_i(\theta ,\varphi |\theta _n,\varphi _n)w_{n1}(\theta ,\varphi ).$$
(28)
From this procedure we clearly see that $`w_n`$ comprises the results of all $`n`$ measurements. Since we have no a priori information in the beginning the initial density distribution is given by
$$w_0(\theta ,\varphi )=\frac{1}{4\pi }.$$
(29)
We recall that so far we have still simply postulated a certain measurement direction $`(\theta _n,\varphi _n)`$ for the $`n`$th step. However, the formulation given above allows us to optimize the choice $`(\theta _n,\varphi _n)`$ in the sense that we gain as much information as possible out of the $`n`$th measurement. The expected information gain of the $`n`$th measurement reads
$$S(\theta _n,\varphi _n)=\underset{i=0}{\overset{1}{}}p_i^{(est)}(\theta _n,\varphi _n)\mathrm{ln}p_i^{(est)}(\theta _n,\varphi _n)$$
(30)
using the estimated probabilities
$`p_1^{(est)}(\theta _n,\varphi _n)`$ $`=`$ $`\theta _n,\varphi _n|\rho _{n1}^{(est)}|\theta _n,\varphi _n`$ (31)
$`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }w_{n1}(\theta ,\varphi )P_1(\theta ,\varphi |\theta _n,\varphi _n)}`$ (32)
and
$$p_0^{(est)}(\theta _n,\varphi _n)=1p_1^{(est)}(\theta _n,\varphi _n)$$
(33)
given by our current knowledge, i.e., given by the density operator $`\rho _{n1}^{(est)}`$. The optimized measurement direction then is the one which maximizes $`S(\theta _n,\varphi _n)`$, Eq. (30).
Eventually, we will have measured all qubits in this way and hence we arrive at $`w_N(\theta ,\varphi )`$. The point $`(\theta ^{(est)},\varphi ^{(est)})`$ at which $`w_N`$ takes on its maximum defines the pure state $`|\theta ^{(est)},\varphi ^{(est)}`$ which is our final estimate for the original qubit $`|1_\stackrel{}{n}`$. The corresponding overlap
$$F=\left|\theta ^{(est)},\varphi ^{(est)}|1_\stackrel{}{n}\right|^2$$
(34)
then measures the fidelity of our estimation.
## IV Numerical simulations and results
Using numerical simulations we will now show that the combination of purification and adaptive measurements leads to an improved estimation. Even though we have more and separable qubits at our disposal, if we apply no purification, the case with purification is superior.
In order to demonstrate this we have calculated the mean fidelity
$$F=<|\theta ^{(est)},\varphi ^{(est)}|1_\stackrel{}{n}|^2>$$
(35)
where $`\mathrm{}`$ means choosing $`4\times 10^4`$ initial states $`|1_\stackrel{}{n}`$ equally distributed over the Bloch sphere, to ensure that we average over sufficiently many histories of measurement results. In the case of the purification we also have to average over the possible results of the purification process. So for the case of $`N`$ input qubits all estimations for $`M=0,2,4,\mathrm{},N`$ must be carried out and weighted by the appropriate probabilities $`p_M`$, Eq. (7).
In Fig.1 the average fidelity obtained by the adaptive estimation of $`N=6`$ purified qubits is compared to the same estimation procedure with unpurified qubits. The results are plotted versus the probability $`c_1`$ of the initial noisy qubits, Eq. (2). Over the whole range from pure states with $`c_1=1`$ to totally mixed states with $`c_1=1/2`$ there is an clear improvement in the estimation results. The relative gain in the fidelity of the measurement goes up to $`5.5\%`$ of the unpurified value. Even on the average there is an improvement of about $`3.3\%`$. So this shows that purification can increase the accessibility of the originally encoded qubit information. Even though the purification protocol leads to entanglement of the purified qubits we find this increase in the average estimation fidelity. That was not clear from the beginning since via this entanglement our single measurements naturally change the quantum state of the qubits not yet measured.
In Fig.2 we compare the adaptive measurement of $`N=6`$ purified qubits described above to a non-adaptive measurement where the directions $`(\theta _n,\varphi _n)`$ of the measurement are choosen randomly on the Bloch sphere. We find that for purified and non-separable qubits one can also improve the estimation quality by using adaptive measurements originally designed for separable states .
Finally we take a look at the evolution of the estimation fidelity during the estimation process. The plot in Fig.3 shows the fidelity after the $`n`$-th qubit measurement. Here the initial probability $`c_1`$ is set to $`0.75`$. For the unpurified qubits there is a slow increase in this fidelity which keeps up until the last measurement. In contrast to this the fidelity of purified qubits increases only during the first measurements. Then follows a domain of nearly stagnant fidelity. Hence the purification which is basically an increase in the length of the Bloch vector helps the estimation process to rapidly find a good estimation direction. The possible loss of qubits in the purification process does not harm dramatically because of the relatively small gain in the last qubits.
## V Conclusions
Reading of quantum information is a delicate task. The state of a qubit can only be estimated using finite resources. We have shown that adaptive measurements are a useful tool for estimating quantum states. This also holds true for decohered quantum bits which is the more common case in quantum information processing. An advantage of our approach is the reduction to simple spin measurements on single quantum systems which considerably simplifies a possible physical implementation of an estimation procedure.
Qubit purification gives us the possibility to even outrange these results. We have presented results which show that an additional purification step improves the estimation of noisy quantum bits. This means that purification does not only work in the single-qubit case but also increases the accessability of quantum information stored in an ensemble of qubits. This holds true even despite the fact that purification is a probabilistic process and may result in a loss of usable quantum systems.
###### Acknowledgements.
We would like to thank G. Alber, A. Delgado, and M. Mussinger for helpful discussions. We acknowledge financial support by the DFG Schwerpunktprogramm “Quanten-Informationsverarbeitung” and by the European Commission within the IST project QUBITS.
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# I Introduction
## I Introduction
Recently quark matter at very high density has attracted a great flurry of interest . In this region, quark matter is expected to behave as a color superconductor . Possible phenomenological applications are associated with the description of neutron star interiors, neutron star collisions and the physics near the core of collapsing stars. A better understanding of highly squeezed nuclear matter might also shed some light on nuclear matter at low density, i.e. densities close to ordinary nuclear matter where some models already exist. For example, in Ref. a rather complete soliton model at low density is constructed containing vector-bosons, along with the Goldstone bosons.
In a superconductive phase, the color symmetry is spontaneously broken and a hierarchy of scales, for given chemical potential, is generated. Indicating with $`g`$, the underlying coupling constant, the relevant scales are: the chemical potential $`\mu `$ itself, the dynamically generated gluon mass $`m_{gluon}g\mu `$ and the gap parameter $`\mathrm{\Delta }\frac{\mu }{g^5}e^{\frac{\alpha }{g}}`$ with $`\alpha `$ a calculable constant. Since for high $`\mu `$ the coupling constant $`g`$ (evaluated at the fixed scale $`\mu `$) is $`1`$ , we have:
$$\mathrm{\Delta }m_{gluon}\mu .$$
(1)
Massless excitations dominate physical processes at very low energy with respect to the energy gap ($`\mathrm{\Delta }`$). Their spectrum is intimately related to the underlying global symmetries and the way they are realized at low energies. Indeed when the dynamics is such that a continuous global symmetry is spontaneously broken, a Goldstone boson appears in order to compensate for the breaking. Massless excitations obey low energy theorems governing their interactions which can be usefully encoded in effective Lagrangians. A well known example, in the region of cold and non-dense QCD, is the effective Lagrangian for pions and kaons. These Lagrangians are seen to describe well the QCD low energy phenomenology .
Another set of relevant constraints is provided by quantum anomalies. At zero density and temperature, ’t Hooft argued that the underlying continuous global anomalies have to be matched in a given low energy phase by a set of massless fermions associated with the intact global symmetries and a set of massless Goldstone bosons associated with the broken ones. The low energy fermions (composite or elementary) contribute via triangle diagrams while for the Goldstones a Wess-Zumino term should be added to correctly implement the associated global anomalies.
In , the ’t Hooft constraints were seen to hold for QCD at finite density. In particular, it was shown, by reviewing the dynamically favored phases for $`N_f=2,3`$ at high density, that the low energy spectrum displayes the correct quantum numbers to saturate the ’t Hooft global anomalies. It was also observed that QCD at finite density can be envisioned, from a global symmetry and anomaly point of view, as a chiral gauge theory for which at least part of the matter field content is in complex representations of the gauge group. Indeed an important distinction from zero density, vector-like theories is that these theories, when strongly coupled, can exist in the Higgs phase by dynamically breaking their own gauge symmetries. This is also the striking feature of the superconductive phase allowed for QCD at high density. In fact at finite density, vector-like symmetries are no longer protected against spontaneous breaking by the Vafa-Witten theorem .
In this paper we build the effective theory describing the low energy excitations for Quantum Chromodynamics with two flavors at high density. The non linear realization framework is used to properly construct the low energy theory. The light degrees of freedom, as required by ’t Hooft anomaly conditions, contains massless fermions which we properly include in the effective Lagrangian. We prove that the non linearly realized low energy effective theory, in general, contains at the lowest order in a derivative expansion for the Goldstones (i.e. two derivative) two independent terms which are responsible for distinct contributions to the gluon masses when gauging the color degrees of freedom.
We finally provide a discussion of the linearly realized Lagrangian. In this framework we see, as in reference , that the linearly realized effective Lagrangian at the dimension four level actually predicts the squared gluon mass ratio between the eighth gluon and the other massive ones to be $`4/3`$. This result does not agree with the underlying calculations . In order to resolve the issue, in Ref. it was pointed out that one can include in the linear Lagrangian another two derivative term possessing mass dimension 6. This result shows the failure of the standard naive dimensional counting arguments for suppressing higher dimensional terms for the effective linearly realized effective Lagrangian. Interestingly, in the non linear realization case all the terms needed to correctly describe the gluon masses arises at the same derivative order. This fact suggests that non linear realization description leads to a consistent counting scheme.
In Section II we summarize some general features of 2 flavor QCD. The general non linear formalism for the spontaneous breaking of color is presented in Section III. We implement the breaking of the baryon number in Section IV. Then, in medium, fermions are introduced in Section V together with the summary of the low energy theory and its extension to the broken Lorentz group case. A linearly realized Lagrangian is explored in Section VI. In Section VII we conclude.
## II General features of 2 flavor QCD at High Density
The underlying gauge group is $`SU(3)`$ while the quantum flavor group is
$$SU_L(2)\times SU_R(2)\times U_V(1),$$
(2)
and the classical $`U_A(1)`$ symmetry is destroyed at the quantum level by the Adler-Bell-Jackiw anomaly. We indicate with $`q_{Lc,i;\alpha }`$ the two component left spinor where $`\alpha =1,2`$ is the spin index, $`c=1,2,3`$ is the color index while $`i=1,2`$ represents the flavor. $`q_{R}^{}{}_{c,i}{}^{\dot{\alpha }}`$ is the two component right spinor. We summarize the transformation properties in the following table:
| | $`[SU(3)]`$ | $`SU_L(2)`$ | $`SU_R(2)`$ | $`U_V(1)`$ |
| --- | --- | --- | --- | --- |
| $`q_L`$ | | | $`1`$ | $`1`$ |
| $`q_R`$ | | $`1`$ | | $`1`$ |
(3)
where $`SU(3)`$ is the gauge group, indicated by square bracket. The theory is subject to the following global anomalies:
$$SU_{L/R}(2)^2U_V(1)\pm 3,$$
(4)
where we have chosen a convention for the flavor generators in which the quadratic anomaly factor is $`1`$.
At zero density we have two possible phases compatible with the anomaly conditions. One is the ordinary Goldstone phase associated with the spontaneous breaking of the underlying global symmetry to $`SU_V(2)\times U_V(1)`$. The other is the Wigner-Weyl phase where, assuming confinement, the global symmetry at low energy is intact and the needed massless spectrum consists of massless baryons. The Goldstone phase, associated with a non vanishing $`\overline{q}_Lq_R`$ condensate, is the one observed in nature. This fact supports a new idea presented in Ref. . There it is suggested that, for an asymptotically free theory, among multiple infrared phases allowed by ’t Hooft anomaly conditions, the one which minimizes the entropy at the approach to freeze-out is preferred.
What happens when we squeeze nuclear matter? At very low baryon density compared to a fixed intrinsic scale of the theory $`\mathrm{\Lambda }`$, it is reasonable to expect that the Goldstone phase persists. On the other hand at very large densities, it is seen, via dynamical calculations , that the ordinary Goldstone phase is no longer favored compared with a superconductive one associated with the following type of quark condensates:
$$ϵ^{\alpha \beta }ϵ^{abc}ϵ^{ij}<q_{Lb,i;\alpha }q_{Lc,j;\beta }>,ϵ_{\dot{\alpha }\dot{\beta }}ϵ^{abc}ϵ^{ij}<q_{Rb,i}^{\dot{\alpha }}q_{Rc,j}^{\dot{\beta }}>.$$
(5)
We can now introduce two scalar fields which play the role of order parameters:
$$L_{}^{}{}_{}{}^{a}ϵ^{abc}ϵ^{ij}q_{Lb,i}^\alpha q_{Lc,j;\alpha },R_{}^{}{}_{}{}^{a}ϵ^{abc}ϵ^{ij}q_{Rb,i;\dot{\alpha }}q_{Rc,j}^{\dot{\alpha }}.$$
(6)
Under a parity transformation
$$L_aR_a.$$
(7)
If parity is not broken spontaneously, we have
$$L_a=R_a=v\delta _a^3,$$
(8)
where we chose the condensate to be in the 3rd direction of color. This condensate is not allowed at zero density by the Vafa-Witten theorem. The order parameters are singlets under the $`SU_L(2)\times SU_R(2)`$ flavor transformations while possessing baryon charge $`2`$.
The vev breaks the gauge symmetry while leaving intact the following group:
$$\left[SU(2)\right]\times SU_L(2)\times SU_R(2)\times \stackrel{~}{U}_V(1),$$
(9)
where $`\left[SU(2)\right]`$ is the unbroken part of the gauge group. The $`\stackrel{~}{U}_V(1)`$ generator is the following linear combination of the previous $`U_V(1)`$ generator $`Q=\mathrm{diag}(1,1,1)`$ and the broken diagonal generator of the $`SU(3)`$ gauge group $`T^8=\frac{1}{2\sqrt{3}}\mathrm{diag}(1,1,2)`$
$$\stackrel{~}{S}=\frac{1}{3\sqrt{2}}\left[Q2\sqrt{3}T^8\right].$$
(10)
The $`\stackrel{~}{S}`$ charge of the quarks with color $`1`$ and $`2`$ is zero.
The superconductive phase for $`N_f=2`$ possesses the same global symmetry group of the confined Wigner-Weyl phase. This remarkable feature when considering chiral gauge theories at zero density (where the superconductive phase is now a Higgs phase) is referred as complementarity. This idea was introduced in Ref. where it was conjectured that any Higgs phase can be described in terms of confined degrees of freedom and vice versa.
It is convenient to introduce the fields:
$$V_a=\frac{L_a+R_a}{\sqrt{2}},A=\frac{L_aR_a}{\sqrt{2}}.$$
(11)
On the vacuum, via (8), we have:
$$V_a=\sqrt{2}v\delta _a^3A_a=0.$$
(12)
The massless excitations are associated with the fluctuations around the vacuum expectation value for $`V_a`$ while the fields described by $`A_a`$ are massive. In Ref. it is argued, based on a dynamical calculation, that at very high density the $`A_a`$ fields might be very light and possibly relevant for the low energy phenomenology. However we will focus mainly on the truly massless excitations since low energy theorems are valid only for these fields. Low energy theorems are, in general, efficiently encoded in a non linear realization framework which we will soon construct.
## III Low Energy Effective Theory without Baryon number
In order to write an effective Lagrangian for 2 flavor QCD, one could, at least in principle, proceed in different ways. For example one could add a mass term for the strange quark in the $`3`$ flavor Lagrangian. However, since the theory is not supersymmetric, an exact decoupling procedure is not known<sup>§</sup><sup>§</sup>§An attempt to generalize to QCD the Seiberg decoupling procedure at the effective Lagrangian level is provided in . at effective Lagrangian level. A general way, which we will explore here, to directly construct the 2 flavor effective Lagrangian makes use of the non linear realization methods . The latter have been already successfully employed for the 3 flavor case at high density in .
Following reference , we now construct the non linearly realized effective Lagrangian containing the diquark degrees of freedom coupled to gluons.
We postpone the breaking of the baryon number and its consequences to the next section. The group of transformations is $`G=SU(3)`$, while the stability group leaving the vacuum invariant is the proper subgroup $`H=SU(2)G`$. The color generators $`T^m`$ of $`SU(3)`$, with $`m=1,\mathrm{},8,`$ obey the normalization condition $`\mathrm{Tr}\left[T^mT^n\right]=\frac{1}{2}\delta ^{mn}`$. We divide the generators $`\{T\}`$ into two classes, calling the generators of $`H`$ $`\{S^a=T^a\}`$, with $`a=1,\mathrm{},3`$ and the broken generators $`\{X^i=T^{i+3}\}`$ with $`i=1,\mathrm{},5`$. It is worth noticing that the quotient space $`G/H`$, in this case, is not a symmetric space. By symmetric space we mean that, if $`X`$ and $`S`$ represent arbitrary linear combinations of the broken and unbroken generators, their commutators should satisfy the restriction:
$$[X,X]=iS.$$
(13)
It is easy to verify that the previous condition is not obeyed. Clearly, our generators should always satisfy the trivial conditions
$$[S,S]=iS,[X,S]=iX,$$
(14)
expressing the fact that the Lie algebra of $`H`$ closes and that the $`\{X\}`$ form a representation of $`H`$.
The coset space $`G/H`$ is parameterized by the group elements
$$𝒱=\mathrm{exp}(i\xi ^iX^i),$$
(15)
with $`\xi ^i=\mathrm{\Pi }^i/v`$ describing the Goldstone fields and are coordinates of the space $`G/H`$. $`𝒱`$ transforms non linearly under a color transformation, i.e.:
$$𝒱(\xi )g𝒱(\xi )h^{}(\xi ,g),$$
(16)
where $`gG`$ and $`hH`$. It is convenient to define the hermitian (algebra valued) Maurer-Cartan one-form
$$\omega _\mu =i𝒱^{}D_\mu 𝒱,$$
(17)
with $`D_\mu `$ the covariant derivative with respect to $`G=SU(3)`$
$$D_\mu 𝒱=(_\mu iG_\mu )𝒱,$$
(18)
and $`G_\mu ^mT^m`$ the gluon fields. Since $`𝒱`$ transforms with respect to $`G`$ as in Eq. (16) it follows that
$$\omega _\mu h(\xi ,g)\omega _\mu h^{}(\xi ,g)+h(\xi ,g){}_{\mu }{}^{}h_{}^{}(\xi ,g).$$
(19)
We decompose $`\omega _\mu `$ into the part parallel to $`H`$
$$\omega _\mu ^{}=2S^a\mathrm{Tr}\left[S^a\omega _\mu \right],\mathrm{Lie}H,$$
(20)
and into the perpendicular part
$$\omega _\mu ^{}=2X^i\mathrm{Tr}\left[X^i\omega _\mu \right],\mathrm{Lie}G\mathrm{Lie}H.$$
(21)
Summation over repeated indices is assumed.
This yields the following transformation properties:
$$\omega _\mu ^{}h(\xi ,g)\omega _\mu ^{}h^{}(\xi ,g)+h(\xi ,g)_\mu h^{}(\xi ,g),$$
(22)
and
$$\omega _\mu ^{}h(\xi ,g)\omega _\mu ^{}h^{}(\xi ,g).$$
(23)
Then, it turns out that the most general invariant of second order in the derivatives is
$$L=v^2\mathrm{Tr}\left[\omega _\mu ^{}\omega ^\mu \right]=2v^2\mathrm{Tr}\left[X^i𝒱^{}D_\mu 𝒱\right]\mathrm{Tr}\left[X^i𝒱^{}D^\mu 𝒱\right].$$
(24)
We can immediately see, by adopting the unitary gauge (corresponding to $`\xi 0`$), that the previous term provides a mass term for the gluons associated with the 5 coset generators, whereas the gluons $`G^{1,2,3}`$ remain massless. Indeed the mass term reads:
$$L=2v^2\mathrm{Tr}\left[X^iG_\mu \right]\mathrm{Tr}\left[X^iG^\mu \right]=\frac{v^2}{2}G_\mu ^iG^{i\mu }.$$
(25)
The fact that we find the same mass for all of the gluons of the broken subgroup is not surprising, since we have treated them equally. In the next section we address the problem of the baryon number and show how, in the non linear framework, this naturally leads to a distinct mass for the eighth gluon.
The massless gluons of the subgroup $`H`$ confine, leaving, as we shall see, only some of the quarks as massless excitations at low energies.
## IV The Abelian Group $`U_V(1)`$
The previous discussion is not yet complete, since we have omitted the breaking of the baryon number. Indeed the group of transformations should be $`G=SU(3)\times U_V(1)`$, while the unbroken subgroup is $`H=SU(2)\times \stackrel{~}{U}_V(1)`$.
The non linear transformations are now:
$$𝒱(\xi )u_Vg𝒱(\xi )h^{}(\xi ,g,u)h_{\stackrel{~}{V}}^{}(\xi ,g,u),$$
(26)
$$u_VU_V(1),gSU(3),h(\xi ,g,u)SU(2),h_{\stackrel{~}{V}}(\xi ,g,u)\stackrel{~}{U}_V(1).$$
(27)
The $`U_V(1)`$ charge is $`Q=\mathrm{diag}(1,1,1)`$. The generator, for $`\stackrel{~}{U}_V(1)`$, leaving the vev invariant is
$$\stackrel{~}{S}=\frac{1}{3\sqrt{2}}\left[Q2\sqrt{3}T^8\right]=\mathrm{diag}(0,0,\sqrt{2}/2),$$
(28)
where we have chosen to normalize it according to Tr$`\left[\stackrel{~}{S}^2\right]=1/2`$. The coset space $`G/H`$ is parameterized, as before, by the group elements
$$𝒱=\mathrm{exp}(i\xi ^iX^i),$$
(29)
but now we identify the Goldstone bosons as:
$`\xi ^i`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }^i}{v}}i=1,2,3,4,`$ (30)
$`\xi ^5`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }^5}{\stackrel{~}{v}}}.`$ (31)
We still have 5 generators $`\{X^i\}`$ which belong to the coset space $`G/H`$, but one of them needs to be modified. While $`X^{1,2,3,4}`$ are still identified, respectively, with $`T^{4,5,6,7}`$, the coset generator which replaces $`X^5=T^8`$ of the previous section is now
$$X^5=\frac{1}{3}\left[Q+\sqrt{3}T^8\right]=\mathrm{diag}(\frac{1}{2},\frac{1}{2},0).$$
(32)
We used the orthogonality condition
$$\mathrm{Tr}\left[XS\right]=0,$$
(33)
to construct $`X^5`$ with Tr$`\left[X^5X^5\right]=\frac{1}{2}`$. Notice that $`X^5`$ is no longer traceless (i.e. $`\mathrm{Tr}\left[X^5\right]=1`$).
It is amusing to note that $`𝒱`$ transforms on the left as an ordinary quark with respect to color and $`U_V(1)`$ transformations. This will turn to be important when adding the fermions.
As for ordinary QCD, we can construct a field transforming linearly by using only non linearly realized fields via:
$$V_a=\frac{v}{\sqrt{2}}ϵ_{abc}𝒱_i^b𝒱_j^cϵ^{ij3}.$$
(34)
To prove that the previous expression transforms linearly we recall that the $`H`$ subgroup acts on the right of $`𝒱`$. So under a general non linear transformation we have:
$$ϵ_{abc}𝒱_i^b𝒱_j^cϵ^{ij3}u_V^2ϵ_{abc}g_d^bg_e^c𝒱_k^d𝒱_l^eh_{}^{}{}_{i}{}^{k}h_{}^{}{}_{j}{}^{l}h_{\stackrel{~}{V}}^{}{}_{}{}^{2}ϵ^{ij3}u_V^2ϵ_{abc}g_d^bg_e^c𝒱_k^d𝒱_l^eϵ^{kl3}.$$
(35)
Indeed due to Eq. (28) we have that $`h_{\stackrel{~}{V}}`$ acting on the color indices $`1`$ and $`2`$ is equivalent to the identity. We also have
$$h_{}^{}{}_{i}{}^{k}h_{}^{}{}_{j}{}^{l}ϵ^{ij3}=ϵ^{kl3},$$
(36)
since $`h`$ is an $`SU(2)`$ matrix and we have conveniently chosen the indices. Therefore, $`V_a`$ transforms linearly. Using the relation:
$$\mathrm{det}𝒱\frac{1}{3!}ϵ_{abc}ϵ^{kij}𝒱_k^a𝒱_i^b𝒱_j^c=\mathrm{exp}\left[i\frac{\mathrm{\Pi }^5}{\stackrel{~}{v}}\frac{\mathrm{Tr}\left[Q\right]}{3}\right]=\mathrm{exp}\left[i\frac{\mathrm{\Pi }^5}{\stackrel{~}{v}}\right],$$
(37)
we can rewrite Eq. (34) in the following form:
$$V_a=\sqrt{2}ve^{i\frac{\mathrm{\Pi }^5}{\stackrel{~}{v}}}𝒱_{}^{1}{}_{a}{}^{3}.$$
(38)
The latter expression corresponds to the polar decomposition of the vector $`V_a`$ with the massive scalar field frozen to its vev value. Since the field $`\mathrm{\Pi }^5`$ is associated with the $`X^5`$ generator, which is a linear combination of an abelian and of a non abelian $`T^8`$ generator, $`\mathrm{det}𝒱`$ does not parameterize an independent field.
Expanding $`𝒱`$ up to first order in the Goldstones leads to:
$$V_a=\sqrt{2}v\left[1+i\xi ^5\left(\frac{2}{3}Q\frac{1}{\sqrt{3}}T^8\right)i\underset{i=1}{\overset{4}{}}\xi ^iX^i\right]_a^c\delta _c^3+v𝒪\left(\xi ^2\right).$$
(39)
This new field explicitly describes the vev properties. First of all, when considering the limit $`\xi 0`$ we recover the correct vacuum, i.e.
$$V_a=\sqrt{2}v\delta _a^3.$$
(40)
Furthermore, as expected, $`V_a`$ transforms under the underlying gauge transformations as a diquark.
Equation (34) mimics what happens in ordinary QCD. There the vev transformations are encoded in the linearly transforming matrix $`Ug_LUg_R^{}`$ while the non linearly transforming field is simply $`\sqrt{U}g_L\sqrt{U}k^{}(\sqrt{U},g_L,g_R)=k(\sqrt{U},g_L,g_R)\sqrt{U}g_R^{}`$, with $`k(\sqrt{U},g_L,g_R)SU_V(3)`$.
If we insist in gauging the color transformations the general discussion in the previous section remains unchanged. The substantial difference is that now the most generic two derivative kinetic Lagrangian
$$L=v^2a_1\mathrm{Tr}\left[\omega _\mu ^{}\omega ^\mu \right]+v^2a_2\mathrm{Tr}\left[\omega _\mu ^{}\right]\mathrm{Tr}\left[\omega ^\mu \right],$$
(41)
acquires a new term. The presence of a double trace term is due to the absence of the traceless condition for the broken generator $`X^5`$.
The kinetic term for the Goldstones is:
$$L_{kin}=\frac{a_1}{2}\underset{i=1}{\overset{4}{}}_\mu \mathrm{\Pi }^i^\mu \mathrm{\Pi }^i+\frac{v^2}{\stackrel{~}{v}^2}\left(\frac{a_1+2a_2}{2}\right)_\mu \mathrm{\Pi }^5^\mu \mathrm{\Pi }^5.$$
(42)
Normalizing the kinetic term we have
$$a_1=1,a_2=\frac{1}{2}\left[\frac{\stackrel{~}{v}^2}{v^2}a_1\right].$$
(43)
The square gluon mass ratio, in the unitary gauge (i.e. $`\xi 0`$), between the eighth gluon and any other massive one now reads:
$$\frac{m_8^2}{m_i^2}=\frac{1}{3}\left[1+2\frac{a_2}{a_1}\right]=\frac{1}{3}\frac{\stackrel{~}{v}^2}{v^2},$$
(44)
with $`i=4,5,6`$ or $`7`$ and in the last step we use Eq. (43). Using perturbation theory at one loop $`a_21/2`$ (i.e. $`\stackrel{~}{v}\sqrt{2}v`$). Also, as we shall see, in the linearly realized lagrangian one can write two two-derivative terms in the effective Lagrangian. However one of them appears in the Lagrangian as a higher order in mass dimension, leading to a not straightforward systematic expansion.
## V In Medium Fermions
Now we turn our attention to the Fermi fields. First of all, let us introduce the following field:
$$\stackrel{~}{\psi }=𝒱^{}\psi ,$$
(45)
which under the action of the group $`G=SU(3)\times U_V(1)`$ transforms as
$$\stackrel{~}{\psi }h_{\stackrel{~}{V}}(\xi ,g,u)h(\xi ,g,u)\stackrel{~}{\psi }.$$
(46)
This construction allows us to easily switch to the non linear representations and it is often encountered in the Heavy Quark Effective formalism .
The fermion field in (45) can be thought as a, in medium, quark surrounded by a Goldstone cloud. Remarkably, as also required by ’t Hooft anomaly conditions , the newly defined (in medium) fermion possesses the unbroken $`\stackrel{~}{U}_V(1)`$ charge. The cloud (effectively included in $`𝒱`$) correctly screens the baryonic charge.
In the new variables, one can construct two independent invariant terms
$$L_1=\overline{\stackrel{~}{\psi }}i\gamma ^\mu (_\mu i\omega _\mu ^{})\stackrel{~}{\psi },L_2=\overline{\stackrel{~}{\psi }}\gamma ^\mu \omega _\mu ^{}\stackrel{~}{\psi }.$$
(47)
One can easily verify that the tree level lagrangian term for the quarks
$$L_{\mathrm{Tree}}=\overline{\psi }\gamma ^\mu (_\mu iG_\mu )\psi ,$$
(48)
corresponds to the linear combination of
$$L_{\mathrm{Tree}}=L_1+L_2.$$
(49)
Another invariant that we can add to the effective Lagrangian corresponds to the following Majorana mass term:
$$L_m=m_M\overline{\stackrel{~}{\psi }^C}_i\gamma ^5(iT^2)\stackrel{~}{\psi }_jϵ^{ij}+\mathrm{h}.\mathrm{c}.$$
(50)
with $`\stackrel{~}{\psi }^C=i\gamma ^2\stackrel{~}{\psi }^{}`$, where $`i,j=1,2`$ are flavor indices. The invariance under a $`h(\xi ,g,u)SU(2)`$ transformation is insured by the relation
$$(S^a)^TT^2=T^2S^a,$$
(51)
where
$$T^2=S^2=\frac{1}{2}\left(\begin{array}{cc}\sigma ^2\hfill & 0\hfill \\ 0\hfill & 0\hfill \end{array}\right).$$
(52)
The invariance under a $`\stackrel{~}{U}_V\left(1\right)`$ transformation is due to the fact that its generator $`\stackrel{~}{S}=\mathrm{diag}(0,0,\frac{\sqrt{2}}{2})`$ is such that $`(\stackrel{~}{S})^TT^2=T^2\stackrel{~}{S}=0`$.
We stress that this Lagrangian term respects the underlying color transformations as well as the global transformations $`SU_L(2)\times SU_R(2)\times U_V\left(1\right)`$. The breaking of the underlying symmetries manifests itself only when we evaluate this Lagrangian on the vev (i.e. $`\xi 0`$). To see this one can rewrite the previous terms as an explicit function of $`𝒱`$. Parity is also enforced. This term does not produce a mass term for the quarks with color index $`3`$ while yielding a Majorana mass term for fermions with color index $`1`$ and $`2`$.
The latter Lagrangian term is similar to the meson-fermion interaction we can write in standard QCD (of the type $`\overline{q}Uq`$ which yields the constituent Dirac quark masses when evaluated on the vev).
At this point, for reader’s convenience, we summarize the total non linearly realized effective Lagrangian describing in medium fermions, gluons and their self interactions, up to two derivatives,
$`=`$ $`v^2a_1\mathrm{Tr}\left[\omega _\mu ^{}\omega ^\mu \right]+v^2a_2\mathrm{Tr}\left[\omega _\mu ^{}\right]\mathrm{Tr}\left[\omega ^\mu \right]`$ (53)
$`+`$ $`b_1\overline{\stackrel{~}{\psi }}i\gamma ^\mu (_\mu i\omega _\mu ^{})\stackrel{~}{\psi }+b_2\overline{\stackrel{~}{\psi }}\gamma ^\mu \omega _\mu ^{}\stackrel{~}{\psi }`$ (54)
$`+`$ $`m_M\overline{\stackrel{~}{\psi }^C}\gamma ^5(iT^2)\stackrel{~}{\psi }+\mathrm{h}.\mathrm{c}..`$ (55)
Here $`a_1,a_2,b_1`$ and $`b_2`$ are real coefficients while $`m_M`$ is complex and we omit flavor indices. At the tree level in the underlying theory we have $`b_1=b_2=1`$. The massless degrees of freedom are the in medium fermions $`\stackrel{~}{\psi }_{a=3,i}`$ which possess the correct quantum numbers prescribed by ’t Hooft anomaly conditions . As already mentioned, the $`SU(2)`$ color subgroup remains unbroken and the associated 3 gluons are expected to confine again. To the previous general effective Lagrangian we should also add the gluon kinetic term.
In writing the effective low energy theory we have not yet considered the breaking of Lorentz invariance at finite density. Following Ref. we impose invariance only under the $`O(3)`$ subgroup of the Lorentz transformations. This amounts to have different coefficients for the temporal and spatial indices of the Lagrangian which now becomes:
$`=`$ $`v^2a_1\mathrm{Tr}\left[\omega _0^{}\omega _0^{}\alpha _1\stackrel{}{\omega }^{}\stackrel{}{\omega }^{}\right]+v^2a_2\left[\mathrm{Tr}\left[\omega _0^{}\right]\mathrm{Tr}\left[\omega _0^{}\right]\alpha _2\mathrm{Tr}\left[\stackrel{}{\omega }^{}\right]\mathrm{Tr}\left[\stackrel{}{\omega }^{}\right]\right]`$ (56)
$`+`$ $`b_1\overline{\stackrel{~}{\psi }}i\left[\gamma ^0(_0i\omega _0^{})+\beta _1\stackrel{}{\gamma }\left(\stackrel{}{}i\stackrel{}{\omega }^{}\right)\right]\stackrel{~}{\psi }+b_2\overline{\stackrel{~}{\psi }}\left[\gamma ^0\omega _0^{}+\beta _2\stackrel{}{\gamma }\stackrel{}{\omega }^{}\right]\stackrel{~}{\psi }`$ (57)
$`+`$ $`m_M\overline{\stackrel{~}{\psi }^C}\gamma ^5(iT^2)\stackrel{~}{\psi }+\mathrm{h}.\mathrm{c}.,`$ (58)
where the new coefficients $`\alpha `$s and $`\beta `$s encode the effective breaking of Lorentz invariance. To the previous Lagrangian we can still add the chemical potential type of term $`\overline{\stackrel{~}{\psi }}\gamma ^0\stackrel{~}{\psi }`$.
Clearly, in future, it would be valuable to compute the coefficients of the effective Lagrangian in Eq. (58) at asymptotically high densities. We notice that, due to the invariance under the global $`SU_L(2)\times SU_R(2)`$ symmetry group, no mass term arises for the third colored quarks and that any dynamical calculation (preserving the flavor symmetries) will have to respect this condition. In Ref. , using some dynamical calculation, it is argued that this seems to be the case.
## VI Linear Realizations
For the linearly realized effective Lagrangian we can start directly with the triplet field $`V_a`$. This complex field encodes the five Goldstones plus a massive scalar. The relation with the non linearly realized fields is provided in Eq. (34).
While in the non linear realization case the effective Lagrangian is ordered in number of derivatives, usually in the linear case, the terms are ordered according to their increasing mass dimension. The general dimension four effective Lagrangian is:
$$=D^\mu V^{}D_\mu V+P(V).$$
(59)
$`P(V)`$ is a potential term of the general form:
$$P(V)=M_V^2V^{}V+\lambda _V\left[V^{}V\right]^2,$$
(60)
with $`M_V`$ and $`\lambda _V`$ real parameters and the negative sign mass term has been chosen to provide a non zero vev for $`V`$
The color covariant derivative is defined as follows:
$$D_\mu V=_\mu V+iVG_\mu ,$$
(61)
and $`G_\mu =G_\mu ^mT^m`$ is the gluon field. In constructing the previous lagrangian, we have also assumed parity invariance. The potential leads to a non vanishing vacuum expectation value for $`V`$ which we choose to align in the color direction $`3`$,
$$V_a=\sqrt{2}v\delta _a^3,v=\frac{1}{2}\sqrt{\frac{M_V^2}{\lambda _V}}.$$
(62)
The local $`SU(3)`$ gauge invariance is then spontaneously broken to the $`SU(2)`$ local gauge subgroup. $`5`$ of the $`6`$ independent scalars contained in $`V_a`$ are massless. We choose the unitary gauge and absorbe the massless degrees of freedom in the longitudinal components of the 5 massive gauge fields. In this gauge we can write the fluctuations around the vacuum for $`V_a`$ as:
$$V=\left(\begin{array}{c}0\\ 0\\ \sqrt{2}v+\sigma \end{array}\right).$$
(63)
where the scalar $`\sigma `$ is a massive real field with $`\sigma =0`$. The masses of the $`5`$ gluons are related to the condensate via:
$$m_i=v,m_8=\frac{2}{\sqrt{3}}v,$$
(64)
with $`i=4,5,6,7`$. Clearly no Goldstone boson survives since there is no global symmetry left unbroken. As expected the massless degrees of freedom are the first $`3`$ gluons for the unbroken $`SU(2)`$ gauge symmetry. However the latter are supposed to confine and hence to generate a new confining scale associated with pure gluon-dynamics with two colors but no flavors.
We immediately notice, as in reference , that the linearly realized effective Lagrangian at the dimension four level predicts the squared gluon mass ratio between the eighth gluon and the other massive ones to be $`4/3`$. This result does not agree with the underlying calculations . In order to resolve the issue in Ref. , it was pointed out that one can include in the linear Lagrangian another two derivative term of the form:
$$\frac{b}{v^2}\left[V^{}D_\mu V\right]\left[V^{}D^\mu V\right].$$
(65)
$`b`$ is a real number. This term has mass dimension 6. At this point we cannot use, in general, the standard naive dimensional counting arguments for suppressing other higher dimensional terms for the effective linearly realized Lagrangian. Interestingly, in the non linear realization case, all the terms needed to correctly describe the gluon masses arises at the same derivative order. This fact suggests that non linear realization description leads to the correct counting scheme.
## VII Conclusions
We constructed the low energy effective theory describing Quantum Chromodynamics with two flavors at high density. In order to correctly implement the low energy theorems, we have used the non linear realization framework. An important difference with respect to the 3 flavor case is the presence, guaranteed by ’t Hooft anomaly conditions, of massless fermions. We have hence properly added the, in medium, fermions to the effective Lagrangian. Then we generalized it to the case when Lorentz invariance is broken, by medium effects, to $`O(3)`$. We also show that there are two independent two derivative terms for the Goldstone bosons (which will eventually become longitudinal components of the color gauge bosons). This is shown to be an effect related to the baryon number violation. We finally investigate the linearly realized effective Lagrangian which, in general, is useful and well defined only when describing a phase transition. Here we see, as also noted in Ref. , that the naive (in mass) dimensional counting rule used to naturally order the terms does not hold when confronted with perturbative dynamical calculations . In contrast the non linearly realized theory provided a consistent counting scheme.
Acknowledgments
It is pleasure for us to thank D.H. Rischke and J. Schechter for interesting discussions and encouragement. The work of Z.D. and F.S. has been partially supported by the US DOE under contract DE-FG-02-92ER-40704.
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# Apse Alignment of Narrow Eccentric Planetary Rings
## 1 INTRODUCTION
Each narrow eccentric ring surrounding Uranus is composed of particles moving on nested elliptical orbits. The outer and inner edges of a given ring define ellipses having semi-major axes $`a\pm \mathrm{\Delta }a/2`$ and eccentricities $`e\pm \mathrm{\Delta }e/2`$, where $`\mathrm{\Delta }aa`$, $`\mathrm{\Delta }ee`$, and $`e1`$. Observed values of $`a`$, $`e`$, $`\mathrm{\Delta }a`$, and $`\mathrm{\Delta }e`$ for the Uranian $`ϵ`$, $`\alpha `$, and $`\beta `$ rings are listed in Table 1.
Remarkably, the set of ellipses describing an individual ring share a common line of apsides. Apse alignment is surprising because the oblateness of Uranus causes orbits of particles with different semi-major axes to precess differentially. Timescales for differential precession in the absence of other forces are extremely short; in the case of the Uranian $`ϵ`$ ring, the inner edge would precess a full revolution relative to the outer edge in 175 years. Rigid precession of an eccentric planetary ring has remained a problem in ring dynamics for over 20 years.
Goldreich & Tremaine (1979, hereafter GT) proposed that apse alignment is maintained by self-gravity. Their theory predicts that the eccentricity gradient across the ring,
$$q_ea\frac{e}{a},$$
(1)
must be positive. A positive eccentricity gradient in an apse-aligned ring implies that the ring is narrowest at periapse and widest at apoapse. Gravitational forces between particles are therefore greatest near periapse. Material in the inner half of the ring pulls radially inward on the outer half at periapse, generating a differential precession which exactly cancels that due to planetary oblateness.
Though the prediction that $`q_e>0`$ accords with observations of all known narrow eccentric rings, the standard self-gravity model (hereafter SSG) predicts Uranian ring masses that are too low compared to those inferred from Voyager radio occultations. Ring masses based on observations exceed predictions by factors of at least $``$3 ($`ϵ`$ ring) to $``$50 ($`\alpha `$ and $`\beta `$ rings) (Tyler et al. 1986; Gresh 1990; see also the reviews by Esposito et al. 1991 and French et al. 1991). Low surface densities are particularly problematic for the $`\alpha `$ and $`\beta `$ rings. With SSG surface densities, torques exerted by inner shepherd satellites would be insufficiently strong to confine the $`\alpha `$ and $`\beta `$ rings against drag from the distended exosphere of Uranus (Goldreich & Porco 1987, hereafter GP). In addition, as discussed by Graps et al. (1995), shapes of the $`ϵ`$ ring surface density profiles as derived from occultation light curves do not accord with SSG predictions.
This paper points the way towards resolving these problems. In §2, a theory of collisionally modified self-gravity (hereafter CMSG) is qualitatively described. A simple quantitative model is set forth in §3, in which new surface density profiles are derived for the $`ϵ`$ and $`\alpha `$ rings that are in better agreement with observations. In §4, implications of our solutions for torque balance, the role of planetary oblateness, and the value of $`q_e`$ are discussed. Directions for future research are summarized in §5.
## 2 QUALITATIVE SOLUTION
For simplicity, consider an apse-aligned eccentric ring having constant, positive $`q_e`$ across its width. The ring is filled with spherical particles of internal mass density $`\rho `$ and radius $`r`$, and the ring surface density is given by $`\mathrm{\Sigma }`$. Let $`n`$ and $`\stackrel{~}{\omega }`$ be the mean motion and apsidal angle, respectively, of a ring particle. Subscripts $`i`$ and $`b`$ denote quantities evaluated in the ring interior and near the ring boundary, respectively. Variables subscripted with $`p`$ or $`s`$ are associated with the central planet or shepherd satellite, respectively, and take their usual meanings. The dimensionless strength of the quadrupole moment of the planet is given by $`J_2`$. Numerical estimates in this section are made using parameters appropriate for the $`ϵ`$ ring.
A key ingredient missing in the SSG model is an accounting for interparticle collisions. Since ring optical depths $`\tau `$ measured normal to the orbital plane are typically of order unity, each particle collides, on average, a few times with its neighbors every orbital period. Only modest collisional impulses per unit mass and time, of order $`0.1\mathrm{cm}\mathrm{s}^1`$ per orbit, are required to generate differential precession rates comparable to those induced by planetary oblateness (GT). Velocity dispersions of order $`c_i0.1\mathrm{cm}\mathrm{s}^1`$ in the ring interior are not unreasonable: both the Keplerian shearing velocity across a particle diameter, $`3nr`$, and the escape velocity from the particle surface, $`r\sqrt{8G\rho }`$, are of that order for the meter-sized bodies that plausibly compose the ring.
Although a single collision can impart an impulse of dynamically significant magnitude, multiple collisions experienced by a particle in the ring interior leave its precession rate largely unaltered. A particle in the ring interior is struck by its inside neighbors about as frequently and as forcefully as by its outside ones. Differential precession across the ring induced by smooth internal pressure gradients occurs on timescales of order $`2\pi \mathrm{\Sigma }nae/|P|2\pi nae\mathrm{\Delta }a/c_i^210^6(0.1\mathrm{cm}\mathrm{s}^1/c_i)^2\mathrm{yr}`$, much longer than misalignment timescales set by planetary oblateness (cf. GT). Here the height-integrated pressure $`P\mathrm{\Sigma }c_i^2`$ is taken to vary over a lengthscale $`\mathrm{\Delta }a`$.
Conditions are dramatically different near ring edges. Pressure-induced accelerations are maximal there because (1) velocity dispersions are enhanced by resonant satellite perturbations, and (2) the surface density declines steeply (Borderies, Goldreich, & Tremaine 1982). The velocity dispersion near the ring boundary could be as high as
$$c_b\sqrt{\frac{d}{w_r}}c_i3\frac{c_i}{0.1\mathrm{cm}\mathrm{s}^1}\mathrm{cm}\mathrm{s}^1,$$
(2)
where $`d10^3\mathrm{km}`$ is the ring-satellite separation, and $`w_r`$ is the width of the annulus perturbed by the satellite. To order-of-magnitude, the latter is given by $`w_ra\sqrt{M_s/M_p}1\mathrm{km}`$, the distance from the resonant edge at which nested periodic orbits cross. Equation (2) is derived by equating the rate of energy dissipation by collisions in the perturbed zone, $`\pi \mathrm{\Sigma }aw_rc_b^2n\tau `$, to the rate of energy deposition by the satellite, $`3nTd/2a9\pi \mathrm{\Sigma }nac_i^2d/2`$, where $`T`$ is the satellite-induced confining torque whose magnitude equals that of the viscous torque, $`3\pi \mathrm{\Sigma }c_i^2a^2`$, in steady-state.
A particle on the ring edge experiences a radially directed, collisional acceleration
$$C\frac{P}{\mathrm{\Sigma }}\pm \frac{c_b^2}{\lambda }\widehat{r}\pm c_bn\widehat{r},$$
(3)
where the upper (lower) sign applies to the outer (inner) ring edge. Here $`P\mathrm{\Sigma }c_b^2`$ is taken to vary over a radial lengthscale, $`\lambda `$, of order the local ring thickness, $`c_b/n`$. In a $`q_e>0`$ ring, collision rates are highest near periapse. At the periapsis of a ring boundary, the radial acceleration, $`C`$, generates a differential precession rate, $`\mathrm{\Delta }d\stackrel{~}{\omega }/dt_CC/nae`$, relative to the precession rate at the ring midline. This collision-induced rate is greater than the local differential rate due to planetary oblateness, $`\mathrm{\Delta }d\stackrel{~}{\omega }/dt_O`$, by a substantial factor:
$$\frac{\mathrm{\Delta }d\stackrel{~}{\omega }/dt_C}{\mathrm{\Delta }d\stackrel{~}{\omega }/dt_O}\frac{c_b/ae}{(21\mathrm{\Delta }a/8a)J_2n(R_p/a)^2}40\frac{c_b}{1\mathrm{cm}\mathrm{s}^1}.$$
(4)
Self-gravity maintains apse alignment against differential precession caused by planetary oblateness and interparticle collisions. For self-gravity to enforce rigid precession near ring edges, surface densities there must be higher than those predicted by SSG. At ring boundaries, self-gravitational attraction must balance the extra repulsive acceleration due to collisions. To estimate the surface density near the edge, $`\mathrm{\Sigma }_b`$, equate the collisional acceleration, $`Cc_b^2/\lambda `$, to the gravitational acceleration from a wire of linear mass density $`\mathrm{\Sigma }_b\lambda `$ located a distance $`\lambda `$ away:
$$c_b^2/\lambda 2G\mathrm{\Sigma }_b.$$
(5)
Take $`\lambda =c_b/n50\mathrm{m}`$ to obtain
$$\mathrm{\Sigma }_bc_bn/\mathrm{\hspace{0.17em}2}G10^3\frac{c_b}{1\mathrm{cm}\mathrm{s}^1}\mathrm{g}\mathrm{cm}^2,$$
(6)
which is greater than corresponding SSG predictions by factors $`40`$. Equation (6) is equivalent to the condition that Toomre’s Q be of order unity at the edge.
These endwires of mass $`2\pi a\mathrm{\Sigma }_b\lambda =\pi c_b^2a/G`$ constitute new boundary conditions not found in SSG. Gravitational forces from massive endwires induce substantial differential precession in the ring interior. For self-gravity to maintain apse alignment in the interior, surface densities there must also be greater than those predicted by SSG.
## 3 QUANTITATIVE MODEL
Divide the region occupied by an apse-aligned, constant $`q_e=a\mathrm{\Delta }e/\mathrm{\Delta }a`$ ring into an even number $`N`$ of equally spaced intervals. The center of the $`j^{\mathrm{th}}`$ interval contains an elliptical wire having mass $`m_j`$, semi-major axis $`a_j=a+[j(N+1)/2]\mathrm{\Delta }a/N`$, and eccentricity $`e_j=e+[j(N+1)/2]\mathrm{\Delta }e/N`$. Denote by $`\mathrm{\Delta }_jd\stackrel{~}{\omega }/dt`$ the precession rate of the $`j^{\mathrm{th}}`$ wire relative to the precession rate of a test particle at the ring midline. Uniform precession requires
$$\mathrm{\Delta }_j\frac{d\stackrel{~}{\omega }}{dt}=\mathrm{\Delta }_j(\frac{d\stackrel{~}{\omega }}{dt}_O+\frac{d\stackrel{~}{\omega }}{dt}_G+\frac{d\stackrel{~}{\omega }}{dt}_C)=0.$$
(7)
Subscripts $`O`$, $`G`$, and $`C`$ denote contributions from planetary oblateness, self-gravity, and interparticle collisions, respectively. The first two terms are given by
$$\mathrm{\Delta }_j\frac{d\stackrel{~}{\omega }}{dt}_O=\frac{21}{4}J_2n\left(\frac{R_p}{a}\right)^2\frac{a_ja}{a}$$
(8)
and
$$\mathrm{\Delta }_j\frac{d\stackrel{~}{\omega }}{dt}_G=\frac{q_eH(q_e^2)}{\pi e}n\frac{a}{M_p}\underset{kj}{}\frac{m_k}{a_ja_k},$$
(9)
where
$$H(q_e^2)\frac{1\sqrt{1q_e^2}}{q_e^2\sqrt{1q_e^2}}$$
(cf. GT).
For $`\mathrm{\Delta }_jd\stackrel{~}{\omega }/dt_C`$, the following simplistic prescription is adopted:
$$\mathrm{\Delta }_j\frac{d\stackrel{~}{\omega }}{dt}_C=\{\begin{array}{cc}+[q_eH(q_e^2)c_b^2/\lambda nae](1x/\lambda )\hfill & \text{if }x(j1/2)\mathrm{\Delta }a/N<\lambda \hfill \\ [q_eH(q_e^2)c_b^2/\lambda nae](1y/\lambda )\hfill & \text{if }y(Nj+1/2)\mathrm{\Delta }a/N<\lambda \hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
(10)
Thus, $`\mathrm{\Delta }_jd\stackrel{~}{\omega }/dt_C`$ is non-zero only within intervals $`\delta a=\lambda `$ from each edge; there, its magnitude rises linearly from 0 to $`q_eH(q_e^2)c_b^2/\lambda nae`$. As a crude justification for this maximum value, approximate the collisional acceleration as
$`C`$ $``$ $`{\displaystyle \frac{P}{\mathrm{\Sigma }}}`$ (11)
$``$ $`\pm {\displaystyle \frac{c_b^2}{\lambda (1q_e\mathrm{cos}f)}}\widehat{r}`$
where the upper (lower) sign applies to the outer (inner) edge. Here the pressure gradient is taken to vary inversely as the separation between streamlines. Insert $`C`$ into Gauss’s perturbation equation for $`d\stackrel{~}{\omega }/dt`$ and average over true anomaly:
$`{\displaystyle \frac{d\stackrel{~}{\omega }}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{\pi nae}}{\displaystyle _0^\pi }C\mathrm{cos}fdf`$ (12)
$`=`$ $`q_eH(q_e^2)c_b^2/\lambda nae.`$
For the constant $`q_e`$ ring models presented here,
$$\lambda =c_b/n,$$
(13)
so that the only remaining free parameter is $`c_b`$. Equation (13) is relaxed for §4.2.2 and §4.3.2.
Note that this prescription for $`\mathrm{\Delta }_jd\stackrel{~}{\omega }/dt_C`$ ignores the decrease in velocity dispersion from $`c_b`$ at the ring edge to $`c_i`$ in the ring interior. The decline in velocity dispersion occurs over a length scale of order $`w_r`$. This length scale is large compared to $`c_b/n`$ so that the gradient of velocity dispersion does not give rise to a significant radial acceleration.
For a given value of $`c_b`$, equations (7), (8), (9), (10), and (13) comprise $`N`$ equations in $`N`$ unknowns $`\{m_j\}`$.<sup>1</sup><sup>1</sup>1Reflection symmetry about the ring midline reduces the number of equations necessary to $`N/2`$. Typically $`N2000`$ wires are needed to converge to within 10% of the solution for $`N\mathrm{}`$. Solutions for surface density profiles at quadrature for various values of $`c_b`$ are displayed in Figure 1, for parameters appropriate to the $`ϵ`$ and $`\alpha `$ rings; models for the $`\beta `$ ring are nearly identical to those of the $`\alpha `$ ring. In CMSG models, higher surface densities near ring edges are evident, as are higher total ring masses.
## 4 DISCUSSION
### 4.1 Surface Density Profiles and Torque Balance
Simple CMSG models, while not fully realistic, demonstrate the existence of a new class of self-gravity solution, that obtained by accounting for the modification of ring boundary conditions by interparticle collisions. Remarkably, forces felt by material in the last $``$100 m of a $``$10 km wide ring can increase equilibrium masses by factors up to 100.
Large S-band opacities measured by Voyager, which are incompatible with SSG surface densities (see, e.g., the review by Esposito et al. 1991), can be reconciled with average CMSG surface densities of $``$75–100$`\mathrm{g}\mathrm{cm}^2`$ for the $`ϵ`$, $`\alpha `$, and $`\beta `$ rings. Moreover, CMSG models predict that surface densities near ring edges are higher than those in the interior. This behavior is reminiscent of the “double-dip” structure seen in occultation light curves for the $`ϵ`$ and $`\alpha `$ rings (see, e.g., the review by French et al. 1991).
Greater ring masses as implied by CMSG resolve problems associated with exospheric drag that were pointed out by GP for rings $`\alpha `$ and $`\beta `$. For the remainder of this subsection, numerical estimates will be made for the $`\alpha `$ ring; similar conclusions hold for the $`\beta `$ and $`ϵ`$ rings. Surface densities are scaled to a typical CMSG value in the ring interior of $`\mathrm{\Sigma }=75\mathrm{g}\mathrm{cm}^2`$. An inner shepherd satellite exerts a repulsive, non-linear torque at first-order Lindblad resonances of magnitude
$$T_{\mathrm{nl}}^L\frac{10\rho _sR_s^3\mathrm{\Sigma }^2n^2a^7}{M_p^2d}6\times 10^{17}\left(\frac{\mathrm{\Sigma }}{75\mathrm{g}\mathrm{cm}^2}\right)^2\left(\frac{R_s}{10\mathrm{km}}\right)^3\left(\frac{\rho _s}{1.5\mathrm{g}\mathrm{cm}^3}\right)\left(\frac{500\mathrm{km}}{d}\right)\mathrm{erg},$$
(14)
where the satellite radius, $`R_s`$, is scaled to the Voyager upper limit of 10 km (Smith et al. 1986). The shepherding torque exceeds the magnitude of the drag torque exerted by the Uranian exosphere,
$$T_d4\pi m_Hn_Hv_Tna^3\mathrm{\Delta }a4\times 10^{16}\left(\frac{n_H}{10^3\mathrm{cm}^3}\right)\mathrm{erg}.$$
(15)
Here $`n_H=7\times 10^6e^{32.4R_p/a}\mathrm{cm}^3`$ is the number density of hydrogen atoms of mass $`m_H`$ in the exosphere, and $`v_T1\mathrm{km}\mathrm{s}^1`$ is their thermal speed normal to the orbital plane (Broadfoot et al. 1986). That $`T_{\mathrm{nl}}^L>|T_d|`$ ensures that the inner shepherd prevents ring particles from spiraling in towards Uranus.
Estimates of viscous torques $`T_v`$ also require revision. For a ring undergoing Keplerian shear, with minimum kinematic viscosity $`n(\mathrm{\Sigma }/\rho )^2`$, the viscous torque is given by
$$T_v\frac{3\pi n^2\mathrm{\Sigma }^3a^2}{\rho ^2}2.5\times 10^{18}\left(\frac{\mathrm{\Sigma }}{75\mathrm{g}\mathrm{cm}^2}\right)^3\left(\frac{1.5\mathrm{g}\mathrm{cm}^3}{\rho }\right)^2\mathrm{erg}$$
(16)
(GP). That $`T_v|T_d|`$ ensures that ring particles on the outer edge press against the inner Lindblad resonance established by the outer shepherd.
Conclusions drawn from comparisons between $`T_{\mathrm{nl}}^L`$ and $`T_v`$ are on less sure footing. For the choice of scaling parameters, the latter exceeds the former, contrary to the requirement of the standard theory of shepherding that the torques be equal. This might be construed as evidence that the angular momentum luminosity in the ring interior is reduced below $`T_v`$ by the non-Keplerian shear associated with a non-zero $`q_e`$ (Borderies, Goldreich, & Tremaine 1982; GP). However, the numerical estimates for the two torques differ only by a factor of a few. The shepherding torque should be evaluated using surface densities near the edge, which CMSG predicts are higher than those in the interior; this would increase the estimate of $`T_{\mathrm{nl}}^L`$. Uncertainties in the choice of parameters preclude drawing any conclusion other than that these torques are of the same order of magnitude.
### 4.2 Relative Importance of Planetary Oblateness
#### 4.2.1 $`J_2=0`$ vs. $`J_20`$
What does CMSG predict if $`J_2=0`$? Figure 2a displays the answer for the $`ϵ`$ ring, for $`c_b=2`$ and 3 $`\mathrm{cm}\mathrm{s}^1`$. In contrast to SSG, a non-vanishing equilibrium surface density does not require a finite planetary oblateness; self-gravity can be balanced entirely by collisional pressure gradients. For the $`\alpha `$ and $`\beta `$ rings, solutions with and without $`J_2`$ are practically indistinguishable for $`c_b0.5\mathrm{cm}\mathrm{s}^1`$. The influence of $`J_2`$ on the equilibrium solution diminishes as $`\mathrm{\Delta }a`$ decreases or as $`c_b`$ increases.
#### 4.2.2 Empirical Scaling Relations for $`J_2=0`$
For $`J_2=0`$ and fixed ring geometry, the surface density at quadrature near a given edge scales as
$$\mathrm{\Sigma }_b(0|x|\lambda )=\frac{c_b^2}{G\lambda }f(|x|/\lambda ),$$
(17)
where $`|x|`$ measures distance from the edge, $`c_b`$ and $`\lambda `$ are the same free parameters as in Equation (10), and $`f`$ is a dimensionless function of the similarity variable $`|x|/\lambda `$. Well away from ring edges, the surface density at quadrature scales as
$$\mathrm{\Sigma }_i(|x|\lambda )=\frac{c_b^2}{G\sqrt{\lambda \mathrm{\Delta }a}}g(|x|/\mathrm{\Delta }a),$$
(18)
where $`g`$ is another dimensionless function. The total ring mass scales as
$$M\frac{c_b^2a}{G}\sqrt{\mathrm{\Delta }a/\lambda }.$$
(19)
### 4.3 Value of $`q_e`$
#### 4.3.1 Sign of $`q_e`$
Figure 2b displays a CMSG model for the $`ϵ`$ ring obtained by reversing the sign of $`q_e`$. In contrast to SSG, a positive eccentricity gradient is not necessary in CMSG to obtain an equilibrium solution. This resurrects the problem of why all known eccentric planetary rings, including the Titan and Huygens ringlets around Saturn, are narrowest at periapse and widest at apoapse.
It is possible that equilibria obtained using $`q_e<0`$ are unstable. To address this issue, a preliminary investigation of ring stability for an $`N=4`$ ringlet model has been undertaken. Forces due to pressure gradients are included only for the first and fourth ringlets. Collisional accelerations are treated as if they arise from anti-self-gravity forces (self-gravity with the sign of the acceleration reversed); i.e., collisional shear stresses are ignored. In this crude approximation, equilibria are found to be stable regardless of the sign of $`q_e`$; small deviations from equilibrium masses result in apsidal librations (Borderies, Goldreich & Tremaine 1983). It remains to be seen whether collisional shear stresses alter stability properties.
Another possibility is that initial conditions set the sign of $`q_e`$. If the ring were initially uniform in width as a function of azimuth, then planetary oblateness would determine the initial sense of differential precession within the ring. The resultant narrowing of the ring width near a true anomaly of $`f=\pi /2`$ would cause a positive eccentricity gradient to grow by self-gravity. Under this hypothesis, an $`N=2`$ ringlet model incorporating forces from self-gravity and planetary oblateness yields the following time evolution for the apse and eccentricity differences between outer and inner ringlets:
$`\delta \stackrel{~}{\omega }=A\mathrm{sin}\mathrm{\Omega }_{\mathrm{lib}}t`$ (20)
$`\delta e=Ae(1\mathrm{cos}\mathrm{\Omega }_{\mathrm{lib}}t)`$ (21)
where $`A>0`$ and $`\mathrm{\Omega }_{\mathrm{lib}}`$ are the amplitude and frequency, respectively, of libration (cf. Borderies, Goldreich, & Tremaine 1983). Note that the time-average of $`\delta e`$ is positive. Inelastic collisions would damp librations and the ring would eventually settle into an equilibrium for which $`q_e>0`$.
#### 4.3.2 Magnitude of $`q_e`$ Near Ring Boundaries
It has been assumed that the eccentricity gradient, $`q_e`$, is finite out to the last $`\lambda =c_b/n50`$ meters of ring material. A finite $`q_e`$ is necessary to generate a non-zero azimuthal average of the collisional acceleration \[see equation (12)\]. The simple quantitative model of §3 employed the observed value of $`q_e`$ averaged over the entire ring width. The true value over the last few hundred meters of ring material is unknown.
In the case of the best-studied $`ϵ`$ ring, Graps et al. (1995) combined Voyager photopolarimeter and radio occultation measurements to infer the eccentricity gradient as a function of semi-major axis. They found that $`q_e`$ decreases over the last $``$5 km from its nearly constant value of $``$0.65 in the interior to $``$0.35 near the edge. The radial resolution of their study was between 1 and 2 km.
A decrease in $`q_e`$ towards ring boundaries is theoretically plausible. Distortions in a circular ring can be described by the change in separation, $`\delta r`$, between neighboring streamlines of the form
$$\delta r\mathrm{cos}m(\varphi \mathrm{\Omega }_{\mathrm{pat}}t),$$
(22)
where $`\mathrm{\Omega }_{\mathrm{pat}}`$ is the pattern speed of the distortion and $`m`$ is an integer. A constant $`q_e`$ ring that precesses rigidly in the quadrupole field of the central planet is equivalent to a distorted circular ring for which $`m=1`$ and $`\mathrm{\Omega }_{\mathrm{pat}}=d\stackrel{~}{\omega }/dt_Q`$. Resonant satellite perturbations, which enhance velocity dispersions within a distance $`w_r1\mathrm{km}`$ of ring edges, are characterized by much higher values of $`m=2a/3d1`$ and $`\mathrm{\Omega }_{\mathrm{pat}}=\mathrm{\Omega }_s`$. Satellite-induced disturbances might therefore reduce the local value of $`q_e`$. A decrease in $`q_e`$ over a distance $`w_r`$ near ring boundaries is roughly equivalent to setting $`\lambda =w_r`$ in equation (10). By the scaling relations (18) and (19), this would reduce surface densities and total ring masses shown in Figure 1 by a factor of $`\sqrt{w_rn/c_b}4`$.
## 5 DIRECTIONS FOR FUTURE RESEARCH
This work is primarily a demonstration that interparticle collisions near ring boundaries play a crucial role in determining ring masses under the self-gravity hypothesis. The nature of ring boundary conditions has not been calculated in rigorous detail; instead a prescription motivated by order-of-magnitude arguments is provided for collision-induced precession rates. Numerical simulations incorporating shepherd satellites will help to determine the actual 3-dimensional collisional stress tensor and eccentricity gradient everywhere within the ring.
Why all narrow eccentric rings surrounding Uranus and Saturn are narrowest at periapse and widest at apoapse remains to be understood. Stability analyses incorporating collisional shear stresses may reveal that rings having $`q_e<0`$ are unstable. Alternatively, the sign of $`q_e`$ may be set by initial conditions. Scenarios for ring formation—e.g., the catastrophic disruption of a small moon—require further elucidation.
Viscous damping gives rise to small differences between apsidal angles of neighboring streamlines (Borderies, Goldreich, & Tremaine 1983). For a given apsidal shift of $`\delta \stackrel{~}{\omega }1`$, the difference between the azimuth of maximum streamline separation and the azimuth of apoapse is given by the “pinch angle”, $`\delta \varphi =\mathrm{arctan}(e\delta \stackrel{~}{\omega }/\delta e)\delta \stackrel{~}{\omega }`$. The pinch angles calculated by Borderies et al. (1983) for their $`N=2`$ streamline models of the Uranian and Saturnian ringlets are suspect, however, because they neglect the boundary effects highlighted in the present work. A careful calculation of $`\delta \varphi (a)`$ that incorporates viscous drag and the global effects of resonant forcing by shepherd satellites has yet to be performed. Upcoming observations of narrow Saturnian ringlets by the Cassini Orbiter might test the predictions of such a calculation, thereby furnishing a powerful diagnostic of stresses within ringlets.
Financial support for this research was provided by NSF grant 94-14232 and by a Caltech Kingsley Foundation Fellowship held by E.C.
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# R-parity Violating Contribution to Neutron EDM at One-loop Order
## Abstract
We present the full result for the down squark mass-squared matrix in the complete theory of supersymmetry without R-parity where all kind of R-parity violating terms are admitted without bias. An optimal parametrization, the single-VEV parametrization, is used. The major result is a new contribution to $`LR`$ squark mixing, involving both bilinear and trilinear R-parity violating couplings. Among other things, the latter leads to neutron electric dipole moment at one-loop level. Similiar mechanism leading to electron electric dipole moment at the same level. We present here a short report on major features of neutron electric dipole moment from supersymmetry without R-parity and give the interesting constraints obtained.
preprint: IPAS-HEP-k004 Apr 2000 rev. Sep 2000
Introduction. The minimal supersymmetric standard model (MSSM) is no doubt the most popular candidate theory for physics beyond the Standard Model (SM). The alternative theory with a discrete symmetry called R-parity not imposed deserves no less attention. A complete theory of supersymmetry (SUSY) without R-parity, where all kind of R-parity violating (RPV) terms are admitted without bias, is generally better motivated than ad hoc versions of RPV theories. The large number of new parameters involved, however, makes the theory difficult to analyze. It has been illustrated that an optimal parametrization, called the single-VEV parametrization, can be of great help in making the task manageable. Here in this letter, we use the formulation to present the full result for the down squark mass-squared matrix. The major result is a new contribution to $`LR`$ squark mixing, involving both bilinear and trilinear RPV couplings. The interesting physics implications of this new contribution are discussed. Among such issues, we focus here on the RPV contribution to neutron electric dipole moment (EDM) at one-loop level.
Neutron and electron EDM’s are important topics for new CP violating physics. Within MSSM, studies on the plausible EDM contributions lead to the so called SUSY-CP problem. In the domain of R-parity violation, two recent papers focus on the contributions from the trilinear RPV terms and conclude that there is no contribution at the 1-loop level. Perhaps it has not been emphasized enough in the two papers that they are not studying the complete theory of SUSY without R-parity. It is interesting to see in the latter case that there is in fact contribution at 1-loop level, as discussed below. We would like to emphasize again that the new contribution involves both bilinear and trilinear (RPV) couplings. Since various other RPV scenarios studied in the literature typically admit only one of the two types of couplings, the contribution has not been previously identified.
The most general renormalizable superpotential for the supersymmetric SM (without R-parity) can be written as
$`W`$ $`=`$ $`\epsilon _{ab}[\mu _\alpha \widehat{H}_u^a\widehat{L}_\alpha ^b+h_{ik}^u\widehat{Q}_i^a\widehat{H}_u^b\widehat{U}_k^C+\lambda _{\alpha jk}^{}\widehat{L}_\alpha ^a\widehat{Q}_j^b\widehat{D}_k^C`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2}}\lambda _{\alpha \beta k}\widehat{L}_\alpha ^a\widehat{L}_\beta ^b\widehat{E}_k^C]+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{\prime \prime }\widehat{U}_i^C\widehat{D}_j^C\widehat{D}_k^C,`$ (2)
where $`(a,b)`$ are $`SU(2)`$ indices, $`(i,j,k)`$ are the usual family (flavor) indices, and $`(\alpha ,\beta )`$ are extended flavor index going from $`0`$ to $`3`$. In the limit where $`\lambda _{ijk},\lambda _{ijk}^{},\lambda _{ijk}^{\prime \prime }`$ and $`\mu _i`$ all vanish, one recovers the expression for the R-parity preserving case, with $`\widehat{L}_0`$ identified as $`\widehat{H}_d`$. Without R-parity imposed, the latter is not a priori distinguishable from the $`\widehat{L}_i`$’s. Note that $`\lambda `$ is antisymmetric in the first two indices, as required by the $`SU(2)`$ product rules, as shown explicitly here with $`\epsilon _{12}=\epsilon _{21}=1`$. Similarly, $`\lambda ^{\prime \prime }`$ is antisymmetric in the last two indices, from $`SU(3)_C`$.
R-parity is exactly an ad hoc symmetry put in to make $`\widehat{H}_d`$ stand out from the other $`\widehat{L}_i`$’s. It is defined in terms of baryon number, lepton number, and spin as, explicitly, $`=(1)^{3B+L+2S}`$. The consequence is that the accidental symmetries of baryon number and lepton number in the SM are preserved, at the expense of making particles and superparticles having a categorically different quantum number, R-parity. The latter is actually not the most effective discrete symmetry to control superparticle mediated proton decay, but is most restrictive for term admitted in the superpotential.
Doing phenomenological studies without specifying a choice of flavor bases is ambiguous. It is like doing SM quark physics with 18 complex Yukawa couplings, instead of the 10 real physical parameters. As far as the SM itself is concerned, the extra 26 real parameters are simply redundant, and attempts to relate the full 36 parameters to experimental data will be futile. In SUSY without R-parity, the choice of an optimal parametrization mainly concerns the 4 $`\widehat{L}_\alpha `$ flavors. Under the SVP, flavor bases are chosen such that : 1/ among the $`\widehat{L}_\alpha `$’s, only $`\widehat{L}_0`$, bears a VEV, i.e. $`\widehat{L}_i0`$; 2/ $`h_{jk}^e(\lambda _{0jk})=\frac{\sqrt{2}}{v_0}\mathrm{diag}\{m_1,m_2,m_3\}`$; 3/ $`h_{jk}^d(\lambda _{0jk}^{}=\lambda _{j0k})=\frac{\sqrt{2}}{}v_0\mathrm{diag}\{m_d,m_s,m_b\}`$; 4/ $`h_{ik}^u=\frac{v_u}{\sqrt{2}}V_{\text{CKM}}^T\mathrm{diag}\{m_u,m_c,m_t\}`$, where $`v_0\sqrt{2}\widehat{L}_0`$ and $`v_u\sqrt{2}\widehat{H}_u`$. The big advantage of here is that the (tree-level) mass matrices for all the fermions do not involve any of the trilinear RPV couplings, though the approach makes no assumption on any RPV coupling including even those from soft SUSY breaking; and all the parameters used are uniquely defined, with the exception of some removable phases. In fact, the (color-singlet) charged fermion mass matrix is reduced to the simple form :
$$_C=\left(\begin{array}{ccccc}M_2& \frac{g_2v_0}{\sqrt{2}}& 0& 0& 0\\ \frac{g_2v_u}{\sqrt{2}}& \mu _0& \mu _1& \mu _2& \mu _3\\ 0& 0& m_1& 0& 0\\ 0& 0& 0& m_2& 0\\ 0& 0& 0& 0& m_3\end{array}\right).$$
(3)
Readers are referred to Ref. for details conerning the RPV effects on the leptons.
Squark mixing and EDM. The soft SUSY breaking part of the Lagrangian can be written as follows :
$`V_{\mathrm{soft}}`$ $`=`$ $`ϵ_{ab}B_\alpha H_u^a\stackrel{~}{L}_\alpha ^b+ϵ_{ab}\left[A_{ij}^U\stackrel{~}{Q}_i^aH_u^b\stackrel{~}{U}_j^C+A_{ij}^DH_d^a\stackrel{~}{Q}_i^b\stackrel{~}{D}_j^C+A_{ij}^EH_d^a\stackrel{~}{L}_i^b\stackrel{~}{E}_j^C\right]+\mathrm{h}.\mathrm{c}.`$ (4)
$`+`$ $`ϵ_{ab}\left[A_{ijk}^\lambda ^{}\stackrel{~}{L}_i^a\stackrel{~}{Q}_j^b\stackrel{~}{D}_k^C+{\displaystyle \frac{1}{2}}A_{ijk}^\lambda \stackrel{~}{L}_i^a\stackrel{~}{L}_j^b\stackrel{~}{E}_k^C\right]+{\displaystyle \frac{1}{2}}A_{ijk}^{\lambda ^{\prime \prime }}\stackrel{~}{U}_i^C\stackrel{~}{D}_j^C\stackrel{~}{D}_k^C+\mathrm{h}.\mathrm{c}.`$ (5)
$`+`$ $`\stackrel{~}{Q}^{}\stackrel{~}{m}_Q^2\stackrel{~}{Q}+\stackrel{~}{U}^{}\stackrel{~}{m}_U^2\stackrel{~}{U}+\stackrel{~}{D}^{}\stackrel{~}{m}_D^2\stackrel{~}{D}+\stackrel{~}{L}^{}\stackrel{~}{m}_L^2\stackrel{~}{L}+\stackrel{~}{E}^{}\stackrel{~}{m}_E^2\stackrel{~}{E}+\stackrel{~}{m}_{H_u}^2|H_u|^2`$ (7)
$`+{\displaystyle \frac{M_1}{2}}\stackrel{~}{B}\stackrel{~}{B}+{\displaystyle \frac{M_2}{2}}\stackrel{~}{W}\stackrel{~}{W}+{\displaystyle \frac{M_3}{2}}\stackrel{~}{g}\stackrel{~}{g}+\mathrm{h}.\mathrm{c}.,`$
where we have separated the R-parity conserving ones from the RPV ones ($`H_d\widehat{L}_0`$) for the $`A`$-terms. Note that $`\stackrel{~}{L}^{}\stackrel{~}{m}_{\stackrel{~}{L}}^2\stackrel{~}{L}`$, unlike the other soft mass terms, is given by a $`4\times 4`$ matrix. Explicitly, $`\stackrel{~}{m}_{L_{00}}^2`$ is $`\stackrel{~}{m}_{H_d}^2`$ of the MSSM case while $`\stackrel{~}{m}_{L_{0k}}^2`$’s give RPV mass mixings.
We have illustrated above how the SVP keeps the expressions for the down-quark and color-singlet charged fermion mass matrices simple. The SVP performs the same trick to the corresponding scalar sectors as well. Here, we concentrate on the down-squarks. We have the mass-squared matrix as follows :
$$_D^2=\left(\begin{array}{cc}_{LL}^2& _{RL}^2\\ _{RL}^2& _{RR}^2\end{array}\right),$$
(8)
where $`_{LL}^2`$ and $`_{RR}^2`$ are the same as in MSSM while
$$(_{RL}^2)^T=A^D\frac{v_0}{\sqrt{2}}m_D\mu _0^{}\mathrm{tan}\beta (\mu _i^{}\lambda _{ijk}^{})_{}\frac{v_u}{\sqrt{2}}.$$
(9)
Here, $`m_D`$ is the down-quark mass matrix, which is diagonal under the parametrization adopted; $`(\mu _i^{}\lambda _{ijk}^{})_{}`$ denotes the $`3\times 3`$ matrix $`()_{jk}`$ with elements listed; and $`\mathrm{tan}\beta =\frac{v_u}{v_0}`$. Note that in the equation for $`(_{RL}^2)^T`$, we can write the first, $`A`$-term, contribution as
$$A^D\frac{v_0}{\sqrt{2}}=A_dm_D+\delta A^D\frac{v_0}{\sqrt{2}}$$
(10)
with $`A_d`$ being a constant (mass) parameter representing the “proportional” part. The remaining terms in $`(_{RL}^2)^T`$ are $`F`$-term contributions; in particular, the last term is a “SUSY conserving” but R-parity violating contributions given here for the first time. In fact, contributions to $`LR`$ scalar mixing of this type, for the sleptons, is first identified in a recent paper where their role in the SUSY analog of the Zee neutrino mass model is discussed. In a parallel paper by one of the authors (O.K.), a systematic analysis of the full squark and slepton masses as well as their contributions, through $`LR`$ mixings, to 1-loop neutrino masses are also presented. Here, we focus only on the down-quark sector. Note that the full $`F`$-term part in the above equation can actually be written together as $`(\mu _\alpha ^{}\lambda _{\alpha jk}^{})_{}\frac{v_u}{\sqrt{2}}`$ where the $`\alpha =0`$ term gives the second term in RHS of Eq.(9), which is the usual $`\mu `$-term contribution in the MSSM case. The latter is, however, diagonal, i.e. vanishes for $`jk`$. We would like to emphasize that the above result is complete — all RPV contributions are included. The simplicity of the result is a consequence of the SVP. Explicitly, the RPV $`A`$-terms contributions \[c.f. Eq.(10)\] vanish as $`v_i\sqrt{2}\widehat{L}_i=0`$
The $`(\mu _i^{}\lambda _{ijk}^{})_{}`$ term is very interesting. It involves only parameters in the superpotential and has nothing to do with soft SUSY breaking. Without an underlining flavor theory, there is no reason to expect any specific structure among different terms of the matrix. In particular, the off-diagonal terms ($`jk`$) may have an important role to play. They contribute to flavor changing neutral current (FCNC) processes such as $`bs\gamma `$, a topic to be addressed in a later publication. Moreover, both the $`\mu _i`$’s and the $`\lambda _{ijk}^{}`$’s are complex parameters. Hence, diagonal terms in $`(\mu _i^{}\lambda _{ijk}^{})_{}`$ also bear CP-violating phases and contribute to electric dipole moments (EDM’s) of the corresponding quarks. In particular, $`\mu _i^{}\lambda _{i11}^{}`$ gives contribution to neutron EDM at 1-loop level, in exactly the same fashion as the $`A`$-term in MSSM does. The similar term in $`LR`$ slepton mixing gives rise to electron EDM. This result is in direct contrary to the impression one may get from reading the two recent papers on the subject. One should bear in mind that the two papers do not put together both the bilinear and the trilinear RPV terms. Our treatment here, bases on the SVP, gives, for the first time, the result of squark masses for the complete theory of SUSY without R-parity. Going from here, obtaining the EDM contributions is straight forward.
Contribution to EDM of the $`d`$ quark at 1-loop level, from a gaugino loop with $`LR`$-squark mixing in particular (see Fig. 1), has been widely studied within MSSM. With the squark mixings in the down-sector parametrized by $`\delta _{jk}^D`$ (normalized by average squark mass as explicitly shown below), we have the neutron EDM result given by
$$d_n=\frac{8}{27}\frac{e\alpha _s}{\pi }\frac{M_{\stackrel{~}{g}}}{M_{\stackrel{~}{d}}^2}\text{Im}(\delta _{11}^D)F_1\left(\frac{M_{\stackrel{~}{g}}^2}{M_{\stackrel{~}{d}}^2}\right)$$
(11)
where $`M_{\stackrel{~}{g}}`$ and $`M_{\stackrel{~}{d}}`$ are the gluino and down squark masses respectively, and
$$F_1(x)=\frac{1}{(1x)^3}\left(\frac{1+5x}{2}+\frac{2+x}{1x}\mathrm{ln}x\right).$$
(12)
Contribution of $`\mu _i^{}\lambda _{i11}^{}`$ to $`\delta _{11}^D`$ is to be given as
$`\mu _i^{}\lambda _{i11}^{}{\displaystyle \frac{v_u}{\sqrt{2}}}{\displaystyle \frac{1}{M_{\stackrel{~}{d}}^2}}.`$Requiring the contribution alone not to upset the experimental bound on neutron EDM : $`(d_n)^{\text{exp}}<6.310^{26}e\text{cm}`$, a bound can be obtained for the RPV parameters. Note that going from $`d`$ quark EDM to neutron EDM, we assume the simple valence quark model. Taking $`M_{\stackrel{~}{d}}=100\text{GeV}`$ and $`M_{\stackrel{~}{g}}=300\text{GeV}`$ gives the bound
$$\text{Im}(\mu _i^{}\lambda _{i11}^{})10^6\text{GeV},$$
(13)
(with $`v_u200\text{GeV}`$). This result is interesting. Let us first concentrate on the $`i=3`$ part, assuming the $`i=1`$ and $`2`$ contribution to be subdominating. Imposing the $`18.2\text{MeV}`$ experimental bound for the mass of $`\nu _\tau `$ still admits a relatively large $`\mu _3`$, especially for a large $`\mathrm{tan}\beta `$. Reading from the results in Ref., the bound is $`7\text{GeV}`$ at $`\mathrm{tan}\beta =2`$ and $`300\text{GeV}`$ at $`\mathrm{tan}\beta =45`$, while the best bound on the corresponding $`\lambda _{311}^{}`$ (from $`\tau \pi \nu `$) is around $`0.050.1`$.
Here, an explicit comparison with the corresponding R-parity conserving contribution is of interest. From Eqs.(9) and (10), it is obvious that we are talking about $`(A_d\mu _0^{}\mathrm{tan}\beta )m_d`$ verses $`\mu _i^{}\lambda _{i11}^{}\frac{v_u}{\sqrt{2}}`$. Both $`A_d`$ and $`\mu _0`$ are expected to be roughly at the same order as $`v_u`$, i.e. at electroweak scale. We are hence left to compare $`m_d`$ ($`10^3\text{GeV}`$) with $`\mu _i^{}\lambda _{i11}^{}`$. The above discussion then leads to the conclusion that the RPV part could easily be larger by one or even two orders of magnitude.
On the other hand, if one insists on a sub-eV mass for $`\nu _\tau `$ as suggested, but far from mandated, by the result from the Super-Kamiokande (super-K) experiment, we would have $`\mu _3\mathrm{cos}\beta 10^4\text{GeV}`$. This means that at least for the large $`\mathrm{tan}\beta `$ case, the EDM bound as given by Eq.(13) still worths notification, even under this most limiting scenario.
The $`\mu _i^{}\lambda _{i11}^{}`$ contribution to squark mixing, as well as $`\lambda _{i11}^{}`$ in itself together with an $`A`$-term mixing, also gives rise to neutrino mass at 1-loop. Hence, to consistently impose the super-K sub-eV neutrino mass scenario, one should also check the corresponding bound obtained. We are interested here in whether these will further weaken the implication of the EDM bound discussed here. Fig. 2 shows a familiar quark-squark loop neutrino mass diagram. We are interested here in the case where both the $`\lambda ^{}`$-couplings are $`\lambda _{311}^{}`$. We have, for the R-parity conserving $`LR`$ squark mixing, the familiar result
$$m_{\nu _\tau }\frac{3}{8\pi ^2}m_d^2\frac{(A_d\mu _0^{}\mathrm{tan}\beta )}{M_{\stackrel{~}{d}}^2}\lambda _{311}^{}.$$
(14)
However, with the full $`LR`$ mixing result as given in Eq.(9), there is an extra contribution to be given as
$$\frac{3}{8\pi ^2}m_d\frac{v_u}{\sqrt{2}}\frac{\mu _i^{}\lambda _{i11}^{}}{M_{\stackrel{~}{d}}^2}\lambda _{311}^{\mathrm{\hspace{0.17em}2}}.$$
(15)
The latter type of RPV contribution to neutrino masses has not been identified before (see however Refs and ).
From Eq.(14), one can easily see that the requirement for $`m_{\nu _\tau }`$ to be at the super-K atmospheric neutrino oscillation scale only gives a bound for $`\lambda _{311}^{}`$ of about the same magnitude as one quoted above, from the other sources. As for the contribution \[Eq.(15)\], the bound given by Eq.(13) itself says the contribution is smaller than the previous one. Hence, neutrino mass contributions from Fig. 2 do not change our conclusion above.
Note that the EDM bound given by Eq.(13) actually involves a summation over index $`i`$. Results from Ref. indicated that while $`\mu _1`$ is very strongly bounded, the bound on $`\mu _2`$ could be not very strong. Moreover, the bound on $`\lambda _{211}^{}`$ is no better than that on $`\lambda _{311}^{\mathrm{\hspace{0.17em}2}}`$. Hence, the EDM bound may still be of interest there too. The story for imposing the super-K constraint is obviously the same as the above discussion for the $`i=3`$ case.
One should bear in mind that the EDM and the neutrino mass bounds involve different combinations of RPV parameters as well as with the other SUSY parameters. An exact comparison for bounds obtained from the two sources is hence difficult. Our discussion above is aimed at illustrating the fact that the EDM bound is not completely overshadowed by the super-K neutrino mass bound. In other word, even requiring the magnitudes for the RPV parameters to satisfy the most stringently interpreted super-K bounds does not make them so small that the above discussed contribution to neutron EDM will always be satisfied.
Beyond the gluino diagram. Similar RPV contributions on the neutron and electron EDM’s are obtained through neutralino exchange diagram. One simply has to replace the gluino in the diagram with the other neutral gauginos. In the neutron case, the gluino diagram contribution discussed here no doubt dominate, due to the much stronger QCD coupling.
There are other 1-loop contributions. In the case of MSSM, the chargino contribution is known to be competitive or even dominates over the gluino one in some regions of the parameter space. The major part of the chargino contribution comes from a diagram with a gauge and a Yukawa coupling for the loop vertices, with pure $`L`$-squark running in the loop. Here we give the corresponding formula generalized to the case of SUSY without R-parity. This is given by
$$\left(\frac{d_f}{e}\right)_{\chi ^{\text{ -}}}=\frac{\alpha _{\text{em}}}{4\pi \mathrm{sin}^2\theta _W}\underset{\stackrel{~}{f}^{}}{}\underset{n=1}{\overset{5}{}}\text{Im}(𝒞_{fn})\frac{M_{\chi _n^\text{-}}}{M_{\stackrel{~}{f}^{}}^2}\left[𝒬_{\stackrel{~}{f}^{}}B\left(\frac{M_{\chi _n^\text{-}}^2}{M_{\stackrel{~}{f}^{}}^2}\right)+(𝒬_f𝒬_{\stackrel{~}{f}^{}})A\left(\frac{M_{\chi _n^\text{-}}^2}{M_{\stackrel{~}{f}^{}}^2}\right)\right],$$
(16)
for $`f`$ being $`u`$ ($`d`$) quark and $`f^{}`$ being $`d`$ ($`u`$), where
$`𝒞_{un}`$ $`=`$ $`{\displaystyle \frac{y_u}{g_2}}𝑽_{2n}^{}𝒟_{d1}\left(𝑼_{1n}𝒟_{d1}^{}{\displaystyle \frac{y_d}{g_2}}𝑼_{2n}𝒟_{d2}^{}{\displaystyle \frac{\lambda _{i11}^{}}{g_2}}𝑼_{(i+2)n}𝒟_{d2}^{}\right),`$ (17)
$`𝒞_{dn}`$ $`=`$ $`({\displaystyle \frac{y_d}{g_2}}𝑼_{2n}+{\displaystyle \frac{\lambda _{i11}^{}}{g_2}}𝑼_{(i+2)n})𝒟_{u1}\left(𝑽_{1n}^{}𝒟_{u1}^{}{\displaystyle \frac{y_u}{g_2}}𝑽_{2n}^{}𝒟_{u2}^{}\right).`$ (18)
The terms in $`𝒞_{dn}`$ with only one factor of $`\frac{1}{g_2}`$ and a $`\lambda _{i11}^{}`$ gives the RPV analog of the dominating MSSM chargino contribution. The term is described by a diagram, which at first order requires a $`l_{L_i}^{\text{ -}}`$$`\stackrel{~}{W}^+`$ mass mixing. The latter vanishes, as shown in Eq.(3). From the full formula above, it is easy to see that the exact mass eigenstate result would deviate from zero only to the extent that the mass dependence of the $`B`$ and $`A`$ functions spoils the GIM like cancellation in the sum. The resultant contribution, however, is shown by our exact numerical calculation to be substantial. What is most interesting here is that an analysis through perturbational approximations illustrates that the contribution is proportional to, basically, the same combination of RPV parameters, i.e. $`\mu _i^{}\lambda _{i11}^{}`$. While we cannot give much of the details here (see Ref.), let us list numbers from a sample point for illustration : with $`A_u=A_d=500\text{GeV}`$, $`\mu _0=300\text{GeV}`$, $`\mathrm{tan}\beta =3`$, a common gaugino masses at $`300\text{GeV}`$, $`\stackrel{~}{m}_Q=200\text{GeV}`$, $`\stackrel{~}{m}_u=\stackrel{~}{m}_d=100\text{GeV}`$, $`\mu _3=1\times 10^4\text{GeV}`$, and $`\lambda _{311}^{}=0.1\times \text{exp}(i\pi /6)`$ (being the only complex parameter), we have the resulted neutron EDM contributions from gluino, chargino(-like), and neutralino(-like) 1-loop diagrams given by $`2.49`$, $`0.56`$, and $`0.056`$ times $`10^{27}e\text{cm}`$, respectively.
Summary. In summary, we have presented the complete result for LR squark mixing and analyzed its contribution to neutron EDM through the gluino diagram. The result provide interesting new bounds on RPV parameters. A brief discussion for the chargino(-like) 1-loop contribution is also given, together with a sample result from exact numerical calculations, including also the neutralino(-like) loop. The full details will be report in an publication.
We would also like to mention that there is the analogous case for the slepton mixing and electron EDM. The latter contributions, while having a similar structure, has potential complications from mixings with charged Higgs. The issue is under investigation.
Acknowledgment : Y.Y.K. wishes to thank M. Kobayashi and H.Y.Cheng for their hositality. His work was in part supported by the National Science Council of R.O.C. under the Grant No. NSC-89-2811-M-001-0053. O.K. wants to thank D. Chang for discussions.
Figure captions :
Fig. 1 — EDM for $`d`$ quark at 1-loop.
Fig. 2 — Neutrino mass at 1-loop.
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# 1 Introduction
## 1 Introduction
There are some proposals as constructive definitions of string theory. One of them is called “M-theory” which is limited to eleven-dimensional supergravity at low energy, and that corresponds to the strong coupling limit of the type-IIA superstring . More concretely, M-theory is described by the IIA matrix model, which governs the dynamics of D0-branes. One of the others is the IIB matrix model, which has been proposed by Ishibashi, Kawai, Kitazawa and Tsuchiya . Recently, the authors of have shown that the twisted reduced model can be interpreted as noncommutative Yang-Mills theory. They have obtained noncommutative Yang-Mills theory with D-brane backgrounds in the IIB matrix model. Furthermore, some of the same authors have calculated the Wilson loop in noncommutative Yang-Mills, and have found an open string-like object , whose length is $`|C^{\mu \nu }k_\nu |`$. $`C^{\mu \nu }`$ delineates the noncommutative scale. In , a non-local basis, called bi-local basis, is introduced to provide a simple description for high-momentum $`|k^\mu |>\lambda `$, where $`\lambda `$ is the spacing quanta and $`k^\mu `$ is the eigenvalue of adjoint $`\widehat{P}^\mu `$. As can be recognized above, some new mathematical methods are required to more exactly investigate superstring theory. One of new mathematics is noncommutative geometry, which was introduced to physicists by a mathematician, Alain Connes .
There are many papers in which authors apply noncommutative geometry to rebuilding ordinary physics. The initiative is as follows. A trial for constructing the standard model, $`U(1)\times SU(2)\times SU(3)`$, has been partially successed in the framework of noncommutative geometry . In the base manifold is $`E_4(M)\times Z_2`$ and the algebra is $`𝒜`$=$`C^{\mathrm{}}(M)\{M_2(C)M_1(C)\}`$. The acting Hilbert space is given by that of spinors of $`\left(\begin{array}{c}l\{\mathrm{doublet}\}\\ e\{\mathrm{singlet}\}\end{array}\right)`$. To introduce an $`SU(3)`$ gauge group for the quarks color, the authors have taken notice of the bimodule structure relating two algebras $`(𝒜,)`$, where $`𝒜`$=$`\mathrm{algebra}\mathrm{of}\mathrm{quaternions}C`$ and $``$=$`CM_3(C)`$. Finally, the Higgs field appears as a gauge field between two points represented by $`Z_2`$. From the unification of the standard model and gravitational theory, this work is very meaningful, because it gives a geometrical interpretation of the standard model.
However, the Connes-Lott formulation could not be beyond the standard model, as they have already described in the literature. For example, the generation matrix is given by $`C^{N_G}`$; namely, the authors could not introduce the inside structure of the generation. In A.Connes represents ultimate noncommutative space including the color symmetry $`SU(3)`$ as the 4-dimensional continuum times the finite space ‘F’. What is the quantum group of the finite space ‘F’ ? This is one of problems which he discusses on the last page of , which one should solve when trying to make up a beyond standard model on noncommutative geometry. We have already recognized that we could not excel A.Connes in mathematical abilities. Hence, we attack the open problem from the opposite side to mathematicians.
We expect that noncommutative geometry may better extract the potential abilities of the existing physical models, for instance the standard model, though there are still a few papers which have suggested that noncommutative geometry can expose a new thing in ordinary things. In the future we would like to determine the masses of quarks based on geometric quantities on some noncommutative algebraic structure. Our purpose is to propose a gravitational model coupled with an $`SU(2)`$ doublet complex Higgs-like field. In this paper we artificially elevate the ordinary Higgs couplings to some components of one-form basis for our noncommutative space. We then impose the unitarity condition and the torsion-less condition for the basis of our noncommutative differential geometry as the authors of have performed. Finally, in virtue of imposing these conditions, we will find that our Higgs-like couplings are represented by the Wilson loops of the spin connections of the noncommutative space and $`U(1)_{EM}`$ gauge fields. These gauge fields are expressed by the linear combinations of the $`U(1)_Y`$ and $`SU(2)_W`$ gauge fields. In section two we briefly review . In section three we will introduce our vielbein while comparing it to Chamseddine, Fr$`\ddot{o}`$lich, Grandjean’s metric. Our discussion is given in the last section. The unitarity conditions and the components of torsion are expressed in the appendices.
## 2 Review of Chamseddine, Fr$`\ddot{o}`$hlich, Grandjean metric
In this section we introduce in order to describe the underlying geometric structure of our proposal. The authors have studied the gravitational action of the noncommutative geometry underlying the above Connes-Lott construction of the standard model . Ultimately, they have found the Einstein-Hilbert action in the leptonic sector and one in the quark sector. In these actions they have shown that the distance function between the two points is determined by a real scalar field, $`\sigma `$, whose vacuum expectation value sets the weak scale.
They set up the noncomutative differential geometry on two copies of the four-dimensional continuum, $`M_4`$, defined by
$`M_4\times Z_2.`$ (1)
They choose the algebra $`𝒜`$ describing the above noncommutative space as
$`𝒜=(𝒜_1𝒜_2)C^{\mathrm{}}(M_4),`$ (2)
where $`𝒜_1=M_2(C)`$ and $`𝒜_2=C`$. They firstly consider the leptonic part and the Higgs sector of the standard model.
$`a𝒜`$ is represented by
$`a=\left(\begin{array}{cc}a_1& 0\\ 0& a_2\end{array}\right),`$ (5)
where $`a_1`$ (for $`𝒜_1`$),$`a_2`$ ( for $`𝒜_2`$) are the $`C^{\mathrm{}}`$ function on $`M_4`$. Notice that $`a_1`$ is a two times two matrix.
They define the representation space for $`𝒜`$ and the Hilbert space as follows:
$`=L^2(S_1,dv_1)L^2(S_2,dv_2),`$ (6)
where $`S_i=S_0V_i`$. $`S_0`$ is the bundle of Dirac spinors on $`M_4`$ and $`V_i`$ is a representation space for $`𝒜_i`$. They give the representation space as $`V_1=C^2,V_2=C`$. $`dv_i`$ is the volume form corresponding to a Riemannian metric on $`M_4`$.
They express the Dirac operator on the noncommutative space as
$`D=\left(\begin{array}{cc}_1\mathrm{𝟏}_2\mathrm{𝟏}_3& \gamma ^5𝐌_{12}𝐤\\ \gamma ^5𝐌_{12}^{}𝐤^{}& _2\mathrm{𝟏}_3\end{array}\right),`$ (9)
where $`_i`$ is defined by
$`_i=e_{ia}^\mu \gamma ^a(_\mu +i\omega _{i\mu }).`$ (10)
$`e_{ia}^\mu `$ is a vielbein of $`M_4`$. The index $`i=1,2`$ distinguishes $`𝒜_i`$ and $`\mu `$ represents the vector index on $`M_4`$; $`\mu `$=0,1,2,3; $`a`$ is the coordinate of the tangent bundle $`M_4`$; $`a`$=1,2,3,4. $`\omega _{i\mu }`$ is the spin connection in (1). $`\{\gamma ^a\}_{a=1}^4`$ are the anti-hermitian Euclidean Dirac matrices, with $`\{\gamma ^a,\gamma ^b\}=\gamma ^a\gamma ^b+\gamma ^b\gamma ^a=2\delta ^{ab}`$, $`\gamma ^5=\gamma ^1\gamma ^2\gamma ^3\gamma ^4`$, $`\gamma ^5=(\gamma ^5)^{}`$, and $`\gamma ^a=(\gamma ^a)^{}`$. $`\gamma ^{}`$ means the hermitian conjugate. We follow their notation in the section three. The components of the spin connection are chosen based on the Cartan structure equation of (1) as the Riemannian geometry. $`k`$ is a $`3\times 3`$ family mixing matrix. In the leptonic sector the detailed structure inside of $`k`$ is not necessary. $`M_{12}`$, called a doublet, is written as $`M_{12}=\left(\begin{array}{c}\alpha (x)\\ \beta (x)\end{array}\right).`$ Functions $`\alpha (x)`$ and $`\beta (x)`$ are restricted by the requirements the consistency of $`\mathrm{\Omega }_D^2(𝒜)`$, which is isomorphism to $`\pi (\mathrm{\Omega }^2(𝒜))/\mathrm{Aux}^2`$, where
$`\mathrm{Aux}^2=\{{\displaystyle \underset{i}{}}[D,a_i][D,b_i]=0:{\displaystyle \underset{i}{}}a_i[D,b_i]=0\}.`$ (11)
In other words, we should pull the following elements $`a_i,b_i`$ from $`𝒜`$:
$`\mathrm{If}\rho `$ $`=`$ $`{\displaystyle \underset{i}{}}a_idb_i\mathrm{ker}\pi \pi (\rho )=0,`$
$`\mathrm{then}\pi (d\rho )`$ $`=`$ $`{\displaystyle \underset{i}{}}[D,a_i][D,b_i]=0.`$ (12)
The authors find three possibilities which can satisfy Eq.(12), and in detail compute one of these in ,
$`M_{12}=\mathrm{exp}(\sigma )\left(\begin{array}{c}0\\ 1\end{array}\right),`$ (15)
where $`\sigma (x)`$ is a real scaler field.
They introduce a system of generators of $`\stackrel{~}{\mathrm{\Omega }}_D^1(𝒜)`$, $`\{E^A\}`$, which are suitable for the Hilbert space (6) and the representation space defined by themselves:
$`E^a`$ $`=`$ $`\gamma ^a\left(\begin{array}{cc}\mathrm{𝟏}_2& 0\\ 0& 1\end{array}\right),a=1,2,3,4,`$ (18)
$`E^r`$ $`=`$ $`\gamma ^5\left(\begin{array}{cc}\mathrm{𝟎}_2& ke_r\\ k^{}e_r^{}& 0\end{array}\right),r=5,6,`$ (21)
where $`e_5=\left(\begin{array}{c}1\\ 0\end{array}\right)`$, $`e_6=\left(\begin{array}{c}0\\ 1\end{array}\right),`$ and $``$ means the transposed matrix. <sup>1</sup><sup>1</sup>1The authors of define $`\stackrel{~}{\mathrm{\Omega }}(𝒜):=\stackrel{~}{\pi }(\mathrm{\Omega }^.(𝒜))`$, $`\stackrel{~}{\mathrm{\Omega }}_D^n(𝒜):=\stackrel{~}{\pi }(\mathrm{\Omega }^n(𝒜))/\stackrel{~}{\pi }(dJ_{n1})`$. Here $`J_n`$ is the intersection of the kernel which is given by Eq.(12). $`\stackrel{~}{\pi }`$ is a \*-representation of $`\mathrm{\Omega }^{}(𝒜)`$.
They list up the conditions of vanishing torsion $`T^A,A=1,\mathrm{},6`$, where $`T^A=dE^A+\mathrm{\Omega }^A{}_{B}{}^{}E_{}^{B}`$ and the unitarity conditions on (18) and (21). The components $`\mathrm{\Omega }^A_B`$ which are defined by
$`E^A=\mathrm{\Omega }^A{}_{B}{}^{}E^B,`$ (22)
are the connection coefficients on $`\stackrel{~}{\mathrm{\Omega }}_D^1(𝒜)`$ . After imposing these conditions, they declare that the field $`\sigma `$ becomes non-dynamical. They then weaken the condition of vanishing torsion as
$`Tr_kT^A=0`$ (23)
where $`Tr_k`$ is the trace over the family mixing matrix. As the result of Eq.(23), the $`\sigma `$ field behaves as a dynamical field in the Hilbert-Einstein action of the leptonic sector.
They also construct the gravitational action of the quark sector. It is the same form of action of the leptonic sector, but with different coefficients and with dependence on the generation-mixing matrices of the quarks masses. The action is given by
$`{\displaystyle d^4x\left[\frac{1}{2}(3c_l+4c_q)R+\alpha (_a\sigma )^2+\beta e^{2\sigma }\right]},`$ (24)
where $`c_l,c_q`$ are arbitrary constants and $`\alpha `$,$`\beta `$ depends on the structure of the generation matrices of the quark masses and $`c_l,c_q`$ . However, as the authors of described, since the gravitational action is non-renormalizable and nobody understands the quantum noncommutative geometry, these coefficients do not have any physical significance. $`R`$ is the scalar curvature in $`M_4\times Z_2`$.
In string theory we have known a theory whose action is almost the same form as (24). It is the Liouville field theory, where the action contains an exponential tachyon field, $`\mathrm{exp}(\alpha _{}x^1)`$. Here, $`x^1`$ is one-direction of the target space, and the effective string coupling is diverging at large $`x^1`$. $`\alpha _{}`$ is given by $`(\frac{26D}{6\alpha ^{}})^{1/2}(\frac{2D}{6\alpha ^{}})^{1/2}`$. For $`D>2`$ this is complex, hence this term oscillates, but for $`D2`$, $`\alpha _{}>0`$, we have a real exponential. <sup>2</sup><sup>2</sup>2In the discussion of the Liouville field theory, we have referred to the explanation on pp.323-327 of . The difference between (24) and the Liouville action is that the latter contains the term $`R\times x^1`$. The difference between the sign of the exponential interactions can be eliminated by redefining of the field $`\sigma `$ to $`\sigma `$. What we want to stress here is that in string theory we may express the degree of freedom, $`x^1`$, as a differential operator for some discrete space, since the authors of can gain (24), which is similar to the Liouville action, from the gravitational theory on $`M_4\times Z_2`$.
## 3 Our proposal- geometrized Higgs mechanism
An outline of our idea is that we would like to finally generalize an ordinary Higgs mechanism of the standard model to a gravitational theory on a noncommutative geometry. We attempt to raise the Higgs coupling from a constant to one of components of vielbein on noncommutative geometry, $`M_4\times Z_2`$. Furthermore, we make our Higgs couplings <sup>3</sup><sup>3</sup>3We call our Higgs couplings Higgs-like couplings. depend on the local coordinates of $`M_4`$. In the second subsection we impose the torsion-less condition and the unitarity condition on the vielbein and the spin connection of $`M_4\times Z_2`$, following . As a result, we find that the Higgs-like couplings can be calculated by the Wilson loops of the spin connection and the $`U(1)_{EM}`$ gauge.
### 3.1 Dynamical vielbein and localized Higgs couplings
We choose the algebra $`𝒜`$, defining the noncommutative space underlying our model as
$`𝒜=(𝒜_1𝒜_2𝒜_3)C^{\mathrm{}}(M_4),`$ (25)
where $`M_4`$ is a smooth, compact, four-dimensional Riemannian spin manifold; $`𝒜_1`$ is the algebra of complex $`2\times 2`$ matrices. $`𝒜_2`$ and $`𝒜_3`$ are $`C`$. Elements, $`a`$, of $`𝒜`$ are written as
$`a=\left(\begin{array}{ccc}a_1& 0& 0\\ 0& a_2& 0\\ 0& 0& a_3\end{array}\right),`$ (29)
where $`a_i`$ is a $`C^{\mathrm{}}`$-function on $`M_4`$ with values in $`𝒜_i`$, $`i=1,2,3`$.
The Hilbert space is defined as the spinors of the form $`L=\left(\begin{array}{c}u_L\\ d_L\\ u_R\\ d_R\end{array}\right)`$, where $`R,L`$ respectively mean the two kinds of chiralities which are defined by the four-dimensional $`\gamma ^5`$. We want to regard $`u_R,d_R`$ as an $`SU(2)`$ singlet and $`u_L,d_L`$ as an $`SU(2)`$ doublet. Here, we image the Hilbert space of one-family quarks.
We introduce the following basis as
$`^a`$ $`=`$ $`\gamma ^a\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),`$ (34)
$`^{\overline{r}}`$ $`=`$ $`\gamma ^5\left(\begin{array}{cc}\mathrm{𝟎}& \overline{𝒜}e_{\overline{r}}\\ e_{\overline{r}}^{}\overline{𝒜}& \mathrm{𝟎}\end{array}\right),`$ (37)
where
$`\overline{𝒜}`$ $`=`$ $`\left(\begin{array}{cc}0& \frac{v}{\sqrt{2}}\\ \frac{v}{\sqrt{2}}& 0\end{array}\right),`$ (40)
$`e_u=e_{\overline{u}}`$ $`=`$ $`\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),e_d=e_{\overline{d}}=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),`$ (45)
$`e_u^{}=e_{\overline{u}}^{}=e_d=e_{\overline{d}},e_d^{}=e_{\overline{d}}^{}=e_u=e_{\overline{u}}.`$ (46)
$`a=1,\mathrm{},4`$ and $`\overline{r}=\overline{u},\overline{d}`$. We call this basis flat basis. (3.25) and (3.26) in inspire us (18) and (37). $`v`$ is a constant. We would like to first treat (34) and (37) as a system of generators of the one-form $`\stackrel{~}{\mathrm{\Omega }}^1(𝒜)`$ though (34) and (37) contain auxiliary parts as we will explain in the section $`4`$.
Let us introduce a local basis $`e_\mu ^a(x)`$, $`\mu =1,\mathrm{},4`$ of orthonormal tangent vectors to $`M_4`$ for each $`a`$. $`a`$ is the local Lorentz index on the flat tangent plane. In curved $`M_4`$ we introduce gamma matrices by $`\gamma ^\mu e_a^\mu \gamma ^a`$ and define a curved base which we would like to regard as one-form of $`\stackrel{~}{\mathrm{\Omega }}^1(𝒜)`$,
$`^\mu =\gamma ^\mu \left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).`$ (51)
Moreover, let us consider remaining components for the curved base. These contain Higgs-like couplings $`f(x)`$, $`\stackrel{~}{f}(x)`$, which depend on the local coordinates of $`M_4`$ as
$`^r=\gamma ^5\left(\begin{array}{cc}\mathrm{𝟎}& 𝒜^{}e_r\\ e_r^{}𝒜& 0\end{array}\right),`$ (54)
where
$`𝒜`$ $`=`$ $`\left(\begin{array}{cc}0& f(x)\frac{v}{\sqrt{2}}\\ \stackrel{~}{f}(x)\frac{v}{\sqrt{2}}& 0\end{array}\right),𝒜^{}=\left(\begin{array}{cc}0& \stackrel{~}{f}^{}(x)\frac{v}{\sqrt{2}}\\ f^{}(x)\frac{v}{\sqrt{2}}& 0\end{array}\right).`$ (59)
We express the dual basis, following (3.29) of ,
$`\omega =\left(\begin{array}{cccc}\gamma ^\mu \omega _{1\mu ,ij}(x)& & \omega _{}^{}{}_{2}{}^{}(x)& 0\\ & & 0& \omega _2(x)\\ \omega _{}^{}{}_{1}{}^{}(x)& 0& \gamma ^\mu \omega _{2\mu ,ij}(x),& \\ 0& \omega _1(x)& & \end{array}\right),`$ (64)
where $`\omega _{1\mu ,ij}(x)`$ and $`\omega _{2\mu ,ij}(x)`$ are two-times-two matrices, $`i,j=1,2`$.
Following (3.28) and (3.29) in , we write down the following connection coefficients of (51) and (54), which have been defined in (22), as
$`\mathrm{\Omega }^A{}_{B}{}^{}=\left(\begin{array}{cc}\gamma ^\mu \omega _{1\mu }{}_{}{}^{A}{}_{B}{}^{},_{ij}& 𝒜\gamma ^5e^\sigma \left(\begin{array}{cc}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{A}_B& 0\\ 0& \omega _2{}_{}{}^{A}_B\end{array}\right)\\ \left(\begin{array}{cc}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{A}_B& 0\\ 0& \omega _1{}_{}{}^{A}_B\end{array}\right)\gamma ^5e^\sigma 𝒜^{}& \gamma ^\mu \omega _{2\mu }{}_{}{}^{A}{}_{B}{}^{},_{ij}\end{array}\right),`$ (71)
$`\mathrm{\Omega }^A{}_{B}{}^{}{}_{}{}^{}=\left(\begin{array}{cc}\gamma ^\mu \omega _{1\mu }{}_{}{}^{A}{}_{B}{}^{}^{}& 𝒜\gamma ^5e^\sigma \left(\begin{array}{cc}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{A}_B& 0\\ 0& \omega _1{}_{}{}^{A}_B\end{array}\right)\\ \left(\begin{array}{cc}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{A}_B& 0\\ 0& \omega _2{}_{}{}^{A}_B\end{array}\right)\gamma ^5e^\sigma 𝒜^{}& \gamma ^\mu \omega _{2\mu }{}_{}{}^{A}{}_{B}{}^{}^{}\end{array}\right),`$ (78)
where $`A,B,C=1,2,3,4,5,6`$ and $`5,6`$ respectively, refers to $`\overline{u},\overline{d}`$.
We define the Dirac operator on the curved space, which is described by (51) and (54), as
$`\widehat{D}=\left(\begin{array}{cccc}i_R& 0& & \gamma ^5𝐌^{}\\ 0& i_R& & \\ & & & \\ \gamma ^5𝐌& & & i_L\end{array}\right),`$ (83)
$`𝐌^{}=\left(\begin{array}{c}f^{}(x)(\varphi ^0,\varphi ^+)\\ \stackrel{~}{f}^{}(x)(\varphi ^{},\varphi ^0)\end{array}\right),𝐌=\left(\begin{array}{cc}f(x)\left(\begin{array}{c}\varphi ^0\\ \varphi ^{}\end{array}\right)& \stackrel{~}{f}(x)\left(\begin{array}{c}\varphi ^+\\ \varphi ^0\end{array}\right)\end{array}\right),`$ (91)
where $`\varphi ^0`$ is the neutral Higgs-like field and $`\varphi ^+=\varphi ^{}`$.
$`i_R=\gamma ^\mu (i_\mu {\displaystyle \frac{g^{}}{2}}YB_\mu ),`$ (92)
$`i_L=\gamma ^\mu (i_\mu {\displaystyle \frac{g}{2}}\tau _aA_{a\mu }{\displaystyle \frac{g^{}}{2}}YB_\mu ),`$ (93)
where $`A_{a\mu }(x)`$ is the $`SU(2)_W`$ gauge field and $`B_\mu (x)`$ is the $`U(1)_Y`$ gauge field. $`\tau _a`$ are the Pauli matrices. <sup>4</sup><sup>4</sup>4 $`\tau _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, $`\tau _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)`$, $`\tau _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$.
### 3.2 On-shell Higgs couplings
Following we impose the unitarity condition and the torsion-less condition on these gravitational components given in the previous subsection $`3.1`$. First, we would like to impose the unitarity condition for the components. There are four possibilities from $`\alpha `$ to $`\delta `$: $`\alpha .`$ (18) and (37), $`\beta .`$ (18) and (54), $`\gamma .`$ (51) and (37), $`\delta .`$ (51) and (54). We here treat the second case, $`\beta `$. In the curved quark space, $`u`$, $`d`$, the generalized Higgs couplings depend on the local coordinate, $`f=f(x),\stackrel{~}{f}=\stackrel{~}{f}(x)`$.
In order to discuss the suitable Dirac operator, which we will calculate the unitarity condition, let us remember here the spontaneous symmetry breaking down in the standard model. In $`SU(2)_W\times U(1)_Y`$ non-abelian gauge theory, if one gives the following vacuum expectation value of the Higgs doublet field:
$`0|\mathrm{\Phi }(x)|0={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}0\\ v\end{array}\right),`$ (96)
then the mass terms of the gauge fields appear as $`M_W^2W_\mu ^{}W^\mu +\frac{1}{2}M_Z^2Z_\mu Z^\mu `$ in the Lagrangian. Here, $`W_\mu =\frac{1}{\sqrt{2}}(A_\mu ^1iA_\mu ^2)`$ corresponds to the W boson and $`Z_\mu =\mathrm{cos}\theta _WA_\mu ^3\mathrm{sin}\theta _WB_\mu `$ is the Z boson. From the action one can easily recognize that a linear combination of $`SU(2)_W\times U(1)_Y`$ gauge fields, $`A_\mu =\mathrm{sin}\theta _WA_\mu ^3+\mathrm{cos}\theta _WB_\mu `$, is still a mass-less gauge field.
Next, let us recall the leptonic sector after spontaneous symmetry breaking. The left-handed sector is $`SU(2)_W`$ doublet and the right-handed sector is $`SU(2)_W`$ singlet, because the charged current weak interaction of leptonic fields is the V-A interaction. When the vacuum expectation value of the Higgs doublet is (96), the masses of leptonic fields, $`m_j`$, are in proportion to the Higgs coupling as $`f_j=\sqrt{2m_j/v}`$. The masses of quarks are in proportion to the Higgs couplings, similar to the leptonic sector. The only different point between the leptonic sector and the quarks sector is that the lower quarks of the left handed ($`SU(2)_W`$ doublet) are not mass eigenstates.
We apply the following Dirac operator on our Hilbert space and we will impose the unitarity condition and the torsion-less condition on our bases,
$`\widehat{D}_{EM}=\left(\begin{array}{cc}\begin{array}{cc}i_{EM}& 0\\ 0& i_{EM}\end{array}& \frac{v}{\sqrt{2}}\gamma ^5\left(\begin{array}{cc}f^{}(x)& 0\\ 0& \stackrel{~}{f}^{}(x)\end{array}\right)\\ \frac{v}{\sqrt{2}}\gamma ^5\left(\begin{array}{cc}f(x)& 0\\ 0& \stackrel{~}{f}(x)\end{array}\right)& \begin{array}{cc}i_{EM}& 0\\ 0& i_{EM}\end{array}\end{array}\right),`$ (107)
where
$`i_{EM}`$ $`=`$ $`\gamma ^\mu (i_\mu LA_\mu ^{EM}),`$
$`A_\mu ^{EM}`$ $`=`$ $`\mathrm{sin}\theta _WA_\mu ^3+\mathrm{cos}\theta _WB_\mu .`$ (108)
$`L`$ is the $`U(1)_{EM}`$ charge and $`\theta _W`$ is the Weinberg angle.
We summarize the unitarity condition and the components of torsion in appendices. From (160) and (161) in appendix A, we ultimately acquire
$`\stackrel{~}{f}^2(x)=C_1\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{1\nu }{}_{}{}^{u}{}_{u}{}^{},{}_{11}{}^{}LA_\nu ^{EM})`$ $`(y),`$ (109)
$`\stackrel{~}{f}^2(x)=C_2\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{2\nu }{}_{}{}^{u}{}_{u}{}^{},{}_{11}{}^{}LA_\nu ^{EM})(y).`$ (110)
From (166) and (167) in the appendix, we obtain
$`f^2(x)=C_3\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{1\nu }{}_{}{}^{d}{}_{d}{}^{},{}_{22}{}^{}LA_\nu ^{EM})`$ $`(y),`$ (111)
$`f^2(x)=C_4\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{2\nu }{}_{}{}^{d}{}_{d}{}^{},{}_{22}{}^{}LA_\nu ^{EM})(y).`$ (112)
In the left-hand side, $`f^2`$ and $`\stackrel{~}{f}^2`$, are real, we can actually observe these. On the other hand, we can calculate the right-hand sides of (168), (169), (170) and (171) as the Wilson loops of the spin connections, which is a kind of gauge field and the $`U(1)_{EM}`$ gauge field. In ordinary gauge theories and string theory it is well-known that the Wilson loops are one of the observables. Therefore, we may expect that these equations are geometrically meaningful for our noncommutative geometry, whose algebra is $`\{M_2(C)CC\}C^{\mathrm{}}(M_4)`$.
In the list of the components of the torsion given by the appendix B we find some solvable differential equations:
$`T_{42}^d={\displaystyle \frac{v}{\sqrt{2}}}(i\gamma ^\mu \gamma ^5_\mu f[\gamma ^5,\gamma ^\mu ]_{}fi_{EM\mu }){\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{d}{}_{d,22}{}^{}=0,`$ (113)
$`T_{24}^d=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{d}{}_{d,22}{}^{}+{\displaystyle \frac{v}{\sqrt{2}}}(i\gamma ^\mu \gamma ^5_\mu f^{}[\gamma ^5,\gamma ^\mu ]_{}f^{}i_{EM\mu })=0,`$ (114)
$`T_{13}^u=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}\omega _{1\mu }{}_{}{}^{u}{}_{u,11}{}^{}\stackrel{~}{f}_{}^{}+i{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5_\mu \stackrel{~}{f}^{}+i{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}^{}[\gamma ^\mu ,\gamma ^5]_{}_{EM\mu }=0,`$ (115)
$`T_{31}^u=i{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5_\mu \stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}[\gamma ^\mu ,\gamma ^5]_{}i_{EM\mu }\gamma ^\mu \gamma ^5\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{2\mu }{}_{}{}^{u}{}_{u,11}{}^{}=0,`$ (116)
where ‘$``$’ means that the differential operator, $`_\mu `$, only acts on the next $`f`$.
If at a quantum level we imposed the torsion-less condition on the following ground state
$`_\mu |0>=0,`$ (117)
Solutions of Eq.(113)-Eq.(116) would be given by:
$`f^{}(x)=Q_1:\mathrm{exp}i{\displaystyle ^x}dy^\nu (\widehat{\omega }_{1\nu }{}_{}{}^{d}{}_{d,22}{}^{}2L\widehat{A}_\nu ^{EM})(y):,`$ (118)
$`f(x)=Q_2:\mathrm{exp}i{\displaystyle ^x}dy^\nu (\widehat{\omega }_{2\nu }{}_{}{}^{d}{}_{d,22}{}^{}2L\widehat{A}_\nu ^{EM})(y):,`$ (119)
$`\stackrel{~}{f}^{}(x)=Q_3:\mathrm{exp}i{\displaystyle ^x}dy^\nu (\widehat{\omega }_{1\nu }{}_{}{}^{u}{}_{u,11}{}^{}2L\widehat{A}_\nu ^{EM})(y):,`$ (120)
$`\stackrel{~}{f}(x)=Q_4:\mathrm{exp}i{\displaystyle ^x}dy^\nu (\widehat{\omega }_{2\nu }{}_{}{}^{u}{}_{u,11}{}^{}2L\widehat{A}_\nu ^{EM})(y):,`$ (121)
where $`Q_i`$, $`i=1,\mathrm{},4`$ are arbitrary constants. $`\widehat{\omega }`$, $`\widehat{A}_\nu ^{EM}`$ denote operators at the quantum level. $`:`$ expresses a normal-ordering. This normal ordering should be defined which can be consistent with (168), (169), (170) and (171).
## 4 Conclusion and discussion
We have proposed a gravitational model coupled with $`SU(2)`$ doublet complex Higgs-like fields whose couplings depend on the local coordinates of $`M_4`$. As the results of some on-shell conditions (the unitarity condition and the torsion-less condition)in the noncommutative space, these Higgs couplings have been represented by the Wilson loops of the connections and the gauge fields, (168), (169), (170) and (171). Moreover if we use these equations, we can geometrically present the conditions that the Higgs couplings become zero. The choices of the topological properties for the four dimensional continuum $`M_4`$ are not free, but are restricted by the other conditions,
$`|\stackrel{~}{f}|^2={\displaystyle \frac{2}{v^2}}{\displaystyle \frac{\omega _{1\mu }{}_{}{}^{ua}_{,11}}{\omega _{1\mu }{}_{}{}^{a}_{u,11}}}={\displaystyle \frac{2}{v^2}}{\displaystyle \frac{\omega _{1\mu }{}_{}{}^{ua},_{12}}{\omega _{1\mu }{}_{}{}^{a}_{u,21}}}=\mathrm{}`$
Notice that these restrictions contain the connections which the spread over the noncommutative space $`M_4\times Z_2`$. Therefore, we wish to stress in this paper that if we construct the standard model on the non-commutative space and we elevate the Higgs couplings to some geometrical objects on the noncommutative space, we would gain a system beyond the standard model, which could represent the masses of leptons and quarks as geometrical things of noncommutative space.
Lastly, we confess what we should improve in our construction. We have not decomposed the components of our bases containing the local Higgs couplings as the authors have performed in . They have described the decomposition of their one-form as $`E^r=e^\sigma [D,m^r]`$, where $`m^5=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)`$, $`m^6=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right)`$. In other words, we have not exactly presented the one-form of $`\stackrel{~}{\mathrm{\Omega }}_D^1(𝒜)`$ for our algebra (25). The one-form $`\stackrel{~}{\mathrm{\Omega }}_D^1(𝒜)`$ should be represented by a zero-form times a commutator between a zero-form and a Dirac operator. After excluding the auxiliary part of our base given in this paper, a final description of these conditions with respect to the Higgs-like couplings would be changed. However if we add a one-form of some algebra, we may revive the expressions of the above Wilson loops.
We hope that our attempt will be one of the advantages to found a formulation beyond the standard model.
Acknowledgments
We would like to thank Professor Watamura for giving his seminar at our university and for stimulating discussion. One of the authors (R.K) would also like to thank him for discussing about and . After submitting version 2 of this paper to arXiv, R.K. had a chance to talk about and to some members of theoretical physics group in KEK. They commented on the renormalizability of noncommutative field theories. R.K. would like to thank them for the valuable discussions and comments. R.K. is supported by the Research Fellowships of the Japan Society for the Promotion of Science.
## Appendix A Unitarity conditions
We calculate the unitarity condition by using
$`d_{EM}^A,^B_D=\mathrm{\Omega }^A{}_{C}{}^{}^C,^B_{D}^{}+^A,^C_D(\mathrm{\Omega }^B{}_{C}{}^{})^{}.`$ (122)
In order to calculate the unitarity condition of the bases and the spin connection, we need to estimate the generalized hermitian inner product between the above one-form bases as (3.30) and (3.31) in ,
$`^A,^B_D={\displaystyle \frac{1}{2}}\{\gamma (A),\gamma (B)\}_+\widehat{}^A\widehat{}^B,`$ (123)
where $`\{\gamma (A),\gamma (B)\}_+`$ means the anti-commutator between two gamma matrixes which are respectively containing $`^A`$ and $`^B`$. $`\widehat{}^B`$ means $`^B`$ removing the gamma matrix, $`\gamma ^B`$. The inner products between each (18), (37) (51) and (54) are respectively expressed by
$`^a,^b_D`$ $`=`$ $`\delta ^{ab},`$
$`^{\overline{u}},^{\overline{u}}_D`$ $`=`$ $`{\displaystyle \frac{v^2}{2}}\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\end{array}\right),`$ (128)
$`^{\overline{d}},^{\overline{d}}_D`$ $`=`$ $`{\displaystyle \frac{v^2}{2}}\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right),`$ (133)
$`^{\overline{u}},^{\overline{d}}_D`$ $`=`$ $`^{\overline{d}},^{\overline{u}}_D=0.`$ (134)
For the bases of the curved space:
$`^\mu ,^\nu _D=g^{\mu \nu }(x),g^{\mu \nu }=e_a^\mu e_b^\nu \delta ^{ab},`$ (135)
$`^u,^u_D`$ $`=`$ $`\stackrel{~}{f}(x)^2{\displaystyle \frac{v^2}{2}}\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\end{array}\right),`$ (140)
$`^d,^d_D`$ $`=`$ $`f(x)^2{\displaystyle \frac{v^2}{2}}\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right),`$ (145)
$`^u,^d_D`$ $`=`$ $`^d,^u_D=0.`$ (146)
$`(𝐀,𝐁)=(a,b)`$
$`\omega _{1\mu }{}_{}{}^{ab},{}_{11}{}^{}+\omega _{1\mu }{}_{}{}^{ba},_{11}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{ab},{}_{12}{}^{}+\omega _{1\mu }{}_{}{}^{ba},_{21}`$ $`=`$ $`0,`$
$`\omega _2{}_{}{}^{ab}\omega _1{}_{}{}^{ba}=0,`$ (147)
$`\omega _{1\mu }{}_{}{}^{ab},{}_{21}{}^{}+\omega _{1\mu }{}_{}{}^{ba},_{12}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{ab},{}_{22}{}^{}+\omega _{1\mu }{}_{}{}^{ba},_{22}`$ $`=`$ $`0,`$
$`\omega _{}^{}{}_{2}{}^{}{}_{}{}^{ab}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{ab}=0,`$ (148)
$`\omega _{}^{}{}_{1}{}^{}{}_{}{}^{ab}\omega _{}^{}{}_{2}{}^{}^{ba}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ab},{}_{11}{}^{}+\omega _{2\mu }{}_{}{}^{ba},_{11}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ab},{}_{12}{}^{}+\omega _{2\mu }{}_{}{}^{ba},{}_{21}{}^{}=0,`$ (149)
$`\omega _1{}_{}{}^{ab}\omega _2^{ba}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ab},{}_{21}{}^{}+\omega _{2\mu }{}_{}{}^{ba},_{12}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ab},{}_{22}{}^{}+\omega _{2\mu }{}_{}{}^{ba},{}_{22}{}^{}=0.`$ (150)
$`(𝐀,𝐁)=(a,u)`$
$`\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{a}{}_{u}{}^{},{}_{11}{}^{}+\omega _{1\mu }{}_{}{}^{ua},_{11}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{ua},{}_{21}{}^{}=0,\omega _1{}_{}{}^{ua}=0,`$
$`\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{a}{}_{u}{}^{},{}_{21}{}^{}+\omega _{1\mu }{}_{}{}^{ua},_{12}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{a}{}_{u}{}^{}+\omega _{}^{}{}_{1}{}^{}^{ua}`$ $`=`$ $`0,`$
$`\omega _{}^{}{}_{2}{}^{}^{ua}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{a}{}_{u}{}^{},{}_{11}{}^{}+\omega _{2\mu }{}_{}{}^{ua},_{11}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ua},_{21}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _1{}_{}{}^{a}{}_{u}{}^{}+\omega _2^{ua}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{a}{}_{u}{}^{},{}_{21}{}^{}+\omega _{2\mu }{}_{}{}^{ua},_{12}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ua},{}_{22}{}^{}=0.`$ (151)
$`(𝐀,𝐁)=(a,d)`$
$`\omega _{1\mu }{}_{}{}^{da},_{11}`$ $`=`$ $`0,`$
$`f^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{a}{}_{d}{}^{},{}_{12}{}^{}+\omega _{1\mu }{}_{}{}^{6a},_{21}`$ $`=`$ $`0,`$
$`f^2{\displaystyle \frac{v^2}{2}}\omega _2{}_{}{}^{a}{}_{d}{}^{}+\omega _1{}_{}{}^{da}=0,`$ (152)
$`\omega _{1\mu }{}_{}{}^{da},_{12}`$ $`=`$ $`0,`$
$`f^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{a}{}_{d}{}^{},{}_{22}{}^{}+\omega _{1\mu }{}_{}{}^{da},_{22}`$ $`=`$ $`0,`$
$`\omega _{}^{}{}_{1}{}^{}{}_{}{}^{da}=0,`$ (153)
$`f^2{\displaystyle \frac{v^2}{2}}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{a}{}_{d}{}^{}+\omega _{}^{}{}_{2}{}^{}^{da}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{da},_{11}`$ $`=`$ $`0,`$
$`f^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{a}{}_{d}{}^{},{}_{12}{}^{}+\omega _{2\mu }{}_{}{}^{da},{}_{21}{}^{}=0,`$ (154)
$`\omega _2{}_{}{}^{da}=0,\omega _{2\mu }{}_{}{}^{da},{}_{12}{}^{}=0,`$
$`f^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{a}{}_{d}{}^{},{}_{22}{}^{}+\omega _{2\mu }{}_{}{}^{da},{}_{22}{}^{}=0.`$ (155)
$`(𝐀,𝐁)=(u,b)`$
$`\omega _{1\mu }{}_{}{}^{ub},{}_{11}{}^{}+\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{b}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{ub},{}_{12}{}^{}+\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{b}{}_{u}{}^{},_{21}`$ $`=`$ $`0,`$
$`\omega _2{}_{}{}^{ub}+\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _1{}_{}{}^{b}{}_{u}{}^{}=0,`$ (156)
$`\omega _{1\mu }{}_{}{}^{ub},_{21}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{ub},_{22}`$ $`=`$ $`0,`$
$`\omega _{}^{}{}_{2}{}^{}{}_{}{}^{ub}=0,`$ (157)
$`\omega _{}^{}{}_{1}{}^{}{}_{}{}^{ub}+\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{b}_u`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ub},{}_{11}{}^{}\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{b}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ub},{}_{12}{}^{}+\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{b}{}_{u}{}^{},{}_{21}{}^{}=0,`$ (158)
$`\omega _1^{ub}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ub},_{21}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{ub},{}_{22}{}^{}=0.`$ (159)
$`(𝐀,𝐁)=(u,u)`$
$`i_\mu ^{EM}\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}+v^2\stackrel{~}{f}^2\omega _{1\mu }{}_{}{}^{u}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$ (160)
$`\omega _{1\mu }{}_{}{}^{u}{}_{u}{}^{},_{21}`$ $`=`$ $`0,`$
$`\omega _1{}_{}{}^{u}_u`$ $`=`$ $`0,`$
$`\omega _{}^{}{}_{2}{}^{}{}_{}{}^{u}_u`$ $`=`$ $`0,`$
$`i_\mu ^{EM}\stackrel{~}{f}^2{\displaystyle \frac{v^2}{2}}+v^2\stackrel{~}{f}^2\omega _{2\mu }{}_{}{}^{u}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$ (161)
$`\omega _{2\mu }{}_{}{}^{u}{}_{u}{}^{},{}_{21}{}^{}=0.`$ (162)
$`(𝐀,𝐁)=(u,d)`$
$`\omega _{1\mu }{}_{}{}^{d}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$
$`f^2\omega _{1\mu }{}_{}{}^{u}{}_{d}{}^{},{}_{12}{}^{}+\stackrel{~}{f}^2\omega _{1\mu }{}_{}{}^{d}{}_{u}{}^{},_{21}`$ $`=`$ $`0,`$
$`f^2\omega _2{}_{}{}^{u}{}_{d}{}^{}+\stackrel{~}{f}^2\omega _1{}_{}{}^{d}_u`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{u}{}_{d}{}^{},_{22}`$ $`=`$ $`0,`$
$`f^2\omega _{}^{}{}_{1}{}^{}{}_{}{}^{u}{}_{d}{}^{}+\stackrel{~}{f}^2\omega _{}^{}{}_{2}{}^{}{}_{}{}^{d}_u`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{d}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$
$`f^2\omega _{2\mu }{}_{}{}^{u}{}_{d}{}^{},{}_{12}{}^{}+\stackrel{~}{f}^2\omega _{2\mu }{}_{}{}^{d}{}_{u}{}^{},_{21}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{u}{}_{d}{}^{},{}_{22}{}^{}=0.`$ (163)
$`(𝐀,𝐁)=(d,b)`$
$`\omega _{1\mu }{}_{}{}^{db},{}_{11}{}^{}=\omega _{1\mu }{}_{}{}^{db},_{12}`$ $`=`$ $`\omega _2{}_{}{}^{db}=0,`$
$`\omega _{1\mu }{}_{}{}^{db},{}_{21}{}^{}+f^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{b}{}_{d}{}^{},_{12}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{db},{}_{22}{}^{}+f^2{\displaystyle \frac{v^2}{2}}\omega _{1\mu }{}_{}{}^{b}{}_{d}{}^{},_{22}`$ $`=`$ $`0,`$
$`\stackrel{´}{\omega }_2{}_{}{}^{db}+f^2{\displaystyle \frac{v^2}{2}}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{b}_d`$ $`=`$ $`0,`$
$`\omega _{}^{}{}_{1}{}^{}{}_{}{}^{db}=\omega _{2\mu }^{db},{}_{11}{}^{}=\omega _{2\mu }{}_{}{}^{db},_{12}`$ $`=`$ $`0,`$
$`\omega _1{}_{}{}^{db}+f^2{\displaystyle \frac{v^2}{2}}\omega _2{}_{}{}^{b}_d`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{db},{}_{21}{}^{}+f^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{b}{}_{d}{}^{},_{12}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{db},{}_{22}{}^{}+f^2{\displaystyle \frac{v^2}{2}}\omega _{2\mu }{}_{}{}^{b}{}_{d}{}^{},{}_{22}{}^{}=0.`$ (164)
$`(𝐀,𝐁)=(d,u)`$
$`\omega _{1\mu }{}_{}{}^{d}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2\omega _{1\mu }{}_{}{}^{d}{}_{u}{}^{},{}_{21}{}^{}+f^2\omega _{1\mu }{}_{}{}^{u}{}_{d}{}^{},_{12}`$ $`=`$ $`0,`$
$`\omega _{1\mu }{}_{}{}^{u}{}_{d}{}^{},_{22}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2\omega _{}^{}{}_{2}{}^{}{}_{}{}^{d}{}_{u}{}^{}+f^2\omega _{}^{}{}_{1}{}^{}{}_{}{}^{u}_d`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{d}{}_{u}{}^{},_{11}`$ $`=`$ $`0,`$
$`\stackrel{~}{f}\omega _1{}_{}{}^{d}{}_{u}{}^{}+f^2\omega _2{}_{}{}^{u}_d`$ $`=`$ $`0,`$
$`\stackrel{~}{f}^2\omega _{2\mu }{}_{}{}^{d}{}_{u}{}^{},{}_{21}{}^{}+f^2\omega _{2\mu }{}_{}{}^{u}{}_{d}{}^{},_{12}`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{u}{}_{d}{}^{},{}_{22}{}^{}=0.`$ (165)
$`(𝐀,𝐁)=(d,d)`$
$`\omega _{1\mu }{}_{}{}^{d}{}_{d}{}^{},{}_{12}{}^{}=\omega _2{}_{}{}^{d}_d`$ $`=`$ $`0,`$
$`i_\mu ^{EM}f^2{\displaystyle \frac{v^2}{2}}+f^2v^2\omega _{1\mu }{}_{}{}^{d}{}_{d}{}^{},_{22}`$ $`=`$ $`0,`$ (166)
$`\omega _{}^{}{}_{1}{}^{}{}_{}{}^{d}_d`$ $`=`$ $`0,`$
$`\omega _{2\mu }{}_{}{}^{d}{}_{d}{}^{},_{12}`$ $`=`$ $`0,`$
$`i_\mu ^{EM}f^2{\displaystyle \frac{v^2}{2}}+f^2v^2\omega _{2\mu }{}_{}{}^{d}{}_{d}{}^{},{}_{22}{}^{}=0.`$ (167)
From (160) and (161), we obtain
$`\stackrel{~}{f}^2(x)=C_1\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{1\nu }{}_{}{}^{u}{}_{u}{}^{},{}_{11}{}^{}LA_\nu ^{EM})`$ $`(y),`$ (168)
$`\stackrel{~}{f}^2(x)=C_2\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{2\nu }{}_{}{}^{u}{}_{u}{}^{},{}_{11}{}^{}LA_\nu ^{EM})(y).`$ (169)
From (166) and (167), we obtain
$`f^2(x)=C_3\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{1\nu }{}_{}{}^{d}{}_{d}{}^{},{}_{22}{}^{}LA_\nu ^{EM})`$ $`(y),`$ (170)
$`f^2(x)=C_4\mathrm{exp}i{\displaystyle ^x}dy^\nu (2\omega _{2\nu }{}_{}{}^{d}{}_{d}{}^{},{}_{22}{}^{}LA_\nu ^{EM})(y).`$ (171)
## Appendix B Components of torsion
We have calculated the torsion $`T^A=dE^A+\mathrm{\Omega }^A{}_{B}{}^{}E_{}^{B}`$ where $`dE^A=[\widehat{D}_{EM},E^A]`$ and (54) <sup>5</sup><sup>5</sup>5As we described in the discussion, we did not eliminate the auxiliary part which has been defined by (11) and (12) from (51). The each component of $`T^A`$ is given by
$`T_{11}^a=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{a}{}_{b,11}{}^{}+i\gamma ^\mu _\mu \gamma ^a+[\gamma ^\mu ,\gamma ^a]_{}i_\mu ^{EM}`$
$`T_{12}^a=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{a}{}_{b,12}{}^{}\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma )f^2{\displaystyle \frac{v^2}{2}}\omega _2{}_{}{}^{a}_d`$
$`T_{13}^a=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}\omega _{1\mu }{}_{}{}^{a}{}_{u,11}{}^{}\stackrel{~}{f}_{}^{}+{\displaystyle \frac{v}{\sqrt{2}}}[\gamma ^5,\gamma ^a]_{}f^{}`$
$`T_{14}^a=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{a}{}_{d,12}{}^{}+{\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^5\gamma ^b\mathrm{exp}(\sigma )\omega _2{}_{}{}^{a}_b`$
$`T_{21}^a=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{a}{}_{b,21}{}^{}\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma ){\displaystyle \frac{v^2}{2}}\stackrel{~}{f}^2\omega _{}^{}{}_{2}{}^{}{}_{}{}^{a}_u`$
$`T_{22}^a=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{a}{}_{b,22}{}^{}+\gamma ^\mu i_\mu \gamma ^a+[\gamma ^\mu ,\gamma ^a]_{}i_\mu ^{EM}`$
$`T_{23}^a=\gamma ^5\gamma ^b\mathrm{exp}(\sigma )\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{a}{}_{b}{}^{}+\gamma ^\mu \gamma ^5\omega _{1\mu }{}_{}{}^{a}{}_{u,21}{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}^{}`$
$`T_{24}^a=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{a}{}_{d,22}{}^{}+{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}^{}[\gamma ^5,\gamma ^a]_{}`$
$`T_{31}^a=\gamma ^\mu \gamma ^5\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{2\mu }{}_{}{}^{a}{}_{u,11}{}^{}+{\displaystyle \frac{v}{\sqrt{2}}}[\gamma ^5,\gamma ^a]_{}f`$
$`T_{32}^a=\omega _{}^{}{}_{1}{}^{}{}_{}{}^{a}{}_{b}{}^{}\stackrel{~}{f}_{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma ){\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{a}_{d,12}`$
$`T_{33}^a=i\gamma ^\mu _\mu \gamma ^a+[\gamma ^\mu ,\gamma ^a]_{}i_\mu ^{EM}+\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{a}_{b,11}`$
$`T_{34}^a=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{a}{}_{b,12}{}^{}+\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma ){\displaystyle \frac{v^2}{2}}f^{}\stackrel{~}{f}^{}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{a}_d`$
$`T_{41}^a=\omega _1{}_{}{}^{a}{}_{b}{}^{}f_{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma ){\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5\stackrel{~}{f}\omega _{2\mu }{}_{}{}^{a}_{u,21}`$
$`T_{42}^a={\displaystyle \frac{v}{\sqrt{2}}}[\gamma ^5,\gamma ^a]_{}\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{a}_{d,22}`$
$`T_{43}^a=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{a}{}_{b,21}{}^{}+\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma )\stackrel{~}{f}^{}f^{}{\displaystyle \frac{v^2}{2}}\omega _1{}_{}{}^{a}_u`$
$`T_{44}^a=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{a}{}_{b,22}{}^{}+\gamma ^\mu i_\mu \gamma ^a+[\gamma ^\mu ,\gamma ^a]_{}i_\mu ^{EM}`$ (172)
$`T_{11}^d=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{d}_{b,11}`$
$`T_{12}^d=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{d}{}_{b,12}{}^{}\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma )f^2{\displaystyle \frac{v^2}{2}}\omega _2{}_{}{}^{d}_d`$
$`T_{13}^d=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}\omega _{1\mu }{}_{}{}^{d}{}_{u,11}{}^{}\stackrel{~}{f}_{}^{}`$
$`T_{14}^d=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{d}_{d,12}`$
$`T_{21}^d=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{d}{}_{b,21}{}^{}\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma ){\displaystyle \frac{v^2}{2}}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{d}_u`$
$`T_{22}^d={\displaystyle \frac{v^2}{2}}(\stackrel{~}{f}^{}f+f^{}\stackrel{~}{f})\mathrm{𝟏}_{4\times 4}`$
$`T_{23}^d=\gamma ^5\gamma ^b\mathrm{exp}(\sigma )\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{d}{}_{b}{}^{}+\gamma ^\mu \gamma ^5\omega _{1\mu }{}_{}{}^{d}{}_{u,21}{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}^{}`$
$`T_{24}^d=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{d}{}_{d,22}{}^{}+{\displaystyle \frac{v}{\sqrt{2}}}(i\gamma ^\mu \gamma ^5_\mu f^{}[\gamma ^5,\gamma ^\mu ]_{}f^{}i_{EM\mu })`$
$`T_{31}^d=\gamma ^\mu \gamma ^5\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{2\mu }{}_{}{}^{d}_{u,11}`$
$`T_{32}^d=\omega _{}^{}{}_{1}{}^{}{}_{}{}^{d}{}_{b}{}^{}\stackrel{~}{f}_{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma ){\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{d}_{d,12}`$
$`T_{33}^d=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{d}_{b,11}`$
$`T_{34}^d=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{d}{}_{b,12}{}^{}+\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma ){\displaystyle \frac{v^2}{2}}f^{}\stackrel{~}{f}^{}\stackrel{´}{\omega }_1{}_{}{}^{d}_d`$
$`T_{41}^d=\omega _1{}_{}{}^{d}{}_{b}{}^{}f_{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma ){\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5\stackrel{~}{f}\omega _{2\mu }{}_{}{}^{d}_{u,21}`$
$`T_{42}^d={\displaystyle \frac{v}{\sqrt{2}}}(i\gamma ^\mu \gamma ^5_\mu f[\gamma ^5,\gamma ^\mu ]_{}fi_{EM\mu }){\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{d}_{d,22}`$
$`T_{43}^d=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{d}{}_{b,21}{}^{}+\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma )\stackrel{~}{f}^{}f^{}{\displaystyle \frac{v^2}{2}}\omega _1{}_{}{}^{d}_u`$
$`T_{44}^d=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{d}{}_{b,22}{}^{}+{\displaystyle \frac{v^2}{2}}(\stackrel{~}{f}f^{}+f\stackrel{~}{f}^{})\mathrm{𝟏}_{4\times 4}`$ (173)
$`T_{11}^u={\displaystyle \frac{v^2}{2}}\mathrm{𝟏}_{4\times 4}(f^{}\stackrel{~}{f}+\stackrel{~}{f}^{}f)+\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{u}_{b,11}`$
$`T_{12}^u=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{u}{}_{b,12}{}^{}\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma )f^2{\displaystyle \frac{v^2}{2}}\omega _2{}_{}{}^{u}_d`$
$`T_{13}^u=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}\omega _{1\mu }{}_{}{}^{u}{}_{u,11}{}^{}\stackrel{~}{f}_{}^{}+i{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5_\mu \stackrel{~}{f}^{}+i{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}^{}[\gamma ^\mu ,\gamma ^5]_{}_{EM\mu }`$
$`T_{14}^u=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{u}{}_{d,12}{}^{}+f{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma )\omega _2{}_{}{}^{u}_b`$
$`T_{21}^u=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{u}{}_{b,21}{}^{}\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma ){\displaystyle \frac{v^2}{2}}\stackrel{~}{f}^2\omega _{}^{}{}_{2}{}^{}{}_{}{}^{u}_u`$
$`T_{22}^u=\gamma ^\mu \gamma ^b\omega _{1\mu }{}_{}{}^{u}_{b,22}`$
$`T_{23}^u=\gamma ^5\gamma ^b\mathrm{exp}(\sigma )\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{}^{}{}_{2}{}^{}{}_{}{}^{u}{}_{b}{}^{}+\gamma ^\mu \gamma ^5\omega _{1\mu }{}_{}{}^{u}{}_{u,21}{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}^{}`$
$`T_{24}^u=\gamma ^\mu \gamma ^5{\displaystyle \frac{v}{\sqrt{2}}}f^{}\omega _{1\mu }{}_{}{}^{u}_{d,22}`$
$`T_{31}^u=i{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5_\mu \stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\stackrel{~}{f}[\gamma ^\mu ,\gamma ^5]_{}i_{EM\mu }\gamma ^\mu \gamma ^5\stackrel{~}{f}{\displaystyle \frac{v}{\sqrt{2}}}\omega _{2\mu }{}_{}{}^{u}_{u,11}`$
$`T_{32}^u=\omega _{}^{}{}_{1}{}^{}{}_{}{}^{u}{}_{b}{}^{}\stackrel{~}{f}_{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma ){\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{u}_{d,12}`$
$`T_{33}^u=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{u}{}_{b,11}{}^{}+{\displaystyle \frac{v^2}{2}}\mathrm{𝟏}_{4\times 4}(f\stackrel{~}{f}^{}+\stackrel{~}{f}f^{})`$
$`T_{34}^u=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{u}{}_{b,12}{}^{}+\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma ){\displaystyle \frac{v^2}{2}}f^{}\stackrel{~}{f}^{}\omega _{}^{}{}_{1}{}^{}{}_{}{}^{u}_d`$
$`T_{41}^u=\omega _1{}_{}{}^{u}{}_{b}{}^{}f_{}^{}{\displaystyle \frac{v}{\sqrt{2}}}\gamma ^5\gamma ^b\mathrm{exp}(\sigma ){\displaystyle \frac{v}{\sqrt{2}}}\gamma ^\mu \gamma ^5\stackrel{~}{f}\omega _{2\mu }{}_{}{}^{u}_{u,21}`$
$`T_{42}^u={\displaystyle \frac{v}{\sqrt{2}}}f\gamma ^\mu \gamma ^5\omega _{2\mu }{}_{}{}^{u}_{d,22}`$
$`T_{43}^u=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{u}{}_{b,21}{}^{}+\mathrm{𝟏}_{4\times 4}\mathrm{exp}(\sigma )\stackrel{~}{f}^{}f^{}{\displaystyle \frac{v^2}{2}}\omega _2{}_{}{}^{u}_u`$
$`T_{44}^u=\gamma ^\mu \gamma ^b\omega _{2\mu }{}_{}{}^{u}_{b,22}`$ (174)
In section three we have imposed the torsion-less condition for a part of these components.
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# References
UCSD-PTH-00-08
TECHNION-PH-00-27
hep-ph/yymmddd
April 2000
EXTRAPOLATING SU(3) BREAKING EFFECTS
FROM $`D`$ TO $`B`$ DECAYS
Michael Gronau
Physics Department, Technion – Israel Institute of Technology
32000 Haifa, Israel
and
Dan Pirjol
Department of Physics, University of California at San Diego
La Jolla, CA 92093
ABSTRACT
> We consider two SU(3) breaking parameters, $`R_1(m_B)`$ and $`R_2(m_B)`$, appearing in a relation between $`B^+K\pi `$ and $`B^+\pi \pi `$ amplitudes, which plays an important role in determining the weak phase $`\gamma `$. We identify an isospin-related quantity $`R_2(m_D)`$ measured in $`D`$ decays, exhibiting large SU(3) breaking which is likely due to nonfactorizable effects. With a cautious remark about possible nonfactorizable SU(3) breaking in $`B`$ decays, we proceed to calculate factorizable SU(3) breaking corrections. Applying heavy quark symmetry to semileptonic $`D`$ and $`B`$ decay form factors, we find that SU(3) breaking in $`R_2(m_B)/R_1(m_B)`$ may be significantly larger than estimated from certain model calculations of form factors.
PACS codes: 12.15.Hh, 12.15.Ji, 13.25.Hw, 14.40.Nd
Weak nonleptonic decays of $`B`$ mesons provide an important source of information about the elements of the Cabibbo-Kobayashi-Maskawa (CKM) matrix. Flavor SU(3) symmetry of strong interactions plays an essential role in some of the methods proposed to determine the weak phases . First order SU(3) breaking effects in hadronioc $`B`$ decays may be parametrized in a completely general way in terms of several unknown parameters some of which can be determined from experiments. In certain hadronic amplitudes, such as in $`B\overline{D}\pi `$, experimental evidence exists for factorization in terms of products of two current matrix elements . In these cases, the corresponding SU(3) breaking parameters are given by ratios of $`K`$ and $`\pi `$ decay constants and ratios of $`B/B_s`$ to $`D/D_s`$ form factors. In decays to two charmless pseudoscalar mesons, which are useful for weak phase determinations , experimental evidence for factorization of hadronic matrix elements is still lacking. It was argued recently that within QCD nonfactorizable corrections due to hard gluon exchange are calculable and those which are due to soft exchanges are suppressed by $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ in a heavy quark expansion. Actual calculations of these corrections, controlling the former in a model-independent manner and showing that the latter are indeed small, are both desirable and challenging. Furthermore, in order to treat SU(3) breaking within the factorization approximation, one still needs the values of certain ratios of unmeasured form factors, for which one oftens relies on theoretical models.
The purpose of this Letter is to learn about SU(3) breaking in $`B`$ decays from the corresponding measured effects in $`D`$ decays. SU(3) breaking does not necessarily decrease monotonously with the decaying heavy quark mass. We will try to address the two relevant questions, of factorizable and nonfactorizable SU(3) violating corrections to hadronic decays, and of SU(3) breaking in semileptonic form factors which are used in the factorization approximation.
In general, soft final state interactions which spoil factorization are expected to affect $`D`$ and $`B`$ decays differently. It was often argued , and it has recently been shown by an actual calculation , that $`D`$ decay amplitudes involve large contributions from nearby light $`q\overline{q}`$ resonances which induce large SU(3) breaking effects. Such effects are not expected in $`B`$ decays. To avoid resonance effects, and thus study $`D`$ and $`B`$ decays on common grounds, we will consider only decays to “exotic” final states involving $`\pi \pi `$ in $`I=2`$ and $`K\pi `$ in $`I=3/2`$.
We will find very large SU(3) breaking in hadronic $`D`$ decays, in the absence of resonant terms, implying in the most likely scenario large nonfactorizable corrections. This should serve as a warning for what may be the case also in $`B`$ decays. In the factorization approximation, we then proceed to calculate SU(3) breaking in hadronic $`B`$ decays, where $`B`$ meson form factors are obtained from those measured for $`D`$ by applying a heavy quark symmetry scaling law. Our result will be compared with a model-dependent calculation.
As our test case, we consider an SU(3) relation between the isospin $`I=3/2`$ amplitude in $`BK\pi `$ and the $`I=2`$ amplitude in $`B\pi \pi `$
$`A(B^+K^0\pi ^+)`$ $`+`$ $`\sqrt{2}A(B^+K^+\pi ^0)=\sqrt{2}\mathrm{tan}\theta _c(R_1\delta _+e^{i\gamma }R_2)A(B^+\pi ^0\pi ^+),`$
$`\delta _+`$ $``$ $`[3/(2\lambda |V_{ub}/V_{cb}|)][(c_9+c_{10})/(c_1+c_2)]=0.66\pm 0.15.`$ (1)
This SU(3) relation generalizes a triangle relation proposed in by including, in addition to the current-current (“tree”) contributions, also the effects of dominant electroweak penguin (EWP) amplitudes given by the second term on the right-hand-side. Eq. (S0.Ex1) and its charge-conjugate were proposed as a way for determining the weak phase $`\gamma \mathrm{Arg}V_{ub}^{}`$. Other suggestions for using $`BK\pi `$ decays to study $`\gamma `$ were discussed in .
The coefficients $`R_{1,2}`$ in Eq. (S0.Ex1) parametrize SU(3) breaking effects and are in general complex numbers. In the SU(3) limit they are both equal to 1. Knowledge of the precise values of $`R_1`$ and $`R_2/R_1`$, in the presence of SU(3) breaking, is crucial for an accurate determination of $`\gamma `$ . Using the factorization approximation, it is customary to apply the value $`R_1f_K/f_\pi =1.22`$ to the tree part. SU(3) breaking corrections to the EWP-to-tree ratio $`R_2/R_1`$ were estimated in the generalized factorization approximation, assuming a certain model-dependent value for the ratio of $`B`$ to $`K`$ and $`B`$ to $`\pi `$ form factors, and were found to amount to a few percent . Numerically, this follows from an accidental cancellation between the contributions of the color-allowed and color-suppressed amplitudes. In addition, nonfactorizable SU(3) breaking corrections can in principle be significant .
Our main concern will be the SU(3) breaking parameter $`R_2`$. We will show that there exists a corresponding quantity $`R_2(m_D)`$, which has already been measured in $`D`$ decays exhibiting large SU(3) breaking. It will be argued that this effect is likely due to nonfactorizable corrections. It is not obvious why such effects should be much suppressed in $`B`$ decays. With this cautious remark, we nevertheless assume factorization in order to calculate $`R_1(m_B)`$ and $`R_2(m_B)`$ in this approximation. Using heavy quark symmetry to extrapolate form factors from measured semileptonic $`D`$ decays to $`B`$ decays, we calculate factorizable SU(3) breaking in $`R_2(m_B)/R_1(m_B)`$ and compare with estimates based on certain models for form factors.
For completeness, and in order to define $`R_1`$ and $`R_2`$ in broken SU(3) and to prove Eq. (S0.Ex1), we start by quickly reviewing the SU(3) structure of the amplitudes entering Eq. (S0.Ex1). The tree and electroweak penguin four-quark operators describing charmless decays transform under flavor SU(3) as a sum of $`\overline{\mathrm{𝟑}},\mathrm{𝟔}`$ and $`\overline{\mathrm{𝟏𝟓}}`$
$`_T^{\mathrm{\Delta }S=1}+_T^{\mathrm{\Delta }S=0}+_{EWP}^{\mathrm{\Delta }S=1}=`$
$`{\displaystyle \frac{G_F}{\sqrt{2}}}\lambda _u^{(s)}[{\displaystyle \frac{1}{2}}(c_1c_2)(\overline{\mathrm{𝟑}}_{I=0}^{(a)}\mathrm{𝟔}_{I=1})+{\displaystyle \frac{1}{2}}(c_1+c_2)(\overline{\mathrm{𝟏𝟓}}_{I=1}{\displaystyle \frac{1}{\sqrt{2}}}\overline{\mathrm{𝟏𝟓}}_{I=0}+{\displaystyle \frac{1}{\sqrt{2}}}\overline{\mathrm{𝟑}}_{I=0}^{(s)})]`$
$`+`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\lambda _u^{(d)}[{\displaystyle \frac{1}{2}}(c_1c_2)(\mathrm{𝟔}_{I=\frac{1}{2}}\overline{\mathrm{𝟑}}_{I=\frac{1}{2}}^{(a)})+{\displaystyle \frac{1}{2}}(c_1+c_2)({\displaystyle \frac{2}{\sqrt{3}}}\overline{\mathrm{𝟏𝟓}}_{I=\frac{3}{2}}{\displaystyle \frac{1}{\sqrt{6}}}\overline{\mathrm{𝟏𝟓}}_{I=\frac{1}{2}}+{\displaystyle \frac{1}{\sqrt{2}}}\overline{\mathrm{𝟑}}_{I=\frac{1}{2}}^{(s)})]`$
$``$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \frac{\lambda _t^{(s)}}{2}}\left({\displaystyle \frac{c_9c_{10}}{2}}(3\mathrm{𝟔}_{I=1}+\overline{\mathrm{𝟑}}_{I=0}^{(a)})+{\displaystyle \frac{c_9+c_{10}}{2}}(3\overline{\mathrm{𝟏𝟓}}_{I=1}{\displaystyle \frac{3}{\sqrt{2}}}\overline{\mathrm{𝟏𝟓}}_{I=0}{\displaystyle \frac{1}{\sqrt{2}}}\overline{\mathrm{𝟑}}_{I=0}^{(s)})\right),`$
where $`\lambda _q^{(q^{})}=V_{qb}^{}V_{qq^{}}`$. The explicit expressions of the four-quark operators appearing in the Hamiltonian can be found in .
The left-hand-side of Eq. (S0.Ex1) receives only contributions from the $`\mathrm{\Delta }S=1,I=1`$ terms in the weak Hamiltonian, which transform as 6 and $`\overline{\mathrm{𝟏𝟓}}`$, (QCD penguin operators are pure $`I=0`$ and do not contribute)
$`A(B^+K^0\pi ^+)+\sqrt{2}A(B^+K^+\pi ^0)=`$
$`\lambda _u^{(s)}(C_{\overline{15}_{I=1}}+C_{6_{I=1}})+\lambda _t^{(s)}\left({\displaystyle \frac{3}{2}}{\displaystyle \frac{c_9+c_{10}}{c_1+c_2}}C_{\overline{15}_{I=1}}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{c_9c_{10}}{c_1c_2}}C_{6_{I=1}}\right).`$ (3)
Here
$$C_{\overline{15}_{I=1}}(m_B)=\frac{G_F}{\sqrt{2}}\frac{1}{2}(c_1+c_2)(K^0\pi ^+|\overline{\mathrm{𝟏𝟓}}_{I=1}|B^++\sqrt{2}K^+\pi ^0|\overline{\mathrm{𝟏𝟓}}_{I=1}|B^+)$$
(4)
and
$$C_{6_{I=1}}(m_B)=\frac{G_F}{\sqrt{2}}\frac{1}{2}(c_1c_2)(K^0\pi ^+|\mathrm{𝟔}_{I=1}|B^++\sqrt{2}K^+\pi ^0|\mathrm{𝟔}_{I=1}|B^+)$$
(5)
are hadronic matrix elements of operators transforming as $`\overline{\mathrm{𝟏𝟓}}`$ and $`\mathrm{𝟔}`$. Using the approximate equality
$$\frac{c_9+c_{10}}{c_1+c_2}\frac{c_9c_{10}}{c_1c_2}1.12\alpha ,$$
(6)
which holds to better than $`3\%`$ , one finds
$$A(B^+K^0\pi ^+)+\sqrt{2}A(B^+K^+\pi ^0)=\lambda _u^{(s)}[(C_{\overline{15}_{I=1}}+C_{6_{I=1}})\delta _+e^{i\gamma }(C_{\overline{15}_{I=1}}C_{6_{I=1}})].$$
(7)
On the other hand, the amplitude on the right-hand-side of (S0.Ex1) is given by the matrix element of the $`\mathrm{\Delta }S=0,I=3/2`$ term in the weak Hamiltonian, (we neglect a very small EWP contribution )
$`\sqrt{2}A(B^+\pi ^+\pi ^0)=\lambda _u^{(d)}C_{\overline{15}_{I=3/2}},`$ (8)
where
$`C_{\overline{15}_{I=3/2}}(m_B)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}(c_1+c_2)\sqrt{{\displaystyle \frac{2}{3}}}\pi ^+\pi ^0|\overline{\mathrm{𝟏𝟓}}_{I=3/2}|B^+.`$ (9)
Taking the ratio of (3) and (8) reproduces the factor on the right-hand-side of (S0.Ex1) with
$`R_1(m_B)={\displaystyle \frac{C_{\overline{15}_{I=1}}+C_{6_{I=1}}}{C_{\overline{15}_{I=3/2}}}},R_2(m_B)={\displaystyle \frac{C_{\overline{15}_{I=1}}C_{6_{I=1}}}{C_{\overline{15}_{I=3/2}}}}.`$ (10)
Both final states on the left-hand-side of (3) and (8) belong to a 27 multiplet of SU(3), such that the matrix elements of $`\overline{\mathrm{𝟏𝟓}}_{I=1}`$ and $`\overline{\mathrm{𝟏𝟓}}_{I=3/2}`$ are related in the SU(3) limit, $`C_{\overline{15}_{I=1}}=C_{\overline{15}_{I=3/2}}`$. (The different numerical factors defining these amplitudes in Eqs. (4) and (9) are related to the different isospins involved). Furthermore, the matrix element of 6 in (3) vanishes in the same limit, such that $`R_1=R_2=1`$ in the SU(3) symmetric case. However, in broken SU(3) $`C_{\overline{15}_{I=1}}C_{\overline{15}_{I=3/2}}`$ and $`C_{6_{I=1}}0`$, which causes both $`R_1`$ and $`R_2`$ to differ from unity.
Whereas $`R_1(m_B)`$ and $`R_2(m_B)`$ are purely theoretical quantities, which cannot be directly measured, we prove now that another SU(3) breaking parameter,
$`R_2(m_D)={\displaystyle \frac{V_{us}}{V_{ud}}}{\displaystyle \frac{A(D^{}K^0\pi ^{})}{\sqrt{2}A(D^{}\pi ^{}\pi ^0)}},`$ (11)
measured in $`D`$ decays, is related to $`R_2(m_B)`$ by isospin in a fictitious heavy quark limit $`m_c=m_b`$.
The final states in the numerator and denominator of $`R_2(m_D)`$ have quantum numbers $`|I=\frac{3}{2},I_3=\frac{3}{2}`$ and $`|I=2,,I_3=1`$, respectively, and belong to the same isospin multiplets as the states $`|K^0\pi ^++\sqrt{2}|K^+\pi ^0`$ and $`|\pi ^+\pi ^0`$ in (S0.Ex1). The initial states $`D^{}`$ and $`B^+`$ are related to each other by isospin in the (fictitious) limit of identical heavy quarks. The weak Hamiltonian responsible for the relevant $`\overline{D}`$ decays is
$`_W`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{ud}^{}V_{cs}[{\displaystyle \frac{1}{2}}(c_1c_2)\sqrt{2}\mathbf{\hspace{0.17em}\hspace{0.17em}6}_{I=1}{\displaystyle \frac{1}{2}}(c_1+c_2)\sqrt{2}\overline{\mathrm{𝟏𝟓}}_{I=1}]`$
$`+`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}V_{us}^{}V_{cs}[(c_1+c_2)({\displaystyle \frac{1}{\sqrt{3}}}\overline{\mathrm{𝟏𝟓}}_{I=3/2}\sqrt{{\displaystyle \frac{2}{3}}}\overline{\mathrm{𝟏𝟓}}_{I=1/2})+(c_1c_2)\mathbf{\hspace{0.17em}6}_{I=1/2}],`$
where we neglect a small CP-violating contribution proportional to $`\frac{1}{2}(V_{us}^{}V_{cs}+V_{ud}^{}V_{cd})=𝒪(\lambda ^5)`$ in the Cabibbo-suppressed part and very small contributions of penguin operators .
The $`\mathrm{\Delta }S=1`$ $`(\mathrm{\Delta }S=0)`$ $`I=1,I_3=1`$ $`(I=\frac{3}{2},I_3=\frac{1}{2})`$ operators in (S0.Ex6) are the isospin partners of the $`I=1,I_3=0`$ $`(I=\frac{3}{2},I_3=\frac{1}{2})`$ operators in the $`B`$ decay Hamiltonian (S0.Ex4). (This can also be shown in terms of their quark structure). Therefore, in the limit of identical heavy quarks, isospin symmetry of strong interactions relates the amplitudes for $`D^{}`$ decays in (11) to those in $`B^+`$ decays
$`A(D^{}K^0\pi ^{})`$ $`=`$ $`V_{ud}^{}V_{cs}(C_{\overline{15}_{I=1}}(m_D)C_{6_{I=1}}(m_D)),`$ (13)
$`\sqrt{2}A(D^{}\pi ^{}\pi ^0)`$ $`=`$ $`V_{us}^{}V_{cs}C_{\overline{15}_{I=3/2}}(m_D).`$ (14)
Taking the ratio of these amplitudes yields the SU(3) breaking parameter $`R_2(m_D)`$ given in Eq. (11).
The experimental value of the ratio of amplitudes (11) is
$`|R_2(m_D)|=0.56\pm 0.08.`$ (15)
The large SU(3) breaking effect in $`R_2(m_D)`$ is somewhat surprising since the relevant final states are exotic, $`I=\frac{3}{2}`$ and $`2`$, and thus receive no resonant contributions . The large deviation of the ratio $`|R_2(m_D)|`$ from 1 raises the concern of a similar large SU(3) breaking effect in the $`B`$ case. In view of this possibility, let us review previous attempts and difficulties in explaining the numerical value (15).
A common way of studying SU(3) breaking in hardonic $`D`$ (and $`B`$) decays is by using the generalized factorization approach . In this approach one finds
$`R_2(m_D)={\displaystyle \frac{a_2^{(DK\pi )}}{a_1^{(D\pi \pi )}+a_2^{(D\pi \pi )}}}{\displaystyle \frac{f_K}{f_\pi }}{\displaystyle \frac{F_0^{D\pi }(m_K^2)}{F_0^{D\pi }(m_\pi ^2)}}+{\displaystyle \frac{a_1^{(DK\pi )}}{a_1^{(D\pi \pi )}+a_2^{(D\pi \pi )}}}{\displaystyle \frac{m_D^2m_K^2}{m_D^2m_\pi ^2}}{\displaystyle \frac{F_0^{DK}(m_\pi ^2)}{F_0^{D\pi }(m_\pi ^2)}}.`$ (16)
The phenomenological parameters $`a_{1,2}`$, describing the external and internal $`W`$-emission amplitudes respectively, are related to corresponding Wilson coefficients through $`a_{1,2}=c_{1,2}+\zeta c_{2,1}`$. The parameter $`\zeta `$ is process- and scale-dependent and is determined from experiments. When fitting nonleptonic two-body $`DK\pi `$ decays, using $`F_0^{DK}(m_\pi ^2)=0.77`$ and $`F_0^{D\pi }(m_\pi ^2)=0.7`$ , one obtains $`a_1^{(DK\pi )}=1.26`$ and $`a_2^{(DK\pi )}=0.51`$, corresponding to $`\zeta (m_c)=0`$. This is compatible with neglecting $`1/N_c`$ contributions in charm decays . This fit neglects, however, resonance contributions in nonexotic channels which, when included in an appropriate way, modifies the extracted values of $`a_{1,2}`$ to become $`a_1^{(DK\pi )}=1.06,a_2^{(DK\pi )}=0.64`$ .
An attempt was made to explain the large SU(3) breaking in (15) by using Eq. (16). In this attempt one faces three kinds of problems. First, there is an uncertainty in the values of $`a_i^{(DK\pi )}`$ due to resonance contributions in fitted nonexotic $`D`$ decays. Second, the values of $`a_i^{(D\pi \pi )}`$ may differ from those of $`a_i^{(DK\pi )}`$ which causes another uncertainty. In fact, a determination of $`a_i^{(D\pi \pi )}`$ from the corresponding Cabibbo suppressed decays (neglecting resonance contributions) gives very different results for $`a_2`$ compared with the $`DK\pi `$ case
$`a_1^{(D\pi \pi )}=1.05\left({\displaystyle \frac{0.7}{F_0^{D\pi }(m_\pi ^2)}}\right),a_2^{(D\pi \pi )}=0.07\left({\displaystyle \frac{0.7}{F_0^{D\pi }(m_\pi ^2)}}\right),`$ (17)
where $`F_0^{D\pi }(m_\pi ^2)=0.7`$ is used for normalization. This large deviation was blamed on inelastic hadronic rescattering .
Finally, a third uncertainty in evaluating $`R_2(m_D)`$ using (16) is due to the present experimental error in the ratio of form factors $`F_0^{DK}(0)/F_0^{D\pi }(0)`$. It was noted in that the value of $`R_2(m_D)`$ is very sensitive to this ratio. In Table 1 we list the results of four experiments for which the average value is $`F_0^{DK}(0)/F_0^{D\pi }(0)=1.00\pm 0.08`$.
| | Mark III | CLEO | CLEO | E687 |
| --- | --- | --- | --- | --- |
| $`\frac{F_0^{DK}(0)}{F_0^{D\pi }(0)}`$ | $`0.951\pm 0.214`$ | $`1.054\pm 0.246`$ | $`0.990\pm 0.230`$ | $`1.000\pm 0.110`$ |
> Table 1. Experimental results for the ratio of $`D\pi (K)`$ form factors at $`q^2=0`$. In quoting the numbers we used $`|V_{cd}/V_{cs}|=0.226`$.
We conclude that it is difficult to evaluate $`R_2(m_D)`$ and to explain its experimental value in a reliable manner within the generalized factorization approach. It is not entirely impossible that the failure to account for this large SU(3) breaking is due to resonant contributions in other $`D`$ decay processes which modify the extracted values of $`a_i`$. Assuming, for instance, $`a_2^{(DK\pi )}/a_1^{(DK\pi )}=0.6`$ , $`a_i^{(D\pi \pi )}=a_i^{(DK\pi )},F_0^{DK}(0)/F_0^{D\pi }(0)=1.1`$, one finds using Eq. (16) the value $`R_2(m_D)=0.64`$ consistent with (15). Still, a probable explanation for this failure is the presence of signficant nonfactorizable nonresonant contributions.
In view of the situation of $`R_2`$ at the $`D`$ mass, one should be aware of the possible presence of uncalculable nonfactorizable SU(3) breaking terms at the $`B`$ mass. We will disregard such terms for the rest of the discussion and study $`R_1(m_B)`$ and $`R_2(m_B)`$ in the generalized factorization approximation, keeping in mind that larger SU(3) breaking may be caused by nonfactorizable contributions.
In the factorization approximation one has
$`R_{1,2}(m_B)={\displaystyle \frac{a_{1,2}^{(BK\pi )}}{a_1^{(B\pi \pi )}+a_2^{(B\pi \pi )}}}{\displaystyle \frac{f_K}{f_\pi }}{\displaystyle \frac{F_0^{B\pi }(m_K^2)}{F_0^{B\pi }(m_\pi ^2)}}+{\displaystyle \frac{a_{2,1}^{(BK\pi )}}{a_1^{(B\pi \pi )}+a_2^{(B\pi \pi )}}}{\displaystyle \frac{m_B^2m_K^2}{m_B^2m_\pi ^2}}{\displaystyle \frac{F_0^{BK}(m_\pi ^2)}{F_0^{B\pi }(m_\pi ^2)}},`$ (18)
where $`R_2(m_B)`$ is given by an expression analogous to (16). The parameters $`a_i^{(BK\pi )}`$ and $`a_i^{(B\pi \pi )}`$ cannot be determined direcly from experiments. The closest one can get empirically is to measure these parameters at a different scale, the scale of hadronic $`bc`$ decays. An analysis of measured rates for $`BD^{()}\pi (\rho )`$ and $`BJ/\psi K`$, yields values $`a_1^{BD\pi }1`$ and $`a_2^{BD\pi }=0.20.3`$. A recent perturbative QCD calculation of $`B\pi \pi `$ decays , including nonfactorizable contributions due to hard gluon exchange, suggests that the corresponding value of the effective $`a_2`$ for two light pions could be even smaller, around $`a_2^{(B\pi \pi )}=0.1`$. Actually, $`a_2`$ acquires a sizable complex phase. This calculation does not include nonfactorizable terms due to soft exchanges, which are argued to be power suppressed in the heavy quark limit. A precise calculation of these soft corrections is a challenging task. In our estimate below of the SU(3) breaking parameters $`R_{1,2}`$ we will use the range $`a_2=0.10.3`$, assuming for simplicity $`a_i^{(BK\pi )}=a_i^{(B\pi \pi )}`$. Note that in general $`a_i`$ acquire complex phases and therefore $`R_i`$ become complex. Neglecting complex phases has a small effect on our estimates.
Under these assumptions, it is convenient to introduce the sum and difference of $`R_1`$ and $`R_2`$, in which a dependence on $`(a_1a_2)/(a_1+a_2)`$ is restricted to the difference
$`R_1+R_2`$ $`=`$ $`{\displaystyle \frac{f_K}{f_\pi }}{\displaystyle \frac{F_0^{B\pi }(m_K^2)}{F_0^{B\pi }(m_\pi ^2)}}+{\displaystyle \frac{m_B^2m_K^2}{m_B^2m_\pi ^2}}{\displaystyle \frac{F_0^{BK}(m_\pi ^2)}{F_0^{B\pi }(m_\pi ^2)}},`$ (19)
$`R_1R_2`$ $`=`$ $`{\displaystyle \frac{a_1a_2}{a_1+a_2}}\left({\displaystyle \frac{f_K}{f_\pi }}{\displaystyle \frac{F_0^{B\pi }(m_K^2)}{F_0^{B\pi }(m_\pi ^2)}}{\displaystyle \frac{m_B^2m_K^2}{m_B^2m_\pi ^2}}{\displaystyle \frac{F_0^{BK}(m_\pi ^2)}{F_0^{B\pi }(m_\pi ^2)}}\right).`$ (20)
The sum $`R_1+R_2`$ can be estimated more reliably than the difference $`R_1R_2`$, since the former does not depend on the poorly known coefficients $`a_{1,2}`$.
Important ingredients entering the factorization expressions (18) are the hadronic form factors $`F_0^{BP}(q^2)`$, defined in the usual way
$`P(p_P)|\overline{b}\gamma _\mu q|B(p_B)=\left((p_B+p_P)_\mu {\displaystyle \frac{m_B^2m_P^2}{q^2}}q_\mu \right)F_1^{BP}(q^2)+{\displaystyle \frac{m_B^2m_P^2}{q^2}}q_\mu F_0^{BP}(q^2),`$ (21)
where $`q=p_Bp_P`$. The form factors $`F_0^{B\pi (K)}(0)`$ were computed in a variety of quark models , light front model , MIT bag model , QCD sum rules and lattice QCD . The results obtained for these form factors at $`q^2=0`$ are presented in Table 2.
| | BSW | QCDSR | LCSR | RQM | LFM | BM | Lattice QCD |
| --- | --- | --- | --- | --- | --- | --- | --- |
| $`F_0^{B\pi }(0)`$ | 0.33 | 0.24 | $`0.30\pm 0.04`$ | $`0.37\pm 0.12`$ | 0.26 | 0.33 | $`0.27\pm 0.11`$ |
| $`F_0^{BK}(0)`$ | 0.38 | 0.25 | $`0.35\pm 0.05`$ | $`0.26\pm 0.08`$ | 0.34 | $``$ | $``$ |
> Table 2. Theory predictions for semileptonic $`B\pi (K)`$ form factors at $`q^2=0`$.
The ratio of form factors $`F_0^{B\pi }(m_K^2)/F_0^{B\pi }(m_\pi ^2)`$ is expected to differ from 1 by less than one percent; this difference will be neglected in the following discussion. Using the numerical values in Table 1 gives a typical value for the form factor ratio appearing in the second term of (18)
$$\frac{F_0^{BK}(m_\pi ^2)}{F_0^{B\pi }(m_\pi ^2)}=1.16.$$
(22)
It is hard to assign a theoretical uncertainty to this value, considering the large spread of model-predictions, some of which even involve values smaller than one. The particular value (22) implies a near cancellation of the two terms in (20) , giving $`R_1R_2=0.06(a_1a_2)/(a_1+a_2)=0.05(0.03)`$, corresponding to $`a_2=0.1(0.3)`$. Together with the sum $`R_1+R_2=2.37`$, this predicts $`R_1=1.21(1.20)`$ and $`R_2=1.16(1.17)`$. Thus, with the particular choice (22), SU(3) breaking in $`R_2/R_1`$ is at most about $`4\%`$.
In view of the wide range of model-dependent results for $`F_0^{BK}(0)/F_0^{B\pi }(0)`$ (see Table 2), and in order to perhaps narrow this range, we propose an alternative calculation of this ratio, which is based on the measured ratio of corresponding form-factors in $`D`$ decays, $`F_0^{DK}(0)/F_0^{D\pi }(0)=1.00\pm 0.08`$. Semileptonic $`B`$ and $`D`$ decay form factors, at points of equal $`\pi (K)`$ energy in the rest frame of the decaying meson, are related by a heavy quark symmetry scaling law
$`F_0^{B\pi }(q_{}^2)=\left({\displaystyle \frac{\alpha _s(m_b)}{\alpha _s(m_c)}}\right)^{6/25}\sqrt{{\displaystyle \frac{m_D}{m_B}}}F_0^{D\pi }(0),F_0^{BK}(q_{}^2)=\left({\displaystyle \frac{\alpha _s(m_b)}{\alpha _s(m_c)}}\right)^{6/25}\sqrt{{\displaystyle \frac{m_D}{m_B}}}F_0^{DK}(0).`$ (23)
The momentum transfer for $`B`$ form factors corresponding to $`q^2=0`$ in $`D`$ decays is $`q_{}^2=18.0`$ GeV<sup>2</sup>, for $`K`$ in the final state, and $`q_{}^2=17.6`$ GeV<sup>2</sup> for $`\pi `$. Taking the double ratio of $`B`$ and $`D`$ form factors cancels the leading $`𝒪(1/m_Q)`$ and $`𝒪(m_s/\mathrm{\Lambda }_{\chi SB})`$ corrections to the scaling laws of the individual form factors
$$\frac{F_0^{BK}(q_{}^2)/F_0^{B\pi }(q_{}^2)}{F_0^{DK}(0)/F_0^{D\pi }(0)}=1+𝒪(m_s/m_cm_s/m_b).$$
(24)
We use this relation to predict the ratio of $`B`$ form factors (22) in terms of the corresponding ratio for $`D`$ decays. The extrapolation of the former from $`q_{}^2`$ down to $`q^2=0`$ is made by assuming dominance by the $`0^+`$ states $`B_{0(s)}`$ for which we take $`m_{B_0}=5.75.8`$ GeV, $`m_{B_{s0}}=5.85.9`$ GeV. This gives
$$\frac{F_0^{BK}(0)}{F_0^{B\pi }(0)}=(1.013\pm 0.002)\frac{F_0^{BK}(q_{}^2)}{F_0^{B\pi }(q_{}^2)}1.01\pm 0.11,$$
(25)
where we introduced an error of 7% associated with the $`𝒪(m_s/m_c)`$ term in (24) . The rest of the uncertainty is due to the error in $`F_0^{DK}(0)/F_0^{D\pi }(0)`$. This uncertainty is expected to be reduced in future experiments of semileptonic $`D`$ decays. The relation between ratios of form factors in $`D`$ and $`B`$ decays can be tested by measuring $`B\pi \mathrm{}\nu `$ and $`BK\mathrm{}^+\mathrm{}^{}`$.
The value (25) is somewhat lower than the result (22) taken from certain models. Inserting this value into the relations (19) and (20), we find $`R_1+R_2=2.22\pm 0.11`$ and $`R_1R_2=(0.21\pm 0.11)(a_1a_2)/(a_1+a_2)`$. The central values yield $`R_1=1.20(1.17)`$ and $`R_2=1.02(1.05)`$ for $`a_2=0.1(0.3)`$. This implies very small SU(3) breaking in $`R_2`$ and larger SU(3) breaking in $`R_2/R_1`$, at a level of $`15\%(10\%)`$. This is significantly higher than the $`4\%`$ effect estimated from Eq. (22). An even larger SU(3) breaking in $`R_2/R_1`$ is obtained in the factorization approximation for values of $`F_0^{BK}(0)/F_0^{B\pi }(0)`$ which are smaller than 1.
We conclude with an interesting observation. Our discussion of the large measured SU(3) breaking in hadronic $`D`$ decays indicates the likely need for a significant nonfactorizable nonresonant contribution. Such effects may be smaller in $`B`$ decays but ought to be considered with care. In spite of this warning, one may argue from rather simple grounds that in the generalized factorization approximation SU(3) breaking in $`R_2(m_D)`$ is expected to be much larger than in $`R_2(m_B)`$. Assuming universal values for $`a_i`$, separately for $`B`$ and $`D`$ decays, both $`R_2(m_B)`$ in Eq. (18) and $`R_2(m_D)`$ in Eq. (16) consist of two SU(3) breaking contributions weighed by $`a_2/(a_1+a_2)`$ and $`a_1/(a_1+a_2)`$. In $`B`$ decays, where $`a_2/a_10.10.3`$, the dominant $`a_1`$ term involves SU(3) breaking given by $`F_0^{BK}(0)/F_0^{B\pi }(0)1`$ which is expected to be at a level of $`10\%`$. On the other hand, in $`D`$ decays in which $`a_2/a_1(0.6)(0.4)`$ is large and negative, the $`22\%`$ SU(3) breaking of $`f_K/f_\pi `$ in the $`a_2`$ term may be effectively roughly doubled by the destructive interference of this term with the $`a_1`$ term.
Acknowledgement: We thank H. Y. Cheng, R. Fleischer, B. Grinstein and J. L. Rosner for useful discussions. This work is supported in part by the National Science Foundation, by the United States - Israel Binational Science Foundation under Research Grant Agreement 98-00237, and by the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities.
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# Electron-Phonon Interactions in Correlated Systems: Adiabatic Expansion of Dynamical Mean Field Theory
## Abstract
We use the dynamical mean field theory to develop a systematic and computationally tractable method for studying electron-phonon interactions in systems with arbitrary electronic correlations. The method is formulated as an adiabatic expansion around the limit of static phonons. No specific electronic groundstate is assumed. We derive an effective low-frequency phonon action whose coefficients are static local correlation functions of the underlying electron system. We identify the correct expansion parameters. At a critical electron-phonon interaction strength the system undergoes a transition to a polaronic state. We determine the location of this polaronic instability in the presence of electron-electron interactions, doping, and quantum lattice fluctuations and present the formalism needed for study of the electron self-energy and effective mass.
The development of a quantitative theory of electronic properties of “correlated materials” (i.e. those for which the local density approximation plus Boltzmann transport theory is inadequate) is an important goal of materials science. An important feature of real materials is the electron-phonon interaction. The conventional theory of electron-phonon interactions in metals is due to Migdal and Eliashberg (ME) and is based on two assumptions: that the underlying electronic state is well described by Landau’s Fermi liquid theory and that the typical phonon frequency $`\omega _0`$ is small compared to the electronic Fermi energy $`E_F`$, so that an expansion is possible in the “adiabatic parameter” $`\gamma =\omega _0/E_F`$. However, in many materials of current interest electron correlations are strong, so a Fermi liquid description may not be appropriate and the relevant expansion parameter is unclear.
The introduction and more recent improvements of the dynamical mean field (DMF) method have opened an important avenue for progress, by showing how a good approximation to the correlation physics can be obtained from the solution of a quantum impurity problem plus a self-consistency condition. Unfortunately the straightforward inclusion of the electron-phonon coupling in the DMF formalism presents a difficult technical problem: the mismatch between the typical phonon frequency scale $`\omega _00.1\mathrm{eV}`$ and electron energy scale $`t1\mathrm{e}\mathrm{V}`$ renders conventional numerical approaches to the impurity problem prohibitively expensive, except in the “anti-adiabatic” limit $`\gamma 1`$ relevant to rather few materials.
In this paper we show how to turn this apparent difficulty to an advantage, by developing a systematic adiabatic expansion of the DMF formalism about the limit $`\gamma =0`$ of static (classical) phonons which was shown to be easily tractable in Refs. . We derive an effective low-energy phonon action which, at leading order, reproduces ME theory, appropriately generalized to a ground state which is not necessarily a Fermi liquid. The simplifications inherent in the DMF theory allow us to go beyond leading nontrivial order and to calculate the effect of dynamic terms representing quantum lattice fluctuations. The vertices of the action are higher-order static local density correlation functions of the electron system in the absence of phonons. We identify the correctly renormalized expansion parameters and show that the expansion breaks down in the vicinity of a polaronic instability, occurring at a critical value of the electron-phonon coupling. We determine the value of this critical coupling using analytical and Quantum Monte Carlo methods. We also present the formalism needed to calculate the electron self-energy.
We write a general tight-binding based Hamiltonian as follows: $`H=H_{\mathrm{el}}+H_{\mathrm{ph}}+H_{\mathrm{el}\mathrm{ph}}`$. The electronic part is
$$H_{\mathrm{el}}=\underset{ij\sigma }{}t_{ij}(c_{i\sigma }^{}c_{j\sigma }+c_{j\sigma }^{}c_{i\sigma })N\mu n+H_{\mathrm{ee}}.$$
(1)
Here $`H_{\mathrm{ee}}`$ represents electron-electron interactions, not explicitly written. The operator $`c_{i\sigma }^{}`$ creates an electron with spin $`\sigma `$ on lattice site $`i`$. The chemical potential is $`\mu `$ and the mean density is $`n=(1/N)_{i\sigma }c_{i\sigma }^{}c_{i\sigma }`$, where $`N`$ is the number of sites in the lattice.
We model the phonons as quantum oscillators:
$$H_{\mathrm{ph}}=\frac{1}{2}\underset{i}{}(M\dot{x}_i^2+Kx_i^2).$$
(2)
The operator $`x_i`$ measures the ionic displacement at site $`i`$, $`M`$ is the ion mass, and $`K`$ is the spring constant. We have taken the oscillator frequency $`\omega _0=(K/M)^{1/2}`$ to be dispersionless (Einstein model). To apply the method to more realistic situations one should interpret $`K`$ and $`M`$ as averages over the appropriate phonon bands or use the extended DMF theory of Ref. . Anharmonic terms can be trivially added and will be seen to be generated by the electron-phonon interaction.
We take the electron-phonon interaction to be local:
$$H_{\mathrm{el}\mathrm{ph}}=g\underset{i}{}x_i(n_in).$$
(3)
Here $`n_i=_\sigma c_{i\sigma }^{}c_{i\sigma }`$, and $`x_i=0`$ is the equilibrium phonon displacement for a uniform electron distribution.
In DMF theory the properties of $`H`$ may be obtained from the solution of an impurity model specified by the action $`S[c,\overline{c},x,a]=S_0[x]+S_{\mathrm{ee}}[c,\overline{c},a]+S_1[c,\overline{c},x,a]`$, with
$`S_0[x]`$ $`=`$ $`{\displaystyle \frac{1}{2T}}{\displaystyle \underset{k}{}}x_k\left(K+M\omega _k^2\right)x_k,`$ (4)
$`S_1[c,\overline{c},x,a]`$ $`=`$ $`{\displaystyle \underset{n\sigma }{}}\overline{c}_{n\sigma }c_{n\sigma }a_n+g{\displaystyle \underset{nk\sigma }{}}\overline{c}_{n\sigma }c_{n+k,\sigma }x_k.`$ (5)
Here $`S_{\mathrm{ee}}`$ arises from $`H_{\mathrm{ee}}`$ and is not explicitly written. The impurity electron and phonon are represented by $`c,\overline{c}`$ and $`x`$, respectively, and have been Fourier transformed according to $`x(\tau )=_k\mathrm{exp}(i\omega _k\tau )x_k`$ etc. Since the problem is local the fields depend on frequency but not on momentum. We use the compact notation $`x_k=x(i\omega _k)`$, $`c_n=c(i\omega _n)`$, with bosonic (fermionic) Matsubara frequencies $`\omega _k=2k\pi T`$ ($`\omega _n=(2n+1)\pi T`$) indexed by integers $`k`$ ($`n`$). The local Green function $`G_{\mathrm{loc}}`$ is calculated from the impurity partition function $`Z[a]=[dcd\overline{c}dx]\mathrm{exp}S[c,\overline{c},x,a]`$,
$$G_{\mathrm{loc}}[a]_n=\frac{\delta \mathrm{ln}Z[a]}{\delta a_n}\frac{1}{a_n\mathrm{\Sigma }[a]_n},$$
(6)
which defines the self-energy $`\mathrm{\Sigma }`$. Notice that $`Z`$, $`G_{\mathrm{loc}}`$, and $`\mathrm{\Sigma }`$ are functionals of the mean field function $`a`$ which is fixed by the condition that $`G_{\mathrm{loc}}[a]`$ agrees with the momentum integral of the full Green function, using the same self-energy, i.e.
$$G_{\mathrm{loc}}[a]_n=𝑑ϵ_k\frac{\rho (ϵ_k)}{i\omega _n+\mu \mathrm{\Sigma }[a]_nϵ_k},$$
(7)
which depends on the lattice density of states $`\rho (ϵ_k)`$.
The foregoing is general. To analyze the phonon problem it is convenient to formally integrate out the electron fields
$$\mathrm{exp}(S[x,a])=[dcd\overline{c}]\mathrm{exp}(S[c,\overline{c},x,a])$$
(8)
and work with an effective phonon action $`S[x,a]=S_0[x]+S_{\mathrm{ee}}[a]+S_1[x,a]`$. We proceed by formally expanding $`S_1`$ about the values $`\overline{x}`$ which extremize $`S`$. The crucial fact (of course already well-known to Migdal ) is that the characteristic phonon frequency scale is small. Therefore, as we shall show, the frequency sums involving phonon fields are dominated by frequencies of order $`\omega _0=(K/M)^{1/2}`$, thus $`S_1`$ may be evaluated by an expansion about the adiabatic limit.
For small $`g`$, $`\overline{x}=0`$; this is the conventional metallic state with no lattice distortion. As the coupling increases, eventually a state with $`\overline{x}0`$ becomes preferred. This corresponds to a polaronic instability at which the ground state is fundamentally reconstructed and has been extensively discussed for models involving only electron-phonon interactions. For the rest of this paper we assume $`\overline{x}=0`$; i.e. we expand about the undistorted ground state. We thus have:
$`S_1[x,a]`$ $`=`$ $`\mathrm{tr}\mathrm{ln}a`$ (10)
$`{\displaystyle \underset{n=2}{}}{\displaystyle \frac{g^n}{n}}{\displaystyle \underset{k_1,\mathrm{},k_n}{\overset{}{}}}\mathrm{\Gamma }_n[a]_{k_1,\mathrm{},k_n}x_{k_1}\mathrm{}x_{k_n},`$
where the prime denotes the restriction to $`k_1+\mathrm{}+k_n=0`$. The $`\mathrm{\Gamma }_n[a]_{k_1,\mathrm{},k_n}\mathrm{\Gamma }_n[a](i\omega _{k_1},\mathrm{},i\omega _{k_n})`$ are the connected $`n`$-point local density correlation functions of the electronic action. Their explicit form depends on the non-phonon physics contained in $`S_{\mathrm{ee}}`$. It is convenient to introduce a scale $`t`$ on which electronic properties vary (this could be the bandwidth or some interaction scale) and to define the parameters
$$\gamma =\frac{(K/M)^{1/2}}{t}=\frac{\omega _0}{t},\lambda =\frac{g^2}{Kt}.$$
(11)
We measure all frequencies and temperatures on the phonon frequency scale $`\gamma t`$, writing $`\stackrel{~}{\omega }_n=\omega _n/(\gamma t)`$ and $`\stackrel{~}{T}=T/(\gamma t)`$, and we rescale the phonon fields to $`\stackrel{~}{x}_k\stackrel{~}{x}(i\stackrel{~}{\omega }_k)=(K/T)^{1/2}x_k`$ so the free action becomes $`S_0[\stackrel{~}{x}]=\frac{1}{2}_k\stackrel{~}{x}_k(1+\stackrel{~}{\omega }_k^2)\stackrel{~}{x}_k`$, and the interaction part becomes
$`S_1[\stackrel{~}{x},a]`$ $`=`$ $`\mathrm{tr}\mathrm{ln}a{\displaystyle \underset{n=2}{}}{\displaystyle \frac{\lambda ^{n/2}(\stackrel{~}{T}\gamma )^{n/21}}{n}}`$ (13)
$`{\displaystyle \underset{k_1,\mathrm{},k_n}{\overset{}{}}}\stackrel{~}{\mathrm{\Gamma }}_n[a]_{k_1,\mathrm{},k_n}\stackrel{~}{x}_{k_1}\mathrm{}\stackrel{~}{x}_{k_n}.`$
The parameter $`\gamma `$ now appears in the prefactor and in the frequency arguments of the rescaled vertex functions $`\stackrel{~}{\mathrm{\Gamma }}_n[a]_{k_1,\mathrm{},k_n}\stackrel{~}{\mathrm{\Gamma }}_n[a](i\gamma \stackrel{~}{\omega }_{k_1},\mathrm{},i\gamma \stackrel{~}{\omega }_{k_n})=t^{n1}T\mathrm{\Gamma }_n[a](i\omega _{k_1},\mathrm{},i\omega _{k_n})`$. Provided (as is usually the case) that the dominant contributions to $`\stackrel{~}{\mathrm{\Gamma }}_n`$ come from frequencies of order $`t`$ we may evaluate the $`\stackrel{~}{\mathrm{\Gamma }}_n`$ via a low-frequency expansion and may for most purposes set the external frequencies to zero and define coupling constants $`\stackrel{~}{\mathrm{\Gamma }}_n[a]:=\stackrel{~}{\mathrm{\Gamma }}_n[a]_{0,\mathrm{},0}`$. Notice that it is not necessary that the electron Green function has no structure on the scale of $`\omega _0`$; merely that the dominant contribution to $`\stackrel{~}{\mathrm{\Gamma }}_n`$ comes from the scale of $`t`$. With $`S`$ given as in Eq. (13), $`\mathrm{\Sigma }`$, $`G_{\mathrm{loc}}`$, and $`a`$ may be computed from Eqs. (6), (7) as formal power series in $`\gamma `$.
The complete action, up to terms of order $`\gamma ^{3/2}`$, is
$`S[\stackrel{~}{x},a]=\mathrm{tr}\mathrm{ln}a+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k}{}}\stackrel{~}{x}_k\stackrel{~}{D}_k^1\stackrel{~}{x}_k`$ (14)
$`{\displaystyle \frac{1}{3}}\lambda ^{3/2}\gamma ^{1/2}\stackrel{~}{T}^{1/2}\stackrel{~}{\mathrm{\Gamma }}_3[a]{\displaystyle \underset{k_1,k_2}{}}\stackrel{~}{x}_{k_1}\stackrel{~}{x}_{k_2}\stackrel{~}{x}_{k_1k_2}`$ (15)
$`{\displaystyle \frac{1}{4}}\lambda ^2\gamma \stackrel{~}{T}\stackrel{~}{\mathrm{\Gamma }}_4[a]{\displaystyle \underset{k_1,k_2,k_3}{}}\stackrel{~}{x}_{k_1}\stackrel{~}{x}_{k_2}\stackrel{~}{x}_{k_3}\stackrel{~}{x}_{k_1k_2k_3}.`$ (16)
Eq. (16) can be regarded as an effective action on the low-energy scale $`\gamma t`$. The phonon Green function is $`\stackrel{~}{D}_k^1=1+\stackrel{~}{\omega }_k^2\lambda \stackrel{~}{\mathrm{\Gamma }}_2[a]_{k,k}1\lambda /\lambda _c+\stackrel{~}{\omega }_k^2+\lambda \gamma \alpha _p|\stackrel{~}{\omega }_k|+𝒪(\gamma ^2)`$, where $`\lambda _c=\stackrel{~}{\mathrm{\Gamma }}_2[a]^1`$ is a critical interaction strength and $`\alpha _p`$ is a damping parameter. The inclusion of static particle-hole diagrams $`\stackrel{~}{\mathrm{\Gamma }}_2`$ in the phonon self energy reduces the phonon frequency scale but increases the electron-phonon interaction strength. The renormalized expansion parameters are
$$\overline{\gamma }=\gamma (1\lambda /\lambda _c)^{1/2},\overline{\lambda }=\frac{\lambda }{1\lambda /\lambda _c}.$$
(17)
Physical properties $`\overline{\lambda }^m\overline{\gamma }^n`$ (where $`m`$ and $`n`$ are the number of phonon lines and loops, respectively, in the corresponding Feynman diagram) diverge like $`(1\lambda /\lambda _c)^{n/2m}`$ at $`\lambda \lambda _c`$ (notice that $`mn`$ always), and the uniform metallic state $`\overline{x}=0`$ becomes unstable to local distortions. Here we consider only $`\lambda <\lambda _c`$.
The coefficients $`\stackrel{~}{\mathrm{\Gamma }}_n`$ in (16) are higher order static susceptibilities of the interacting electron problem and may be computed with the techniques presently available. Here as an example we study the Holstein-Hubbard model with spin-$`\frac{1}{2}`$ electrons for which
$$S_{\mathrm{ee}}=U_0^\beta 𝑑\tau [n_{}(\tau )1/2][n_{}(\tau )1/2],$$
(18)
where $`n_\sigma (\tau )=\overline{c}_\sigma (\tau )c_\sigma (\tau )`$. We specialize to a Bethe lattice for which $`\rho (ϵ_k)=(4t^2ϵ_k^2)^{1/2}\theta (4t^2ϵ_k^2)/(2\pi t^2)`$ and the self-consistency equation (7) becomes $`a_n=i\omega _n+\mu t^2G_n`$. Here we focus on $`\stackrel{~}{\mathrm{\Gamma }}_2`$, but most considerations remain valid for $`\stackrel{~}{\mathrm{\Gamma }}_3`$ and $`\stackrel{~}{\mathrm{\Gamma }}_4`$. Combining Eqs. (8) and (13) we can show that $`\stackrel{~}{\mathrm{\Gamma }}_2=tT(nn_{\mathrm{ee}})^2_{\mathrm{ee}}`$, the static local density-density correlation function, where the average $`\mathrm{}_{\mathrm{ee}}`$ is taken with respect to $`S_{\mathrm{ee}}`$. All calculations will be done in the paramagnetic (PM) state.
In the weakly interacting limit $`u:=U/(2t)1`$ we use standard perturbation theory. To second order in $`u`$ we can replace the mean-field function $`a`$ entering $`\stackrel{~}{\mathrm{\Gamma }}_2[a]`$ by its noninteracting form $`a_0=i\omega _n+\mu t^2a_0^1`$. We find
$`\stackrel{~}{\mathrm{\Gamma }}_2[a](u,\beta ,\mu )=\stackrel{~}{\mathrm{\Gamma }}_2[a_0](0,\beta ,\mu )u\stackrel{~}{\mathrm{\Gamma }}_2^2[a_0](0,\beta ,\mu )`$ (19)
$`+u^2\left({\displaystyle \frac{1}{2}}\stackrel{~}{\mathrm{\Gamma }}_2^3[a_0](0,\beta ,\mu )+C_2(\beta ,\mu )\right)+𝒪(u^3),`$ (20)
where $`\stackrel{~}{\mathrm{\Gamma }}_2[a_0](0,\beta ,\mu )=(8/\pi )_1^1𝑑xf(x\mu )\sqrt{1x^2}x`$ is the noninteracting particle-hole bubble and $`C_2(\beta ,\mu )`$ is given as a sum of integrals which can be evaluated numerically. We set $`\beta =2t/T`$ and measure $`\mu `$ on the scale $`2t`$. Details will be given in a separate publication.
In the strongly interacting limit $`u1`$ we use Hartree-Fock (HF) theory which approximates the local Green function by $`G_n=\frac{1}{2}[(a_nU/2)^1+(a_n+U/2)^1]`$. The resulting cubic equation for $`G_n`$ may be evaluated numerically and the result used to calculate $`\stackrel{~}{\mathrm{\Gamma }}_2`$. At half filling we find $`\stackrel{~}{\mathrm{\Gamma }}_2\frac{1}{4}u^3`$ for $`u\mathrm{}`$.
In order to study finite doping and the crossover from weak to strong electron-electron interactions we use the Quantum Monte Carlo (QMC) method as described in Ref. . The imaginary time interval is divided into $`L=\beta /\mathrm{\Delta }\tau `$ time slices. At each time slice an Ising spin is introduced to decouple the quartic $`U`$ interaction. The partition function and all observables can then be sampled stochastically over all $`2^L`$ auxiliary spin configurations. Computer time grows with $`L^3`$. The QMC algorithm has to be iterated self-consistently \[using Eq. (7)\] until a convergent solution for $`a`$ is obtained. We found 15-20 iterations to be sufficient. The starting guess for $`a`$ was computed using finite-temperature Iterated Perturbation Theory.
Fig. 1 shows the results of QMC simulations of $`\stackrel{~}{\mathrm{\Gamma }}_2`$ as a function of $`u`$. Physical parameters are $`\beta =6`$ and $`\mu =0`$. The main graph shows the crossover from weak to strong coupling. In the insets we compare QMC (circles) to analytical results (lines), with excellent agreement. Upper inset: weak-coupling perturbation theory as discussed above, with $`\stackrel{~}{\mathrm{\Gamma }}_2[a_0](0,6,0)=0.745702`$ and $`C_2(6,0)=0.032789`$. Lower inset: $`\mathrm{log}\stackrel{~}{\mathrm{\Gamma }}_2`$ vs. $`\mathrm{log}u`$ which converges to the asymptotic solution $`\mathrm{log}\stackrel{~}{\mathrm{\Gamma }}_2=2\mathrm{log}23\mathrm{log}u`$ found above.
In Fig. 2 we plot $`\stackrel{~}{\mathrm{\Gamma }}_2`$ as a function of electron density $`n`$, for $`\beta =6`$ and various $`u`$, computed by QMC with $`L=32`$. For $`u=0`$ single-particle density fluctuations decrease when the system is doped away from half filling, as expected from the Pauli principle. At $`u=\mathrm{}`$, density fluctuations are forbidden at half filling, so $`\stackrel{~}{\mathrm{\Gamma }}_2(n=1)=0`$. The minimum at $`n=1`$ develops as pairs (for $`n1`$) or holes (for $`n1`$) become the dominant density fluctuations at low doping.
In general we expect $`\lambda _c`$ also to depend on lattice vibrations, i.e. non-zero $`\gamma `$. Integrating out cubic and quartic terms in Eq. (16) leads to an effective quadratic action $`S[\stackrel{~}{x},a]=\frac{1}{2}_k\stackrel{~}{x}_k(1\lambda /\lambda _c(\gamma )+\stackrel{~}{\omega }_k^2)\stackrel{~}{x}_k`$ where $`\lambda _c^1(\gamma )=\stackrel{~}{\mathrm{\Gamma }}_2+\frac{3}{4}\overline{\lambda }\overline{\gamma }\stackrel{~}{\mathrm{\Gamma }}_4+𝒪(\gamma ^2)`$ at half filling and $`T=0`$. The vertex $`\stackrel{~}{\mathrm{\Gamma }}_4[a_0]=16(1\mu ^2)^{3/2}(16\mu ^2)/(15\pi )`$ is negative at $`u=0`$. A $`u\mathrm{}`$ HF calculation gives $`\stackrel{~}{\mathrm{\Gamma }}_4u^3`$ and no sign change. Thus $`\stackrel{~}{\mathrm{\Gamma }}_2`$ (classical phonon self-energy) and $`\stackrel{~}{\mathrm{\Gamma }}_4`$ (quantum lattice fluctuations) have competing effects: the former increases the coupling $`\lambda \overline{\lambda }`$ while the latter decreases $`\overline{\lambda }`$ by renormalizing $`\lambda _c\lambda _c(\gamma )`$. For superconductivity, a similar competition between “phonon dressing” and “vertex corrections” was observed in Ref. .
We finally consider the electron self-energy $`\mathrm{\Sigma }=\mathrm{\Sigma }^{\mathrm{ee}}+\mathrm{\Sigma }^{\mathrm{ph}}=\delta \varphi /\delta G`$, where the Luttinger-Ward functional $`\varphi [G]`$ is the sum of all vacuum-to-vacuum skeleton diagrams. The electron part $`\mathrm{\Sigma }^{\mathrm{ee}}`$ comes from $`\varphi `$ arising from $`S_{\mathrm{ee}}`$, the phonon part $`\mathrm{\Sigma }^{\mathrm{ph}}`$ from $`S_1`$ which may be expanded in powers of $`\gamma `$ as in Eq. (16). To leading order we find
$$\mathrm{\Sigma }_{n\sigma }^{\mathrm{ph}}=\frac{\lambda }{2}\underset{k}{}D_k\frac{\delta \stackrel{~}{\mathrm{\Gamma }}_2[a]_{k,k}}{\delta G_{n\sigma }}+𝒪(\gamma ^2).$$
(21)
In general $`\mathrm{\Sigma }^{\mathrm{ph}}`$ is of order $`\gamma `$ and its frequency dependence is on the scale of $`t`$ so that it may be neglected compared to either the bare frequency dependence or $`\mathrm{\Sigma }^{\mathrm{ee}}`$. However, if $`\delta \stackrel{~}{\mathrm{\Gamma }}/\delta G`$ is singular at low frequencies, as in Fermi liquids, then $`\mathrm{\Sigma }/\omega `$ may be of order unity. To investigate the possibility of such a singularity we note that
$`\stackrel{~}{\mathrm{\Gamma }}_2[a]_{k,k}`$ $`=`$ $`tT{\displaystyle \underset{n\sigma }{}}G_{n\sigma }G_{n+k,\sigma }\times `$ (23)
$`\left(1+tT{\displaystyle \underset{n^{}\sigma ^{}}{}}\mathrm{\Lambda }_{nn^{}k}^{\sigma \sigma ^{}}G_{n^{}\sigma ^{}}G_{n^{}+k,\sigma ^{}}\right),`$
where the vertex $`\mathrm{\Lambda }`$ is given in terms of the particle-hole irreducible vertex $`\mathrm{\Lambda }^I`$ via
$$\mathrm{\Lambda }_{nn^{}k}^{\sigma \sigma ^{}}=\mathrm{\Lambda }_{nn^{}k}^{I\sigma \sigma ^{}}+tT\underset{n^{\prime \prime }\sigma ^{\prime \prime }}{}\mathrm{\Lambda }_{nn^{\prime \prime }k}^{I\sigma \sigma ^{\prime \prime }}G_{n^{\prime \prime }\sigma ^{\prime \prime }}G_{n^{\prime \prime }+k,\sigma ^{\prime \prime }}\mathrm{\Lambda }_{n^{\prime \prime }n^{}k}^{\sigma ^{\prime \prime }\sigma ^{}}.$$
Following the usual Fermi liquid arguments we observe that $`\mathrm{\Lambda }^I`$ is a smooth function of its arguments, so the required singular behavior can only occur if we differentiate on one of the explicit $`G`$ factors, leading to
$$\mathrm{\Sigma }_{n\sigma }^{\mathrm{ph}}=\lambda tT\underset{k}{}D_kG_{n+k,\sigma }\mathrm{\Lambda }_{k\sigma }^2+\mathrm{\Sigma }_{n\sigma }^{\mathrm{ph}\mathrm{reg}}$$
(24)
with $`\mathrm{\Lambda }_{k\sigma }=1+tT_{n^{}\sigma ^{}}\mathrm{\Lambda }_{0n^{}0}^{\sigma \sigma ^{}}G_{n^{}\sigma ^{}}G_{n^{}+k,\sigma ^{}}`$ and $`\mathrm{\Sigma }^{\mathrm{ph}\mathrm{reg}}`$ a function which varies on the scale of $`t`$ or $`U`$. If the ground state is a Fermi liquid then $`G_{n\sigma }=i\pi \mathrm{sign}(\omega _n)\rho (\mu )+G^{\mathrm{inc}}`$ and we obtain for the phonon contribution to the mass enhancement
$$\frac{m^{}}{m}|_{\mathrm{ph}}=1+\overline{\lambda }t\rho (\mu )\mathrm{\Lambda }_{0\sigma }^2.$$
(25)
We have calculated $`\mathrm{\Lambda }`$ and find it decreases rapidly as correlations increase, so the density-coupled electron-phonon interaction is “turned off” as the Mott insulator is approached. Details will be presented elsewhere.
To summarize, we have combined a small-$`\gamma `$ expansion with DMF theory to obtain a general formalism for studying electron-phonon effects in correlated electron systems. We have identified a local polaronic instability at $`\lambda \lambda _c`$ where physical properties diverge. We have shown that $`\lambda _c`$ can be increased by electron-electron interactions and quantum fluctuations and decreased by doping. We have obtained formal expressions for the one-phonon electron self-energy and effective mass. Future papers will apply the method to studies of conductivity, isotope effects, and dynamical consequences of electron-phonon interactions near a Mott transition.
We thank S. Blawid and R. L. Greene for useful discussions. We acknowledge NSF Grant No. DMR00081075 and the University of Maryland-Rutgers MRSEC for support.
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# Coherent backscattering effect on wave dynamics in a random medium
\[
## Abstract
A dynamical effect of coherent backscattering is predicted theoretically and supported by computer simulations: The distribution of single-mode delay times of waves reflected by a disordered waveguide depends on whether the incident and detected modes are the same or not. The change amounts to a rescaling of the distribution by a factor close to $`\sqrt{2}`$. This effect appears only if the length of the waveguide exceeds the localization length; there is no effect of coherent backscattering on the delay times in the diffusive regime.
\]
Coherent backscattering refers to the systematic constructive interference of waves reflected from a medium with randomly located scatterers. The constructive interference occurs in a narrow cone around the angle of incidence, and is a fundamental consequence of time-reversal symmetry . The resulting peak in the angular dependence of the reflected intensity is a generic wave effect: It has been observed using light waves and acoustic waves , for classical and quantum scatterers , in passive and active media .
These studies mainly addressed static properties. Dynamic aspects of wave propagation in random media are now entering the focus of attention , and the work on acoustic waves has started to study the connection with the coherent backscattering effect. The key observable in the dynamic experiments is the derivative $`\varphi ^{}=\mathrm{d}\varphi /\mathrm{d}\omega `$ of the phase $`\varphi `$ of the wave amplitude with respect to the frequency $`\omega `$. The quantity $`\varphi ^{}`$ has the dimension of a time and is interpreted as a delay time. Van Tiggelen et al. have developed a statistical theory for the distribution of the delay time $`\varphi ^{}`$ and the intensity $`I`$ in a waveguide geometry (where angles of incidence are discretized as modes). Although the theory was worked out mainly for the case of transmission, the implications for reflection are that the distribution $`P(\varphi ^{})`$ does not depend on whether the detected mode $`n`$ is the same as the incident mode $`m`$ or not. This is in contrast with $`P(I)`$, which is rescaled by a factor of $`1/2`$ when $`n`$ becomes equal to $`m`$—so that the mean $`\overline{I}`$ becomes twice as large. Hence it appears that no coherent backscattering effect exists for $`P(\varphi ^{})`$.
What we will demonstrate here is that this is true only if wave localization may be disregarded. Previous studies dealt with the diffusive regime of waveguide lengths $`L`$ below the localization length $`\xi `$. Here we consider the localized regime $`L>\xi `$ (assuming that also the absorption length $`\xi _\mathrm{a}>\xi `$). The distribution of reflected intensity is insensitive to the presence or absence of localization, being given in both regimes by Rayleigh’s law:
$$P(I)=\{\begin{array}{cc}Ne^{NI}& \text{ if }nm,\\ \frac{1}{2}Ne^{NI/2}& \text{ if }n=m\end{array}$$
(1)
(for unit incident intensity). In contrast, we find that the delay-time distribution changes markedly as one enters the localized regime, decaying more slowly for large $`|\varphi ^{}|`$. Moreover, a coherent backscattering effect appears: For $`L>\xi `$ the peak of $`P(\varphi ^{})`$ is higher for $`n=m`$ than for $`nm`$ by a factor which is close to $`\sqrt{2}`$. We present a complete analytical theory, compare it with numerical simulations, and offer a qualitative argument for this unexpected dynamical effect of coherent backscattering.
Let us begin with a more precise formulation of the problem. We consider a disordered medium (mean free path $`l`$) in a waveguide geometry (length $`L`$, with $`N1`$ propagating modes at frequency $`\omega `$) and study the correlator $`\rho `$ of the reflected wave amplitudes at two nearby frequencies $`\omega \pm \frac{1}{2}\delta \omega `$,
$$\rho =r_{nm}(\omega +\frac{1}{2}\delta \omega )r_{nm}^{}(\omega \frac{1}{2}\delta \omega ).$$
(2)
The indices $`n`$ and $`m`$ specify the detected and incident mode, respectively. (We assume single-mode excitation and detection.) The amplitudes $`r_{nm}`$ form the $`N\times N`$ reflection matrix $`r`$. In the localized regime (localization length $`\xi Nl`$ smaller than both $`L`$ and the absorption length $`\xi _\mathrm{a}`$), the matrix $`r`$ is approximately unitary because transmission is negligibly small. We assume time-reversal symmetry (no magneto-optical effects), so that $`r`$ is also symmetric. Following Genack et al. we define the single-mode (or single-channel) delay time $`\varphi ^{}`$ as
$$\varphi ^{}=\underset{\delta \omega 0}{lim}\frac{\mathrm{Im}\rho }{\delta \omega I},$$
(3)
where $`I=|r_{nm}(\omega )|^2`$ is the intensity of the reflected wave in the detected mode for unit incident intensity. We seek the joint distribution function $`P(I,\varphi ^{})`$ in an ensemble of different realizations of disorder.
The single-mode delay time $`\varphi ^{}`$ is a linear combination of the Wigner-Smith delay times $`\tau _i`$ ($`i=1,2,\mathrm{},N`$), which are the eigenvalues of the matrix
$$ir^{}\frac{\mathrm{d}r}{\mathrm{d}\omega }=U^{}\mathrm{diag}(\tau _1,\mathrm{},\tau _N)U.$$
(4)
(The matrix of eigenvectors $`U`$ is unitary for a unitary reflection matrix.) For small $`\delta \omega `$ we can expand
$$r(\omega \pm \frac{1}{2}\delta \omega )=U^\mathrm{T}U\pm \frac{1}{2}i\delta \omega U^\mathrm{T}\mathrm{diag}(\tau _1,\mathrm{},\tau _N)U,$$
(5)
hence the relations
$$\varphi ^{}=\mathrm{Re}\frac{A_1}{A_0},I=|A_0|^2,A_k=\underset{i}{}\tau _i^ku_iv_i.$$
(6)
We have abbreviated $`u_i=U_{im}`$, $`v_i=U_{in}`$.
The distribution of the Wigner-Smith delay times for this problem was determined recently . In terms of the rates $`\mu _i=1/\tau _i`$ it has the form of the Laguerre ensemble of random-matrix theory,
$$P(\{\mu _i\})\underset{i<j}{}|\mu _i\mu _j|\underset{k}{}\mathrm{\Theta }(\mu _k)e^{\gamma (N+1)\mu _k},$$
(7)
where $`\mathrm{\Theta }(x)=1`$ for $`x>0`$ and $`0`$ for $`x<0`$. The parameter $`\gamma =\alpha l/c`$ (with wave velocity $`c`$) equals the scattering rate times a numerical coefficient ($`\alpha =\pi ^2/4`$, $`8/3`$ for two, three-dimensional scattering). Eq. (7) extends the single-mode ($`N=1`$) result of Refs. to any $`N`$. The matrix $`U`$ is uniformly distributed in the unitary group. We consider first the typical case $`nm`$ of different incident and detected modes. (The special case $`n=m`$ is addressed later.) For $`nm`$ the vectors $`𝐮`$ and $`𝐯`$ become uncorrelated in the large-$`N`$ limit, and their elements become independent Gaussian random numbers with vanishing mean and variance $`|u_i^2|=|v_i^2|=N^1`$.
It is convenient to work momentarily with the weighted delay time $`W=\varphi ^{}I`$ and to recover $`P(I,\varphi ^{})`$ from $`P(I,W)`$ at the end. The characteristic function $`\chi (p,q)=e^{ipIiqW}`$ is the Fourier transform of $`P(I,W)`$. The average $`\mathrm{}`$ is over the vectors $`𝐮`$ and $`𝐯`$ and over the set of eigenvalues $`\{\tau _i\}`$. The average over one of the vectors, say $`𝐯`$, is easily carried out, because it is a Gaussian integration. The result is a determinant,
$`\chi (p,q)=det(1+iH/N)^1,`$ (8)
$`H=p𝐮^{}𝐮^\mathrm{T}+{\displaystyle \frac{1}{2}}q(\overline{𝐮}^{}𝐮^\mathrm{T}+𝐮^{}\overline{𝐮}^\mathrm{T}).`$ (9)
The Hermitian matrix $`H`$ is a sum of dyadic products of the vectors $`𝐮`$ and $`\overline{𝐮}`$, with $`\overline{u}_i=u_i\tau _i`$, and hence has only two non-vanishing eigenvalues $`\lambda _+`$ and $`\lambda _{}`$. Some straightforward linear algebra gives
$$\lambda _\pm =\frac{1}{2}\left(qB_1+p\pm \sqrt{2pqB_1+q^2B_2+p^2}\right),$$
(10)
where we have defined the spectral moments
$$B_k=\underset{i}{}|u_i|^2\tau _i^k.$$
(11)
The resulting determinant is $`det(1+H/N)^1=(1+\lambda _+/N)^1(1+\lambda _{}/N)^1`$, hence
$`\chi (p,q)=\left[1+{\displaystyle \frac{ip}{N}}+{\displaystyle \frac{iq}{N}}B_1+{\displaystyle \frac{q^2}{4N^2}}(B_2B_1^2)\right]^1.`$ (12)
An inverse Fourier transform, followed by a change of variables from $`I`$, $`W`$ to $`I`$, $`\varphi ^{}`$, gives
$`P(I,\varphi ^{})=\mathrm{\Theta }(I)(N^3I/\pi )^{1/2}e^{NI}`$ (13)
$`\times (B_2B_1^2)^{1/2}\mathrm{exp}\left(NI{\displaystyle \frac{(\varphi ^{}B_1)^2}{B_2B_1^2}}\right).`$ (14)
The average is over the spectral moments $`B_1`$ and $`B_2`$, which depend on the $`u_i`$’s and $`\tau _i`$’s via Eq. (11).
This result in the localized regime is to be compared with the result of diffusion theory ,
$`P_{\mathrm{diff}}(I,\varphi ^{})=\mathrm{\Theta }(I)(N^3I/\pi )^{1/2}e^{NI}`$ (15)
$`\times (Q\overline{\varphi ^{}}^2)^{1/2}\mathrm{exp}\left(NI{\displaystyle \frac{(\varphi ^{}\overline{\varphi ^{}})^2}{Q\overline{\varphi ^{}}^2}}\right).`$ (16)
The constants are given by $`QL/l`$ and $`\overline{\varphi ^{}}L/c`$ up to numerical coefficients of order unity . Comparison of Eqs. (14) and (16) shows that the two distributions would be identical if statistical fluctuations in the spectral moments $`B_1`$, $`B_2`$ could be ignored. However, as we shall see shortly, the distribution $`P(B_1,B_2)`$ is very broad, so that fluctuations can not be ignored. The large fluctuations are a consequence of the high density of anomalously large Wigner-Smith delay times $`\tau _i`$ in the Laguerre ensemble (7), and are related to the penetration of the wave deep into the localized regions. The large $`\tau _i`$’s are eliminated in the diffusive regime $`L\xi `$. Then $`B_1`$ and $`B_2`$ can be replaced by their ensemble averages, and the Gaussian theory is recovered. (The same applies if the absorption length $`\xi _\mathrm{a}\xi `$.)
To determine how the statistical fluctuations in the spectral moments alter $`P(I,\varphi ^{})`$, we need the joint distribution $`P(B_1,B_2)`$. This can be calculated by applying the random-matrix technique of Refs. to the Laguerre ensemble. The result is
$`P(B_1,B_2)=\mathrm{\Theta }(B_1)\mathrm{\Theta }(B_2)\mathrm{exp}\left({\displaystyle \frac{NB_1^2}{B_2}}\right)`$ (17)
$`\times [{\displaystyle \frac{B_1^2\gamma N^3}{B_2^4}}(B_2+\gamma N^2B_1)\mathrm{exp}({\displaystyle \frac{2\gamma N}{B_1}})`$ (18)
$`{\displaystyle \frac{\gamma ^3N^5}{4B_2^5}}(2B_2^24B_1^2B_2N+B_1^4N^2)\mathrm{Ei}({\displaystyle \frac{2\gamma N}{B_1}})],`$ (19)
where $`\mathrm{Ei}(x)`$ is the exponential-integral function. The most probable values are $`B_1\gamma N`$, $`B_2\gamma ^2N^3`$, while the mean values $`B_1`$, $`B_2`$ diverge—demonstrating the presence of large fluctuations. The distribution $`P(I,\varphi ^{})`$ follows from Eq. (14) by integrating over $`B_1`$ and $`B_2`$ with weight given by Eq. (19). This is an exact result in the large-$`N`$ limit.
For the discussion we concentrate on the distribution $`P(\varphi ^{})=_0^{\mathrm{}}dIP(I,\varphi ^{})`$ of the single-mode delay time by itself, which takes the form
$$P(\varphi ^{})=\underset{0}{\overset{\mathrm{}}{}}\underset{0}{\overset{\mathrm{}}{}}dB_1dB_2\frac{P(B_1,B_2)(B_2B_1^2)}{2(B_2+\varphi ^22B_1\varphi ^{})^{3/2}}.$$
(20)
We compare this distribution in the localized regime with the result of diffusion theory ,
$$P_{\mathrm{diff}}(\varphi ^{})=(Q/2\overline{\varphi ^{}})[Q+(\varphi ^{}/\overline{\varphi ^{}}1)^2]^{3/2}.$$
(21)
In the localized regime the value $`\varphi _{\mathrm{peak}}^{}\gamma N`$ at the center of the peak of $`P(\varphi ^{})`$ is much smaller than the width of the peak $`\mathrm{\Delta }\varphi ^{}\gamma N^{3/2}\varphi _{\mathrm{peak}}^{}(\xi /l)^{1/2}`$. This holds also in the diffusive regime, where $`\varphi _{\mathrm{peak}}^{}=\overline{\varphi ^{}}`$ and $`\mathrm{\Delta }\varphi ^{}\varphi _{\mathrm{peak}}^{}(L/l)^{1/2}`$. However, the mean $`\varphi ^{}=B_1`$ diverges for $`P`$, but is finite (equal to $`\overline{\varphi ^{}}`$) for $`P_{\mathrm{diff}}`$. In the tails $`P`$ decays $`|\varphi ^{}|^2`$, while $`P_{\mathrm{diff}}|\varphi ^{}|^3`$.
The transition from the diffusive to the localized regime with increasing $`L`$ is illustrated in Fig. 1. The data points are obtained from the numerical simulation of scattering of a scalar wave by a two-dimensional random medium. The reflection matrices $`r(\omega \pm \frac{1}{2}\delta \omega )`$ are computed by applying the method of recursive Green functions to the Helmholtz equation on a square lattice (lattice constant $`a`$). The width $`W=100a`$ and the frequency $`\omega =1.4c/a`$ are chosen such that there are $`N=50`$ propagating modes. The mean free path $`l=14.0a`$ is found from the formula $`T=(1+s)^1`$ for the transmission probability in the diffusive regime, where $`s=2L/\pi l`$ for two-dimensional scattering. The corresponding localization length $`\xi =NL/s=1100a`$. The parameter $`\gamma =46.3a/c`$ is found from $`\overline{\varphi ^{}}`$ in the diffusive regime . The relationship between the parameters $`\gamma `$, $`\overline{\varphi ^{}}`$, and $`Q`$ appearing in $`P`$ and $`P_{\mathrm{diff}}`$ is given by
$$\overline{\varphi ^{}}=\gamma \frac{s(3+2s)}{3(1+s)},Q=\frac{8s^3+28s^2+30s+15}{5(2s+3)^2}.$$
(22)
In Fig. 1, the same set of parameters is used for all lengths to plot the distributions $`P`$ (solid curve) and $`P_{\mathrm{diff}}`$ (dashed). The numerical data agrees very well with the analytical predictions in their respective regimes of validity.
We now turn to the case $`n=m`$ of equal-mode excitation and detection. The vectors $`𝐮`$ and $`𝐯`$ in Eq. (6) are then identical, and we can write
$$\varphi ^{}=\mathrm{Re}\frac{C_1}{C_0},I=|C_0|^2,C_k=\underset{i}{}\tau _i^ku_i^2.$$
(23)
The joint distribution function of the complex numbers $`C_0`$ and $`C_1`$ can be calculated in the same way as $`P(B_1,B_2)`$. We find
$`P(C_0,C_1)\mathrm{exp}(N|C_0|^2/2){\displaystyle _0^{\mathrm{}}}dxx^2e^x`$ (24)
$`\times \left(1+{\displaystyle \frac{|C_1|^2x^2}{\gamma ^2N^2}}{\displaystyle \frac{2x}{\gamma N}}\mathrm{Re}C_0C_1^{}\right)^{5/2}.`$ (25)
The maximal value $`P(\varphi _{\mathrm{peak}}^{})=\sqrt{2/\pi N^3\gamma ^2}`$ for $`n=m`$ is larger than the maximum of $`P(\varphi ^{})`$ for $`nm`$ by a factor $`\sqrt{2}\times \frac{4096}{1371\pi }=1.35`$ in the large-$`N`$ limit. This is in contrast to the diffusive regime, where there is no difference in the distributions of single-mode delay times for $`n=m`$ and $`nm`$. Our analytical expectations are again in excellent agreement with the numerical simulations, presented in Fig. 2.
In order to explain the coherent backscattering enhancement of the peak of $`P(\varphi ^{})`$ in qualitative terms, we compare Eq. (23) for $`n=m`$ with the corresponding relation (6) for $`nm`$. The quantities $`A_0`$ and $`A_1`$, as well as the quantities $`C_0`$ and $`C_1`$, become mutually independent in the large-$`N`$ limit. \[The cross-term $`(\gamma N)^1\mathrm{Re}C_0C_1^{}`$ in Eq. (25) is of order $`N^{1/2}`$ because $`C_0N^{1/2}`$ and $`C_1\gamma N`$.\] The main contribution to the enhancement of the peak height, namely the factor of $`\sqrt{2}`$, has the same origin as the factor-of-two enhancement of the mean intensity $`\overline{I}`$. More precisely, the relation $`P(A_0)=\sqrt{2}P(\sqrt{2}C_0)`$ leads to a rescaling of $`P(I)`$ for $`n=m`$ by a factor of $`1/2`$ \[see Eq. (1)\] and to a rescaling of $`P(\varphi ^{})`$ by a factor of $`\sqrt{2}`$. The remaining factor of $`\frac{4096}{1371\pi }=0.95`$ comes from the difference in the distributions $`P(A_1)`$ and $`P(C_1)`$. These distributions turn out to be very similar, hence the factor is close to unity. The asymptotic independence of $`A_0`$ and $`A_1`$ (as well as of $`C_0`$ and $`C_1`$) is another consequence of the strong fluctuations originating from the high density of anomalously large Wigner-Smith delay times $`\tau _i`$. In the diffusive regime the corresponding quantities are strongly correlated, and the coherent backscattering enhancement of the intensity affects both in the same way. Because only their ratio features in $`\varphi ^{}`$, this effect cancels and no difference is observed in $`P_{\mathrm{diff}}(\varphi ^{})`$ for $`n=m`$ and $`nm`$.
In conclusion, we have discovered a dynamical effect of coherent backscattering, that requires localization for its existence. Computer simulations confirm our prediction, which now awaits experimental observation.
We thank P. W. Brouwer for valuable advice. This work was supported by the Dutch Science Foundation NWO/FOM.
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# Simple stochastic models showing strong anomalous diffusion
## I Introduction
Anomalous diffusion, i.e., when the scaling of the moments of the position $`x(t)`$ is $`x^2(t)t^{2\nu }`$ with $`\nu >1/2`$, has been observed in a rather wide class of dynamical systems, e.g., intermittent maps , $`2D`$ symplectic maps and random velocity field as well in $`2D`$ time-dependent flow and Hamiltonian systems (e.g. the egg-crate potential) . In highly nontrivial systems, as those described in , the existence of anomalous diffusion has been established only numerically. On the other hand, for random shears it is possible to give an analytical criterium both for the existence of anomalous diffusion and for the computation of $`\nu `$ . As far as we know, the simplest nontrivial system showing anomalous diffusion is the continuous-time random walk (CTRW), sometimes also called Lévy walk. The CTRW is entirely specified by the probability density function (pdf) $`\psi (r,\tau )`$ to move a distance $`r`$ in a time $`\tau `$ in a single motion event. Let us assume, as in ,
$$\psi (r,\tau )=P(\tau )P(r\tau ),$$
(1)
where $`P(\tau )`$ is the pdf of having a flight of duration $`\tau `$ and $`P(r\tau )`$ is the conditional pdf of a displacement $`r`$ given the flight time $`\tau `$. The cases corresponding to $`P(\tau )\tau ^g`$ and $`P(r\tau )=\delta (r\tau ^\alpha )/2`$ can be treated analytically .
If the scaling of all moments can be described by just one exponent, i.e. $`x^{2n}(t)t^{2n\nu }`$, a collapse of the pdf’s at different times is obtained exploiting the rescaling
$$P(x,t)=t^\nu F(x/t^\nu ).$$
(2)
Then it also becomes clear that the value of $`\nu `$, in general, does not completely characterize the statistical properties of the diffusion process, as the function $`F(\xi )`$ needs to be specified.
In many cases, the use of just one exponent is not enough to describe all the moments, i.e., we have the so-called strong anomalous diffusion . For which $`|x(t)|^qt^{q\nu (q)}`$, where $`q\nu (q)`$ is a nonlinear function of $`q`$. The existence of a non-unique scaling exponent implies the failure of the data collapse for the pdf in the form given by Eq. 2. The best known case of a process showing strong anomalous diffusion is the advection of a passive scalar by a turbulent velocity field. In many cases it has been observed that the function $`q\nu (q)`$ is piecewise linear, i.e.,
$$q\nu (q)\{\begin{array}{cc}\nu _1q\hfill & q<q_c\hfill \\ qc\hfill & q>q_c.\hfill \end{array}$$
(3)
This is basically due to the existence of two mechanisms: a weak (i.e., with a unique exponent $`\nu _1>1/2`$) anomalous diffusion for the typical events, and a ballistic transport for the rare excursions (i.e., excursions much larger than $`x_{typ}(t)\mathrm{exp}\mathrm{ln}x(t)`$). The behavior (3) suggests the validity of the data collapse (2) for the pdf core, i.e. $`x/t^{\nu _1}`$ not too large, and two peaks at $`xt`$ ,i.e., the footprint of ballistic events. The cases where strong anomalous diffusion for which (3) holds are relatively simple and surely different from the cases of the relative dispersion in the fully developed turbulence .
By using only elementary techniques, in this paper we show that the bi-linear behavior for the scaling of the moments (3), which is present in the special case of the CTRW commonly found in the literature, i.e. $`P(r\tau )=\delta (r\tau ^\alpha )/2`$, does not hold in the general case. To show this point, we shall consider a generalized CTRW of the form:
$$P(r\tau ^{1+h}\tau )\tau ^{S(h)}$$
(4)
The inspiration for this choice comes from the multifractal description of turbulence .
The paper is organized as follows. In section II A we present the “standard” CTRW model. We then present a simple method to find the scaling for the even order moments (section II B). In section II C the model is generalized, and the same method is again applied to find the scaling. Numerical analysis to corroborate the analytical results of the general model are presented in Section III together with some discussions related to the shape of $`P(x,t)`$. Discussions and conclusions can finally be found in Section IV.
## II CTRW models
Anomalous diffusion occurs when some, or all, of the hypothesis of the central limit theorem break down. More specifically, the system has to violate at least one of the two following conditions:
1. Finite variance of the velocity.
2. Fast enough decay of the auto-correlation function of the Lagrangian velocities.
The paradigmatic model for anomalous diffusion, namely the Lévy flights violate the first condition. In the one dimensional case a Lévy flight corresponds to the evolution in discrete time
$$x(t_{i+1})=x(t_i)+v_i\mathrm{\Delta }t$$
(5)
of the particle position $`x`$ with $`t_{i+1}=t_i+\mathrm{\Delta }t`$ and $`v_i`$ being independent stochastic variables identically distributed according to a Lévy-stable distribution such that
$$P(v)v^g\mathrm{for}\mathrm{large}v$$
(6)
where $`1<g3`$. It is easy to show that $`x^2=\mathrm{}`$ for $`g<3`$ and that this stochastic process shows anomalous diffusion being $`x_{typ}t^{1/(g1)}>t^{1/2}`$.
### A The “standard” CTRW model
The existence of an infinite variance is not very pleasing from a physical point of view. This has lead to the introduction of the CTRW (also called Lévy walks). The idea is to relax the condition of a fixed, discrete time step in such a way that the process still has anomalous diffusion, but finite variance of the velocity.
Firstly, we introduce the particle trajectory
$$x(t)=x(t\tau _i)+v_i\tau _i$$
(7)
where $`x(t)`$ denotes the position of the particle at time $`t`$. During the random intervals $`\tau _i`$, the particles move a distance $`r_i`$ with constant random velocity $`v_i`$ independent of $`\tau _i`$. After each interval they choose new random values for $`\tau _i`$ and $`v_i`$.
The relevant quantity to characterize the motion of the particle is the pdf $`\psi (r,\tau )`$ of having a displacement $`r`$ in time $`\tau `$ in a single motion event. This pdf is chosen in the form (1). In the simple case where $`v_i=\pm v`$ we have
$$P(r\tau )=\frac{1}{2}\delta (rv\tau ).$$
(8)
Taking
$$P(\tau )\tau ^g\mathrm{for}\mathrm{large}\tau $$
(9)
we can determine the pdf $`P(x,t)`$ to be in $`x`$ at time $`t`$ . Actually, by introducing the probability density $`\mathrm{\Psi }(x,t)`$ to pass at location $`x`$ at time $`t`$ in a single motion event (and not necessarily to stop at $`x`$)
$$\mathrm{\Psi }(x,t)=P(xt)_t^{\mathrm{}}𝑑\tau _x^{\mathrm{}}𝑑r\psi (r,\tau )=\frac{1}{2}\delta (xvt)_t^{\mathrm{}}𝑑\tau P(\tau )$$
(10)
the pdf $`P(x,t)`$ can be written in the following way
$$P(x,t)=\mathrm{\Psi }(x,t)+_{\mathrm{}}^{\mathrm{}}𝑑x^{}_0^t𝑑\tau \psi (x^{},\tau )\mathrm{\Psi }(xx^{},t\tau )+\mathrm{}$$
(11)
the first term denotes the probability density to reach the position $`x`$ at time $`t`$ in a single motion event, the second term is the probability density to reach $`x`$ at time $`t`$ with one stop in $`x^{}`$ and so on to include all the combinations of motion events. In the Fourier-Laplace space $`(xk,tu)`$ the series in Eq. (11) assumes the closed form
$$\widehat{P}(k,u)=\frac{\widehat{\mathrm{\Psi }}(k,u)}{1\widehat{\psi }(k,u)}$$
(12)
and the behavior of $`x(t)^2`$ can be calculated analytically by using the relation
$$x(t)^2=^1\left[\frac{^2}{^2k}\widehat{P}(k,u)_{k=0}\right],$$
where $`^1`$ denotes the inverse Laplace transform $`(ut)`$ .
The results is that
$$x(t)^2\{\begin{array}{ccccc}t^2& 1& <g<& 2& \mathrm{b}allisticmotion\hfill \\ t^{4g}& 2& <g<& 3& \mathrm{a}nomalousdiffusion\hfill \\ t& 3& <g& & \mathrm{n}ormaldiffusion.\hfill \end{array}$$
(13)
Thus, enhanced anomalous diffusion occurs for $`2<g<3`$. For the moments of small order, which describe the core of the pdf, it has been shown that the asymptotic behavior gives $`|x(t)|^qt^{q\nu }`$ with $`\nu =1/(g1)`$ . Thus the core of the pdf can be scaled as in (2) using $`\nu `$. The ballistic motions which are responsible for the different scaling of higher order moments, show up as wings on the pdf, which does not scale using $`\nu `$.
The previous approach can be generalized to the case where
$$P(rt)=\frac{1}{2}\delta (r\tau ^\alpha ),v=\pm \tau ^{\alpha 1},P(\tau )\tau ^g\mathrm{with}g>1.$$
(14)
In addition one can treat more complicated situations by considering that the particle can move ballistically but it can be also trapped in some structures as vortices or chaotic islands (standard map or “egg-crate” potential) .
### B Finding the scaling of the moments
The usual method used in to find the scaling of the moments in the simple CTRW model is not elementary. We now present an alternative easier way to calculate the displacement moments $`x^q(t)`$ and thus to characterize the anomalous diffusion.
Let us consider a particle moving ballistically with velocity $`v_i`$ during the interval times $`\tau _i`$. The velocities $`v_i`$ are identically distributed, independent, random variables assuming the values $`\pm 1`$ alternatingly. The intervals times $`\tau _i`$ are identically distributed, independent, random variables and assuming the value $`\tau `$ with probability
$$P(\tau )\tau ^g\mathrm{with}g>1\mathrm{and}\tau [t_{min},T]$$
(15)
where $`t_{min}`$ and $`T`$ are the lowest and highest cutoffs, respectively. The reason for which we need to introduce such cutoffs will be clear later. The particle position at the time $`t`$ can be written as
$$x(t)=\underset{i=1}{\overset{n}{}}v_i\tau _i+v_{n+1}ϵ_{n+1}\mathrm{with}ϵ_{n+1}=t\underset{i=1}{\overset{n}{}}\tau _i$$
(16)
$`n`$ being the (random) integer value for which $`t_nt`$ and $`t_{n+1}>t`$. The total time $`t`$ can be rewritten as $`t=_{i=1}^n\tau _i+ϵ_{n+1}`$. Denoting by $`Nn`$ the average value of the number of time steps necessary to reach the time $`t`$, one has for large times $`tN\tau `$ and thus $`x(t)=_{i=1}^nv_i\tau _i`$.
From simple considerations related to the symmetry of the velocity pdf under the transformation $`vv`$, it immediately follows that the odd-order moments $`x(t)^q`$ are trivially zero. Conversely, even-order moments are nonzero and can be evaluated exploiting the following properties: $`v_i\tau _j=0`$, $`v_iv_j\delta _{ij},\tau _i\tau _j\delta _{ij}`$. For times large enough, the mean squared displacement thus reads:
$$x(t)^2=\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{n}{}}(v_i\tau _i)(v_j\tau _j)=N(v\tau )^2.$$
(17)
Similarly, for the fourth and sixth order moments the limit of large times yields
$`x^4(t)`$ $`=`$ $`N(v\tau )^4+3N^2(v\tau )^2^2`$ (18)
$`x^6(t)`$ $`=`$ $`N(v\tau )^6+15N^2(v\tau )^2(v\tau )^4+`$ (20)
$`15N^3(v\tau )^2^3.`$
and so on.
Our attention being focused on the behavior of $`x(t)^q`$ as a function of $`t`$, we make the substitution $`N=t/\tau `$ in previous expressions (17) and (20) and again exploit the facts that $`v_i`$ and $`\tau _i`$ are uncorrelated and $`v_i^2=1`$ to get:
$`x^2(t)`$ $`=`$ $`{\displaystyle \frac{t}{\tau }}\tau ^2`$ (21)
$`x^4(t)`$ $`=`$ $`{\displaystyle \frac{t}{\tau }}\tau ^4+3\left({\displaystyle \frac{t}{\tau }}\right)^2\tau ^2^2`$ (22)
$`x^6(t)`$ $`=`$ $`{\displaystyle \frac{t}{\tau }}\tau ^6+15\left({\displaystyle \frac{t}{\tau }}\right)^2\tau ^2\tau ^4+`$ (24)
$`15\left({\displaystyle \frac{t}{\tau }}\right)^3\tau ^2^3.`$
For times $`tT`$ the leading term for $`x^{2n}(t)`$ is that one proportional to $`t^n`$, therefore one has:
$$x(t)^{2n}t^n\left(\frac{\tau ^2}{\tau }\right)^n\mathrm{w}herenisinteger.$$
(25)
which is just ordinary diffusive behaviour. Diffusive behavior in such limit is actually expected from general considerations. Indeed, when $`tT`$ the particle position at the time $`t`$ can be rearranged in the form of a sum of almost independent displacements. If the number of the latter is large enough, central limit arguments apply, with the immediate consequence is that particles undergo diffusive motion. Such result can also be rigorously proved exploiting multiscale perturbative expansions in the (small) parameter $`T/t`$ as done, e.g., in Ref. .
In the opposite regime, where $`tT`$, the system is strongly correlated, central limit arguments does not apply, and the final result is that non-diffusive (i.e. anomalous) regimes can occur where $`x^q(t)t^{q\nu (q)}`$ with $`\nu (q)1/2`$.
The possible emergence of anomalous behaviors in the limit $`t\mathrm{}`$ can be investigated by looking at the dependence of moments on the cutoff $`T`$. For times shorter than $`T`$ they behave as $`x^q(t)t^{q\nu (q)}`$, but around $`tT`$ the moments have a crossover to diffusive behaviour, $`x^q(t)(t/\tau )^{q/2}\tau ^2^{q/2}`$.
By matching the two different regimes at $`t=T`$, and using the results $`\tau ^qT^{g+q+1}`$ for $`g+q+1>0`$ and $`\tau ^q=O(1)`$ for $`g+q+1<0`$, the following expressions for $`q\nu (q)`$ are found as a function of the exponent $`g`$:
$$\begin{array}{cccc}g(1,2]\hfill & q\nu (q)=q\hfill & q=2,4,6,\mathrm{}\hfill & \\ & & & \\ g(2,3]\hfill & q\nu (q)=q+2g\hfill & q=2,4,6,\mathrm{}\hfill & \\ & & & \\ g(3,4]\hfill & q\nu (q)=q/2\hfill & q=2\hfill & \\ & q\nu (q)=q+2g\hfill & q=4,6,8,\mathrm{}\hfill & \end{array}$$
(26)
and so on for higher values of $`g`$.
From (26) it follows that the anomalous diffusion phase takes place for higher and higher moments as $`g`$ increases. For $`q`$ large enough one has $`q\nu (q)=q+\mathrm{c}onstant`$ and noting that $`2\nu (2)2`$ one can conclude that $`\nu (q)`$ can not be constant and therefore a strong anomalous diffusion is present.
The matching argument had been used in Ref. in the context of the multiscale method for the anomalous diffusion in random shear flows .
### C Generalized CTRW model
We now present a generalization of the previous model showing a strong anomalous diffusion regime characterized by a non-piecewise linear behavior as a function of the order $`q`$.
As in the previous case the particle moves ballistically with random velocity $`v_i`$ during the random interval times $`\tau _i=t_{i+1}t_i`$. We assume that $`\tau _i`$ has the same pdf as in (15) but now the velocities $`v_i`$ assume the value $`\pm \tau _i^h`$ where $`h`$ is a random positive variable conditioned on $`\tau _i`$. Specifically, the conditional pdf $`P(h\tau )`$ is
$$P(h\tau )\tau ^{S(h)}$$
(27)
where $`S(h)`$ is a positive smooth concave function. This has the effect of giving a larger variance to the velocity, the larger the time $`\tau _i`$. The function $`S(h)`$ can be taylored to a special need, i.e., to mimic the intermittent turbulent velocity field. Here we just use a generic, simple function $`S(h)=h^2/(2\sigma ^2)`$, as we are mainly interested in the generic properties of the model.
The moments $`(v\tau )^{2q}`$ can now be calculated:
$`(v\tau )^{2q}`$ $`=`$ $`{\displaystyle _0^+\mathrm{}}𝑑h{\displaystyle _{t_{min}}^T}𝑑\tau (v\tau )^{2q}P(\tau )P(h\tau )`$ (30)
$`{\displaystyle _{t_{min}}^T}𝑑\tau \tau ^{g+2q}{\displaystyle _0^+\mathrm{}}𝑑h\tau ^{2qhS(h)}=`$
$`{\displaystyle _{t_{min}}^T}𝑑\tau \tau ^{g+2q+y(2q)}T^{g+2q+y(2q)+1}.`$
where the integral $`𝑑h\tau ^{2qhS(h)}`$ has been evaluated exploiting the steepest descent method and we have defined
$$y(2q)\underset{h}{\mathrm{max}}[2qhS(h)].$$
(31)
Now using the specific shape of $`S(h)`$, expression (31) takes the form:
$$y(2q)=q^2y(2)\mathrm{with}y(2)=2\sigma ^2.$$
(32)
Exploiting the same matching arguments as in section II B, the following expressions for the exponents $`q\nu (q)`$ are obtained:
$$\begin{array}{cccc}g(1,2]\hfill & q\nu (q)=\frac{q^2}{2}\sigma ^2+q\hfill & q=2,4,6,\mathrm{}\hfill & \\ & & & \\ g(2,3+y(2)]\hfill & q\nu (q)=\frac{q^2}{2}\sigma ^2+q+2g\hfill & q=2,4,6,\mathrm{}\hfill & \\ & & & \\ g(3+y(2),4+y(4)]\hfill & q\nu (q)=q/2\hfill & q=2\hfill & \\ & q\nu (q)=\frac{q^2}{2}\sigma ^2+q+2g\hfill & q=4,6,8,\mathrm{}\hfill & \end{array}$$
(33)
and so on, for higher values of $`g`$.
We now see that in the general case of the CTRW, $`q\nu (q)`$ is not just a bi-linear function, and the exact form depends upon the shape of $`S(h)`$. The generalization to an arbitrary shape of $`S(h)`$ is straightforward.
As far as we know it is not simple to obtain the results in Eq.(26) and Eq.(33) with the method discussed in section II A.
## III Numerical simulations
We present results of numerical simulations of the general model defined by Eqs. (15) and (27). The goal of the numerical simulations has been to verify the validity of the theoretical expectations (26) and (33) and to study the pdf of $`x(t)`$ in more detail.
In practice the pdf for the length of the jumps $`P(\tau )`$ has to be supplemented with a lower cutoff. As it is the tail of the distribution which governs the scaling of the moments of $`x(t)`$ the results are independent upon how the cutoff is made. We used a hard cutoff at $`\tau =1`$, such that $`P(\tau )=0`$ for $`\tau <1`$ and
$$P(\tau )=(g1)\tau ^g\mathrm{f}or\tau >1.$$
(34)
The moments have been calculated as an ensemble average of many realizations of the process for times $`0<t<\stackrel{~}{T}`$. For all the simulations presented here, $`\stackrel{~}{T}`$ has been set to $`10^6`$.
In Fig. 1 the scaling of the second moment $`x^2(t)`$ with $`t`$ for different values of $`g`$ and $`\sigma =0.2`$ is shown. Examples have been chosen where diffusive, anomalous and ballistic behavior is expected. After an initial ballistic motion for short times, a transient towards scaling is taking place. For long times, clean scaling with an exponent corresponding to Eqs. (26) is evident.
To show how the introduction of the velocity field from section II C influences the moments, we have calculated the higher order moments of $`x(t)`$ for $`\sigma =0`$,$`0.2`$ and $`0.4`$ for $`g=2.5`$ (Fig. 2). For the case with $`\sigma =0`$ we see a clear bi-linear behavior as expected, with a cross over at $`q_c1.5`$. For low order moments ($`q<1`$) the core of the pdf is important, and the behavior approches $`1/(g1)`$. As is seen in the inset, this behavior is only to be strictly valid in the limit $`q0`$. The scaling of the moments changes smoothly into the prediction $`\nu (q)=q+2g`$ for $`q2`$. Using the fact that $`q\nu (q)`$ is a concave increasing function and, for $`\sigma =0`$, the slope can not be larger than 1, one obtains that the prediction given by Eq. 26, obtained only for even order moments, surely is valid for any moment larger than 2. For the generalized model ($`\sigma >0`$), the fluctuations become much stronger, and it was not possible to get a clear convergence for large orders. However, for $`q=2`$ the prediction in Eq. 33 and the numerical results are in perfect agreement.
The characterization of the process by the scaling exponents of the moments can be seen as a way to probe the pdf $`p(x,t)`$ of the process. Roughly speaking, the low order structure functions characterize the core of the pdf, while higher orders characterize the tails. In the case of “ordinary” anomalous diffusion, the pdfs can be rescaled according to Eq. 2, due to the fact that the whole process can be characterized by just one scaling exponent. However for the bi-linear or the strong diffusive cases, no such renormalization can be done, as more than one exponent is needed to characterize the process. As all the cases seem to have the same limiting behavior for $`q0`$, it might be possible to make a collapse of the core of the pdf. In Fig. 3 the pdf for $`g=2.5`$ and $`\sigma =0.2`$ has been rescaled using the typical value $`x_{typ}(t)=\mathrm{exp}(\mathrm{ln}(|x(t)|)`$. As expected the core of the rescaled pdfs show a good overlap, while the tails diverge. For higher values of $`\sigma `$ the area of the core with overlap becomes smaller and smaller.
## IV Discussions and conclusions
By using only elementary techniques, we have shown in this paper that the strong anomalous diffusion appears in CTRW. Our approach, in which one computes the even moments $`x(t)^q`$ for a system with a cutoff $`T`$ on the pdf $`P(\tau )`$, and then the matching of the behaviors at $`tT`$ and $`tT`$, allows us to treat rather general cases. This is a relevant advantage with respect to the approach commonly used for CTRW which seems to us to be difficult to apply for more general cases. In the case of generalized CTRW, i.e. with Eq. 27, one has a nontrivial (non bi-linear) shape of $`q\nu (q)`$. This implies that the pdf $`P(x,t)`$ cannot be written in the form (2). On the other hand we found that the rescaling (2), with $`\nu =\nu (0)`$, is valid in a limited range of $`x/x_{typ}`$. It is rather natural to wonder if, at least in this limited range, the function $`F(\xi )`$ is determined by the exponent $`\nu (0)`$. From previous works and the results discussed in section III, it seems to us that one has a negative answer: $`\nu `$ does not determine the function $`F(\xi )`$.
As the shape of the function describing the scaling of the moments looks similar to what one finds for the relative dispersion of a passive scalar by a turbulent velocity field, one might like to use the the CTRW model to describe this process. There is however some differences that has to be taken into consideration. In fully developed turbulence, a flight of duration $`\tau `$ can be considered as associated to an eddy of size $`l\tau ^{3/2}`$. In a process of relative diffusion, one typically has an increase of the relative distance with the time. This implies that a flight of duration $`\tau `$ is preferentially followed by another flight with a larger duration than the previous. The easiest way for a realistic description of relative dispersion in fully developed turbulence should thus consist in a further generalization of the CTRW where a Markov ingredient, i.e. a dependence of the flight duration at the time $`t`$ on the duration at the previous time, is introduced. It seems to us that this is a nontrivial task; a first attempt in this direction can be found in the recent work .
## Acknowledgments
We are particularly grateful to Paolo Muratore-Ginanneschi for very stimulating discussions and suggestions and J. Klafter for useful correspondence. The authors thank M. H. Jensen for fruitful discussions and all the CATS staff for the nice hospitality at Niels Bohr Institute in Copenhagen. A.M. was partially supported by the INFM PA GEPAIGG01. P.C. and A.V. are partially supported by INFM (Progetto Ricerca Avanzata–TURBO) and by MURST (program no. 9702265437). P.C. and K.H.A. acknowledge the support of the European Commission’s TMR Programme, Contract no.ERBFMRXCT980175, “Intermittency in turbulent systems”. Simulations were performed at CINECA (INFM Parallel Computing Initiative).
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# Next-to-leading Order Calculation of the Color-Octet ³𝑆₁ Gluon Fragmentation Function for Heavy Quarkonium
## I Introduction
The cross sections for heavy quarkonium states probe the production of heavy-quark-antiquark pairs with small relative momenta. Many of the theoretical uncertainties in quarkonium production decrease at large transverse momentum. Factorization theorems for inclusive single-hadron production guarantee that the dominant mechanism for producing heavy quarkonia with high $`p_T`$ is fragmentation, the production of a parton which subsequently decays into the quarkonium state and other partons. This process is described by a fragmentation function $`D(z,\mu )`$, where $`z`$ is the longitudinal momentum fraction of the quarkonium state and $`\mu `$ is a factorization scale.
The NRQCD factorization formalism can be used to factor the fragmentation functions $`D(z,\mu )`$ for quarkonium into NRQCD matrix elements, which can be regarded as phenomenological parameters, and short-distance factors, which depend on $`z`$ and are calculable in perturbation theory. Most of the phenomenologically relevant short-distance factors begin at order $`\alpha _s^2`$ and have been calculated to leading order. However there is one short-distance factor that begins at order $`\alpha _s`$. It is the color-octet $`{}_{}{}^{3}S_{1}^{}`$ term in the gluon fragmentation function, whose NRQCD matrix element is denoted $`O_8(^3S_1)`$. This term is of particular phenomenological importance. Braaten and Yuan showed that it must be included in the gluon fragmentation function for triplet P-wave states in order to avoid an infrared divergence in the short-distance coefficient of the color-singlet matrix element $`O_1(^3P_J)`$ . Braaten and Fleming argued that the $`O_8(^3S_1)`$ term is also phenomenologically necessary in the gluon fragmentation function for spin-triplet S-wave states in order to explain the production rate of direct $`J/\psi `$ and $`\psi ^{}`$ at large $`p_T`$ at the Tevatron . This led to the remarkable prediction by Cho and Wise that $`J/\psi `$ and $`\psi ^{}`$ at large $`p_T`$ should be transversely polarized .
In the earliest calculations of fragmentation functions for heavy quarkonium , the short-distance factors were deduced by comparing the cross sections for quarkonium production with the form predicted by the factorization theorems for inclusive single-hadron production. However the fragmentation functions can also be defined formally as matrix elements of bilocal operators in a light-cone gauge or, more generally, as matrix elements of non-local gauge-invariant operators . The gauge-invariant definition of Collins and Soper was first applied to calculations of the fragmentation functions for heavy quarkonium by Ma . The definition is particularly convenient for carrying out calculations beyond leading order in $`\alpha _s`$. It was used by Ma to calculate the short-distance factor of the color-octet $`{}_{}{}^{3}S_{1}^{}`$ term in the gluon fragmentation function to next-to-leading order in $`\alpha _s`$ .
In this paper, we calculate the short-distance coefficients for the color-octet $`{}_{}{}^{3}S_{1}^{}`$ term in the fragmentation function for a gluon to split into heavy quarkonium states to order $`\alpha _s^2`$. We use the gauge-invariant definition of the fragmentation function given by Collins and Soper, and we remove ultraviolet divergences using the $`\overline{\mathrm{MS}}`$ renormalization procedure. Our result for the longitudinal term in the fragmentation function agrees with a previous calculation by Beneke and Rothstein . Our result for the next-to-leading order correction in the transverse term disagrees with a previous calculation by Ma .
## II Gauge-invariant Definition
The fragmentation function $`D_{gH}(z,\mu )`$ gives the probability that a gluon produced in a hard-scattering process involving momentum transfer of order $`\mu `$ decays into a hadron $`H`$ carrying a fraction $`z`$ of the gluon’s longitudinal momentum. This function can be defined in terms of the matrix element of a bilocal operator involving two gluon field strengths in a light-cone gauge . In Ref. , Collins and Soper introduced a gauge-invariant definition of the gluon fragmentation function that involves the matrix element of a nonlocal operator consisting of two gluon field strengths and eikonal operators. One advantage of this definition is that it avoids subtleties associated with products of singular distributions. The gauge-invariant definition is also advantageous for explicit perturbative calculations, because it allows the calculation of radiative corrections to be simplified by using Feynman gauge.
The gauge-invariant definition of Collins and Soper is
$`D_{gH}(z,\mu )`$ $`=`$ $`{\displaystyle \frac{(g_{\mu \nu })z^{N2}}{16\pi (N1)k^+}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x^{}e^{ik^+x^{}}`$ (2)
$`0|G_c^{+\mu }(0)^{}(0^{})_{cb}𝒫_{H(zk^+,0_{})}(x^{})_{ba}G_a^{+\nu }(0^+,x^{},0_{})|0.`$
The operator $`(x^{})`$ in (2) is an eikonal operator that involves a path-ordered exponential of gluon field operators along a light-like path:
$$(x^{})_{ba}=\mathrm{P}\mathrm{exp}[+ig_x^{}^{\mathrm{}}dz^{}A^+(0^+,z^{},0_{})]_{ba},$$
(3)
where $`A^\mu (x)`$ is the matrix-valued gluon field in the adjoint representation: $`[A^\mu (x)]_{ac}=if^{abc}A_b^\mu (x)`$. The operator $`𝒫_{H(p^+,p_{})}`$ in (2) is a projection onto states that in the asymptotic future contain a hadron $`H`$ with momentum $`p=(p^+,p^{}=(m_H^2+p_{}^2)/p^+,p_{})`$, where $`m_H`$ is the mass of the hadron. In the definition (2), the hard-scattering scale $`\mu `$ can be identified with the renormalization scale of the nonlocal operator. In perturbative calculations, it is convenient to use dimensional regularization to regularize ultraviolet divergences. The prefactor in the definition (2) has therefore been expressed as a function of the number of spatial dimensions $`N=32ϵ`$.
For any state $`H`$ that can be defined in perturbation theory, the definition (2) can be used to calculate the fragmentation function $`D_{gH}(z,\mu )`$ as a power series in $`\alpha _s`$. A convenient set of Feynman rules for this perturbative expansion is given in Ref.. If the state consists of a $`Q\overline{Q}`$ pair with invariant mass $`p^2`$, the lowest order diagram is shown in Fig. 1. The circles connected by the double pair of lines represent the nonlocal operator consisting of the gluon field strengths and the eikonal operators. The momentum $`k=(k^+,k^{},k_{})`$ flows into the circle on the left and out the circle on the right. The cutting line represents the projection onto states that in the asymptotic future include a $`Q\overline{Q}`$ pair with total momentum $`p=(zk^+,p^2/(zk^+),0_{})`$. If the heavy quark has relative momentum $`𝐪`$ in the $`Q\overline{Q}`$ rest frame, the invariant mass is $`p^2=4(m_Q^2+𝐪^2)`$.
The fragmentation of a gluon into heavy quarkonium state $`H`$ involves many momentum scales, ranging from the hard-scattering scale $`\mu `$, which we will assume to be larger than $`m_Q`$, to momenta much smaller than $`m_Q`$ where nonperturbative effects are large. The NRQCD factorization formalism allows the systematic separation of momentum scales of order $`m_Q`$ and larger from scales of order $`m_Qv`$ or smaller, where $`v`$ is the typical relative velocity of the heavy quark in the hadron. The factorization formula has the form
$`D_{gH}(z,\mu )={\displaystyle \underset{mn}{}}d_{mn}(z,\mu )𝒪_{nm}^H.`$ (4)
The short-distance coefficients $`d_{mn}(z,\mu )`$ are independent of the quarkonium state $`H`$ and can be calculated as a perturbation series in $`\alpha _s(m_Q)`$. All long-distance effects are factored into the NRQCD matrix elements $`𝒪_{nm}^H`$, which can be expressed as matrix elements in an effective field theory. They have the general form
$`𝒪_{nm}^H`$ $`=`$ $`0|\chi ^{}𝒦_n\psi 𝒫_H\psi ^{}𝒦_m\chi |0,`$ (5)
where $`𝒦_m`$ and $`𝒦_n`$ are constructed out of color matrices, spin matrices, and covariant derivatives and the operator $`𝒫_H`$ projects onto states that in the asymptotic future contain a quarkonium state $`H`$. The NRQCD matrix elements are nonperturbative but they are universal, with the same matrix elements describing inclusive production in other high energy processes.
The threshold expansion method is a general prescription for determining the short-distance coefficients $`d_{nm}(z,\mu )`$ in the NRQCD factorization formula (4). The diagrams for the fragmentation function of a $`Q\overline{Q}`$ pair are computed in perturbation theory, except that the $`Q\overline{Q}`$ pair is allowed to be in a different state on the two sides of the final-state cut. To the left of the cutting line, the $`Q`$ and $`\overline{Q}`$ have relative momentum $`𝐪`$ in the $`Q\overline{Q}`$ rest frame and color and spin states specified by Pauli spinors $`\eta `$ and $`\xi `$. To the right of the cutting line, the $`Q`$ and $`\overline{Q}`$ have relative momentum $`𝐪^{}`$ and color and spin states specified by $`\eta ^{}`$ and $`\xi ^{}`$. After expanding around the threshold $`𝐪=𝐪^{}=0`$, the resulting expression for the diagrams has the form
$`{\displaystyle \underset{mn}{}}d_{mn}(z,\mu )\eta _{}^{}{}_{}{}^{}\kappa _n\xi ^{}\xi ^{}\kappa _m\eta ,`$ (6)
where the color and spin matrices $`\kappa _m`$ and $`\kappa _n`$ are polynomials in the relative momenta $`𝐪`$ and $`𝐪^{}`$. The short-distance coefficients $`d_{nm}(z,\mu )`$ in the factorization formula (4) can then be read off from this expression. For example, if the sum of the diagrams includes the color-octet $`{}_{}{}^{3}S_{1}^{}`$ terms
$`Nm_Q\left[d_T(z)(\delta ^{ij}\widehat{z}^i\widehat{z}^j)+d_L(z)\widehat{z}^i\widehat{z}^j\right]\eta _{}^{}{}_{}{}^{}\sigma ^jT^a\xi ^{}\xi ^{}\sigma ^iT^a\eta ,`$ (7)
then the fragmentation function includes the terms
$`D_{gH}(z)`$ $`=`$ $`{\displaystyle \frac{N}{4m_Q}}\left[d_T(z)(\delta ^{ij}\widehat{z}^i\widehat{z}^j)+d_L(z)\widehat{z}^i\widehat{z}^j\right]0|\chi ^{}\sigma ^jT^a\psi 𝒫_H\psi ^{}\sigma ^iT^a\chi |0.`$ (8)
If we sum over the spin states of $`H`$, the matrix element in (8) is proportional to $`\delta ^{ij}`$ and (8) reduces to
$`D_{gH}(z)=[(N1)d_T(z)+d_L(z)]𝒪_8^H(^3S_1),`$ (9)
where $`𝒪_8^H(^3S_1)`$ is the color-octet $`{}_{}{}^{3}S_{1}^{}`$ matrix element defined in Ref. :
$`𝒪_8^H(^3S_1)={\displaystyle \frac{1}{4m_Q}}0|\chi ^{}\sigma ^iT^a\psi 𝒫_H\psi ^{}\sigma ^iT^a\chi |0.`$ (10)
The factor of $`1/(4m_Q)`$ accounts for the relativistic normalization of the projection operator used in Refs. .
## III Leading order
The only Feynman diagram of order $`\alpha _s`$ for the fragmentation process $`gQ\overline{Q}`$ is shown in Fig. 1. Using the Feynman rules of Ref. , the expression for the diagram can be easily written down in terms of spinors $`u`$ and $`v`$ that describe the color, spin, and relative momentum states of the $`Q`$ and $`\overline{Q}`$. We can allow for the $`Q`$ and $`\overline{Q}`$ on the right side of the final-state cut to have different color and spin states and different relative momentum $`𝐪^{}`$ by replacing their spinors by $`u^{}`$ and $`v^{}`$. The resulting expression is
$`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{2(N1)(p^2)^2}}\delta (1z)\left(g^{\alpha \beta }{\displaystyle \frac{p^2}{(kn)^2}}n^\alpha n^\beta \right)\overline{v}^{}\gamma _\beta T^au^{}\overline{u}\gamma _\alpha T^av.`$ (11)
Setting $`𝐪=𝐪^{}=0`$, we can replace the Dirac spinors with Pauli spinors by substituting
$`\overline{u}\gamma _\alpha T^av=\mathrm{\hspace{0.33em}2}m_QL_{\alpha i}\xi ^{}\sigma ^iT^a\eta ,`$ (12)
where $`L_{\alpha i}`$ is a boost matrix that satisfies the identities
$`g^{\alpha \beta }L_{\alpha i}L_{\beta j}`$ $`=`$ $`\delta ^{ij},`$ (13)
$`n^\alpha L_{\alpha i}`$ $`=`$ $`{\displaystyle \frac{pn}{\sqrt{p^2}}}\widehat{z}^i.`$ (14)
The expression (11) then reduces to
$`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{8(N1)m_Q^2}}\delta (1z)(\delta ^{ij}\widehat{z}^i\widehat{z}^j)\eta _{}^{}{}_{}{}^{}\sigma ^jT^a\xi ^{}\xi ^{}\sigma ^iT^a\eta .`$ (15)
Comparing with (7), we can read off the order-$`\alpha _s`$ terms in the short-distance functions $`d_T(z)`$ and $`d_L(z)`$ defined in (8):
$`d_T^{(\mathrm{LO})}(z)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{8N(N1)m_Q^3}}\delta (1z),`$ (16)
$`d_L^{(\mathrm{LO})}(z)`$ $`=`$ $`0.`$ (17)
The dependence on the number of spatial dimensions $`N`$ agrees with that in Ref. .
## IV Virtual Corrections
The Feynman diagrams for the fragmentation function for $`gQ\overline{Q}`$ at order $`\alpha _s^2`$ consist of virtual corrections, for which the final state is $`Q\overline{Q}`$, and real-gluon corrections, for which the final state is $`Q\overline{Q}g`$. The diagrams with virtual-gluon corrections to the left of the cutting line are shown in Fig. 2. The black blob in Fig. 2(a) includes the vertex corrections and propagator corrections shown in Fig. 3.
We calculate the diagrams using Feynman gauge. In this case, the diagram Fig. 2(b) vanishes, because the gluon attaching to the eikonal line gives a factor of $`n^\mu `$. When contracted with the factor $`kng^{\nu \beta }p^\nu n^\beta `$ from the circle on the left side of the cut, it gives a factor $`(kp)n=0`$. In the limit $`𝐪=𝐪^{}=0`$, all the other diagrams can be reduced to the leading order diagram in Fig. 1 times a multiplicative factor. For the diagrams in Fig. 2(a), this follows from the Dirac equation for the heavy-quark spinors. For the remaining diagrams, we must also use the fact that the contraction of $`n^\mu `$ with the right side of the diagram vanishes.
The virtual corrections contribute only to the transverse short-distance function $`d_T(z)`$ defined in (8). We will express the various contributions in the form of the leading-order result (16) times a multiplicative factor. The sum of the diagrams in Fig. 2(a), together with their complex conjugates, is
$`d_T^{(2a)}(z)`$ $`=`$ $`d_T^{(\mathrm{LO})}(z)\times 2\mathrm{Re}\left[\mathrm{\Lambda }+\mathrm{\Pi }+(Z_Q1)\right],`$ (18)
where $`\mathrm{\Lambda }`$ is the vertex correction factor from the diagrams in Fig. 3(b), $`\mathrm{\Pi }`$ is the propagator correction factor for a gluon with invariant mass $`4m_Q^2`$ from the diagrams in Fig. 3(d), and $`Z_Q1`$ comes from the wavefunction renormalization factors $`Z_Q^{1/2}`$ for the heavy quark from the diagrams in Fig. 3(c). These correction factors are given by
$`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\left[{\displaystyle \frac{13}{12}}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}}}{\displaystyle \frac{1}{12}}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{IR}}}+{\displaystyle \frac{13}{6}}+3\mathrm{ln}2i{\displaystyle \frac{\pi }{2}}\right],`$ (19)
$`Z_Q1`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\left[{\displaystyle \frac{1}{3}}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{IR}}}{\displaystyle \frac{4}{3}}2\mathrm{ln}2\right],`$ (20)
$`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}e^{i\pi }\right)^ϵ\left[{\displaystyle \frac{152n_f}{12}}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}}}+{\displaystyle \frac{9310n_f}{36}}\right],`$ (21)
where $`n_f`$ is the number of light quark flavors. The subscripts on the poles in $`ϵ`$ indicate whether the divergences are of ultraviolet or infrared origin. The sum of the virtual corrections given in (18) is
$`d_T^{(2a)}(z)`$ $`=`$ $`d_T^{(\mathrm{LO})}(z)\times {\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\mathrm{\Gamma }(1+ϵ)\left[{\displaystyle \frac{12n_f}{3ϵ_{UV}}}{\displaystyle \frac{3}{2ϵ_{IR}}}+{\displaystyle \frac{12310n_f}{18}}+2\mathrm{ln}2\right],`$ (22)
Ma’s result for these diagrams differs by an additive constant $`5.487`$ inside the square brackets and by setting $`N(N1)3(N1)`$ in $`d_T^{(\mathrm{LO})}(z)`$.
The contribution from the diagrams in Fig. 2(c) with its complex conjugate is
$`d_T^{(2c)}(z)`$ $`=`$ $`d_T^{(\mathrm{LO})}(z)\times 24\pi \alpha _s\mu ^{2ϵ}\mathrm{Re}\left[i\left(2knI_{ABC}I_{AB}\right)\right],`$ (23)
where the scalar integrals $`I_{AB\mathrm{}}`$ are given in appendix A. The contribution from the sum of the diagrams in Figs. 2(d) and 2(e) with their complex conjugates is
$`d_T^{(2d,e)}(z)`$ $`=`$ $`d_T^{(\mathrm{LO})}(z)\times 96\pi \alpha _s\mu ^{2ϵ}m_Q^2\mathrm{Re}\left[i\left(knI_{ABCD}I_{ABD}\right)\right].`$ (24)
The sum of the virtual corrections given in (23) and (24) is
$`d_T^{(2c,d,e)}(z)`$ $`=`$ $`d_T^{(\mathrm{LO})}(z)\times {\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\mathrm{\Gamma }(1+ϵ)[{\displaystyle \frac{3(1ϵ)}{2ϵ_{UV}ϵ_{IR}}}+{\displaystyle \frac{3}{2ϵ_{UV}}}`$ (26)
$`+3{\displaystyle \frac{\pi ^2}{2}}+6\mathrm{ln}2+6\mathrm{ln}^22].`$
Ma’s result for these diagrams differs by an additive constant $`3+\frac{3}{2}\pi ^26\mathrm{ln}^22`$ inside the square brackets and by setting $`N(N1)3(N1)`$ in $`d_T^{(\mathrm{LO})}(z)`$.
The total virtual correction is the sum of (22) and (26):
$`d_T^{(\mathrm{virtual})}(z)`$ $`=`$ $`d_T^{(\mathrm{LO})}(z)\times {\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ[{\displaystyle \frac{3(1ϵ)}{2}}{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}ϵ_{IR}}}+\beta _0{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}}}`$ (28)
$`+{\displaystyle \frac{17710n_f}{18}}{\displaystyle \frac{\pi ^2}{2}}+8\mathrm{ln}2+6\mathrm{ln}^22],`$
where $`\beta _0=(332n_f)/6`$.
## V Real Gluon Corrections
The Feynman diagrams for the real-gluon corrections to the fragmentation function for $`gQ\overline{Q}`$ pair are shown in Fig. 4. We draw the 5 left-half diagrams only, but they must be multiplied by their complex conjugates to give a total of 25 diagrams. Only diagrams 4(a), 4(b) and 4(c) contribute to the color-octet $`{}_{}{}^{3}S_{1}^{}`$ term, which reduces the total number of diagrams to 9. We calculate the diagrams using Feynman gauge. After a considerable amount of algebra, they reduce to
$`d_T^{(\mathrm{real})}(z)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{8N(N1)m_Q^3}}\times {\displaystyle \frac{3\alpha _s}{\pi \mathrm{\Gamma }(1ϵ)}}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\left(1{\displaystyle \frac{1}{z(1z)}}\right)^2{\displaystyle _{(1z)/z}^{\mathrm{}}}𝑑x{\displaystyle \frac{t^{1ϵ}}{x^2}},`$ (29)
$`d_L^{(\mathrm{real})}(z)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{8Nm_Q^3}}\times {\displaystyle \frac{3\alpha _s}{\pi \mathrm{\Gamma }(1ϵ)}}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\left({\displaystyle \frac{1z}{z}}\right)^2{\displaystyle _{(1z)/z}^{\mathrm{}}}𝑑x{\displaystyle \frac{t^ϵ}{x^2}},`$ (30)
where $`t=(1z)(zx+z1)`$, $`x=2qp/p^2`$, $`q`$ is the final-state gluon momentum, and $`p`$ is the $`Q\overline{Q}`$ momentum. Integration over $`x`$ can be done using the following identities:
$`{\displaystyle _{(1z)/z}^{\mathrm{}}}𝑑x{\displaystyle \frac{t^{1ϵ}}{x^2}}`$ $`=`$ $`z(1z)^{12ϵ}{\displaystyle \frac{1ϵ}{ϵ_{UV}}}B(1+ϵ,1ϵ),`$ (31)
$`{\displaystyle _{(1z)/z}^{\mathrm{}}}𝑑x{\displaystyle \frac{t^ϵ}{x^2}}`$ $`=`$ $`z(1z)^{12ϵ}B(1+ϵ,1ϵ).`$ (32)
The subscript on $`ϵ`$ in (31) indicates that the pole has an ultraviolet origin. In (29), there is also an infrared divergence associated with the limit $`z1`$. It can be made explicit by using the expansion:
$`{\displaystyle \frac{1}{(1z)^{1+2ϵ}}}`$ $`=`$ $`{\displaystyle \frac{1}{2ϵ_{IR}}}\delta (1z)+{\displaystyle \frac{1}{(1z)_+}}2ϵ\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_++O(ϵ^2).`$ (33)
The final results for the real-gluon corrections are
$`d_T^{(\mathrm{real})}(z)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{8N(N1)m_Q^3}}\times {\displaystyle \frac{\alpha _s}{\pi }}\left({\displaystyle \frac{\pi \mu ^2}{m_Q^2}}\right)^ϵ\mathrm{\Gamma }(1+ϵ)[{\displaystyle \frac{3(1ϵ)}{2ϵ_{UV}ϵ_{IR}}}\delta (1z)`$ (36)
$`+{\displaystyle \frac{3(1ϵ)}{ϵ_{UV}}}\left({\displaystyle \frac{z}{(1z)_+}}+{\displaystyle \frac{1z}{z}}+z(1z)\right)`$
$`{\displaystyle \frac{6}{z}}\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_++6(2z+z^2)\mathrm{ln}(1z)],`$
$`d_L^{(\mathrm{real})}(z)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s}{8Nm_Q^3}}\times {\displaystyle \frac{3\alpha _s}{\pi }}{\displaystyle \frac{1z}{z}}.`$ (37)
Note that the double-pole term proportional to $`1/(ϵ_{UV}ϵ_{IR})`$ in (36) exactly cancels its counter part in the virtual correction (28). The sum of (28) and (36) is free of infrared divergences. Ma’s result for the real-gluon corrections is completely different from the sum of (36) and (37). In particular, he found that the $`[\mathrm{ln}(1z)/(1z)]_+`$ terms cancelled. In the way we organized the calculation, there is no possibility of such a cancellation.
## VI Renormalization
The sum of the order-$`\alpha _s^2`$ corrections to $`d_T(z)`$ in (28) and (36) still contains ultraviolet divergences in the form of single poles in $`ϵ`$. These divergences are cancelled by the renormalization of the coupling constant $`\alpha _s`$ in the leading-order expression (17) and by the renormalization of the nonlocal operator in (2). The renormalization of $`\alpha _s`$ in the $`\overline{\mathrm{MS}}`$ scheme can be carried out by making the following substitution in (17):
$`\alpha _s\alpha _s\left[1{\displaystyle \frac{\alpha _s}{2\pi }}\beta _0\left(4\pi e^\gamma \right)^ϵ{\displaystyle \frac{1}{ϵ_{UV}}}\right],`$ (38)
where $`\beta _0=(332n_f)/6`$. The operator renormalization in the $`\overline{\mathrm{MS}}`$ scheme can be carried out by making the following substitution
$`d_T^{(\mathrm{LO})}(z)d_T^{(\mathrm{LO})}(z){\displaystyle \frac{\alpha _s}{2\pi }}\left(4\pi e^\gamma \right)^ϵ{\displaystyle \frac{1}{ϵ_{UV}}}{\displaystyle _z^1}{\displaystyle \frac{dy}{y}}P_{gg}(y)d_T^{(\mathrm{LO})}(z/y),`$ (39)
where $`P_{gg}(y)`$ is the gluon splitting function:
$`P_{gg}(z)=6\left[{\displaystyle \frac{z}{(1z)_+}}+{\displaystyle \frac{1z}{z}}+z(1z)+{\displaystyle \frac{\beta _0}{6}}\delta (1z)\right].`$ (40)
The sum of the two contributions of order $`\alpha _s^2`$ from renormalization is
$`d_T^{(\mathrm{ren})}(z)`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s\mu ^{2ϵ}}{8N(N1)m_Q^3}}`$ (41)
$`\times `$ $`{\displaystyle \frac{\alpha _s}{\pi }}\left(4\pi e^\gamma \right)^ϵ{\displaystyle \frac{(3)}{ϵ_{UV}}}\left[{\displaystyle \frac{z}{(1z)_+}}+{\displaystyle \frac{1z}{z}}+z(1z)+{\displaystyle \frac{\beta _0}{3}}\delta (1z)\right].`$ (42)
This cancels the single ultraviolet poles in the sum of (28) and (36). Our final result for the transverse fragmentation function is obtained by adding the order-$`\alpha _s^2`$ corrections to $`d_T(z)`$ from (28), (36), and (42) and taking $`ϵ0`$:
$`d_T(z,\mu )={\displaystyle \frac{\pi \alpha _s(\mu )}{48m_Q^3}}`$ $`\{`$ $`\delta (1z)+{\displaystyle \frac{\alpha _s(\mu )}{\pi }}[A(\mu )\delta (1z)+(\mathrm{ln}{\displaystyle \frac{\mu }{2m_Q}}{\displaystyle \frac{1}{2}})P_{gg}(z)`$ (44)
$`+6(2z+z^2)\mathrm{ln}(1z){\displaystyle \frac{6}{z}}\left({\displaystyle \frac{\mathrm{ln}(1z)}{1z}}\right)_+]\},`$
where the coefficient $`A(\mu )`$ is
$`A(\mu )=\beta _0\left(\mathrm{ln}{\displaystyle \frac{\mu }{2m_Q}}+{\displaystyle \frac{13}{6}}\right)+{\displaystyle \frac{2}{3}}{\displaystyle \frac{\pi ^2}{2}}+8\mathrm{ln}2+6\mathrm{ln}^22.`$ (45)
The result (44) disagrees with the final result obtained by Ma . Our final result for the longitudinal fragmentation function is obtained by setting $`ϵ0`$ in (37):
$`d_L(z,\mu )={\displaystyle \frac{\alpha _s^2(\mu )}{8m_Q^3}}{\displaystyle \frac{1z}{z}}.`$ (46)
This agrees with the result of Beneke and Rothstein. The fragmentation probabilities obtained by integrating $`d_T(z)`$ and $`d_L(z)`$ diverge, because these functions behave like $`1/z`$ as $`z0`$. The higher moments of the fragmentation functions however are well-defined.
## VII Discussion
The color-octet $`{}_{}{}^{3}S_{1}^{}`$ term in the fragmentation function is important for calculating the production at large $`p_T`$ of spin-singlet S-wave states, like the $`J/\psi `$, and spin-triplet P-wave states, like the $`\chi _{cJ}`$. We can deduce the fragmentation functions for each of their spin states by using the approximate spin symmetry of NRQCD to simplify the expression (8). The functions $`d_T(z)`$ in (44) and $`d_L(z)`$ (46) give the fragmentation functions for the transverse and longitudinal spin states of the $`J/\psi `$, respectively:
$`D_{gJ/\psi (\pm 1)}(z)`$ $`=`$ $`d_T(z)𝒪_8^{J/\psi }(^3S_1),`$ (47)
$`D_{gJ/\psi (0)}(z)`$ $`=`$ $`d_L(z)𝒪_8^{J/\psi }(^3S_1).`$ (48)
The sum over spin states is
$`D_{gJ/\psi }(z)`$ $`=`$ $`[2d_T(z)+d_L(z)]𝒪_8^{J/\psi }(^3S_1).`$ (49)
The fragmentation functions for each of the spin states of the $`\chi _{cJ}`$ are
$`D_{g\chi _{c0}}(z)`$ $`=`$ $`[2d_T(z)+d_L(z)]𝒪_8^{\chi _{c0}}(^3S_1),`$ (50)
$`D_{g\chi _{c1}(0)}(z)`$ $`=`$ $`3d_T(z)𝒪_8^{\chi _{c0}}(^3S_1),`$ (51)
$`D_{g\chi _{c1}(\pm 1)}(z)`$ $`=`$ $`{\displaystyle \frac{3}{2}}[d_T(z)+d_L(z)]𝒪_8^{\chi _{c0}}(^3S_1),`$ (52)
$`D_{g\chi _{c2}(0)}(z)`$ $`=`$ $`[d_T(z)+2d_L(z)]𝒪_8^{\chi _{c0}}(^3S_1),`$ (53)
$`D_{g\chi _{c2}(\pm 1)}(z)`$ $`=`$ $`{\displaystyle \frac{3}{2}}[d_T(z)+d_L(z)]𝒪_8^{\chi _{c0}}(^3S_1),`$ (54)
$`D_{g\chi _{c2}(\pm 2)}(z)`$ $`=`$ $`3d_T(z)𝒪_8^{\chi _{c0}}(^3S_1).`$ (55)
The sums over spin states are
$`D_{g\chi _{cJ}}(z)`$ $`=`$ $`(2J+1)[2d_T(z)+d_L(z)]𝒪_8^{\chi _{c0}}(^3S_1).`$ (56)
In order to give accurate predictions for the production of quarkonium at large $`p_T`$, it is important to know the next-to-leading order correction to the color-octet $`{}_{}{}^{3}S_{1}^{}`$ term in the gluon fragmentation function. Unfortunately our calculation of the short-distance coefficient disagrees with the previous calculation by Ma. An independent calculation of this important function is therefore essential.
This work was supported in part by the U.S. Department of Energy Division of High Energy Physics under grant DE-FG02-91-ER40690, by the Alexander von Humboldt Foundation, and by the Korea Institute for Advanced Study. J.L. would like to thank the OSU theory group for its hospitality during his stay in Columbus.
## A Integral Table
In this appendix, we present the explicit values of the integrals encountered in evaluating the virtual-gluon corrections. These integrals have the form
$$I_{AB\mathrm{}}=\frac{d^{N+1}l}{(2\pi )^{N+1}}\frac{1}{AB\mathrm{}},$$
(A1)
where the denominator $`AB\mathrm{}`$ can be a product of 1, 2, 3, or 4 of the following factors:
$`A`$ $`=`$ $`l^2+iϵ,`$ (A2)
$`B`$ $`=`$ $`(lp)^2+iϵ=l^22lp+4m_Q^2+iϵ,`$ (A3)
$`C`$ $`=`$ $`(pl)n+iϵ,`$ (A4)
$`D`$ $`=`$ $`(lp/2)^2m_Q^2+iϵ=l^2lp+iϵ.`$ (A5)
The momentum $`p`$ is that of a $`Q\overline{Q}`$ pair with zero relative momentum ($`p^2=4m_Q^2`$) and $`n`$ is light-like ($`n^2=0`$). The integrals $`I_A`$ and $`I_B`$ vanish in dimensional regularization. By symmetry under $`plp`$, we have $`I_{AD}=I_{BD}`$. Some of the integrals can be reduced to ones with fewer denominators by using the identity $`A+B2D=4m_Q^2`$:
$`4m_Q^2I_{ABD}`$ $`=`$ $`2(I_{AD}I_{AB}),`$ (A6)
$`4m_Q^2I_{ABCD}`$ $`=`$ $`I_{ACD}+I_{BCD}2I_{ABC}.`$ (A7)
The independent integrals that need to be evaluated are therefore
$`I_{AB}`$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^2}}\left({\displaystyle \frac{\pi e^{i\pi }}{m_Q^2}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }^2(1ϵ)}{ϵ_{UV}\mathrm{\Gamma }(22ϵ)}},`$ (A8)
$`I_{AD}`$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^2}}\left({\displaystyle \frac{4\pi }{m_Q^2}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}(12ϵ)}},`$ (A9)
$`I_{ABC}`$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^2pn}}\left({\displaystyle \frac{\pi e^{i\pi }}{m_Q^2}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }^2(1ϵ)}{ϵ_{UV}ϵ_{IR}\mathrm{\Gamma }(12ϵ)}},`$ (A10)
$`I_{ACD}`$ $`=`$ $`{\displaystyle \frac{+i}{(4\pi )^2pn}}\left({\displaystyle \frac{4\pi }{m_Q^2}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}}}\left[2\mathrm{ln}2+ϵ\left({\displaystyle \frac{\pi ^2}{3}}6\mathrm{ln}^22\right)+O(ϵ^2)\right],`$ (A11)
$`I_{BCD}`$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^2pn}}\left({\displaystyle \frac{4\pi }{m_Q^2}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)}{ϵ_{UV}ϵ_{IR}}}.`$ (A12)
The subscripts on the poles in $`ϵ`$ indicate whether the divergences are of ultraviolet or infrared origin.
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# Introduction
## Introduction
Poisson brackets play a key role in Hamiltonian mechanics. According to the theorem of Darboux (Olver, 1986), degenerated Poisson manifolds are stratified to symplectic manifolds (leaves). A reduction of a Hamiltonian system on common level of all Casimir functions leads to usual Hamiltonian mechanics. The reduction is especially algebraic problem and can be made without referring to a concrete physical problem. In the present paper an algebraic method of finding of symplectic coordinates of Lie—Poisson brackets in some important for physical application cases is demonstrated.
## 1 A symplectic basis of the Lie—Poisson brackets with an algebra $`e(3)`$
The equations of motion of a rigid body with a fixed point in a homogeneous gravitation field are traditionally written in variables $`𝐌,𝜸,`$ where $`𝐌`$ is a vector of an angular momentum of the body and $`𝜸`$ is unit vector parallel to the vector of gravitational field strength. The coordinate system is toughly connected with the body.
The equations of motion are Hamiltonian with a degenerated Poisson brackets (Olver, 1986)
$$\{M_i,M_j\}=ϵ_{ijk}M_k,\{\gamma _i,\gamma _j\}=0,\{M_i,\gamma _j\}=ϵ_{ijk}\gamma _k$$
(1.1)
The algebra (1.1) is a semi–direct sum of an algebra of rotation $`so(3)`$ and an algebra of translation R$`{}_{}{}^{3}:`$ $`e(3)=so(3)_s`$R$`^3.`$ Under reduction of the algebra (1.1) to Casimir functions $`H=(𝐌,𝜸)`$ and $`(𝜸,𝜸)=1,`$ the Poisson brackets (1.1) become non–degenerated and the equations of motion can be written in canonical form. In the capacity of canonical coordinates, most convenient for qualitative analysis, the canonical variables of Andoyer—Deprit $`(l,L;g,G;h,H)`$ are used.
For describing of mechanical sense of these variables we denote by $`OXYZ`$ a stationary trihedron with the origin in the point of fixation of the body, $`Oxyz`$ is connected with the body the frame of coordinates, $`\mathrm{\Sigma }`$ is a plane passing through the point of fixation perpendicular to the vector of angular momentum $`𝐌.`$ At last, in the admitted notations
$`L`$ is a projection of the angular momentum to the moving axis $`Oz`$;
$`G`$ is a value of the angular momentum;
$`H`$ is a projection of the angular momentum to the fixed axis $`OZ`$;
$`l`$ is an angle between the axis $`Ox`$ and a line of intersection of $`\mathrm{\Sigma }`$ with $`Oxy`$;
$`g`$ is an angle between the line of intersection of $`\mathrm{\Sigma }`$ with planes $`Oxy`$ and $`OXY`$;
$`h`$ is an angle between the axis $`OX`$ and the line of intersection of $`\mathrm{\Sigma }`$ with the plane $`OXY.`$
The essence of the method consists in extracting from (1.1) of the subalgebra $`so(3)`$ and getting of canonical variables there. The canonical variables $`l,`$ $`L`$
$$M_1=\sqrt{G^2L^2}\mathrm{sin}l,M_2=\sqrt{G^2L^2}\mathrm{cos}l,M_3=L,$$
(1.2)
are cylindrical coordinates on two–dimensional concentric spheres (orbits of the $`so(3))`$ (Olver, 1986). The angular momentum $`G=\sqrt{M_1^2+M_2^2+M_3^2}`$ is the Casimir function of the considered algebra. The Poisson brackets between $`l,L,G,\gamma _1,\gamma _2,\gamma _3`$ are following:
$$\{l,L\}=1,\{G,L\}=\{G,l\}=0,\{\gamma _i,\gamma _j\}=0,$$
(1.3)
$$\{L,\gamma _1\}=\gamma _2,\{L,\gamma _2\}=\gamma _1,\{L,\gamma _3\}=0,$$
(1.4)
$$\{l,\gamma _1\}=\frac{\mathrm{sin}l\gamma _3}{\sqrt{G^2L^2}},\{l,\gamma _2\}=\frac{\mathrm{cos}l\gamma _3}{\sqrt{G^2L^2}},$$
$$\{l,\gamma _3\}=\frac{HL\gamma _3}{G^2L^2},$$
(1.5)
$$\{G,\gamma _1\}=\frac{1}{G}(\sqrt{G^2L^2}\mathrm{cos}l\gamma _3L\gamma _2),$$
$$\{G,\gamma _2\}=\frac{1}{G}(L\gamma _1\sqrt{G^2L^2}\mathrm{sin}l\gamma _3),$$
(1.6)
$$\{G,\gamma _3\}=\frac{1}{G}\sqrt{G^2L^2}(\mathrm{sin}l\gamma _2\mathrm{cos}l\gamma _1),$$
where $`H`$ is a projection of the angular momentum to the fixed axis. In addition, $`H=(𝐌,𝜸)`$ is the Casimir function of the algebra $`e(3).`$
Now we solve step by step the systems of partial differential equations (1.2)–(1.6), assuming $`\gamma _i`$ as functions of $`(l,L;g,G;H).`$ As a result, we get the required relations with the Andoyer—Deprit variables:
$$\gamma _1=\left(\frac{H}{G}\sqrt{1\left(\frac{L}{G}\right)^2}+\frac{L}{G}\sqrt{1\left(\frac{H}{G}\right)^2}\mathrm{cos}g\right)\mathrm{sin}l+\sqrt{1\left(\frac{H}{G}\right)^2}\mathrm{sin}g\mathrm{cos}l,$$
(1.7)
$$\gamma _2=\left(\frac{H}{G}\sqrt{1\left(\frac{L}{G}\right)^2}+\frac{L}{G}\sqrt{1\left(\frac{H}{G}\right)^2}\mathrm{cos}g\right)\mathrm{cos}l\sqrt{1\left(\frac{H}{G}\right)^2}\mathrm{sin}g\mathrm{sin}l,$$
$$\gamma _3=\left(\frac{H}{G}\right)\left(\frac{L}{G}\right)\sqrt{1\left(\frac{L}{G}\right)^2}\sqrt{1\left(\frac{H}{G}\right)^2}\mathrm{cos}g.$$
The canonical variables of the four–dimensional symplectic leaf of the Poisson brackets of the algebra $`e(3)`$ are $`(l,L;g,G).`$ They are enumerated by the variable $`H.`$
## 2 A symplectic basis of the Lie—Poisson brackets with an algebra $`l(7)`$
The very convenient variables for description of the rigid body dynamics instead of the angles $`(𝜶,𝜷,𝜸)`$ that form the orthogonal matrix of rotations are quaternions (Koshlyakov, 1985). The real variables $`\lambda _0,\lambda _1,\lambda _2,\lambda _3`$ normalized by the condition
$$\lambda _0^2+\lambda _1^2+\lambda _2^2+\lambda _3^2=1,$$
(2.1)
and connected with the vector $`𝜸`$ by the formulae (Koshlyakov, 1985):
$$\gamma _1=2(\lambda _1\lambda _3\lambda _0\lambda _2),$$
$$\gamma _2=2(\lambda _0\lambda _1+\lambda _2\lambda _3),$$
(2.2)
$$\gamma _3=\lambda _0^2\lambda _1^2\lambda _2^2+\lambda _3^2,$$
are named the Hamilton—Rodrigues parameters.
A point M(r) under some rotation around a vector $`𝐞(\alpha ^{},\beta ^{},\gamma ^{})`$ is displaced to $`M^{}(𝐫^{})`$ on an angle $`\chi .`$ The vector $`𝜽=2tan(\chi /2)𝐞`$ is a vector of finite rotation. Instead of projections $`\theta _1,\theta _2,\theta _3`$ one can introduce variables $`\lambda _k,`$ by $`\lambda _k=1/2\lambda _0\theta _k,`$ $`k=1,2,3`$ are subjected to the condition (2.1). So we get
$$\lambda _0=\mathrm{cos}\chi /2,\lambda _1=\mathrm{cos}\alpha ^{}\mathrm{sin}\chi /2,\lambda _2=\mathrm{cos}\beta ^{}\mathrm{sin}\chi /2,\lambda _3=\mathrm{cos}\gamma ^{}\mathrm{sin}\chi /2.$$
(2.3)
Now, it is not difficult to obtain the algebra of Poisson brackets of variables $`𝐌,\lambda _0,\lambda _1,\lambda _2,\lambda _3.`$ The Lagrangian of a free rigid body with a diagonal tensor of inertion $`𝐈=diag(A,B,C)`$ is
$$L=\frac{1}{2}(𝐈𝝎,𝝎),$$
(2.4)
where $`𝝎`$ is a vector of angular velocity. In capacity of quasimomenta we take components of the angular momenta
$$𝐌=L/𝝎.$$
(2.5)
The Hamiltonian $``$ is defined by the Legandre transformation
$$=(𝝎,\frac{L}{𝝎})L_{𝝎𝐌}.$$
(2.6)
Using (2.5), (2.6) and kinematic Euler relations (Koshlyakov, 1985):
$$\omega _1=\dot{\psi }\mathrm{sin}\theta \mathrm{sin}\phi +\dot{\theta }\mathrm{cos}\phi ,$$
$$\omega _2=\dot{\psi }\mathrm{sin}\theta \mathrm{cos}\phi \dot{\theta }\mathrm{sin}\phi ,$$
(2.7)
$$\omega _3=\dot{\psi }\mathrm{cos}\theta +\dot{\phi },$$
where $`\theta ,\phi ,\psi `$ are the Euler angles, we obtain the formulae:
$$M_1=\frac{\mathrm{sin}\phi }{\mathrm{sin}\theta }(p_\psi p_\phi \mathrm{cos}\theta )+p_\theta \mathrm{cos}\phi ,$$
$$M_2=\frac{\mathrm{cos}\phi }{\mathrm{sin}\theta }(p_\psi p_\phi \mathrm{cos}\theta )p_\theta \mathrm{sin}\phi ,$$
(2.8)
$$M_3=p_\phi .$$
Taking under consideration canonical brackets between generalized coordinates —Euler angles and corresponding to them canonical momenta $`p_\theta ,p_\phi ,p_\psi ,`$ it is possible to obtain the algebra
$$\{M_i,M_j\}=ϵ_{ijk}M_k,\{M_i,\lambda _0\}=\frac{1}{2}\lambda _i,$$
(2.9)
$$\{M_i,\lambda _j\}=\frac{1}{2}(ϵ_{ijk}\lambda _k+\delta _{ij}\lambda _0),\{\lambda _\mu ,\lambda _\nu \}=0.$$
Let us obtain by the developed method canonical coordinates of the six–dimensional symplectic leaf of the Poisson brackets with a seven–dimensional Lie algebra $`l(7)=so(3)_s`$R$`^4.`$ Also, as in the first section, we solve step by step the following systems of partial differential equations corresponding to the algebras.
$$1.\{L,\lambda _0\}=\frac{1}{2}\lambda _3,\{L,\lambda _1\}=\frac{1}{2}\lambda _2,$$
(2.10)
$$\{L,\lambda _2\}=\frac{1}{2}\lambda _1,\{L,\lambda _3\}=\frac{1}{2}\lambda _0;$$
$$2.\{H,\lambda _0\}=\frac{1}{2}\lambda _3,\{H,\lambda _1\}=\frac{1}{2}\lambda _2,$$
(2.11)
$$\{H,\lambda _2\}=\frac{1}{2}\lambda _1,\{H,\lambda _3\}=\frac{1}{2}\lambda _0;$$
$$3.\{l,\lambda _0\}=\frac{\lambda _1\mathrm{cos}l\lambda _2\mathrm{sin}l}{2\sqrt{G^2L^2}},\{l,\lambda _1\}=\frac{\lambda _3\mathrm{cos}l+\lambda _0\mathrm{sin}l}{2\sqrt{G^2L^2}},$$
(2.12)
$$\{l,\lambda _2\}=\frac{\lambda _0\mathrm{sin}l\lambda _3\mathrm{cos}l}{2\sqrt{G^2L^2}},\{l,\lambda _3\}=\frac{\lambda _1\mathrm{sin}l+\lambda _2\mathrm{cos}l}{2\sqrt{G^2L^2}};$$
$$4.\{G,\lambda _0\}=\frac{\sqrt{G^2L^2}}{2G}(\lambda _1\mathrm{sin}l+\lambda _2\mathrm{cos}l)+\frac{L}{2G}\lambda _3,$$
$$\{G,\lambda _1\}=\frac{\sqrt{G^2L^2}}{2G}(\lambda _0\mathrm{sin}l+\lambda _3\mathrm{cos}l)\frac{L}{2G}\lambda _2,$$
(2.13)
$$\{G,\lambda _2\}=\frac{\sqrt{G^2L^2}}{2G}(\lambda _0\mathrm{cos}l+\lambda _3\mathrm{sin}l)+\frac{L}{2G}\lambda _1,$$
$$\{G,\lambda _3\}=\frac{\sqrt{G^2L^2}}{2G}(\lambda _2\mathrm{sin}l\lambda _1\mathrm{cos}l)\frac{L}{2G}\lambda _0.$$
Using the condition of the norm (2.1) and relations (2.2) we get the following solutions:
$$\lambda _0=\frac{1}{\sqrt{2}}(\mathrm{sin}(g/2)\mathrm{sin}(y_+)\mathrm{cos}(x_{})+\mathrm{sin}(g/2)\mathrm{cos}(y_+)\mathrm{cos}(x_{})+$$
$$\mathrm{cos}(g/2)\mathrm{sin}(y_+)\mathrm{sin}(x_+)\mathrm{cos}(g/2)\mathrm{cos}(y_+)\mathrm{sin}(x_+)),$$
$$\lambda _1=\frac{1}{\sqrt{2}}(\mathrm{sin}(g/2)\mathrm{cos}(y_{})\mathrm{sin}(x_{})\mathrm{sin}(g/2)\mathrm{sin}(y_{})\mathrm{sin}(x_{})$$
$$\mathrm{cos}(g/2)\mathrm{sin}(y_{})\mathrm{cos}(x_+)\mathrm{cos}(g/2)\mathrm{cos}(y_{})\mathrm{cos}(x_+)),$$
$$\lambda _2=\frac{1}{\sqrt{2}}(\mathrm{sin}(g/2)\mathrm{sin}(y_{})\mathrm{sin}(x_{})\mathrm{sin}(g/2)\mathrm{cos}(y_{})\mathrm{sin}(x_{})+$$
(2.14)
$$\mathrm{cos}(g/2)\mathrm{sin}(y_{})\mathrm{cos}(x_+)\mathrm{cos}(g/2)\mathrm{cos}(y_{})\mathrm{cos}(x_+)),$$
$$\lambda _3=\frac{1}{\sqrt{2}}(\mathrm{sin}(g/2)\mathrm{sin}(y_+)\mathrm{cos}(x_{})\mathrm{sin}(g/2)\mathrm{cos}(y_+)\mathrm{cos}(x_{})$$
$$\mathrm{cos}(g/2)\mathrm{sin}(y_+)\mathrm{sin}(x_+)\mathrm{cos}(g/2)\mathrm{cos}(y_+)\mathrm{sin}(x_+)).$$
In the formulae (2.14) there are introduced the angles $`\zeta ,\tau :`$
$$\zeta =\mathrm{arcsin}H/G,\tau =\mathrm{arcsin}L/G$$
and the combinations:
$$x_+\frac{1}{2}(\zeta +\tau ),x_{}\frac{1}{2}(\zeta \tau ),y_+\frac{1}{2}(l+h),y_{}\frac{1}{2}(lh).$$
The obtained formulae (2.14) can be used for applications of methods of the theory of perturbation to the problem of the rigid body rotations in superposition of some potential strength fields.
## References
Arhangelski Yu.A. (1977). Analytical dynamics of a rigid body, Nauka, Moscow \[in Russian\].
Koshlyakov V.N. (1985). Problems of dynamics of rigid body and applied theory of gyroscopes: Analytical methods, Nauka, Moscow \[in Russian\].
Olver P. (1986). Applications of Lie groups to differential equations, Springer–Verlag, New York.
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# Electromagnetic form factors of nucleons and 𝑝→Δ 11footnote 1Published in Physica Sinica, Vol.24, No.2, 124, March 1975 (Chinese, translated by Jun Gao)
## 1 Introduction
The relativistic quark model is successful in studying the electromagnetic and weak interactions of ground-state baryons and mesons. However, some results are inconsistent with experimental data, for instance, theoretical ratio $`\mu _pG_E^p(q^2)/G_M^p(q^2)`$ drops faster with $`q^2`$ than experimental values. On the other hand, recently there has been some experimental results about the electromagnetic properties of ground-state baryons, for example, the measurement of the magnetic transition form factor $`G_{M1+}(q^2)`$ for $`p\mathrm{\Delta }^+(1236)`$. These results need theoretical explanation. In Ref., a new set of baryon wave functions have been constructed by requiring $`SU(6)`$ symmetry in the frame of center of mess. For example, in the wave function of $`\frac{1}{2}^+`$ baryon there are additional terms
$$\{(1+\frac{i}{m}\widehat{p})\gamma _5C\}_{\alpha \beta }u_\lambda (p)_\gamma +2C_{\alpha \beta }\{\gamma _5u_\lambda (p)\}_\gamma .$$
(1)
p is the momentum of the baryon, m is the mass, and $`C=i\gamma _2\gamma _4`$ which is the charge-conjugation operator. The new wavefunctions are still s-wave and satisfy $`SU(6)`$ symmetry in the frame of center of mass. It is interesting to point out that the original wave function
$$\{(1\frac{i}{m}\widehat{p})\gamma _5C\}_{\alpha \beta }u_\lambda (p)_\gamma $$
is constructed by the spinors of quarks with zero momentum and the new terms(1) are constructed by the spinors of antiquarks with zero momentum. Both satisfy $`SU(6)`$ in the frame of center of mass.
In this paper, the method proposed in Ref. and the wave functions constructed in Ref. are used to study the electromagnetic properties of nucleons and $`p\mathrm{\Delta }^+(1236)`$.
## 2 Matrix element of electric currents
The effective Hamiltonian of electromagnetic interaction in quark model is
$$H_i(x)=ie\overline{\psi }(x)Q\{\widehat{A}(x)\frac{i\kappa }{4m_p}\sigma _{\mu \nu }F_{\mu \nu }(x)\}\psi (x).$$
(2)
The wave function of $`\frac{1}{2}^+`$ baryon is
$`B_{\alpha \beta \gamma ,ijk}^{\frac{1}{2}\lambda }(x_1,x_2,x_3)_l^l^{^{}}`$ $`=`$ $`{\displaystyle \frac{1}{6\sqrt{2}}}\sqrt{{\displaystyle \frac{m}{E}}}\epsilon _{i^{^{}}j^{^{}}k^{^{}}}\{\mathrm{\Gamma }_{\alpha \beta ,\gamma }(p)_\lambda (\epsilon _{ijl^{^{}}}\delta _{kl}+\epsilon _{ikl^{^{}}}\delta _{jl})`$
$`+\mathrm{\Gamma }_{\beta \gamma ,\alpha }(p)_\lambda (\epsilon _{jkl^{^{}}}\delta _{il}+\epsilon _{ikl^{^{}}}\delta _{jl})\},`$
$`\overline{B}_{\alpha \beta \gamma ,ijk}^{\frac{1}{2}\lambda }(x_1,x_2,x_3)_l^l^{^{}}`$ $`=`$ $`{\displaystyle \frac{1}{6\sqrt{2}}}\sqrt{{\displaystyle \frac{m}{E}}}\epsilon _{i^{^{}}j^{^{}}k^{^{}}}\{\overline{\mathrm{\Gamma }}_{\alpha \beta ,\gamma }(p)_\lambda (\epsilon _{ijl^{^{}}}\delta _{kl}+\epsilon _{ikl^{^{}}}\delta _{jl})`$ (3)
$`+\overline{\mathrm{\Gamma }}_{\beta \gamma ,\alpha }(p)_\lambda (\epsilon _{jkl^{^{}}}\delta _{il}+\epsilon _{ikl^{^{}}}\delta _{jl})\}.`$
$`\mathrm{\Gamma }_{\alpha \beta ,\gamma }(p)_\lambda `$ $`=`$ $`\{[f_1(x_1,x_2,x_3){\displaystyle \frac{i}{m}}\widehat{p}f_2(x_1,x_2,x_3)]\gamma _5C\}_{\alpha \beta }u_\lambda (p)_\gamma `$
$`+\{f_1(x_1,x_2,x_3)f_2(x_1,x_2,x_3)\}C_{\alpha \beta }\{\gamma _5u_\lambda (p)\}_\gamma ,`$
$`\overline{\mathrm{\Gamma }}_{\alpha \beta ,\gamma }(p)_\lambda `$ $`=`$ $`\{C[f_1(x_1,x_2,x_3)+{\displaystyle \frac{i}{m}}\widehat{p}f_2(x_1,x_2,x_3)]\gamma _5\}_{\alpha \beta }\overline{u}_\lambda (p)_\gamma `$ (4)
$`+\{f_1(x_1,x_2,x_3)f_2(x_1,x_2,x_3)\}C_{\alpha \beta }\{\overline{u}_\lambda (p)\gamma _5\}_\gamma .`$
The wave function of $`\frac{3}{2}^+`$ baryon is
$`B_{\alpha \beta \gamma ,ijk}^{\frac{3}{2}\lambda ,lmn}(x_1,x_2,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\sqrt{{\displaystyle \frac{m}{E}}}\epsilon _{i^{^{}}j^{^{}}k^{^{}}}d_{ijk}^{lmn}\mathrm{\Gamma }_{\alpha \beta \gamma }(p)_\lambda ,`$
$`\overline{B}_{\alpha \beta \gamma ,ijk}^{\frac{3}{2}\lambda ,lmn}(x_1,x_2,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}\sqrt{{\displaystyle \frac{m}{E}}}\epsilon _{i^{^{}}j^{^{}}k^{^{}}}d_{ijk}^{lmn}\overline{\mathrm{\Gamma }}_{\alpha \beta \gamma }(p)_\lambda .`$ (5)
$`\mathrm{\Gamma }_{\alpha \beta \gamma }(p)_\lambda `$ $`=`$ $`\{[f_2(x_1,x_2,x_3){\displaystyle \frac{i}{m}}\widehat{p}f_1(x_1,x_2,x_3)]\gamma _\mu C\}_{\alpha \beta }\psi _\mu ^\lambda (p)_\gamma `$
$`+{\displaystyle \frac{i}{m}}\{f_1(x_1,x_2,x_3)f_2(x_1,x_2,x_3)\}\{\gamma _\mu \widehat{p}\gamma _5C\}_{\alpha \beta }\{\gamma _5\psi _\mu ^\lambda (p)\}_\gamma ,`$
$`\overline{\mathrm{\Gamma }}_{\alpha \beta \gamma }(p)_\lambda `$ $`=`$ $`\{C[f_2(x_1,x_2,x_3)+{\displaystyle \frac{i}{m}}\widehat{p}f_1(x_1,x_2,x_3)]\gamma _\mu \}_{\alpha \beta }\overline{\psi }_\mu ^\lambda (p)_\gamma `$ (6)
$`+{\displaystyle \frac{i}{m}}\{f_1(x_1,x_2,x_3)f_2(x_1,x_2,x_3)\}`$
$`\times \{C\widehat{p}\gamma _\mu \gamma _5\}_{\alpha \beta }\{\overline{\psi }_\mu ^\lambda (p)\gamma _5\}_\gamma ,`$
where $`i^{^{}}j^{^{}}k^{^{}}`$ are color indices, $`f_{1,2}(x_1,x_2,x_3)`$ are two Lorentz-invariant spacial functions, p is the momentum of baryon. The wave functions(3-6) satisfy $`SU(6)`$ symmetry in the frame of center of mass. They are s-wave in the rest frame.
For the model with three degenerate states , the electric charge operator can be written as
$$Q_{k_1k_2}^{k_1^{^{}}k_2^{^{}}}=\delta _{k_11}\delta _{k_21}\delta _{k_1^{^{}}k_2^{^{}}}\delta _{k_1k_2}\delta _{k_1^{^{}}3}\delta _{k_2^{^{}}3}$$
(7)
and the following relationship is obtained
$$\frac{1}{6}\epsilon _{k_1^{^{}}j^{^{}}i^{^{}}}\epsilon _{k_2^{^{}}j^{^{}}i^{^{}}}Q_{k_1k_2}^{k_1^{^{}}k_2^{^{}}}=Q_{k_1k_2}.$$
(8)
$`Q_{k_1k_2}`$ is the electric charge operator in fractional-charge scheme. Therefore, the matrix elements of electric currents should be the same for both schemes of interger charge and fractional charge of quarks.
By using Eqs.(2,3,5,8) and the method of Ref. , the matrix elements of electric currents of $`\frac{1}{2}^+`$ baryon are obtained
$`<`$ $`B_\lambda ^{^{}}^{\frac{1}{2}}(p_f)_{l_1}^{l_1^{^{}}}J_\mu (0)B_\lambda ^{\frac{1}{2}}(p_i)_{l_2}^{l_2^{^{}}}>={\displaystyle \frac{ie}{2}}M^2Q_{k_1k_1^{^{}}}(\gamma _\mu +{\displaystyle \frac{\kappa }{2m_p}}q_\nu \sigma _{\mu \nu })_{\gamma \gamma ^{^{}}}`$ (9)
$`\times d^4x_1d^4x_2\overline{B}_{\alpha \beta \gamma ,ijk_1}^{\frac{1}{2}\lambda ^{^{}}}(x_1,x_2,0)_{l_1}^{l_1^{^{}}}B_{\gamma ^{^{}}\beta \alpha ,k_1^{^{}}ji}^{\frac{1}{2}\lambda }(0,x_2,x_1)_{l_2}^{l_2^{^{}}}`$
$`=`$ $`{\displaystyle \frac{ie}{24}}({\displaystyle \frac{mm^{^{}}}{EE^{^{}}}})^{\frac{1}{2}}\{A_1I_1A_2I_2\},`$
where M is the rest mass of the quark, m, m’ and E, E’ are the initial and final mass and energy of the baryon respectively,
$$A_1=S_p\overline{B}QB,A_2=S_p\overline{B}BQ.$$
(10)
B, $`\overline{B}`$ are the $`SU_3`$ matrices of the initial and final baryon,
$`I_1`$ $`=`$ $`20\{D_2(q^2)(1{\displaystyle \frac{m_+}{5m}})+D_2^{^{}}(q^2)(1{\displaystyle \frac{m_+}{5m^{^{}}}})`$
$`+{\displaystyle \frac{1}{2mm^{^{}}}}(m_{}^2+q^2+{\displaystyle \frac{\kappa m_+}{5m_p}}q^2)D_3(q^2)\}\overline{u}_\lambda ^{^{}}(p_f)\gamma _\mu u_\lambda (p_i)`$
$`20\{2D_1(q^2)(1{\displaystyle \frac{2m_p}{5\kappa m}})D_2(q^2)(1{\displaystyle \frac{2m_p}{5\kappa m^{^{}}}})D_2^{^{}}(q^2)`$
$`+{\displaystyle \frac{1}{2mm^{^{}}}}(m_+^2+{\displaystyle \frac{3}{5}}q^2)D_3(q^2)\}{\displaystyle \frac{\kappa }{2m_p}}\overline{u}_\lambda ^{^{}}(p_f)q_\nu \sigma _{\mu \nu }u_\lambda (p_i)`$
$`4i\{{\displaystyle \frac{1}{m}}D_2(q^2){\displaystyle \frac{1}{m^{^{}}}}D_2^{^{}}(q^2)+{\displaystyle \frac{\kappa }{2mm^{^{}}m_p}}(m^{}_{}{}^{}2m^2)D_3(q^2)\}`$
$`\times q_\mu \overline{u}_\lambda ^{^{}}(p_f)u_\lambda (p_i),`$
$`I_2`$ $`=`$ $`4\{(12{\displaystyle \frac{m_+}{m}})D_2(q^2)+(12{\displaystyle \frac{m_+}{m^{^{}}}})D_2^{^{}}(q^2)`$ (11)
$`+{\displaystyle \frac{1}{2mm^{^{}}}}(m_{}^2+q^2+2{\displaystyle \frac{\kappa m_+}{m_p}}q^2)D_3(q^2)\}\overline{u}_\lambda ^{^{}}(p_f)\gamma _\mu u_\lambda (p_i)`$
$`4\{2D_1(q^2)(1{\displaystyle \frac{4m_p}{\kappa m}})D_2(q^2)(1{\displaystyle \frac{4m_p}{\kappa m^{^{}}}})D_2^{^{}}(q^2)`$
$`+{\displaystyle \frac{1}{2mm^{^{}}}}(m_+^23q^2)D_3(q^2)\}{\displaystyle \frac{\kappa }{2m_p}}\overline{u}_\lambda ^{^{}}(p_f)q_\nu \sigma _{\nu \mu }u_\lambda (p_i)`$
$`+8i\{{\displaystyle \frac{1}{m}}D_2(q^2){\displaystyle \frac{1}{m^{^{}}}}D_2^{^{}}(q^2)+{\displaystyle \frac{\kappa (m^{}_{}{}^{}2m^2)}{2mm^{^{}}m_p}}D_3(q^2)\}`$
$`\times q_\mu \overline{u}_\lambda ^{^{}}(p_f)u_\lambda (p_i),`$
where
$$q_\mu =p_{i\mu }p_{f\mu },m_+=m+m^{^{}},m_{}=m^{^{}}m,$$
(12)
$`D_1(q^2)`$ $`=`$ $`M^2f_1^{^{}}(x_1,x_2,0)f_1(0,x_2,x_1)d^4x_1d^4x_2,`$
$`D_2(q^2)`$ $`=`$ $`M^2f_1^{^{}}(x_1,x_2,0)f_2(0,x_2,x_1)d^4x_1d^4x_2,`$
$`D_2^{^{}}(q^2)`$ $`=`$ $`M^2f_2^{^{}}(x_1,x_2,0)f_1(0,x_2,x_1)d^4x_1d^4x_2,`$
$`D_3(q^2)`$ $`=`$ $`M^2f_2^{^{}}(x_1,x_2,0)f_2(0,x_2,x_1)d^4x_1d^4x_2.`$ (13)
$`m,m^{^{}}`$ are the rest mass of the initial and final baryon. $`f_j^{^{}}(x_1,x_2,0)`$, $`f_j(0,x_2,x_1)`$ are the spacial part of the initial and final wave function respectively. Eq. (13) shows that when $`p_{f\mu }p_{i\mu }`$ is taken, we have
$$D_2(q^2)D_2^{^{}}(q^2),$$
(14)
therefore when m = m’
$$D_2(q^2)=D_2^{^{}}(q^2).$$
(15)
Similarly, the current matrix elements of $`\frac{1}{2}^+`$ baryon$`>\frac{2}{3}^+`$ baryon are obtained
$`<B_\lambda ^{^{}}^{\frac{3}{2}}(p_f)^{lmn}J_\mu (0)B_\lambda ^{\frac{1}{2}}(p_i)_{l_1}^{l_1^{^{}}}>={\displaystyle \frac{ie}{2}}M^2Q_{kk^{^{}}}(\gamma _\mu +{\displaystyle \frac{\kappa }{2m_p}}q_\nu \sigma _{\mu \nu })_{\gamma \gamma ^{^{}}}`$
$`\times d^4x_1d^4x_2\overline{B}_{\alpha \beta \gamma ,ijk}^{\frac{3}{2}\lambda ^{^{}},lmn}(x_1,x_2,0)B_{\gamma ^{^{}}\beta \alpha ,k^{^{}}ji}^{\frac{1}{2}\lambda }(0,x_2,x_1)_{l_1}^{l_1^{^{}}}`$
$`=`$ $`{\displaystyle \frac{ie}{4}}({\displaystyle \frac{mm^{^{}}}{EE^{^{}}}})d_{l_1jk}^{lmn}\epsilon _{jk^{^{}}l_1^{^{}}}Q_{kk^{^{}}}\{2D_2(q^2)+\kappa [{\displaystyle \frac{m_+}{m_p}}D_3(q^2)+2{\displaystyle \frac{m}{m_p}}D_1(q^2)`$ (16)
$`{\displaystyle \frac{m}{m_p}}D_2(q^2){\displaystyle \frac{m}{m_p}}D_2^{^{}}(q^2)]\}{\displaystyle \frac{1}{mm^{^{}}}}p_\rho q_\sigma \epsilon _{\rho \sigma \nu \mu }\overline{\psi }_\nu ^\lambda ^{^{}}(p_f)u_\lambda (p_i)`$
$`+ie({\displaystyle \frac{mm^{^{}}}{EE^{^{}}}})^{\frac{1}{2}}d_{l_1jk}^{lmn}\epsilon _{jk^{^{}}l_1^{^{}}}Q_{kk^{^{}}}\{D_3(q^2)D_2(q^2)+{\displaystyle \frac{\kappa m}{2m_p}}[D_2(q^2)`$
$`+D_2^{^{}}(q^2)2D_1(q^2)]\}{\displaystyle \frac{1}{mm^{^{}}}}(p_{f\mu }q_\nu p_fq\delta _{\mu \nu })\overline{\psi }_\nu ^\lambda ^{^{}}(p_f)\gamma _5u_\lambda (p_i)`$
$`+ie({\displaystyle \frac{mm^{^{}}}{EE^{^{}}}})^{\frac{1}{2}}d_{l_1jk}^{lmn}\epsilon _{jk^{^{}}l_1^{^{}}}Q_{kk^{^{}}}\{D_2^{^{}}(q^2){\displaystyle \frac{m^{^{}}}{m}}D_2(q^2)+{\displaystyle \frac{m_{}}{m}}D_3(q^2)\}`$
$`\times \overline{\psi }_\mu ^\lambda ^{^{}}(p_f)\gamma _5u_\lambda (p_i).`$
m, m’ are the rest mass of $`\frac{1}{2}^+`$ baryon and $`\frac{2}{3}^+`$ baryon respectively,
$$p_\mu =p_{i\mu }+p_{f\mu }.$$
(17)
In Eq.(11), when m = m’ is taken, the terms in $`I_1`$ and $`I_2`$, which are proportional to $`q_\mu `$ vanish. Thus, when m = m’, the current matrix element of $`\frac{1}{2}^+`$ baryon automatically satisfies current conservation. In general, in order to satisfy current conservation, the following relationship must be satisfied
$$D_2^{^{}}(q^2)\frac{m^{^{}}}{m}D_2(q^2)+\frac{m_{}}{m}D_3(q^2)=0.$$
(18)
For $`\frac{1}{2}^+`$ baryons the only matrix element with $`m^{}m`$ is $`\mathrm{\Sigma }^0>\mathrm{\Lambda }`$. For this process, we have
$$A_1=A_2=\frac{1}{2\sqrt{3}}.$$
(19)
The condition(18) guarantees current conservation.
## 3 Relationship Between $`f_1(x_1,x_2,x_3)`$ and $`f_2(x_1,x_2,x_3)`$
In this section, we study the behavior of two invariant spacial functions $`f_1(x_1,x_2,x_3)`$ and $`f_2(x_1,x_2,x_3)`$ in the frame of center-of-mass. $`\mathrm{\Gamma }_{\alpha \beta ,\gamma }(p)_\lambda `$(4) can be written as
$`\mathrm{\Gamma }_{\alpha \beta ,\gamma }(x_1,x_2,x_3)_\lambda `$ $`=`$ $`g_1(x_1,x_2,x_3)\{(1+\gamma _4)\gamma _5C\}_{\alpha \beta }u_{\lambda ,\gamma }`$
$`+g_2(x_1,x_2,x_3)\{[(1\gamma _4)\gamma _5C]_{\alpha \beta }u_{\lambda ,\gamma }+2C_{\alpha \beta }(\gamma _5u_\lambda )_\gamma \},`$
$`g_1(x_1,x_2,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{f_1(x_1,x_2,x_3)+f_2(x_1,x_2,x_3)\},`$
$`g_2(x_1,x_2,x_3)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{f_1(x_1,x_2,x_3)f_2(x_1,x_2,x_3)\},`$ (20)
$`x_1,x_2,x_3`$ are the time-space coordinates of three quarks. In Ref., in order to guarantee $`SU_6`$ symmetry, it is assumed that strong interaction satisfies $`SU_6`$ symmetry when the speed of the quark is much less than the speed of light. One possibility is that strong interaction takes scalar form.The B-S equation of the baryon is written as
$`(i\widehat{p}_1+M)_{\alpha \alpha ^{^{}}}(i\widehat{p}_2+M)_{\beta \beta ^{^{}}}(i\widehat{p}_3+M)_{\gamma \gamma ^{^{}}}B_{\alpha ^{^{}}\beta ^{^{}}\gamma ^{^{}},ijk}^{\frac{1}{2}\lambda }(p_1,p_2,p_3)_l^m`$ (21)
$`=`$ $`i(i\widehat{p}_3+M)_{\gamma \gamma ^{^{}}}U(q)B_{\alpha \beta \gamma ^{^{}},ijk}^{\frac{1}{2}\lambda }(p_1q,p_2+q,p_3)_l^md^4q`$
$`i(i\widehat{p}_1+M)_{\alpha \alpha ^{^{}}}U(q)B_{\alpha ^{^{}}\beta \gamma ,ijk}^{\frac{1}{2}\lambda }(p_1,p_2q,p_3+q)_l^md^4q`$
$`i(i\widehat{p}_2+M)_{\beta \beta ^{^{}}}U(q)B_{\alpha \beta ^{^{}}\gamma ,ijk}^{\frac{1}{2}\lambda }(p_1+q,p_2,p_3q)_l^md^4q`$
$`V(q_1,q_2,q_3)\delta ^4(q_1+q_2+q_3)B_{\alpha \beta \gamma ,ijk}^{\frac{1}{2}\lambda }(p_1+q_1,p_2+q_2,p_3+q_3)`$
$`\times d^4q_1d^4q_2d^4q_3.`$
$`B_{\alpha \beta \gamma ,ijk}^{\frac{1}{2}\lambda }(p_1,p_2,p_3)_l^m`$ is the wave function for $`\frac{1}{2}^+`$ baryon in the frame of center-of-mass. We assume $`U(q)`$ and $`V(q_1,q_2,q_3)`$ are independent of the momentum of the baryon.
$$p_1+p_2+p_3=p,$$
(22)
where p is the momentum of $`\frac{1}{2}^+`$ baryon.
According to Ref., in order to satisfy $`SU_6`$ symmetry the terms at $`O(\frac{|𝐩_j|}{M}(j=1,2,3))`$ in the wave function are ignored. The same treatment is used in Eq.(21). Substituting the wave function of $`\frac{1}{2}^+`$ into Eq.(21), we obtain
$`(M\gamma _4p_{10})_{\alpha \alpha ^{^{}}}(M\gamma _4p_{20})_{\beta \beta ^{^{}}}(M\gamma _4p_{30})_{\gamma \gamma ^{^{}}}\mathrm{\Gamma }_{\alpha ^{^{}}\beta ^{^{}}\gamma ^{^{}}}(p_1,p_2,p_3)_\lambda `$ (23)
$`=`$ $`i(M\gamma _4p_{30})_{\gamma \gamma ^{^{}}}U(q)\mathrm{\Gamma }_{\alpha \beta ,\gamma ^{^{}}}(p_1q,p_2+q,p_3)_\lambda d^4q`$
$`i(M\gamma _4p_{10})_{\alpha \alpha ^{^{}}}U(q)\mathrm{\Gamma }_{\alpha ^{^{}}\beta ,\gamma }(p_1,p_2q,p_3+q)_\lambda d^4q`$
$`i(M\gamma _4p_{20})_{\beta \beta ^{^{}}}U(q)\mathrm{\Gamma }_{\alpha \beta ^{^{}},\gamma }(p_1+q,p_2,p_3q)_\lambda d^4q`$
$`V(q_1,q_2,q_3)\delta ^4(q_1+q_2+q_3)\mathrm{\Gamma }_{\alpha \beta ,\gamma }(p_1+q_1,p_2+q_2,p_3+q_3)`$
$`\times d^4q_1d^4q_2d^4q_3.`$
where $`\mathrm{\Gamma }_{\alpha \beta ,\gamma }(p_1,p_2,p_3)_\lambda `$ is the expression of $`\mathrm{\Gamma }_{\alpha \beta \gamma }(x_1,x_2,x_3)_\lambda `$(20) in the momentum representation. Calculations lead to
$`(Mp_{10})(Mp_{20})(Mp_{30})g_1(p_1,p_2,p_3)`$ (24)
$`=`$ $`iU(q)\{(Mp_{30})g_1(p_1q,p_2+q,p_3)`$
$`+(Mp_{10})g_1(p_1,p_2q,p_3+q)`$
$`+(Mp_{20})g_1(p_1+q,p_2,p_3q)\}d^4q`$
$`V(q_1,q_2,q_3)\delta ^4(q_1+q_2+q_3)`$
$`\times g_1(p_1+q_1,p_2+q_2,p_3+q_3)d^4q_1d^4q_2d^4q_3,`$
$`(M+p_{10})(M+p_{20})(Mp_{30})g_2(p_1,p_2,p_3)`$ (25)
$`=`$ $`iU(q)\{(Mp_{30})g_2(p_1q,p_2+q,p_3)`$
$`+(M+p_{10})g_2(p_1,p_2q,p_3+q)`$
$`+(M+p_{20})g_2(p_1+q,p_2,p_3q)\}d^4q`$
$`V(q_1,q_2,q_3)\delta ^4(q_1+q_2+q_3)`$
$`\times g_2(p_1+q_1,p_2+q_2,p_3+q_3)d^4q_1d^4q_2d^4q_3,`$
$`(M+p_{10})(Mp_{20})(M+p_{30})g_2(p_1,p_2,p_3)`$ (26)
$`=`$ $`iU(q)\{(M+p_{30})g_2(p_1q,p_2+q,p_3)`$
$`+(M+p_{10})g_2(p_1,p_2q,p_3+q)`$
$`+(Mp_{20})g_2(p_1+q,p_2,p_3q)\}d^4q`$
$`V(q_1,q_2,q_3)\delta ^4(q_1+q_2+q_3)`$
$`\times g_2(p_1+q_1,p_2+q_2,p_3+q_3)d^4q_1d^4q_2d^4q_3,`$
$`(Mp_{10})(M+p_{20})(M+p_{30})g_2(p_1,p_2,p_3)`$ (27)
$`=`$ $`iU(q)\{(M+p_{30})g_2(p_1q,p_2+q,p_3)`$
$`+(Mp_{10})g_2(p_1,p_2q,p_3+q)`$
$`+(M+p_{20})g_2(p_1+q,p_2,p_3q)\}d^4q`$
$`V(q_1,q_2,q_3)\delta ^4(q_1+q_2+q_3)`$
$`\times g_2(p_1+q_1,p_2+q_2,p_3+q_3)d^4q_1d^4q_2d^4q_3.`$
Since $`V(q_1,q_2,q_3)`$ are totally symmetric functions of $`q_1,q_2,q_3`$. $`g_1(p_1,p_2,p_3)`$ are totally symmetric functions of $`p_1,p_2,p_3`$, which is consistent with Ref.. From Eqs.(25-27), we see that $`g_2(p_1,p_2,p_3)`$ have following symmetries: (1) totally symmetric in $`p_1,p_2,p_3`$. (2) since $`U(q)`$ and $`V(q_1,q_2,q_3)`$ are independent of the momentum p , the equation is invariant under the transformations $`p_{20}p_{20}`$, $`p_{30}p_{30}`$; $`p_{10}p_{10}`$, $`p_{30}p_{30}`$; $`p_{10}p_{10}`$, $`p_{20}p_{20}`$. By using the second symmetry of $`g_2(p_1,p_2,p_3)`$, Eq.(3.6) becomes Eq.(3.5) under the transformation $`p_{10}p_{10}`$, $`p_{20}p_{20}`$, thus $`g_1(p_1,p_2,p_3)`$ and $`g_2(p_1,p_2,p_3)`$ satisfy the same equation. $`g_1(p_1,p_2,p_3)`$ is related to $`g_2(p_1,p_2,p_3)`$ by
$$g_1(p_1,p_2,p_3)=bg_2(p_1,p_2,p_3),$$
(28)
where b is a constant. Eq.(28) leads to
$$f_2(x_1,x_2,x_3)=af_1(x_1,x_2,x_3).$$
(29)
Thus, in the wave functions (3,5), there is only one independent spacial function.
Substituting Eq.(29) into Eq.(18), we obtain
$$a=\frac{1}{1\frac{m_0}{m}}ora=1.$$
(30)
$`m_0`$ is a parameter, m is the physical mass of the baryon. Generally $`a1`$, a takes the first expression of Eq.(30).
## 4 Electromagetice Properties of $`\frac{1}{2}^+`$ Baryons
The electromagnetic form factors of $`\frac{1}{2}^+`$ baryon are obtained from the current matrix elements Eq.(9,11)
$`G_E(q^2)`$ $`=`$ $`{\displaystyle \frac{2}{3}}(A_1+2A_2)(1+{\displaystyle \frac{q^2}{4m^2}})\{D_2(q^2){\displaystyle \frac{\kappa q^2}{4mm_p}}D_3(q^2)\}`$ (31)
$`+{\displaystyle \frac{1}{3}}(A_2+5A_1)\{D_2(q^2)+{\displaystyle \frac{q^2}{4m^2}}[D_3(q^2)+{\displaystyle \frac{\kappa m}{m_p}}D_2(q^2)`$
$`{\displaystyle \frac{\kappa m}{m_p}}D_1(q^2)\kappa {\displaystyle \frac{m}{m_p}}(1+{\displaystyle \frac{q^2}{4m^2}})D_3(q^2)]\}.`$
$`G_M(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{3}}(A_2+5A_1)\{D_2(q^2)+{\displaystyle \frac{q^2}{4m^2}}D_3(q^2)+\kappa {\displaystyle \frac{m}{m_p}}[D_1(q^2)`$ (32)
$`D_2(q^2)+(1+{\displaystyle \frac{q^2}{4m^2}})D_3(q^2)]\},`$
where m is the mass of the baryon. From Eq.(31) we obtain
$$D_2(0)=1.$$
(33)
The expression of the magnetic moment of $`\frac{1}{2}^+`$ baryon is obtained from Eq.(32)
$`\mu `$ $`=`$ $`{\displaystyle \frac{1}{3}}(A_2+5A_1)\{{\displaystyle \frac{m_p}{m}}+\kappa [D_1(0)+D_3(0)1]\}.`$ (34)
From Eqs.(40,41,33)
$`\mu `$ $`=`$ $`{\displaystyle \frac{1}{3}}(A_2+5A_1)\{{\displaystyle \frac{m_p}{m}}+\kappa ({\displaystyle \frac{1}{1\frac{m_0}{m}}}{\displaystyle \frac{m_0}{m}})\}`$ (35)
is obtained. The two parameters $`\kappa ,m_0`$ in Eq.(35) are determined to be
$$\kappa =0.481,m_0=0.778m_p$$
(36)
by input the magnetic moments of proton and $`\mathrm{\Sigma }`$ . The magnetic moments of other six baryons are determined to be
| | $`\mu _p`$ | $`\mu _n`$ | $`\mu _\mathrm{\Lambda }`$ | $`\mu _{\mathrm{\Sigma }^+}`$ | $`\mu _{\mathrm{\Sigma }^0}`$ | $`\mu _\mathrm{\Sigma }^{}`$ | $`\mu _{\mathrm{\Xi }^0}`$ | $`\mu _\mathrm{\Xi }^{}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| theory | 2.79 | -.1.86 | -0.64 | 1.74 | 0.58 | -0.57 | -0.97 | -0.51 |
| | (input) | | (input) | | | | | |
| exp | 2.79 | -1.91 | -0.67 | 2.59 | | | | -1.93 |
| | | | $`\pm 0.46`$ | $`\pm 0.46`$ | | | | $`\pm 0.75`$ |
The electromagnetic form factors of proton and neutron are found from Eq.(31,32)
$`G_E^p(q^2)`$ $`=`$ $`D_2(q^2)+{\displaystyle \frac{q^2}{4m_N^2}}\{D_3(q^2)+\kappa [D_2(q^2)D_1(q^2)(1+{\displaystyle \frac{q^2}{4m_N^2}})D_3(q^2)]\},`$
$`G_M^p(q^2)`$ $`=`$ $`D_2(q^2)+\kappa [D_1(q^2)+D_3(q^2)D_2(q^2)]+(1+\kappa ){\displaystyle \frac{q^2}{4m_N^2}}D_3(q^2),`$ (37)
$`G_E^n(q^2)`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{q^2}{4m_N^2}}\{D_3(q^2)D_2(q^2)+\kappa [D_2(q^2)D_1(q^2)]\},`$
$`G_M^n(q^2)`$ $`=`$ $`{\displaystyle \frac{2}{3}}G_M^p(q^2).`$ (38)
By using Eqs.(29,30,36)
$`G_E^p(q^2)`$ $`=`$ $`D_2(q^2)\{1+\tau (2.712.17\tau )\},`$
$`G_M^p(q^2)`$ $`=`$ $`\mu _pD_2(q^2)\{1+2.39\tau \},`$
$`G_E^n(q^2)`$ $`=`$ $`1.39\mu _n\tau D_2(q^2)`$ (39)
are obtained, where $`\tau =\frac{q^2}{4m_N^2}`$. It is seen from Eq.(39) that there is an invariant function D$`{}_{2}{}^{}(q^2)`$ in the three form factors, which can be determined from the experimental data of the magnetic form factor of proton
$$D_2(q^2)=\frac{1}{(1+\frac{q^2}{0.71})^2(1+2.39\tau )}.$$
(40)
The ratio of the electric and magnetic form factor of proton is obtained
$$\frac{\mu _pG_E^p(q^2)}{G_M^p(q^2)}=\frac{1+\tau (2.712.17\tau )}{1+2.39\tau }.$$
(41)
Comparisons with data are shown in Fig.1 and 2. The experimental data of Fig.1 is from Ref., and that for Fig.2 is from Ref..
The expression of the electric form factor of neutron is obtained
$$G_E^n(q^2)=1.39\tau G_M^n(q^2)(1+2.39\tau )^1.$$
(42)
The slope of $`G_E^n(q^2)`$ at $`q^2=0`$ is
$$\frac{dG_E^n(q^2)}{dq^2}_{q^2=0}=1.39\frac{\mu _n}{4m_N^2}=0.73[GeV]^2.$$
(43)
The experimental data are
$$0.579\pm 0.018^{[12]},0.512\pm 0.049^{[13]},0.495\pm 0.010^{[14]}.$$
(44)
Comparisons of Eq.(42) with the experimental data are shown in Fig.3 and Fig.4. The experimental data of Fig.3 comes from Ref. and that for Fig.4 comes from Ref..
At $`q^2=4m_N^2`$, there are
$`G_E^p(4m^2)`$ $`=`$ $`G_M^p(4m^2)=0.18,`$
$`G_E^n(4m^2)`$ $`=`$ $`G_M^p(4m^2)={\displaystyle \frac{2}{3}}G_E^p(4m^2).`$ (45)
The S-matrix element of $`\mathrm{\Sigma }^0\mathrm{\Lambda }+\gamma `$ is studied
$`<\gamma \mathrm{\Lambda }S\mathrm{\Sigma }^0>`$ $`=`$ $`ie(2\pi )^4\delta (p_ip_fq){\displaystyle \frac{e_\mu ^\lambda }{\sqrt{2\omega }}}({\displaystyle \frac{m_\mathrm{\Lambda }}{E_\mathrm{\Lambda }}})^{\frac{1}{2}}\mu _{\mathrm{\Sigma }_0\mathrm{\Lambda }}`$ (46)
$`\times {\displaystyle \frac{\kappa }{2m_p}}\overline{u}_\lambda ^{^{}}(p_f)q_\nu \sigma _{\nu \mu }u_\lambda (p_i),`$
$`\mu _{\mathrm{\Sigma }_0\mathrm{\Lambda }}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{3}}}D_3(0)\{{\displaystyle \frac{2}{a_\mathrm{\Lambda }a_{\mathrm{\Sigma }^0}}}{\displaystyle \frac{1}{a_\mathrm{\Lambda }}}(1{\displaystyle \frac{m_p}{\kappa m_\mathrm{\Sigma }}}){\displaystyle \frac{1}{a_{\mathrm{\Sigma }^0}}}(1{\displaystyle \frac{m_p}{\kappa m_\lambda }})+{\displaystyle \frac{m_+^2}{2m_\mathrm{\Lambda }m_\mathrm{\Sigma }}}\}.`$ (47)
The dependences of $`D_1(0)`$, $`D_2(0)`$, $`D_1^{^{}}(0)`$ and $`D_1^{^{}}(0)`$ on the mass of initial and final baryon need to be found. From Eq.(28) we have
$$\frac{D_2(0)}{D_2^{^{}}(0)}=\frac{a}{a^{^{}}}.$$
(48)
On the other hand, Eq.(14) shows when $`mm^{^{}}`$ is taken, we have
$`D_2(0)`$ $``$ $`D_2^{^{}}(0).`$ (49)
When m = m’, Eqs.(15,33) lead to
$`D_2(0)=D_2^{^{}}(0)=1.`$ (50)
The general expressions of $`D_2(0)`$, $`D_2^{^{}}(0)`$ which satisfy Eqs.(48-50) are found
$`D_2(0)`$ $`=`$ $`({\displaystyle \frac{a}{a^{^{}}}})^{\frac{1}{2}}f(m,m^{^{}}),`$
$`D_2^{^{}}(0)`$ $`=`$ $`({\displaystyle \frac{a^{^{}}}{a}})^{\frac{1}{2}}f(m,m^{^{}}).`$ (51)
$`f(m,m^{^{}})`$ is a symmetric function of m, m’ and
$$f(m,m)=1.$$
(52)
When $`mm^{^{}}`$, the deviation of $`f(m,m^{^{}})`$ from 1 is proportional to $`(mm^{^{}})^2`$. According to Ref., $`f(m,m^{^{}})`$ is the effect of Lorentz contraction. We obtain
$$f(m,m^{^{}})=\frac{4mm^{^{}}}{(m+m^{^{}})^2}.$$
(53)
For $`\mathrm{\Sigma }^0\mathrm{\Lambda }+\gamma `$, the deviation of $`f(m,m^{^{}})`$ from 1 is only $`0.1\%`$. From Eq.(51), the expressions of $`D_1(0)`$ and $`D_3(0)`$ are found
$`D_3(0)`$ $`=`$ $`\sqrt{aa^{^{}}}f(m,m^{^{}}),`$
$`D_1(0)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{aa^{^{}}}}}f(m,m^{^{}}).`$ (54)
The magnetic moment of $`\mathrm{\Sigma }\mathrm{\Lambda }`$ and the decay rate are computed to be
$$\mu _{\mathrm{\Sigma }_0\mathrm{\Lambda }}=1.053$$
(55)
$`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{\alpha }{8}}\mu _{\mathrm{\Sigma }^0\mathrm{\Lambda }}^2{\displaystyle \frac{m_\mathrm{\Sigma }^3}{m_p^2}}(1{\displaystyle \frac{m_\mathrm{\Lambda }^2}{m_\mathrm{\Sigma }^2}})^3=3.79\times 10^3MeV,`$
$`\tau `$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }}}=1.74\times 10^{19}sec.`$ (56)
The experimental upper limit is
$$\tau <1.0\times 10^{14}sec.$$
## 5 Electromagnetic transition of $`p\mathrm{\Delta }(1236)`$
The matrix elements of currents are obtained from Eqs.(18,28,16). Substituting Eqs.(18),(28) into Eq.(16), we derive
$`<`$ $`\mathrm{\Delta }_\lambda ^{^{}}^+(p_f)J_\mu (0)p_\lambda (p_i)>={\displaystyle \frac{ie}{4\sqrt{3}}}({\displaystyle \frac{mm^{^{}}}{EE^{^{}}}})^{\frac{1}{2}}A{\displaystyle \frac{1}{mm^{^{}}}}D_3(q^2)p_\rho q_\sigma `$ (57)
$`\times \epsilon _{\rho \sigma \nu \mu }\overline{\psi }_\nu ^\lambda ^{^{}}(p_f)u_\lambda (p_i){\displaystyle \frac{ie}{\sqrt{3}}}({\displaystyle \frac{mm^{^{}}}{EE^{^{}}}})^{\frac{1}{2}}{\displaystyle \frac{B}{mm^{^{}}}}D_3(q^2)`$
$`\times (p_{f\mu }q_\nu p_fq\delta _{\mu \nu })\overline{\psi }_\nu ^\lambda ^{^{}}(p_f)\gamma _5u_\lambda (p_i),`$
where
$`A`$ $`=`$ $`{\displaystyle \frac{2}{a^{^{}}}}+\kappa \{1+{\displaystyle \frac{m^{^{}}}{m_p}}+{\displaystyle \frac{2}{aa^{^{}}}}{\displaystyle \frac{1}{a}}{\displaystyle \frac{1}{a^{^{}}}}\},`$
$`B`$ $`=`$ $`1{\displaystyle \frac{1}{a^{^{}}}}+{\displaystyle \frac{\kappa }{2}}\{{\displaystyle \frac{1}{a}}+{\displaystyle \frac{1}{a^{^{}}}}{\displaystyle \frac{2}{aa^{^{}}}}\},`$ (58)
and
$$A=1.717,B=0.699.$$
(59)
The S matrix element of $`\gamma p\pi N`$ is written as
$`<\pi NS\gamma p>`$ $`=`$ $`i(2\pi )^4\delta (p_\gamma +p_ip_\pi p_N){\displaystyle \underset{\lambda ^{^{}}}{}}<\pi NU\mathrm{\Delta }_\lambda ^{^{}}^+(p_f)>`$ (60)
$`\times <\mathrm{\Delta }_\lambda ^{^{}}^+(p_f)U\gamma p>{\displaystyle \frac{E_\mathrm{\Delta }}{m_\mathrm{\Delta }}}{\displaystyle \frac{1}{Wm_\mathrm{\Delta }+\frac{i}{2}\mathrm{\Gamma }(W)}},`$
where W is the mass of the final state, $`\mathrm{\Gamma }(W)`$ is the total width of the strong decay of $`\mathrm{\Delta }(1236)`$. The calculation is done in the rest frame of $`\mathrm{\Delta }(1236)`$. $`<\pi NU\mathrm{\Delta }_\lambda ^{^{}}^+(p_f)>`$ is the amplitude of the strong decay of $`\mathrm{\Delta }(1236)`$
$`<\pi NU\mathrm{\Delta }_\lambda ^{^{}}^+(p_f)>=({\displaystyle \frac{m_N}{2E_\pi E_N}})^{\frac{1}{2}}g(W){\displaystyle \frac{p_{\pi \mu }}{m_N}}\overline{u}(p_N)\psi _\mu ^\lambda ^{^{}}.`$ (61)
The electric transition amplitude in Eq.(60) is expressed as
$`<\mathrm{\Delta }_\lambda ^{^{}}U\gamma p>={\displaystyle \frac{1}{\sqrt{2E_\gamma }}}e_\mu <\mathrm{\Delta }_\lambda ^{^{}}J_\mu (0)p>.`$ (62)
By using following equation
$`{\displaystyle \underset{\lambda ^{^{}}}{}}\psi _\mu ^\lambda ^{^{}}\overline{\psi }_\mu ^{^{}}^\lambda ^{^{}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(1+\gamma _4)\{\delta _{\mu \mu ^{^{}}}+{\displaystyle \frac{1}{2}}\gamma _5\gamma _j\epsilon _{j\mu \mu ^{^{}}}\}`$
$`(j,\mu ,\mu ^{^{}}`$ $`=`$ $`1,2,3)`$ (63)
and Eq.(58), we obtain
$`{\displaystyle \underset{\lambda ^{^{}}}{}}\overline{u}_\gamma (p_N)\psi _\mu ^\lambda ^{^{}}`$ $`<`$ $`\mathrm{\Delta }_\lambda ^{^{}}J_\nu (0)p_\lambda >p_{\pi \mu }e_\nu `$ (64)
$`=`$ $`{\displaystyle \frac{eD_3(0)}{24\sqrt{3}m_N^2m_\mathrm{\Delta }}}\{{\displaystyle \frac{m_N(m_N+E_N)}{E_i(m_N+E_i)}}\}^{\frac{1}{2}}\overline{u}_\gamma \{[A(m_N+m_\mathrm{\Delta })^2+B(m_\mathrm{\Delta }^2m_N^2)]`$
$`\times [2𝐤(𝐞𝐩_\pi )+i\sigma \mathrm{𝐞𝐩}_\pi 𝐤i\sigma \mathrm{𝐤𝐩}_\pi 𝐞]`$
$`3iB(m_\mathrm{\Delta }^2m_N^2)(\sigma \mathrm{𝐞𝐩}_\pi 𝐤+\sigma \mathrm{𝐤𝐩}_\pi 𝐞)\}u_\lambda ,`$
where $`E_i`$ is the energy of the initial proton, k is the energy of the photon. The amplitudes of the magnetic and electric transitions are obtained by comparing with the photo production amplitudes in Ref.
$`M1+`$ $`=`$ $`{\displaystyle \frac{eD_3(0)}{96\sqrt{3}\pi m_N^2m_\mathrm{\Delta }}}\{{\displaystyle \frac{m_N+E_N}{m_\mathrm{\Delta }E_i(m+E_i)}}\}^{\frac{1}{2}}{\displaystyle \frac{g(W)p_\pi k}{Wm_\mathrm{\Delta }+\frac{i}{2}\mathrm{\Gamma }(W)}}`$ (65)
$`\times \{A(m_N+m_\mathrm{\Delta })^2+B(m_0^2m_N^2)\},`$
$`E1+`$ $`=`$ $`{\displaystyle \frac{eD_3(0)}{96\sqrt{3}\pi m_N^2m_\mathrm{\Delta }}}\{{\displaystyle \frac{m_N+E_N}{m_\mathrm{\Delta }E_i(m+E_i)}}\}^{\frac{1}{2}}{\displaystyle \frac{g(W)p_\pi k}{Wm_\mathrm{\Delta }+\frac{i}{2}\mathrm{\Gamma }(W)}}`$ (66)
$`\times B(m_\mathrm{\Delta }^2m_N^2),`$
$$\frac{E1+}{M1+}=\frac{B(m_\mathrm{\Delta }m_N)}{A(m_\mathrm{\Delta }+m_N)+B(m_\mathrm{\Delta }m_N)}=5.4\%.$$
(67)
There are several experimental values: -0.045 , -0.073 , -0.024 .
The partial width of $`\mathrm{\Delta }^+(1236)p+\gamma `$ is derived
$`\mathrm{\Gamma }_\gamma `$ $`=`$ $`{\displaystyle \frac{k^2}{2\pi }}{\displaystyle \frac{m_N}{m_\mathrm{\Delta }}}\{A_{\frac{3}{2}}^2+A_{\frac{1}{2}}^2\},`$ (68)
$`A_{\frac{3}{2}}`$ $`=`$ $`{\displaystyle \frac{eD_3(0)(m_\mathrm{\Delta }+m_N)(m_\mathrm{\Delta }^2m_N^2)}{8\sqrt{6}(m_Nm_\mathrm{\Delta })^{3/2}}}\{A+2B{\displaystyle \frac{m_\mathrm{\Delta }m_N}{m_\mathrm{\Delta }+m_N)}}\}`$
$`=`$ $`0.21[GeV]^{\frac{1}{2}},`$
$`A_{\frac{1}{2}}`$ $`=`$ $`{\displaystyle \frac{eD_3(0)(m_\mathrm{\Delta }+m_N)(m_\mathrm{\Delta }^2m_N^2)}{24\sqrt{2}(m_Nm_\mathrm{\Delta })^{3/2}}}\{A2B{\displaystyle \frac{m_\mathrm{\Delta }m_N}{m_\mathrm{\Delta }+m_N)}}\}`$ (69)
$`=`$ $`0.10[GeV]^{\frac{1}{2}}.`$
The experimental data are
$$A_{\frac{3}{2}}=0.24[GeV]^{\frac{1}{2}},A_{\frac{1}{2}}=0.14[GeV]^{\frac{1}{2}}.$$
(70)
The decay width is computed to be
$$\mathrm{\Gamma }_\gamma =0.64MeV,$$
(71)
The experimental data is 0.65MeV.
## 6 Magnetic moment and electromagnetic form factors of $`p\mathrm{\Delta }^+(1236)`$
The differential cross section of the electric production
$`e+pe+\mathrm{\Delta }^+(1236)`$
$`N+\pi `$
is expressed as
$$\frac{1}{\mathrm{\Gamma }_t}\frac{d^2\sigma }{d\mathrm{\Omega }dE^{^{}}}=\sigma _T+\epsilon \sigma _S.$$
(72)
where E’ is the energy of the outgoing electron. Use of the equation
$`{\displaystyle \underset{\lambda }{}}\psi _\mu ^\lambda (p)\overline{\psi }_\mu ^{^{}}^\lambda (p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1{\displaystyle \frac{i}{m^{^{}}}}\widehat{p})\{\delta _{\mu \mu ^{^{}}}+{\displaystyle \frac{2}{3}}{\displaystyle \frac{p_\mu p_\mu ^{^{}}}{m^{}_{}{}^{}2}}{\displaystyle \frac{1}{3}}\gamma _\mu \gamma _\mu ^{^{}}`$ (73)
$`{\displaystyle \frac{i}{3m^{^{}}}}(p_\mu \gamma _\mu ^{^{}}p_\mu ^{^{}}\gamma _\mu )\}`$
and Eq.(58) leads to
$`\sigma _T`$ $`=`$ $`{\displaystyle \frac{m\alpha q^2}{m^{^{}}(W^2m^2)}}{\displaystyle \frac{\mathrm{\Gamma }(W)}{(Wm^{^{}})^2+\frac{1}{4}\mathrm{\Gamma }^2(W)}}{\displaystyle \frac{D_3^2(q^2)}{18m^2}}\{A^2(q^2+m_+^2)`$ (74)
$`+2AB(m^{}_{}{}^{}2m^2q^2)+4B^2(q^2+m_{}^2)(1{\displaystyle \frac{q^2}{q^2}})\},`$
$`\sigma _S`$ $`=`$ $`{\displaystyle \frac{m\alpha q^2}{m^{^{}}(W^2m^2)}}{\displaystyle \frac{\mathrm{\Gamma }(W)}{(Wm^{^{}})^2+\frac{1}{4}\mathrm{\Gamma }^2(W)}}{\displaystyle \frac{2D_3^2(q^2)}{9m^2}}B^2(q^2+m_+^2){\displaystyle \frac{q^2}{q^2}},`$ (75)
where
$$W^2=(p_i+p_ep_e^{^{}})^2,q^2=q^2+\frac{1}{4m^{}_{}{}^{}2}(m^{}_{}{}^{}2m^2q^2)^2.$$
(76)
The ratio of $`\sigma _S`$ and $`\sigma _T`$ is obtained
$`R`$ $`=`$ $`{\displaystyle \frac{\sigma _S}{\sigma _T}}=[4B^2(q^2+m_{}^2){\displaystyle \frac{q^2}{q^2}}]/[(Am_++Bm_{})^2`$ (77)
$`+(AB)^2q^2+3B^2(q^2+m^2)4B^2(q^2+m_{}^2){\displaystyle \frac{q^2}{q^2}}].`$
The behavior of R is released
$`q^2`$ $`=`$ $`0,R=0;`$
$`q^2`$ $``$ $`\mathrm{},R{\displaystyle \frac{1}{q^2}}0.`$ (78)
In the range of $`q^2>3[GeV]^2`$, $`R0.27`$.
According to the definition of mutipoles, the magnetic transition form factor $`G_{M1+}^2(q^2)`$, the electric transition form factor $`G_{E1+}^2(q^2)`$ and $`G_{S1+}^2(q^2)`$ are found
$`G_{M1+}^2(q^2)`$ $`=`$ $`{\displaystyle \frac{D_3^2(q^2)}{18m^2}}\{(Am_++Bm_{})^2+(AB)q^2B^2(q^2+m_{}^2){\displaystyle \frac{q^2}{q^2}}\},`$ (79)
$`G_{E1+}^2(q^2)`$ $`=`$ $`{\displaystyle \frac{D_3^2(q^2)}{18m^2}}B^2(q^2+m_{}^2)(1{\displaystyle \frac{q^2}{q^2}}),`$ (80)
$`G_{S1+}^2(q^2)`$ $`=`$ $`{\displaystyle \frac{D_3^2(q^2)}{18m^2}}B^2(q^2+m_{}^2).`$ (81)
The differential cross section(73) is expressed as
$`{\displaystyle \frac{1}{\mathrm{\Gamma }_t}}{\displaystyle \frac{d^2\sigma }{dE^{^{}}d\mathrm{\Omega }}}`$ $`=`$ $`{\displaystyle \frac{m\alpha q^2}{m^{^{}}(W^2m^2)}}\{G_{M1+}^2(q^2)+3G_{E1+}^2(q^2)`$ (82)
$`+4\epsilon G_{S1+}^2(q^2){\displaystyle \frac{q^2}{q^2}}\}{\displaystyle \frac{\mathrm{\Gamma }(W)}{(Wm^{^{}})^2+\frac{1}{4}\mathrm{\Gamma }^2(W)}}.`$
From Eq.(79), the magnetic moment of $`p\mathrm{\Delta }^+(1236)`$ is derived
$$\mu =G_{M1+}(0)=\frac{D_3(0)}{3\sqrt{2}m}(Am_++Bm_{})^2=1.23\frac{2\sqrt{2}}{3}\mu _p.$$
(83)
The data are
$$1.22\frac{2\sqrt{2}}{3}\mu _p^{[18]},\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1.28}\frac{2\sqrt{2}}{3}\mu _p^{[21]}.$$
Lorentz contraction effect is considered in Eq.(84), for $`p\mathrm{\Delta }^+(1236)`$
$$f(m,m^{^{}})=0.98.$$
(84)
If this effect is ignored, the theoretical value of the magnetic transition moment is $`1.26\frac{2\sqrt{2}}{3}\mu _p`$. The electric multipole moments are obtained from Eq.(81)
$$E1+=\frac{D_3(0)}{3\sqrt{2}m}Bm_{}=0.17,$$
(85)
$`S1+`$ $`=`$ $`E1+,`$
$`{\displaystyle \frac{S1+}{\mu }}`$ $`=`$ $`5.4\%.`$ (86)
The data is
$$\frac{S1+}{\mu }=(5\pm 3)\%.$$
(87)
Theoretical results agree with experimental data.
The expression
$$\sigma _T^R=\frac{m\alpha q^2}{m^{^{}}(W^2m^2)}\frac{\mathrm{\Gamma }(W)}{(Wm^{^{}})^2+\frac{1}{4}\mathrm{\Gamma }^2(W)}G_M^2(q^2)$$
(88)
has been used to determine $`G_M^2(q^2)`$. $`G_M`$ is obtained from Eq.(75) that
$`G_M^2(q^2)`$ $`=`$ $`G_{M1+}^2(q^2)+3G_{E1+}^2(q^2)`$ (89)
$`=`$ $`{\displaystyle \frac{D_3^2(q^2)}{18m^2}}\{(Am_++Bm_{})^2+(AB)^2q^2+B^2(q^2+m_{}^2)(34{\displaystyle \frac{q^2}{q^2}})\}.`$
The mass difference between $`\mathrm{\Delta }^+(1236)`$ and proton is ignored. The integral $`D_3(q^2)`$ for $`p\mathrm{\Delta }^+(1236)`$ is expressed as
$$D_3(q^2)=\frac{4mm^{^{}}\sqrt{aa^{^{}}}}{(m+m^{^{}})^2}(1+2.39\frac{q^2}{4m^2})^1(1+\frac{q^2}{0.71})^1.$$
(90)
Substituting Eq.(92) into Eq.(91), the expression of $`G_M^2(q^2)`$ is obtained. Comparisons with experimental data are shown in Fig.5 and 6. The data for Fig.5 comes from Ref. and that for Fig.6 comes from Ref.. It can be seen from these two figures that as $`q^2`$ increases, the theoretical curve drops a little bit faster than the experimental one. At $`q^2=0.8[GeV]^2`$, the theoretical value is $`10\%`$ less than the experimental value. This difference can be regarded as to be from the ignorance of the mass difference between proton and $`\mathrm{\Delta }^+(1236)`$.
Taking
$$W=m^{^{}}=1.236GeV,\mathrm{\Gamma }(m^{^{}})=0.12GeV,$$
(91)
we obtain
$$\sigma _S=48.4q^2(q^2+0.0888)(1+0.679q^2)^2(1+\frac{q^2}{0.71})^4\times 10^{28}cm^2.$$
(92)
Comparison with the data is shown in Fig.7.
## 7 Discussion
$`SU(6)`$ summetric wave functions of $`\frac{1}{2}^+`$ and $`\frac{3}{2}^+`$ of s- wave(in the frame of center of mass) are applied to study the electromagnetic form factors of nucleons and $`p\mathrm{\Delta }`$. A new expression of $`\frac{\mu _pG_E^p(q^2)}{G_M^p(q^2)}`$ is obtained. Nonzero electric form factor $`G_E^n(q^2)`$ is found. The magnetic form factor of $`p\mathrm{\Delta }`$ decreases faster than $`G_M^p`$. Nonzero multipole moments $`E1+`$ and $`S1+`$ are obtained. They are small and negative. It is interesting to point out that nonzero $`G_E^n`$, $`E1+`$, and $`S1+`$ are resulted in the addtional terms of the wave function, which are constructed by the spinors of antiquarks. The amplitudes and decay rate of $`\mathrm{\Delta }p+\gamma `$ are computed and theory agrees with data.
The magnetic moments of hyperons are calculated under the assumption that the anomalous magnetic moment of strange quark(except for the charge factor) is the same as the one of u and d quarks.
Figure Captions
FIG. 1. Ratio of electric and magnetic form factors of proton.
FIG. 2. Ratio of electric and magnetic form factors of proton.
FIG. 3. Electic form factor of neutron.
FIG. 4. Electric form factor of neutron.
FIG. 5. Magnetic form factor of $`p\mathrm{\Delta }`$.
FIG. 6. Magnetic form factor of $`p\mathrm{\Delta }`$.
FIG. 7. Cross Section of virtual scalar photon.
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# Laterally coupled few-electron quantum dots
## I Introduction
Semiconductor quantum dots are manmade nanoscale structures in which electrons are confined in all three spatial directions similar to the physical situation in atoms. As they show typical atomic properties like discrete energy levels and shell structures they are often referred to as artificial atoms. However, in contrast to natural atoms, in quantum dots the number of electrons $`N`$ is tunable and the characteristic lengths of the system corresponding to external confinement potential, electron-electron interaction, and an applied magnetic field are of comparable size. Therefore, these systems are ideal objects to study interaction effects such as transitions in the ground-state spin configuration as a function of the magnetic field.
Starting from quantum dots as a structure which is well understood by now, more complex systems are conceivable and likely to have perspective in future applications <sup>,</sup><sup>,</sup>. A simple example is analogous to a two-atom molecule consisting of two coupled quantum dots. The more principal properties of coupled quantum dots follow from the simple model of e.g. two disk-like dots side by side in a plane or on top of each other (lateral or vertical quantum-dot molecules (QDMs)), respectively, which are the subject of experimental<sup>,</sup><sup>,</sup> and theoretical<sup>,</sup><sup>,</sup><sup>,</sup> studies. Tunneling and overlap of orbital wavefunctions lead to bonding and antibonding states with defined parity, which are delocalized over the two-dot system as in natural systems. The advantage of QDMs is the tunability of the electron number and of the coupling strength by proper design. The latter allows to investigate the transition between weak and strong coupling.
The increased interest in coupled quantum dots is indicated by a growing number of papers on this topic: Austing et al. measured addition spectra of vertical QDM in different coupling regimes. The characteristic features found in their spectra have been confirmed by Imamura et al. and Rontani et al. applying exact diagonalization techniques. These ground-state calculations were performed for up to 6 electrons as a function of the interdot distance and an external magnetic field. Coulomb blockade effects in lateral QDMs have been measured by Waugh et al.<sup>,</sup>, who found a conductance pattern similar to that of single dots for strong coupling, whereas the weakly coupled systems showed a pairing of conductance peaks. Theoretical work on ground-state properties of lateral QDMs was performed by Yannouleas and Landman using an unrestricted Hartree-Fock approach and by Nagaraja et al. applying SDFT.
In this paper we investigate the electronic structure of the ground state of two laterally coupled (identical) quantum dots as a function of the interdot distance. In our calculations, applying SDFT, we focus on spin-dependent effects in few-electron quantum dot molecules containing 2, 4, and 8 electrons. In the case of 4 (8) electrons we find a change of spin configuration with decreasing coupling strength from $`S_z=1`$ ($`S_z=1`$) to $`S_z=0`$ ($`S_z=2`$). For 2 electrons SDFT yields a transition from spin singlet to spin triplet. It turns out, however, that the latter result is an artifact of the SDFT, as can be shown (at least for our choice of the confinement potential) by using alternative approaches. Furthermore, we show the characteristic Coulomb staircases for the chemical potential in the weak and strong coupling limits, respectively.
For all our calculations we assume that the two dots are coupled in a quantum-mechanically coherent way. This means we use delocalized electron states extended over the whole quantum dot molecule whereas states localized to the left or right dot do not exist in this description. Note, that this model does not contain all aspects of the actual physical situation where an electron occupying a covalent state interacts with other electrons and charges in the surrounding semiconductor material. These interactions may cause a dephasing of the quantum mechanical wavefunction leading to a breakdown of the delocalized state. Up to now no theory is available which allows a reliable approximation for dephasing rates in realistic nanostructures.
The outline of this paper is the following: In the next section we describe our model of the QDM and discuss the single-particle spectrum of the non-interacting system. In Sec. III we sketch the basics of SDFT and of the numerical method for the simultaneous solution of the Kohn-Sham (KS) equations and calculation of the selfconsistent potentials. The following part is devoted to the 2-electron QDM. We compare the SDFT results with exact (analytical) results and with the Heitler-London approach to point out a particular problem of SDFT. Sec. V relates the spin configuration for 4 and 8 electrons to the single-particle spectrum and summarizes the corresponding SDFT results. Finally we show the typical Coulomb staircase of the addition energies for the weak and the strong coupling limits, respectively.
## II The model and its single-particle system
We consider quantum dots prepared by laterally confining two-dimensional electrons in a GaAs/AlGaAs heterostructure (material parameters of GaAs: $`m^{}=0.067m_e`$, $`\epsilon =12.4`$). As the lateral confinement is much weaker than that in growth direction ($`z`$-direction) we adopt the standard situation of the electronic quantum limit and describe the electron density by $`n(x,y,z)=n(x,y)\delta (z)`$. Thus, our system is defined by the effective two-dimensional Hamiltonian
$`H`$ $`=`$ $`T+V+W`$ (1)
$`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}\left({\displaystyle \frac{𝐩_j^2}{2m^{}}}+V(𝐫_j)\right)+{\displaystyle \frac{e^2}{4\pi \epsilon \epsilon _0}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{k=j+1}{\overset{N}{}}}{\displaystyle \frac{1}{\left|𝐫_j𝐫_k\right|}}`$ (2)
with $`𝐩_j`$ and $`𝐫_j`$ being vectors in the $`xy`$-plane. Our choice of the external confinement potential $`V(𝐫)`$ is
$$V(𝐫)=\frac{1}{2}m^{}\omega _0^2\mathrm{min}(\left(𝐫𝐋\right)^2,\left(𝐫+𝐋\right)^2),$$
(3)
where the coupling strength depends on the distance $`d=\left|2𝐋\right|`$ between the potential minima at $`\pm 𝐋`$ in the $`xy`$-plane (Fig. 1). Along the line between these minima there is a barrier of height
$$V_\mathrm{B}(d)=\frac{1}{2}m^{}\omega _0^2\left(\frac{d}{2}\right)^2.$$
(4)
This model potential describes a single parabolic quantum dot for $`d=0`$ as the extreme case of strong coupling and two separate quantum dots of the same shape in the weak-coupling limit for $`d\mathrm{}`$. Due to the broken axial symmetry of $`V(𝐫)`$ for $`d0`$ the numerical requirements are increased in comparison with those for a single circular quantum dot.
In Fig. 2 we show the lowest energy levels of the noninteracting electrons in dependence on the distance between the dot centers for the typical confinement energy $`\mathrm{}\omega _0=3`$ $`\mathrm{meV}`$. For $`d=0`$, when the confinement potential of the QDM degenerates to the simple parabolic potential with axial symmetry, the spectrum shows the typical shell structure of an isotropic harmonic oscillator. For large distances $`d`$ we obtain the same spectrum, now twofold because the system consists of two identical but completely separated quantum dots. By switching on the coupling (or lowering $`d`$ starting from the weak-coupling limit) we recover the anticipated properties of a diatomic molecule: the energies decrease (increase) with respect to the reference level due to formation of bonding (antibonding) states, which have even (odd) parity as indicated by solid (dashed) lines. The intermediate coupling regime is dominated by crossings leading to the rearrangement of the level structure. In the weak-coupling limit the nth bundle of the twofold level structure with energy $`En\mathrm{}\omega _0`$ consists of $`n`$ bonding states and $`n`$ antibonding states.
## III SDFT and KS equations
In order to include the electron-electron interaction we employ the spin-density functional theory (SDFT). It is the generalization of the DFT formalism, originally established by Hohenberg, Kohn, and Sham<sup>,</sup>, to spin-polarized systems by including the coupling of the magnetization to an applied magnetic field. Accordingly, the Hohenberg-Kohn (HK) theorem has to be modified with respect to the spin degrees of freedom. For this case, it states that two different nondegenerate ground-state wavefunctions $`|\mathrm{\Psi }`$ and $`|\mathrm{\Psi }^{}`$ always yield different tupels $`(n(𝐫),𝐦(𝐫))(n^{}(𝐫),𝐦^{}(𝐫))`$ of electron density $`n(𝐫)`$ and magnetization $`𝐦(𝐫)`$. This is sufficient to establish a functional of the total energy with the usual functional properties
$$E_{V_0,𝐁_0}[n,𝐦]=F_{\mathrm{HK}}[n,𝐦]+d𝐫\left[V(𝐫)n(𝐫)𝐁(𝐫)𝐦(𝐫)\right]$$
(5)
and the universal HK functional
$$F_{\mathrm{HK}}[n,𝐦]=\mathrm{\Psi }[n,𝐦]|T+W|\mathrm{\Psi }[n,𝐦].$$
(6)
In the limit $`𝐁0`$ considered here, the SDFT scheme can yield a spin-polarized ground state for even electron numbers due to Hund’s rule<sup>,</sup>. These effects have already attracted much interest and will also be discussed in this paper.
The spin-degree of freedom is considered in the Kohn-Sham (KS) equations by assuming the total spin $`S_z`$ in $`z`$-direction to be a good quantum number
$`\{{\displaystyle \frac{\mathrm{}^2}{2m^{}}}^2+V(𝐫)+{\displaystyle \frac{e^2}{4\pi \epsilon \epsilon _0}}{\displaystyle }\mathrm{d}𝐫^{}{\displaystyle \frac{n(𝐫^{})}{\left|𝐫𝐫^{}\right|}}`$ (7)
$`+V_{\mathrm{XC}}^\sigma ([n_+,n_{}],𝐫]\}\phi _j^\sigma (𝐫)=\epsilon _j^\sigma \phi _j^\sigma (𝐫)`$ (8)
with the spin $`\sigma =\pm `$ in $`z`$-direction and the KS energies $`\epsilon _1^\sigma \epsilon _2^\sigma \mathrm{}`$ . For a system containing $`N`$ particles we calculate the occupation numbers of the KS levels in the ground state due to
$`\gamma _j^\sigma =1`$ $`\epsilon _j<\mu `$ (9)
$`0\gamma _j^\sigma 1`$ $`\epsilon _j=\mu `$ (10)
$`\gamma _j^\sigma =0`$ $`\epsilon _j>\mu `$ (11)
($`\mu `$: chemical potential) with the constraints
$`{\displaystyle \underset{j}{}}\gamma _j^\sigma `$ $`=`$ $`N^\sigma `$ (12)
$`N^++N^{}`$ $`=`$ $`N.`$ (13)
This leads to the spin densities
$$n_\sigma (𝐫)=\underset{j}{}\gamma _j^\sigma \left|\phi _j^\sigma (𝐫)\right|^2,$$
(14)
the total density
$$n(𝐫)=n_+(𝐫)+n_{}(𝐫),$$
(15)
and the magnetization in $`z`$-direction ($`\mu _B`$: Bohr’s magneton)
$$m_z(𝐫)=\mu _\mathrm{B}\left(n_+(𝐫)n_{}(𝐫)\right).$$
(16)
The exchange-correlation (xc) potentials
$$V_{\mathrm{XC}}^\sigma ([n_+,n_{}],𝐫)=\frac{\delta E_{\mathrm{XC}}[n_+,n_{}]}{\delta n_\sigma (𝐫)}$$
(17)
are defined as functional derivatives of the xc energy functional
$`E_{\mathrm{XC}}[n_+,n_{}]`$ $`=`$ $`F_{\mathrm{HK}}[n_+,n_{}]{\displaystyle \frac{1}{2}}{\displaystyle \frac{e^2}{4\pi \epsilon \epsilon _0}}{\displaystyle d𝐫d𝐫^{}\frac{n(𝐫)n(𝐫^{})}{\left|𝐫𝐫^{}\right|}}`$ (19)
$`T_\mathrm{S}[n_+,n_{}].`$
($`T_\mathrm{S}[n_+,n_{}]`$ denotes the kinetic energy functional of the KS system.) The total ground-state energy $`E_0`$ of the interacting system can be calculated from
$`E_0`$ $`=`$ $`{\displaystyle \underset{j,\sigma }{}}\gamma _j^\sigma \epsilon _j^\sigma {\displaystyle \frac{1}{2}}{\displaystyle \frac{e^2}{4\pi \epsilon \epsilon _0}}{\displaystyle d𝐫d𝐫^{}\frac{n(𝐫)n(𝐫^{})}{\left|𝐫𝐫^{}\right|}}`$ (21)
$`+E_{\mathrm{XC}}[n_+,n_{}]{\displaystyle \underset{\sigma }{}}{\displaystyle d𝐫V_{\mathrm{XC}}^\sigma ([n_+,n_{}],𝐫)n_\sigma (𝐫)}.`$
Concerning the xc potentials we apply the local spin density approximation (LSDA)
$$E_{\mathrm{XC}}[n_+,n_{}]d𝐫\left(n_+(𝐫)+n_{}(𝐫)\right)\epsilon _{\mathrm{XC}}(n_+(𝐫),n_{}(𝐫))$$
(22)
and the Padé approximation of $`\epsilon _{\mathrm{XC}}(n_+(𝐫),n_{}(𝐫))`$ in two dimensions following Tanatar and Ceperley.
For the numerical solution of the KS equations we calculate all quantities with spatial dependence on a two-dimensional grid in real space. The damped gradient iteration ensures a simultaneous solution of the KS equations together with the corresponding self-consistent potentials. This method uses the iteration scheme
$$|\mathrm{\Psi }^{(k+1)}=H^1|\mathrm{\Psi }^{(k)}$$
(23)
with the KS Hamiltonian $`H`$ (only positive eigenvalues), which converges to the ground-state wavefunction. The inversion of the Hamiltonian is performed approximately<sup>,</sup> using
$`|\mathrm{\Psi }^{(k+1)}`$ $``$ $`\{[1\stackrel{~}{x}_0({\displaystyle \frac{p_x^2}{2m^{}}}+V_x^{\mathrm{}}+E_{x0})^1`$ (26)
$`({\displaystyle \frac{p_y^2}{2m^{}}}+V_y^{\mathrm{}}+E_{y0})^1`$
$`(H\mathrm{\Psi }^{(k)}|H|\mathrm{\Psi }^{(k)})]\}|\mathrm{\Psi }^{(k)}`$
instead of Eq. (23), where $`\stackrel{~}{x}_0`$, $`E_{x0}`$, $`E_{y0}`$ are iteration parameters and $`V_x^{\mathrm{}}`$, $`V_y^{\mathrm{}}`$ the asymptotic contributions of the external potential for $`|𝐫|\mathrm{}`$. Excited state KS wavefunctions are calculated using the same iteration scheme with an additional orthogonalization routine. The Hartree and xc potentials depending on the densities of the KS wavefunctions are recalculated in each step. Thus, the KS Hamiltonian is modified in each step of the iteration.
## IV Results for the two-electron QDM
Fig. 3 depicts the KS energies for a quantum dot molecule containing two electrons as a function of the distance $`d`$ between the centers of the molecule (we assume $`\mathrm{}\omega _0=3`$ $`\mathrm{meV}`$).
Starting with $`d=0`$ the QDM degenerates into a circular parabolic quantum dot showing the known shell structure of the noninteracting case. Due to the closed shell for 2 electrons the ground state is not spin-polarized. This ground-state spin configuration is stable for increasing $`d`$ as long as the energy gap between the levels (a) and (b) in the single-particle spectrum (Fig. 2) is not too small. At $`d4.5`$ $`\mathrm{a}_0^{}`$ ($`\mathrm{a}_0^{}=9.79`$ $`\mathrm{nm}`$ denotes the effective Bohr’s radius for GaAs) the unpolarized ground state changes into a polarized one. This state with aligned spins persists as SDFT ground state even for $`d\mathrm{}`$, because the energy gap between states (a) and (b) continues to decrease with increasing $`d`$. This result, similar to that obtained e.g. by Nagaraja et al., differs from those of the physically analogous problems of the vertical QDMs<sup>,</sup> and from the hydrogen molecule whose ground states are spin singlets.
The SDFT result is also in contrast to a mathematical theorem stating that the spin configuration of the ground state of a two-electron quantum dot molecule, with our definition of the external confinement potential, should be a spin singlet state: As the total Hamiltonian $`H`$, Eq. (1), does not depend on spin coordinates we can focus on the spatial wavefunctions. Due to Theorem XIII.47 of Ref. the (spatial) ground state wavefunction of the considered two-body-problem is positive and nondegenerate, this means it is symmetric in the spatial coordinates. As a consequence of Pauli’s principle the ground state has to be a spin singlet state.
In view of this rigorous statement, the SDFT results of Fig. 3 for $`d4.5`$ $`\mathrm{a}_0^{}`$ cannot be considered as true ground state values (the same holds for the corresponding SDFT densities).
We test this finding by using a perturbation approach analogous to the Heitler-London model for the hydrogen molecule, which in contrast to SDFT has the advantage of yielding an upper limit for the ground state energy. Similar calculations for other shapes of the external potential were done in Ref. .
For the two-electron QDM we rewrite the total Hamiltonian (1) as
$`H`$ $`=`$ $`{\displaystyle \frac{𝐩_1^2}{2m^{}}}+{\displaystyle \frac{1}{2}}m^{}\omega _0^2\mathrm{min}(\left(𝐫_1𝐋\right)^2,\left(𝐫_1+𝐋\right)^2)`$ (29)
$`{\displaystyle \frac{𝐩_2^2}{2m^{}}}+{\displaystyle \frac{1}{2}}m^{}\omega _0^2\mathrm{min}(\left(𝐫_2𝐋\right)^2,\left(𝐫_2+𝐋\right)^2)`$
$`+{\displaystyle \frac{e^2}{4\pi \epsilon \epsilon _0}}{\displaystyle \frac{1}{\left|𝐫_1𝐫_2\right|}}`$
$`=`$ $`H_1^\mathrm{R}+H_2^\mathrm{L}+U_1^\mathrm{R}+U_2^\mathrm{L}+W`$ (30)
with
$`H_j^{\mathrm{R}/\mathrm{L}}`$ $`=`$ $`{\displaystyle \frac{𝐩_j^2}{2m^{}}}+{\displaystyle \frac{1}{2}}m^{}\omega _0^2\left(𝐫_j𝐋\right)^2`$ (31)
$`U_j^{\mathrm{R}/\mathrm{L}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}m^{}\omega _0^2\mathrm{min}(0,\pm 4𝐫_j𝐋)`$ (32)
$`W`$ $`=`$ $`{\displaystyle \frac{e^2}{4\pi \epsilon \epsilon _0}}{\displaystyle \frac{1}{\left|𝐫_1𝐫_2\right|}}.`$ (33)
The ground-state wavefunctions of the shifted harmonic oscillators $`H^{\mathrm{R}/\mathrm{L}}`$ are denoted by $`|\mathrm{\Psi }^{\mathrm{R}/\mathrm{L}}`$. So we can use the typical Heitler-London ansatz for the (unnormalized) singlet $`|\mathrm{\Psi }^+`$ and triplet $`|\mathrm{\Psi }^{}`$ spatial wavefunctions
$$|\mathrm{\Psi }^\pm =\frac{1}{\sqrt{2}}\left(|\mathrm{\Psi }^\mathrm{R}^{(1)}|\mathrm{\Psi }^\mathrm{L}^{(2)}\pm |\mathrm{\Psi }^\mathrm{R}^{(2)}|\mathrm{\Psi }^\mathrm{L}^{(1)}\right).$$
(34)
to calculate the expectation value of the ground-state energy of the system
$`E_{\mathrm{HL}}^\pm `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }^\pm |H|\mathrm{\Psi }^\pm }{\mathrm{\Psi }^\pm |\mathrm{\Psi }^\pm }}=`$ (35)
$`=`$ $`2\mathrm{}\omega _0+{\displaystyle \frac{\mathrm{\Psi }^\mathrm{R}|U^\mathrm{R}|\mathrm{\Psi }^\mathrm{R}+\mathrm{\Psi }^\mathrm{L}|U^\mathrm{L}|\mathrm{\Psi }^\mathrm{L}}{1\pm \mathrm{e}^{2L^2/l_h^2}}}`$ (39)
$`+{\displaystyle \frac{{}_{}{}^{(1)}\mathrm{\Psi }^\mathrm{R}|{}_{}{}^{(2)}\mathrm{\Psi }^\mathrm{L}|W|\mathrm{\Psi }^\mathrm{R}_{}^{(1)}|\mathrm{\Psi }^\mathrm{L}_{}^{(2)}}{1\pm \mathrm{e}^{2L^2/l_h^2}}}`$
$`\pm {\displaystyle \frac{\mathrm{e}^{L^2/l_h^2}\left(\mathrm{\Psi }^\mathrm{R}|U^\mathrm{R}|\mathrm{\Psi }^\mathrm{R}+\mathrm{\Psi }^\mathrm{L}|U^\mathrm{L}|\mathrm{\Psi }^\mathrm{L}\right)}{1\pm \mathrm{e}^{2L^2/l_h^2}}}`$
$`\pm {\displaystyle \frac{{}_{}{}^{(1)}\mathrm{\Psi }^\mathrm{R}|{}_{}{}^{(2)}\mathrm{\Psi }^\mathrm{L}|W|\mathrm{\Psi }^\mathrm{L}_{}^{(1)}|\mathrm{\Psi }^\mathrm{R}_{}^{(2)}}{1\pm \mathrm{e}^{2L^2/l_h^2}}}.`$
Note that all matrix elements are real quantities. $`l_h=\sqrt{\mathrm{}/m^{}\omega _0}`$ denotes the characteristic oscillator length. The result for the matrix elements is
$`\mathrm{\Psi }^\mathrm{R}|U^\mathrm{R}|\mathrm{\Psi }^\mathrm{R}+\mathrm{\Psi }^\mathrm{L}|U^\mathrm{L}|\mathrm{\Psi }^\mathrm{L}=`$ (40)
$`=2\mathrm{}\omega _0\left({\displaystyle \frac{L}{l_h\sqrt{\pi }}}\mathrm{e}^{L^2/l_h^2}+{\displaystyle \frac{L^2}{l_h^2}}\left(1\mathrm{erf}\left(L/l_h\right)\right)\right)`$ (41)
$`{}_{}{}^{(1)}\mathrm{\Psi }^\mathrm{R}|{}_{}{}^{(2)}\mathrm{\Psi }^\mathrm{L}|W|\mathrm{\Psi }^\mathrm{R}_{}^{(1)}|\mathrm{\Psi }^\mathrm{L}_{}^{(2)}=`$ (42)
$`=\mathrm{}\omega _0\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle \frac{l_h}{\mathrm{a}_0^{}}}\mathrm{e}^{L^2/l_h^2}\mathrm{I}_0(L^2/l_h^2)`$ (43)
$`\mathrm{\Psi }^\mathrm{R}|U^\mathrm{R}|\mathrm{\Psi }^\mathrm{L}+\mathrm{\Psi }^\mathrm{L}|U^\mathrm{L}|\mathrm{\Psi }^\mathrm{R}=2\mathrm{}\omega _0{\displaystyle \frac{L}{l_h\sqrt{\pi }}}\mathrm{e}^{L^2/l_h^2}`$ (44)
$`{}_{}{}^{(1)}\mathrm{\Psi }^\mathrm{R}|{}_{}{}^{(2)}\mathrm{\Psi }^\mathrm{L}|W|\mathrm{\Psi }^\mathrm{L}_{}^{(1)}|\mathrm{\Psi }^\mathrm{R}_{}^{(2)}=\mathrm{}\omega _0\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle \frac{l_h}{\mathrm{a}_0^{}}}\mathrm{e}^{2L^2/l_h^2}`$ (45)
with the modified Bessel’s function $`\mathrm{I}_0(x)`$ and the error function $`\mathrm{erf}\left(x\right)`$ .
In Fig. 4 we compare the ground-state energies resulting from SDFT calculations with the Heitler-London approach. For $`0d4.5`$ $`\mathrm{a}_0^{}`$ the SDFT singlet ground-state energies are lower than the HL energies. As the comparison between SDFT and exact diagonalization for 2-electron quantum dots shows good agreement, we expect this to hold also for the QDM with small $`d`$. Moreover, the HL triplet energy diverges for $`d0`$ and becomes larger than the HL singlet energy. Therefore, we regard the SDFT results for small $`d`$ as a better approximation of the ground state than the HL results.
If, however, $`d`$ is increased beyond $`4.5`$ $`\mathrm{a}_0^{}`$ the HL energies become lower than the SDFT results and replace them as ground-state energies. It can be shown analytically that for sufficiently large $`d`$ the HL energies fulfill the relation $`E_{\mathrm{HL}}^+<E_{\mathrm{HL}}^{}`$. Thus, in the limit $`d\mathrm{}`$ the HL results are consistent with the mathematical theorem quoted above. In addition, in the weak-coupling limit the interaction is reduced to the classical Coulomb repulsion between point charges and the total energy approaches asymptotically $`E=2\mathrm{}\omega _0+e^2/(4\pi \epsilon \epsilon _0d)`$. This is also consistent with the energies of the HL approach.
The energy difference of $`0.36`$ $`\mathrm{meV}`$ between $`E_{\mathrm{SDFT}}^{}`$ and $`E_{\mathrm{HL}}^+`$ in the weak-coupling limit can be ascribed to the self-interaction within the SDFT scheme: The charge density of a weakly coupled quantum dot molecule is mainly localized in its centers (charge 1 $`e^{}`$ per center). However, within the SDFT formalism the energy of each charge in its own Coulomb field is counted additionally leading to an enhanced total energy. To illustrate this, let us to apply the SDFT concept to a single-particle problem. For a circular parabolic quantum dot with $`\mathrm{}\omega _0=3`$ $`\mathrm{meV}`$ this gives a ground-state energy of $`3.18`$ $`\mathrm{meV}`$, i.e. the self-interaction error is $`0.18`$ $`\mathrm{meV}`$. This value has to be doubled for the $`d\mathrm{}`$ limit of the QDM. Thus, we identify the difference $`E_{\mathrm{SDFT}}^{}E_{\mathrm{HL}}^\pm `$ as the self-interaction energy in this limit.
The energy difference between the singlet state and the triplet state in the SDFT approach can be qualitatively explained as an xc effect: Both spin configurations show similar total electron densities. Consequently the contributions from the Hartree energy and the spin-independent xc energy can be taken to be approximately equal for singlet and triplet states. However, the triplet total energy is in addition lowered by spin-dependent xc effects leading to the energy difference in Fig. 4.
To summarize this section: for small $`d`$ ($`0d4.5`$ $`\mathrm{a}_0^{}`$) we find the SDFT (singlet) energies to be the better approximation to the ground state than the HL results, while for larger $`d`$ ($`d4.5`$ $`\mathrm{a}_0^{}`$) the SDFT energies become too large due to the self-interaction error and the HL results provide the better approximation of the ground-state energy.
## V Results for QDMs containing 4 and 8 electrons
In this section we present our results for the ground-state energy and spin-configuration for artificial molecules containing 4 and 8 electrons as a function of the interdot distance $`d`$. Our calculations show that a spin-polarized ground state is accompanied by a (quasi-)degeneracy at the Fermi level in the noninteracting single-particle spectrum whereas unpolarized ground states require a sufficiently large gap between the highest occupied and lowest unoccupied single-particle levels of the noninteracting spectrum.
Single dots with 4 electrons have a spin-polarized ground state ($`S_z=1`$) due to Hund’s rule<sup>,</sup>. This is the case also for the QDM with $`d=0`$ (Fig. 5). The spin-polarized ground state persists with increasing $`d`$ as long as the xc energy overcompensates the energy costs of putting electrons with parallel spins into the bonding and antibonding states ((b) and (c) in Fig. 2) evolving from the second level at $`d=0`$. For $`d>2`$ $`\mathrm{a}_0^{}`$ it becomes energetically more favorable to occupy both the bonding level (a) and antibonding level (b) (Fig. 2) with two particles forming an $`S_z=0`$ ground state. This configuration remains stable even in the weak coupling limit and reflects the noninteracting particle picture (Fig. 2) which does not show level crossings at the Fermi energy for $`N=4`$ and large $`d`$. This result is in agreement with the work of Nagaraja et al. (laterally coupled dots), Rontani et al., and Imamura et al. (vertically coupled dots). On the other hand, Nagaraja et al. do not find a spin-polarized ground state in the strong coupling limit, because their confinement potential obtained by simulating external gates does not approach that of a circular parabolic dot as in our case.
Although even in the limit of large $`d`$ a quantum-mechanical coupling (delocalized bonding and antibonding states) between the dots can be observed within the SDFT scheme the resulting density (Fig. 6) and energy (Fig. 5) tend towards a semiclassical picture: the densities of the dots are well separated and the total energy of the two coupled dots consists of twice the (quantum-mechanical) energy of a single dot containing 2 electrons and the energy of classical repulsion of two point-like charges (2$`e^{}`$) in the distance $`d`$ (see inset of Fig. 5). In contrast to the 2-electron QDM, this test for the asymptotics of the total energy being fulfilled underlines the reliabilty of the SDFT results for 4 electrons.
Turning to QDMs with 8 electrons (Fig. 7) we find for $`d=0`$ the expected spin-polarized ground state ($`S_z=1`$)<sup>,</sup>. Increasing $`d`$ up to $`7.5`$ $`\mathrm{a}_0^{}`$ this state persists with double-occupancy of the levels (a), (b), and (c) of Fig. 2 and single-occupancy of the quasi-degenerate levels (d) and (e) with aligned spins. At $`d7.5`$ $`\mathrm{a}_0^{}`$ the levels (c), (d), (e), and (f) of Fig. 2 get close enough that an alignment of all four spins at the Fermi level can lower the total energy. Therefore, we find a spin-polarized ground state in the weak coupling regime ($`d>8`$ $`\mathrm{a}_0^{}`$) in agreement with Nagaraja et al.. Although the spin-aligned ground state was ruled out in the case of 2 electrons we believe that the $`S_z=2`$ ground state for the 8-electron artificial molecule in the weak-coupling regime is reliable. This is emphasized by the correct asymptotics of the total energy in the limit $`d\mathrm{}`$ (see inset of Fig. 7). Exact diagonalizations<sup>,</sup> for vertically coupled dots performed for up to 6 electrons yield a spin-polarized ground state ($`S_z=1`$) for the 6-electron artificial molecule in the weak coupling regime, too, indicating that spin-polarized ground states can exist in principle for QDMs.
Also for $`N=8`$, in the weak-coupling limit, the total energy and density can be understood in a semiclassical picture: the densities (for sufficiently large $`d`$) are well separated (Fig. 8) and the total energy follows that of two identical dots each containing 4 electrons plus the Coulomb repulsion between two point-like charges (4$`e^{}`$) in distance $`d`$. However, besides the quantum mechanical phenomenon of delocalized bonding and antibonding states we find a (nonclassical) xc energy related effect leading to the spin-polarized ground state.
Yannouleas and Landman find from unrestricted Hartree-Fock calculations electron localization in the form of electron puddles and Wigner supermolecules for a ratio $`\lambda `$ between the oscillator length and Bohr’s radius of about $`1.4`$. In contrast, (exact) Quantum Monte Carlo calculations indicate a transition from the Fermi liquid to the Wigner crystal regime if $`\lambda `$ is increased beyond a critical value of 4. Thus, our calculations performed for $`\lambda 2`$ are obviously in the Fermi liquid regime.
## VI Chemical potential
Although recent experiments succeeded in identifying molecular properties of coupled quantum dots like bonding and antibonding states the most convenient comparison between theoretical and experimental results is based on the chemical potential of the system. The chemical potential is defined as difference of total energies $`\mu (N)=E(N)E(N1)`$ which can be measured by capacitance spectroscopy<sup>,</sup><sup>,</sup>.
As it turns out, the dependence of $`N`$ on $`\mu `$ (Coulomb staircase) exhibits features which depend on the coupling strength<sup>,</sup><sup>,</sup>. While in the strong coupling limit the capacitance spectrum of the QDM is similar to that of the single quantum dot, it exhibits small spacings between the chemical potentials for $`N=1,2`$, $`N=3,4`$, $`N=5,6`$ etc. for the weakly coupled QDM which correspond to the pairing of peaks observed in capacitance spectroscopy<sup>,</sup>. Explanations for this phenomenon provided in the literature are based on states localized at the centers of the molecule and neglect quantum mechanical coupling. Using these assumptions we can easily estimate the dependence of $`\mu `$ on $`N`$ for our system: If the QDM contains an even number $`N`$ of electrons one may assume that $`N/2`$ electrons are localized at each dot. The total energy of the artificial molecule $`E_\mathrm{M}(N)`$ can be estimated by the energy of the single dots containing $`N/2`$ electrons ($`E_\mathrm{D}(N/2)`$) and the Coulomb repulsion between the dots which is approximated as Coulomb energy between point-like charges
$$E_\mathrm{M}(N)=2E_\mathrm{D}\left(\frac{N}{2}\right)+\frac{e^2(N/2)^2}{4\pi \epsilon \epsilon _0d}\text{}N\text{ even.}$$
(46)
For odd $`N`$ one assumes $`(N+1)/2`$ electrons to be in one dot and $`(N1)/2`$ in the other one and the approximation for the total energy yields
$`E_\mathrm{M}(N)`$ $`=`$ $`E_\mathrm{D}\left({\displaystyle \frac{N+1}{2}}\right)+E_\mathrm{D}\left({\displaystyle \frac{N1}{2}}\right)`$ (48)
$`+{\displaystyle \frac{e^2(N+1)(N1)/4}{4\pi \epsilon \epsilon _0d}}\text{}N\text{ odd.}`$
Using these equations we can calculate the chemical potential $`\mu _\mathrm{M}(N)=E_\mathrm{M}(M)E_\mathrm{M}(N1)`$ of the QDM for even and odd $`N`$
$$\mu _\mathrm{M}(N)=E_\mathrm{D}\left(\frac{N}{2}\right)E_\mathrm{D}\left(\frac{N}{2}1\right)+\frac{e^2N/2}{4\pi \epsilon \epsilon _0d}\text{}N\text{ even,}$$
(49)
$$\mu _\mathrm{M}(N)=E_\mathrm{D}\left(\frac{N+1}{2}\right)E_\mathrm{D}\left(\frac{N1}{2}\right)+\frac{e^2(N1)/2}{4\pi \epsilon \epsilon _0d}\text{}N\text{ odd.}$$
(50)
Consequently, the addition energies which are necessary to add a further electron are
$$\mathrm{\Delta }\mu _\mathrm{M}(N)=\frac{e^2}{4\pi \epsilon \epsilon _0d}\text{}N\text{ even,}$$
(51)
$$\mathrm{\Delta }\mu _\mathrm{M}(N)=\mu _\mathrm{D}\left(\frac{N+1}{2}\right)\mu _\mathrm{D}\left(\frac{N1}{2}\right)\text{}N3\text{ odd.}$$
(52)
These equations indicate that in a weakly coupled quantum dot with strongly localized wavefunctions the addition of an electron to a configuration with an even electron number is energetically more expensive as it corresponds to an addition of an electron to a single quantum dot. On the other hand, adding an electron to an odd number of particles is relatively cheap because the addition energy only depends on the interdot Coulomb repulsion which is relatively small for large $`d`$ and gives rise to the typical pairing of conductance peaks. With growing electron numbers this model fails as the increasing Coulomb repulsion caused by more extended charge distributions is not included. This effect destroys the splitting of the paired peaks.
A conceptional difficulty of this model is that it is based on localized states yielding an asymmetric charge distribution for odd electron numbers. In our SDFT calculations, however, the delocalized states also lead to the pairing of the conductance peaks in the weak-coupling regime (Fig. 9) while providing symmetric electron densities for all $`N`$. Therefore, we would like to emphasize, that the pairing of peaks does not depend on electron states localized at one center of the QDM, but can be caused by delocalized wavefunctions as well. In addition, delocalized states provide the advantage that it is not necessary to determine a transition from a coherent to an incoherent regime when the distance $`d`$ between the centers is continuously increased.
## VII Conclusions
We have studied the ground-state properties of QDMs containing 2, 4, and 8 electrons as a function of the distance $`d`$ between the centers of the molecule using SDFT. This concept includes spin effects in addition to the classical Coulomb repulsion. The results obtained for 4 and 8 electrons are reliable over the whole range of $`d`$ and can be predicted from the single-particle spectrum of the noninteracting system for not too weak confinement. For 4 electrons we find a transition from a spin-polarized to a spin-unpolarized configuration with decreasing coupling strength. The ground state of the 8-electron QDM can be characterized by a spin-polarized ground state (2 parallel spins) for strong coupling and a spin-polarized configuration with 4 aligned spins in the weak-coupling limit. In contrast, the two-electron QDM is well described by SDFT only in the strong-coupling regime whereas for large $`d`$ a HL approach yields better results. Finally we have discussed the Coulomb staircase diagram for the weak and the strong coupling regime and identified their characteristic features considering the background of delocalized states in an artificial molecule.
This work was funded by Deutsche Forschungsgemeinschaft (Grants No. SFB 348 and Ro 522/16).
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# An Eulerian-Lagrangian Approach for Incompressible Fluids: Local Theory
## 1 Introduction
The three dimensional Euler equations are evolution equations for the three velocity components $`u(x,t)`$,
$$\frac{u}{t}+uu+p=0,$$
(1)
coupled with a fourth equation, $`u=0`$, expressing incompressibility. In this Eulerian formulation the velocity $`u`$ and pressure $`p`$ are recorded at fixed locations $`x𝐑^3`$. The velocities and pressure vanish at infinity or are periodic. The pressure is determined using incompressibility. The equation is conservative and the total kinetic energy, $`|u|^2𝑑x`$ is a constant of motion.
The Euler equations can be studied in terms of the vorticity (). The vorticity is a vector $`\omega =\times u`$ corresponding to the anti-symmetric part of the gradient matrix $`u`$. It obeys a quadratic equation, whose nature is such that the magnitude of the vorticity may increase in time. If the amplification is not rapid enough then a well-known criterion () guarantees that no blow up can occur: if
$$_0^T\underset{x}{sup}|\omega (x,t)|dt<\mathrm{}$$
and the initial data are smooth, then the solution is smooth on the time interval $`[0,T]`$. The vorticity equation can be interpreted as the vanishing of a commutator
$$[D_t,\mathrm{\Omega }]=0$$
(2)
where
$$D_t=\frac{}{t}+u$$
is the material derivative and
$$\mathrm{\Omega }=\omega .$$
The characteristics of the first order differential operator $`\mathrm{\Omega }`$ are vortex lines; the characteristics of the material derivative $`D_t`$ are Lagrangian particle paths. The Lagrangian variables are the path maps $`aX(a,t)`$. The connection between the Lagrangian description and the Eulerian one is given by the relations
$$u(x,t)=\frac{X(a,t)}{t},x=X(a,t).$$
In this paper we discuss a description of the Euler equations as a system of three coupled active vector equations. The description concerns Lagrangian quantities computed in Eulerian variables and thus combines the physical significance of the Lagrangian description with the analytical advantages of the Eulerian description. The description bears similarities to the Clebsch variable representation. The Clebsch variables are a pair of scalars, $`\theta ,\phi `$ that are constant on particle paths and can be used to re-construct the velocity via
$$u^i(x,t)=\theta (x,t)\frac{\phi (x,t)}{x_i}\frac{n(x,t)}{x_i}.$$
This interesting representation is somewhat restrictive: not all solutions can be represented in this manner. That is because the Clebsch variables impose special constraints on helicity. Helicity is the scalar product of velocity and vorticity $`h=u\omega `$. Although $`h`$ itself is not conserved on particle paths, the integrals
$$_Th(x,t)𝑑x=c$$
are constants of motion, for any vortex tube $`T`$. A vortex tube $`T`$ is a time evolving region in space (not necessarily simply connected) whose boundary is at each point parallel to the vorticity, $`\omega \nu =0`$ where $`\nu `$ is the normal to $`T`$ at $`xT`$. The constants $`c`$ reflect the degree of topological complexity of the flow () and in general are non-trivial but they vanish identically for flows that admit a Clebsch variables representation. Indeed, for such flows the helicity is the divergence of a field that is parallel to the vorticity $`h=(n\omega )`$. Topological properties of streamlines and vortex tubes are relevant to hydrodynamic stability () and turbulence (, ). The description of the flow that allows for arbitrary vortex structures is based on formula (15) (, , , ) that was used for numerical computations. Somewhat related Hamiltonian formulations have been introduced by several authors (, , , ).
## 2 Eulerian-Lagrangian Description
The Lagrangian formulation of the Euler equations describes the flow in terms of a volume preserving diffeomorphism, the map $`aX(a,t)`$. The curve $`tX(a,t)`$ is the Lagrangian path at label $`a`$ and obeys Newton’s law
$$\frac{^2X(a,t)}{t^2}=F_X(a,t).$$
(3)
The incompressibility condition for the map is
$$det\left(_aX\right)=1.$$
(4)
The initial condition sets the labels at the initial time:
$$X(a,0)=a.$$
The forces $`F_X`$ in (3) are
$$F_X(a,t)=(_xp)(X(a,t))=\left[\left(_aX(a,t)\right)^{}\right]^1(_a\stackrel{~}{p})(a,t)$$
(5)
with $`\stackrel{~}{p}(a,t)=p(X(a,t))`$ and where $`p`$ is the Eulerian pressure. The notation $`M^{}`$ means the transposed of the matrix $`M`$, $`(M^{})^1`$ its inverse. Multiplying (3) by $`(_aX)^{}`$ we obtain
$$\frac{^2X(a,t)}{t^2}\left(_aX(a,t)\right)^{}=(_a\stackrel{~}{p})(a,t)$$
(6)
or, on components
$$\frac{^2X^j(a,t)}{t^2}\frac{X^j(a,t)}{a_i}=\frac{\stackrel{~}{p}(a,t)}{a_i}.$$
(7)
Pulling out a time derivative in the left-hand side we obtain
$$\frac{}{t}\left[\frac{X^j(a,t)}{t}\frac{X^j(a,t)}{a_i}\right]=\frac{\stackrel{~}{q}(a,t)}{a_i}$$
(8)
where
$$\stackrel{~}{q}(a,t)=\stackrel{~}{p}(a,t)\frac{1}{2}\left|\frac{X(a,t)}{t}\right|^2$$
(9)
We integrate (8) in time, fixing the label $`a`$:
$$\frac{X^j(a,t)}{t}\frac{X^j(a,t)}{a_i}=u_{(0)}^i(a)\frac{\stackrel{~}{n}(a,t)}{a_i}$$
(10)
where
$$\stackrel{~}{n}(a,t)=_0^t\stackrel{~}{q}(a,s)𝑑s$$
(11)
and
$$u_{(0)}(a)=\frac{X(a,0)}{t}$$
(12)
is the initial velocity. Note that $`\stackrel{~}{n}`$ has dimensions of circulation or of kinematic viscosity (length squared per time). The conservation of circulation
$$_\gamma \frac{X(\gamma ,t)}{t}𝑑\gamma =_\gamma \frac{X(\gamma ,0)}{t}𝑑\gamma $$
follows directly from the form (10). Let us consider
$$A(x,t)=X^1(x,t)$$
(13)
the “back-to-labels” map, and note that it forms a vector of active scalars (an active vector)
$$D_tA=\frac{A}{t}+uA=0.$$
(14)
Turning to (10), multiplying by $`\left[\left(_aX(a,t)\right)^{}\right]^1`$ and reading at $`a=A(x,t)`$ we obtain the formula
$$u^i(x,t)=\left(u_{(0)}^j(A(x,t))\right)\frac{A^j(x,t)}{x_i}\frac{n(x,t)}{x_i}$$
(15)
where
$$n(x,t)=\stackrel{~}{n}(A(x,t))$$
(16)
The equation (15) shows that the general Eulerian velocity can be written in a form that generalizes the Clebsch variable representation:
$$u=(A)^{}Bn$$
(17)
where $`B=u_{(0)}(A(x,t))`$ is also an active vector
$$D_tB=0.$$
(18)
Conversely, and somewhat more generally, if one is given a pair of active vectors $`A=(A^1(x,t),\mathrm{},A^M(x,t))`$ and $`B=(B^1(x,t),\mathrm{},B^M(x,t))`$ of arbitrary dimension $`M`$, such that the active vector equations (14) and (18) hold and if $`u`$ is given by
$$u(x,t)=\underset{k=1}{\overset{M}{}}B^k(x,t)_xA^k(x,t)_xn$$
(19)
with some function $`n`$, then it follows that $`u`$ solves the Euler equations
$$\frac{u}{t}+uu+\pi =0$$
where
$$\pi =D_tn+\frac{1}{2}|u|^2$$
Indeed, the only thing one needs is the kinematic commutation relation
$$D_t_xf=_xD_tf(_xu)^{}_xf$$
(20)
that holds for any scalar function $`f`$. The kinematic commutation relation (20) is a consequence of the chain rule, so it requires no assumption other than smoothness. Differentiating (19) and using the active vector equations (14, 18) it follows that
$$D_t(u)=\underset{k=1}{\overset{M}{}}((_xu)^{}_xA^k)B^k_x(D_tn)+(_xu)^{}n=$$
$$_x(D_tn)(_xu)^{}\left[\underset{k=1}{\overset{M}{}}(_xA^k)B^k_xn\right]=$$
$$_x(D_tn)(_xu)^{}u=_x(\pi ).$$
## 3 The Active Vector Formulation
The previous calculations can be summarized as follows: A function $`u(x,t)`$ solves the incompressible Euler equations if and only if it can be represented in the form $`u=u_A`$ with
$$u_A^i(x,t)=\varphi ^m\left(A(x,t)\right)\frac{A^m(x,t)}{x_i}\frac{n_A(x,t)}{x_i}$$
(21)
and
$$u_A=0$$
(22)
where $`A(x,t)`$ solves the active vector equation
$$\left(_t+u_A\right)A=0,$$
(23)
with initial data
$$A(x,0)=x.$$
The function $`\varphi `$ represents the initial velocity and the function $`n_A(x,t)`$ is determined up to additive constants by the requirement of incompressibility, $`u_A=0`$:
$$\mathrm{\Delta }n_A(x,t)=\frac{}{x_i}\left\{\varphi ^m(A(x,t))\frac{A^m(x,t)}{x_i}\right\}.$$
The periodic boundary conditions are
$$A(x+Le_j,t)=A(x,t)+Le_j;n_A(x+Le_j,t)=n_A(x,t)$$
(24)
with $`e_j`$ the standard basis in $`𝐑^3`$. In this case
$$\delta _A(x,t)=xA(x,t),$$
(25)
$`n_A(x,t)`$, and $`u_A(x,t)`$ are periodic functions in each spatial direction. One may consider also the case of decay at infinity, requiring that $`\delta _A`$, $`u_A`$ and $`n_A`$ vanish sufficiently fast at infinity. The equation of state (21, 22) can be written as
$$u_A=𝐏\left\{\varphi ^m\left(A(,t)\right)A^m(,t)\right\}=𝐏\left\{\left(A\right)^{}\varphi (A)\right\}$$
(26)
where
$$𝐏=\mathrm{𝟏}\mathrm{\Delta }^1$$
(27)
is the Leray-Hodge projector (with appropriate boundary conditions) on divergence free functions.
The Eulerian pressure is determined, up to additive constants by
$$p(x,t)=\frac{n_A(x,t)}{t}+u_A(x,t)n_A(x,t)+\frac{1}{2}|u_A(x,t)|^2.$$
The Jacobian obeys
$$det\left(A(x,t)\right)=1.$$
The vorticity
$$\omega _A(x,t)=\times u_A$$
satisfies the Helmholtz equation
$$D_t^A\omega _A=\omega _Au_A$$
(28)
and is given by the Cauchy formula
$$\omega _A(x,t)=\left[A(x,t)\right]^1\zeta (A(x,t))$$
(29)
where $`\zeta =\times \varphi `$ is the initial vorticity.
The advantage of an active vector formulation is that $`A`$ has conserved distribution, that is, for any function $`\mathrm{\Phi }`$
$$\mathrm{\Phi }(A(x,t))𝑑x=const;$$
in particular $`A(,t)_{L_{loc}^{\mathrm{}}(dx)}`$ is constant in time.
## 4 Local existence
The proof of local existence of solutions to the Euler equations in the active vector formulation is relatively simple and the result can be stated economically. The well-known local existence results in Lagrangian () and Eulerian () variables require more derivatives and thus use more restrictive function spaces.
###### Theorem 1
Let $`\varphi `$ be a divergence free $`C^{1,\mu }`$ periodic vector valued function of three variables. There exists a time interval $`[0,T]`$ and a unique $`C([0,T];C^{1,\mu })`$ spatially periodic vector valued function $`\delta (x,t)`$ such that
$$A(x,t)=x+\delta (x,t)$$
solves the active vector formulation of the Euler equations,
$$\frac{A}{t}+uA=0,$$
$$u=𝐏\left\{(A(x,t))^{}\varphi (A(x,t))\right\}$$
with initial datum $`A(x,0)=x`$.
The same result holds if one replaces periodic boundary conditions with decay at infinity. Differentiating the active vector equation (23) we obtain the equation obeyed by the gradients
$$D_t^A\left(\frac{A^m}{x_i}\right)+\frac{u_A^j}{x_i}\frac{A^m}{x_j}=0.$$
(30)
It is useful to denote
$$𝐏_{jl}=\delta _{jl}_j\mathrm{\Delta }^1_l$$
(31)
the matrix elements of the Leray-Hodge operator. Differentiating in the representation (26) and using the fundamental property
$$𝐏_{jl}\frac{f}{x_l}=0$$
we obtain
$$\frac{u_A^j}{x_i}=𝐏_{jl}\left(Det[\zeta (A);\frac{A}{x_i};\frac{A}{x_l}]\right).$$
(32)
Recall that the function $`\zeta `$ is the curl of $`\varphi `$. This relation shows that the gradient of velocity can be expressed without use of second order derivatives of $`A`$ and is the key to local existence: the equation (30) can be seen as a cubic quasi-local equation on characteristics. Let us make these ideas more precise. We will consider the periodic case first. We write $`C^{j,\mu }`$, $`j=0,1`$ to denote the Hölder spaces of real valued functions that are defined for all $`x𝐑^3`$ and are periodic with period $`L`$ in each direction. We denote by $`f_{0,\mu }`$ the $`C^{0,\mu }`$ norm:
$$f_{0,\mu }=\underset{x}{sup}|f(x)|+\underset{xy}{sup}\left\{|f(x)f(y)|\left(\frac{L}{|xy|}\right)^\mu \right\}$$
(33)
and by $`f_{1,\mu }`$ the $`C^{1,\mu }`$ norm:
$$f_{1,\mu }=f_{0,\mu }+Lf_{0,\mu }$$
(34)
where the notation $`|\mathrm{}|`$ refers to modulus, Euclidean norm, and Euclidean norm for matrices, as appropriate.
We break the solution of the problem in two parts, the map $`\delta u`$ and the map $`u\delta `$. We denote the first one $`W`$.
$$W[\delta ,\varphi ](x,t)=𝐏\left\{(𝐈+\delta (x,t))^{}\varphi (x+\delta (x,t))\right\}$$
(35)
This map is linear in $`\varphi `$ but nonlinear in $`\delta `$.
###### Proposition 1
The map $`W[\delta ,\varphi ]`$ maps
$$W:(C^{1,\mu })^3\times (C^{1,\mu })^3(C^{1,\mu })^3$$
continuously. There exist constants $`C`$ depending on $`\mu `$ alone so that
$$W[\delta ,\varphi ]_{0,\mu }C\varphi _{0,\mu }\{1+\delta _{0,\mu }\}^2$$
and
$$W[\delta ,\varphi ]_{0,\mu }C\times \varphi _{0,\mu }\{1+\delta _{0,\mu }\}^3.$$
hold for any $`\delta \left(C^{1,\mu }\right)^3`$, $`\varphi \left(C^{1,\mu }\right)^3`$.
For the proof we note that $`W`$ is made up from a number of operations. The first operation is the composition $`\varphi (x)\varphi (x+\delta (x))`$. For a fixed $`\delta (C^{1,\mu })^3`$ the map $`xx+\delta `$ is Lipschitz. Composition with a Lipschitz change of variables maps $`C^{0,\mu }`$ into itself continuously (we say that it is a continuous endomorphism). The joint continuity of $`[\varphi ,\delta ]\varphi (x+\delta )`$ in $`C^{1,\mu }`$ follows naturally. The second operation is a sum of products of functions (a matrix applied to a vector). This is a continuous operation because the Hölder spaces $`C^{j,\mu }`$, $`j=0,1`$ we chose are Banach algebras. The third and last operation is the linear operator $`𝐏`$, which is bounded in Hölder spaces. We need to consider also derivatives of $`W`$. We use the formula (32) and note that the expression for the gradient is made of similar operations as above and apply the same kind of reasoning. This finishes the proof.
Time does not play any role in this proposition because the equation of state $`(\delta ,\varphi )W[\delta ,\varphi ]`$ is time independent. The second half of the procedure does depend on time. Let us denote by $`\mathrm{\Theta }`$ the map that associates to two continuous paths $`t\delta (,t)`$ and $`t\varphi (,t)`$ a new path $`t\theta `$; the path $`t\theta =\mathrm{\Theta }[\delta ,\varphi ]`$ is obtained by solving the partial differential equation
$$\frac{\theta }{t}+u\theta +u=0$$
(36)
where
$$u=W[\delta (,t),\varphi (,t)],$$
periodic boundary conditions are imposed on $`\theta `$ and zero initial data
$$\theta (x,0)=0$$
are required. The Euler equation requires only the use of a time independent $`\varphi `$, but allowing time dependent $`\varphi `$ is very useful: one can thus treat more equations, in particular the Navier-Stokes equation. Let us consider the space
$$𝒫_T=C([0,T],(C^{1,\mu })^3)$$
of continuous $`(C^{1,\mu })^3`$ -valued paths defined on a time interval $`[0,T]`$, endowed with the natural norm
$$\theta _{1,𝒫}=\underset{t}{sup}\theta (,t)_{1,\mu }.$$
We will consider also the weaker norm
$$\theta _{0,𝒫}=\underset{t}{sup}\theta (,t)_{0,\mu }.$$
$`\mathrm{\Theta }`$ is nonlinear in both arguments.
###### Proposition 2
The map $`\mathrm{\Theta }[\delta ,\varphi ]`$ maps
$$\mathrm{\Theta }:𝒫_T\times 𝒫_T𝒫_T$$
and is continuous when the topology of the source space $`𝒫_T\times 𝒫_T`$ is the natural product $`C^{1,\mu }`$ topology and the topology of the target space $`𝒫_T`$ is the weaker $`C^{0,\mu }`$ topology. Moreover, there exists a constant $`C`$ depending on $`\mu `$ alone so that
$$\theta (,t)_{0,\mu }\left(_0^tu(,s)_{0,\mu }𝑑s\right)\left\{\mathrm{exp}\{C_0^tu(,s)_{0,\mu }𝑑s\}\right\}$$
holds for each $`tT`$ with $`u=W[\delta ,\varphi ]`$ and $`\theta =\mathrm{\Theta }[\delta ,\varphi ]`$.
Proposition 2 states that the map $`\mathrm{\Theta }`$ is bounded but not that it is continuous in the strong $`C^{1,\mu }`$ topology. The proof follows naturally from the idea to use the classical method of characteristics and ODE Gronwall type arguments. Similar ideas are needed below in the the slightly more difficult proof of Proposition 3 and we will sketch them there and therefore we leave the details of the proof of Proposition 2 to the interested reader.
In order to proceed let us take now a fixed $`\varphi `$, take a small number $`ϵ>0`$ and associate to it the set
$$𝒫_T$$
defined by
$$=\{\delta (x,t);\delta (x,0)=0,\delta (,t)_{0,\mu }ϵ,tT\}.$$
Combining the bounds in the two previous propositions one can choose, for fixed $`\varphi `$, a $`T`$ small enough so that
$$\delta \mathrm{\Theta }[\delta ,\varphi ]=𝒮[\delta ]$$
maps
$$𝒮:.$$
Inspecting the bounds it is clear that it is sufficient to require
$$T\times \varphi _{0,\mu }cϵ$$
with an appropriate $`c`$ depending on $`\mu `$ alone. Leaving $`\varphi `$, $`ϵ`$ and $`T`$ fixed as above, the map $`𝒮`$ is Lipschitz in the weaker norm $`C^{0,\mu }`$:
###### Proposition 3
There exists a constant $`C`$, depending on $`\mu `$ alone, such that, for every $`\delta _1,\delta _2`$, the Lipschitz bound
$$𝒮[\delta _1]𝒮[\delta _1]_{0,𝒫}C\delta _1\delta _2_{0,𝒫}$$
holds.
It is essential that $`\delta _j`$, so that they are smooth and their gradients are small, but nevertheless this is a nontrivial statement. An inequality of the type
$$𝒮[\delta _1]𝒮[\delta _1]_{0,𝒫}C\delta _1\delta _2_{1,𝒫}$$
is easier to obtain, but loses one derivative. This kind of loss of one derivative is a well-known difficulty in general compressible hyperbolic conservation laws. The situation is complicated in addition by the fact that the constitutive law $`W`$ depends on gradients. As we shall see, incompressibility saves one derivative. The heart of the matter is
###### Proposition 4
Let $`\varphi (C^{1,\mu })^3`$ be fixed. There exists a constant depending on $`\mu `$ alone so that
$$W[\delta _1,\varphi ]W[\delta _2,\varphi ]_{0\mu }C\delta _1\delta _2_{0,\mu }\varphi _{1,\mu }$$
holds for any $`\delta _jC^{1,\mu }`$ with $`\delta _j_{1,\mu }1`$.
One could use the condition $`\delta _jC^{1,\mu }`$ with $`\delta _j_{1,\mu }M`$ but then $`C`$ would depend on $`M`$ also.
Proof of Proposition 4. Denoting
$$u=W[\delta _1,\varphi ]W[\delta _2,\varphi ],$$
$$\delta =\delta _1\delta _2,$$
$$\psi (x)=\frac{1}{2}\left(\varphi (x+\delta _1(x))+\varphi (x+\delta _2(x))\right),$$
$$v(x)=\varphi (x+\delta _1(x))\varphi (x+\delta _2(x)),$$
$$\gamma =\frac{1}{2}(\delta _1+\delta _2)$$
we write
$$u=u_1+u_2$$
with
$$u_1=𝐏\left\{(\delta )^{}\psi \right\}$$
and
$$u_2=𝐏\left\{(𝐈+\gamma )^{}v\right\}$$
Now the bound
$$u_2_{0,\mu }C\delta _{0,\mu }\varphi _{1,\mu }$$
is obtained in the same way as the bound in Proposition 1. (Actually $`\varphi `$ Lipschitz is enough here.) The dangerous term is $`u_1`$ because it contains $`\delta `$. But here we can “integrate by parts” and write
$$u_1=𝐏\left\{(\psi )^{}\delta \right\}$$
because of incompressibility. The matrix $`\psi `$ is bounded in $`C^{0,\mu }`$ and the bound follows again easily, as the bounds in Proposition 1. This ends the proof of Proposition 4. We draw the attention to the fact that the presence of the $``$ (transpose) operation is essential for the “integration by parts” to be allowed.
Returning to the proof of Proposition 3 we denote $`\theta _j=𝒮\delta _j`$, $`u_j=W(\delta _j,\varphi )`$, $`u=u_1u_2`$, $`\theta =\theta _1\theta _2`$ and write
$$\frac{\theta }{t}+\frac{u_1+u_2}{2}\theta +u\left(\frac{\theta _1+\theta _2}{2}\right)+u=0$$
We consider the characteristics $`X(a,t)`$ defined by
$$\frac{dX}{dt}=\frac{u_1+u_2}{2}(X,t),X((a,0)=a$$
and note that in view of Proposition 1 and the assumption $`\delta _j`$, the characteristics are well defined for $`0tT`$, their inverse $`A(x,t)=X^1(x,t)`$ (the “back-to-labels” map) is defined too. Moreover,
$$\underset{t,a}{sup}\left|\frac{X}{a}\right|C$$
and
$$\underset{t,x}{sup}\left|\frac{A}{x}\right|C$$
holds with a constant $`C`$ depending on $`\mu `$ alone. Consider now the function
$$F(x,t)=u\left(\frac{\theta _1+\theta _2}{2}\right)+u.$$
Solving by the method of characteristics we obtain
$$\theta (x,t)=_0^tF(X(A(x,t),s),s)𝑑s.$$
Using Proposition 4 in conjunction with the bounds in Propositions 1 and 2 we see that $`F(x,t)`$ is bounded (uniformly in time) in $`C^{0,\mu }`$:
$$\underset{t}{sup}F(,t)_{0,\mu }C\varphi _{1,\mu }\delta _{0,𝒫}$$
Compositions with the uniformly Lipschitz $`X`$ and $`A`$ are harmless and we obtain the desired result
$$\theta _{0,𝒫}C\delta _{0,𝒫}.$$
This ends the proof of Proposition 3. The proof of Theorem 1 follows now using successive approximations. Starting with a first guess $`\delta _1`$ we define inductively
$$\delta _{n+1}=𝒮\delta _n.$$
Proposition 3 implies that the sequence $`\delta _n`$ converges rapidly in the $`C^{0,\mu }`$ topology to a limit $`\delta `$. Because $``$ is convex it contains this weaker limit point, $`\delta `$. Because $`𝒮`$ has the weak Lipschitz property of Proposition 3 it follows that $`𝒮\delta =\delta `$. This actually means that $`A=x+\delta (x,t)`$ solves the active vector formulation of the Euler equations and that $`u=W[\delta ,\varphi ]`$ solves the usual Eulerian formulation.
Now let us consider the case of decay at infinity. This case is instructive to look at this case because it illuminates the difference between $`\varphi ,u,W`$ on the one hand and $`x,\delta ,\mathrm{\Theta }`$ on the other hand; the function spaces need to be modified in a natural fashion to accommodate this difference. The issue of decay at infinity is both a physical one – the total kinetic energy must be defined, and a mathematical one – $`𝐏`$ must be defined. But apart from this, the decay at infinity requirement does not hinder the proof in any respect.
###### Theorem 2
Let $`\varphi `$ be a $`C^{1,\mu }`$ velocity that is square integrable
$$|\varphi (x)|^2𝑑x<\mathrm{}$$
and whose curl is integrable to some power $`1<q<\mathrm{}`$,
$$|\times \varphi (x)|^q𝑑x<\mathrm{}.$$
Then for $`ϵ`$ sufficiently small there exists a time interval $`[0,T]`$ and a $`C^{1,\mu }`$ function $`\delta (x,t)`$ such that
$$\underset{t}{sup}\delta (,t)_{0,\mu }ϵ$$
and such that $`x+\delta (x,t)`$ solves the active vector formulation of the Euler equation. The velocity corresponding to this solution belongs to $`C^{1,\mu }`$, is square integrable and the vorticity is integrable to power $`q`$.
The proof follows along the same lines as above. Because $`\varphi `$ enters linearly in the expression for $`W`$ and because we control $`\delta `$ uniformly, issues of decay at infinity of do not arise. In other words, the function space for velocities does not need to be a Banach algebra, rather a module over the Banach algebra of the $`\delta `$ variables, which need not decay at infinity.
## 5 The blow up issue
Any solution of the Euler equation can be constructed using a sequence of near identity transformations. One starts out with
$$\varphi =u_0$$
and solves for an interval of time $`0tt_1`$ the active vector equation
$$\frac{A}{t}+u_AA=0$$
$$u_A=𝐏\left((A)^{}u_0(A)\right)$$
$$A(x,0)=x.$$
At time $`t=t_1`$ one resets:
$$\varphi =u_1=u_A(,t_1)$$
and solves the system above again, for a new time interval $`t_1tt_2`$ and so one continues the solution. The local existence result guarantees that
$$(t_{n+1}t_n)u_n_{0,\mu }c>0$$
and during this time the solution $`A(x,t)`$ remains close to the identity in the sense that $`\delta =Ax`$ obeys
$$\delta (,t)_{0,\mu }ϵ$$
with a prescribed $`ϵ<<1`$. The formula 32 implies then that
$$u_n_{0,\mu }K^nu_0_{0,\mu }$$
with a fixed $`K>1`$. If the inequalities above would be sharp then, of course, the time steps would have to decrease exponentially and the procedure would diverge in finite time. It is possible that for certain initial data the bounds may be overly pessimistic and the solution may exist for a long time. But with the present knowledge, if one desires long-lived solutions for arbitrary three dimensional data then one needs to smooth either at the end of each step or during each time step. If one applies a smoothing procedure one evidently changes the problem and one introduces an artificial dissipation. There are many ways one could conceivable regularize the Euler equations. The physically correct energy dissipating equation is the Navier-Stokes equation. Unfortunately it is not known in three dimensions if the Navier-Stokes equations have globally defined unique solutions that converge to solutions of the Euler equations. Even in two dimensions, where the existence of smooth solutions is known for both the Euler and Navier-Stokes equations, the situation is not entirely trivial (, ). The two dimensional situation is characterized by the absence of vortex stretching. In the case of the three dimensional Euler equations the vorticity magnitude evolves according to the stretching equation
$$D_t\left(|\omega |\right)=\alpha |\omega |.$$
(37)
The stretching factor $`\alpha `$ is related to the vorticity magnitude through a principal value singular integral ():
$$\alpha (x,t)=P.V.D(\widehat{y},\xi (x,t),\xi (x+y,t))|\omega (x+y,t)|\frac{dy}{|y|^3}.$$
(38)
Here $`\widehat{y}`$ is the unit vector in the direction of $`y`$, $`\xi (x,t)=\frac{\omega }{|\omega |}`$ is the unit vector tangent to the vortex line passing through $`x`$ at time $`t`$ and $`D`$ is a certain geometric factor. The geometric factor is a smooth function of three unit vectors, has zero average on the unit sphere, $`D𝑑S(\widehat{y})=0`$ and vanishes pointwise when $`\xi (x,t)=\pm \xi (x+y,t)`$. Because $`\alpha `$ has the same order of magnitude as $`|\omega |`$, dimensional reasoning suggests blow up of the type one encounters in the ordinary differential equation $`\frac{dm}{dt}=m^2`$,
$$\underset{x}{sup}|\omega (x,t)|\frac{1}{Tt}.$$
But if the vorticity direction $`\xi `$ is smooth then a geometric depletion of $`\alpha `$ occurs; that means that $`\alpha `$ is of the order of magnitude of velocity times the magnitude of the spatial gradient of $`\xi `$ (an inverse length scale, assumed to be finite). The two dimensional Euler equations correspond to the case $`\xi =(0,0,1)`$ and $`\alpha =0`$ identically. If
$$_0^T\underset{x}{sup}|\alpha (x,t)|dt<\mathrm{}$$
then no blow up can occur. This idea of geometric depletion of nonlinearity has been investigated theoretically and numerically for the Euler equations and for a quasi-geostrophic active scalar equation (, , , , , ). In the Eulerian-Lagrangian formulation of the Euler equation the role played by smooth stratifications can be explained in the following manner. Consider functions $`w=w_\psi `$ of the form
$$w_\psi (x,t)=(A(x,t))^{}\psi ((A(x,t))$$
(39)
associated to arbitrary vectors $`\psi `$. Alternately, one might consider solutions of
$$D_t^Aw+\left(u_A\right)^{}w=0$$
(40)
with initial data $`\psi `$. A particular example is provided by choosing $`\psi =\varphi `$, i.e. $`w_\varphi =w_A`$
$$w_A=(A)^{}\varphi (A)$$
(41)
which obeys
$$\times w_A=\times u_A=\omega _A.$$
(42)
Because the vorticity
$$\omega _A=\times u_A$$
(43)
satisfies
$$\omega _A=(A)^1\zeta (A)$$
(44)
it follows that
$$\omega _A(x,t)w_\psi (x,t)=\zeta (A(x,t))\psi (A(x,t)),$$
(45)
holds for any $`\psi `$ or, in other words
$$D_t(\omega _Aw)=0$$
(46)
holds for any solution $`w(x,t)`$ of (40). Global regularity of a solution of the Euler equations would follow from (45) if one could find a sufficient family of vectors $`\psi `$. By a sufficient family for the initial velocity $`\varphi `$ and the time interval $`[0,T]`$ we mean a family of vectors $`\psi `$ such that there exists a non-negative function $`\gamma (t)`$ with $`_0^T\gamma 𝑑t<\mathrm{}`$ such that
$$|\omega _A(x,t)|\gamma (t)\underset{\psi }{sup}|\omega _A(x,t)w(x,t)|$$
holds for $`0tT`$. A sufficient family for all two dimensional flows is provided by just one $`\psi `$, $`\psi =(0,0,1)`$ with $`\gamma =1`$. Generalizations would consist of situations in which one could find sufficient families that depend on the initial data and time and take locally the role played in 2D by the vertical direction.
The blow up issue becomes, in terms of $`A`$, a question of formation of infinite gradients in conserved quantities. This is similar to the case of hyperbolic conservation laws but with the significant difference that the underlying characteristic flow is volume-preserving: $`det(A)=1`$, the matrix $`A`$ is invertible and
$$(\left(A(x,t)\right)^1)_{ij}=\frac{1}{2}ϵ_{imn}Det[e_j;\frac{A}{x_m};\frac{A}{x_n}]$$
(47)
holds, where $`e_j=(\delta _{jk})`$ is the canonical basis in $`𝐑^3`$. Consider the Euler-Lagrange label differentiation
$$L_j^A=\frac{1}{2}\left(ϵ_{imn}ϵ_{jkl}\frac{A_k}{x_m}\frac{A_l}{x_n}\right)\frac{}{x_i}$$
(48)
From the commutation relation (20) and the A equation (23) it follows that
$$[D_t^A,L_j^A]=0$$
(49)
holds for any $`j=1,2,3`$. This commutation relation simply says that in Lagrangian coordinates, time and label derivatives commute. Note, from the formulas (44) and (47) that
$$\frac{1}{2}ϵ_{pil}Det[\zeta (A);\frac{A}{x_i};\frac{A}{x_l}]=\omega _A^p.$$
(50)
It is clear now that $`D_t^A`$ commutes with $`\mathrm{\Omega }_A=\omega _A`$ because it is represented in terms of $`L_j^A`$:
$$\mathrm{\Omega }_A=\zeta _j(A)L_j^A.$$
(51)
Observe that, in view of the definition of the operators $`L_j^A`$,
$$L_j^A=\left(\left(A(x,t)\right)^1\right)_{kj}\frac{}{x_k}$$
(52)
it follows that
$$\left(\left(A(x,t)\right)^1\right)_{ij}=L_j^A[x_i];$$
(53)
on the other hand
$$D_t^A(x_i)=u_A^i$$
(54)
holds, so from the commutation relation (49) we obtain
$$D_t^A\left(\left(A(x,t)\right)^1\right)_{ij}=L_j^A(u_A^i).$$
(55)
This equation, which could have been derived also directly from (30), implies the vorticity equation because of (29):
$$D_t^A\omega _A=\zeta _j(A)L_j^A(u_A)=\mathrm{\Omega }_A(u_A).$$
(56)
Because of the result in () and (29), it is clear that the finiteness of
$$_0^T\left(A(,t)\right)^1_{L^{\mathrm{}}(dx)}𝑑t$$
implies regularity. Or, using (47), we deduce that the finiteness of
$$_0^TA(,t)_{L^{\mathrm{}}(dx)}^2𝑑t$$
implies regularity. Let us introduce now the matrix
$$C_{ij}^A(x,t;z)=\left(A(x+z,t)\right)_{im}\left(\left(A(x,t)\right)^1\right)_{mj}.$$
(57)
and call it the Euler-Lagrange calibrator. The formula
$$C_{ij}^A(x,t;z)=L_j^{A(x,t)}\left(A_i(x+z,t)\right)$$
(58)
shows that the calibrator measures the response of the Eulerian translation to an infinitesimal Lagrangian translation. Note that
$$C_{ij}^A(x,t;0)=\delta _{ij}.$$
(59)
The calibrator is a quotient of gradients at different locations and therefore locally spatially uniform, temporally arbitrary changes like dilations do not affect it. The vorticity equation (56) can be expressed in terms of the Euler-Lagrange calibrator ():
$$D_t^A\omega _A^i=\{\frac{1}{4\pi }P.V.D(\zeta ,C^A\zeta ,C_{.,p}^A)\sigma _{il}(\widehat{z})\frac{dz}{|z|^3}\}\frac{A_p(x,t)}{x_l}$$
(60)
Note that
$$\omega _A^i\frac{A_p(x,t)}{x_i}=\zeta _p(A(x,t))$$
is bounded. It is therefore natural to conjecture that the smoothness of $`C^A`$ prevents finite time blow up for the Euler equations. This conjecture is true for the quasi-geostrophic active scalar. The interested reader is referred to () for details.
The blow up question for the Euler equations remains open. Numerical calculations provide insight and hints, but the answer will have to be analytical. The considerations above point towards a possible incompressible dispersive effect that hinders blow up: as the gradients of $`A`$ become large the resulting rapid (15) and non-uniform (55) motion disperses the large gradients. This might cause instability of blow up or perhaps its suppression.
Acknowledgments. I thank Diego Cordoba, Charles Fefferman and Julian Hunt for helpful comments. This research was supported in part by NSF- DMS9802611. Partial support of AIM and the hospitality of the Princeton Mathematics Department are gratefully acknowledged.
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# The Askey-Wilson function transform
## 1. Introduction.
In several explicit function transforms are realized as Fourier transforms on the non-compact quantum $`SU(1,1)`$ group. In this paper we focus on the function theoretic aspects of the most general Fourier transform on the quantum $`SU(1,1)`$ group. Since the corresponding radial part of the Casimir operator is Askey’s and Wilson’s second order $`q`$-difference operator $`L`$, we are led to the task of studying its spectral properties as a symmetric operator with respect to an explicit measure of unbounded support. This measure naturally arises in as a Haar weight on the quantum $`SU(1,1)`$ group. The kernel of the corresponding generalized Fourier transform is an eigenfunction of $`L`$, and it provides an analytic continuation of the Askey-Wilson polynomial in its degree. Consequently, we name the kernel the Askey-Wilson function.
The Askey-Wilson function satisfies a beautiful symmetry property, which we refer to as duality in this paper. It essentially states that the geometric and the spectral parameter of the Askey-Wilson function are interchangeable (up to a certain involution on the parameters).
For a full exploitation of duality for the development of the $`L^2`$-theory of the corresponding generalized Fourier transform, it is necessary to add another degree of freedom. This can be done in the following natural way. The Askey-Wilson function can be expanded as a linear combination of the asymptotically free solutions of $`L`$ on an arbitrary $`q`$-lattice. The corresponding explicit coefficients, the so-called $`c`$-functions, can then be used to define an explicit measure whose support around infinity lies on the chosen $`q`$-lattice. This gives rise to a one parameter family of measures of unbounded support, in which the extra parameter labels the different $`q`$-lattices. It contains the measure arising from the quantum $`SU(1,1)`$ group as a special case.
The generalized Fourier transform with respect to this one parameter family of measures, and with the Askey-Wilson function as its kernel, is called the Askey-Wilson function transform. In view of its harmonic analytic interpretation on the quantum $`SU(1,1)`$ group, we may think of the Askey-Wilson function transform as a generalization of the classical Jacobi transform, whence in particular of the Mellin-Fock (or Legendre) transform. In fact, it was shown in that the Askey-Wilson function transform is on top of an hierarchy consisting of several different generalizations of the Jacobi and Hankel transform. In particular, the $`L^2`$-theory for the Askey-Wilson function transform unifies the $`L^2`$-theory for the spherical Fourier transform on the (quantum) $`SU(1,1)`$ group as well as on the (quantum) group of plane motions.
The main objective of the paper is to show that the Askey-Wilson function transform extends to an isometric isomorphism in a natural way, and that its inverse is the Askey-Wilson function transform with respect to a different, “dual” choice of parameters.
We expect that the present theory on the Askey-Wilson function transform will turn out to be the special rank one case of a general theory on Macdonald function transforms associated with root systems, in analogy with the polynomial setting (see e.g. , ). From this point of view, it is natural to expect that the structure of the Askey-Wilson function transform, and eventually also the structure of the general Macdonald function transform, will be best understood from an extension of Cherednik’s affine Hecke algebraic approach to the non-polynomial setting. A first indication of the pivotal role of affine Hecke algebras in the structure of the Askey-Wilson function transform is the duality, which in the polynomial setting stems from an involution of the associated double affine Hecke algebra (see for instance , and , for the rank one setting). We hope that future research will shed more light on these matters.
The set-up of the paper is as follows. In §2 we recall the Askey-Wilson polynomials and in §3 we define the Askey-Wilson functions. We emphasize in these sections the important concept of duality. In §4 the asymptotically free solutions of $`L`$ are discussed, as well as the corresponding $`c`$-function expansion of the Askey-Wilson function. We formulate our main result in §5, which states that the Askey-Wilson function transform extends to an isometric isomorphism in a natural way, and that its inverse is given by the Askey-Wilson function transform with respect to dual parameters.
The remaining sections occupy the proof of the main result. Duality plays a pivotal and simplifying role in the proof, since it allows us to avoid the machinery of spectral analysis of (unbounded) self-adjoint operators (in contrast with the analysis in some of its degenerate cases, see e.g. , and ). The other main ingredient of the proof is the computation of the weak limit of the Wronskian of the Askey-Wilson functions. As an interesting by-product of this computation, we obtain in §7 explicit orthogonality relations for “low degree” Askey-Wilson polynomials with respect to the one-parameter family of measures of unbounded support.
Notations and conventions: We assume that $`0<q<1`$ is fixed throughout the paper. We use the standard notations for basic hypergeometric series, see for instance . In particular, we write $`(a_1,\mathrm{},a_r;q)_c=_{i=1}^r(a_i;q)_c`$ with $`(a;q)_c=(a;q)_{\mathrm{}}/(aq^c;q)_{\mathrm{}}`$ and $`(a;q)_{\mathrm{}}=_{i=0}^{\mathrm{}}(1aq^i)`$ for (products) of $`q`$-shifted factorials, and we write
$${}_{r+1}{}^{}\varphi _{r}^{}(\begin{array}{c}a_1,a_2,\mathrm{},a_{r+1}\\ b_1,b_2,\mathrm{},b_r\end{array};q,z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a_1,a_2,\mathrm{},a_{r+1};q)_k}{(q,b_1,\mathrm{},b_r;q)_k}z^k$$
for the $`{}_{r+1}{}^{}\varphi _{r}^{}`$ basic hypergeometric series. The very well poised $`{}_{8}{}^{}\varphi _{7}^{}`$ is defined by
$${}_{8}{}^{}W_{7}^{}(a;b,c,d,e,f;q,z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{1aq^{2k}}{1a}\frac{(a,b,c,d,e,f;q)_kz^k}{(q,qa/b,qa/c,qa/d,qa/e,qa/f;q)_k}.$$
We use the branch of the square root $`\sqrt{}`$ which is positive on $`_{>0}`$ throughout the paper.
Acknowledgements: The second author is supported by a fellowship from the Royal Netherlands Academy of Arts and Sciences (KNAW). Part of the research was done during the second author’s stay at Université Pierre et Marie Curie (Paris VI) and Institut de Recherche Mathématique Avancée (Strasbourg) in France, supported by a NWO-TALENT stipendium of the Netherlands Organization for Scientific Research (NWO) and by the EC TMR network “Algebraic Lie Representations”, grant no. ERB FMRX-CT97-0100.
## 2. The Askey-Wilson polynomials.
In this section we briefly recall the basic properties of the Askey-Wilson polynomials. We formulate these properties in such a way that the connections with the non-polynomial setting (see section 5) will be as transparent as possible. In particular, we emphasize the important concept of duality.
The Askey-Wilson polynomials $`p_n(x)=p_n(x;a,b,c,d;q)`$ ($`n_+`$) are defined by the series expansion
$$p_n(x)={}_{4}{}^{}\varphi _{3}^{}(\begin{array}{c}q^n,q^{n1}abcd,ax,ax^1\\ ab,ac,ad\end{array};q,q),$$
(2.1)
see . If we associate dual parameters $`\stackrel{~}{a}`$, $`\stackrel{~}{b}`$, $`\stackrel{~}{c}`$ and $`\stackrel{~}{d}`$ with $`a,b,c`$ and $`d`$ by the formulas
$$\begin{array}{cc}\hfill \stackrel{~}{a}& =\sqrt{q^1abcd},\stackrel{~}{b}=ab/\stackrel{~}{a}=q\stackrel{~}{a}/cd,\hfill \\ \hfill \stackrel{~}{c}& =ac/\stackrel{~}{a}=q\stackrel{~}{a}/bd,\stackrel{~}{d}=ad/\stackrel{~}{a}=q\stackrel{~}{a}/bc,\hfill \end{array}$$
(2.2)
then it immediately follows from the explicit expression (2.1) that the Askey-Wilson polynomials satisfy the duality relation
$$p_n(aq^m;a,b,c,d;q)=p_m(\stackrel{~}{a}q^n;\stackrel{~}{a},\stackrel{~}{b},\stackrel{~}{c},\stackrel{~}{d};q),m,n_+.$$
(2.3)
The deeper understanding of duality stems from affine Hecke algebraic considerations, see .
The Askey-Wilson polynomials $`\{p_n\}_{n_+}`$ form a basis of the polynomial algebra $`[x+x^1]`$ consisting of eigenfunctions of the Askey-Wilson second order $`q`$-difference operator
$$\begin{array}{cc}\hfill L& =\alpha (x)(T_q1)+\alpha (x^1)(T_q^11),\hfill \\ \hfill \alpha (x)& =\frac{(1ax)(1bx)(1cx)(1dx)}{(1x^2)(1qx^2)},\hfill \end{array}$$
(2.4)
where $`(T_q^{\pm 1}f)(x)=f(q^{\pm 1}x)`$. The eigenvalue of $`L`$ corresponding to the Askey-Wilson polynomial $`p_n`$ is $`\mu (\gamma _n)`$, where $`\gamma _n=\stackrel{~}{a}q^n`$ and
$$\mu (\gamma )=1\stackrel{~}{a}^2+\stackrel{~}{a}(\gamma +\gamma ^1).$$
(2.5)
We remark that the three term recurrence relations for the Askey-Wilson polynomials $`p_m(;\stackrel{~}{a},\stackrel{~}{b},\stackrel{~}{c},\stackrel{~}{d};q)`$ ($`m_+`$) can be derived from the eigenvalue equations $`Lp_n=\mu (\gamma _n)p_n`$ ($`n_+`$) by applying the duality (2.3), see e.g. and .
For the remainder of this section, we restrict our attention to the case that the parameters $`a,b,c`$ and $`d`$ are positive and less than one. Then Askey and Wilson proved the orthogonality relations
$$\frac{1}{2\pi iC_0}_{x𝕋}p_n(x)p_m(x)\mathrm{\Delta }(x)\frac{dx}{x}=\delta _{m,n}\frac{\underset{x=\gamma _0}{\text{Res}}\left(\frac{\stackrel{~}{\mathrm{\Delta }}(x)}{x}\right)}{\underset{x=\gamma _n}{\text{Res}}\left(\frac{\stackrel{~}{\mathrm{\Delta }}(x)}{x}\right)}$$
(2.6)
where $`\delta _{m,n}`$ is the Kronecker delta and $`𝕋`$ is the counterclockwise oriented unit circle in the complex plane, with the weight function given by
$$\mathrm{\Delta }(x)=\frac{(x^2,1/x^2;q)_{\mathrm{}}}{(ax,a/x,bx,b/x,cx,c/x,dx,d/x;q)_{\mathrm{}}},$$
(2.7)
and with $`\stackrel{~}{\mathrm{\Delta }}(x)`$ the weight function $`\mathrm{\Delta }(x)`$ with respect to dual parameters. Here the positive normalization constant $`C_0`$ is given by the Askey-Wilson integral
$$C_0=\frac{1}{2\pi i}_{x𝕋}\mathrm{\Delta }(x)\frac{dx}{x}=\frac{2(abcd;q)_{\mathrm{}}}{(q,ab,ac,ad,bc,bd,cd;q)_{\mathrm{}}}.$$
The orthogonality relations written in the form (2.6) exhibit the duality (2.3) of the Askey-Wilson polynomials on the level of the orthogonality relations, since it expresses the quadratic norms explicitly in terms of the dual weight function $`\stackrel{~}{\mathrm{\Delta }}`$. This description of the quadratic norms was proved in using affine Hecke algebraic techniques.
The above results describe the spectral properties of the Askey-Wilson second order $`q`$-difference operator $`L`$ as a unbounded symmetric operator on the weighted $`L^2`$-space with respect to the measure $`\mathrm{\Delta }(x)\frac{dx}{x}`$ on the unit circle $`𝕋`$. For suitable parameter values, this is directly related to harmonic analysis on the compact quantum $`SU(2)`$ group, see e.g. .
In this paper we study the spectral analysis of $`L`$ as a symmetric operator with respect to a one-parameter family of measures of unbounded support. This setting naturally arises in the study of harmonic analysis on the non-compact quantum $`SU(1,1)`$ group, see . Duality, which we have emphasized in this section on the polynomial level, will play a crucial role in the development of the corresponding $`L^2`$-theory.
## 3. The Askey-Wilson function.
In this section we consider a solution of the eigenvalue equation
$$\left(Lf\right)(x)=\mu (\gamma )f(x),$$
(3.1)
which reduces to the Askey-Wilson polynomial for $`\gamma =\gamma _n`$ ($`n_+`$), and which enjoys the same symmetry properties as the Askey-Wilson polynomials.
Two linearly independent solutions of the eigenvalue equation (3.1) can be derived from Ismail’s and Rahman’s \[5, (1.11)–(1.16)\] solutions for the three term recurrence relation of the associated Askey-Wilson polynomials. The solutions are given in terms of very well poised $`{}_{8}{}^{}\varphi _{7}^{}`$ series. In this section we consider the solution $`\varphi _\gamma (x)=\varphi _\gamma (x;a;b,c;d;q)`$ of (3.1) given by
$$\begin{array}{cc}\hfill \varphi _\gamma (x)=& \frac{(qax\gamma /\stackrel{~}{d},qa\gamma /\stackrel{~}{d}x;q)_{\mathrm{}}}{(\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}\gamma ,q\gamma /\stackrel{~}{d},q\stackrel{~}{a}/\stackrel{~}{d},qx/d,q/dx;q)_{\mathrm{}}}\hfill \\ & \times {}_{8}{}^{}W_{7}^{}(\stackrel{~}{a}\stackrel{~}{b}\stackrel{~}{c}\gamma /q;ax,a/x,\stackrel{~}{a}\gamma ,\stackrel{~}{b}\gamma ,\stackrel{~}{c}\gamma ;q,q/\stackrel{~}{d}\gamma ),|q/\stackrel{~}{d}\gamma |<1.\hfill \end{array}$$
(3.2)
This solution of (3.1) is also the subject of study in Suslov’s papers , in which Fourier-Bessel type orthogonality relations are derived, see also remark 2.
Applying Bailey’s formula \[4, (III.36)\] shows that
$$\begin{array}{cc}\hfill \varphi _\gamma (x)=& \frac{1}{(bc,qa/d,q/ad;q)_{\mathrm{}}}{}_{4}{}^{}\varphi _{3}^{}(\begin{array}{c}ax,a/x,\stackrel{~}{a}\gamma ,\stackrel{~}{a}/\gamma \\ ab,ac,ad\end{array};q,q)\hfill \\ \hfill +& \frac{(ax,a/x,\stackrel{~}{a}\gamma ,\stackrel{~}{a}/\gamma ,qb/d,qc/d;q)_{\mathrm{}}}{(qx/d,q/dx,q\gamma /\stackrel{~}{d},q/\stackrel{~}{d}\gamma ,ab,ac,bc,qa/d,ad/q;q)_{\mathrm{}}}\hfill \\ & \times {}_{4}{}^{}\varphi _{3}^{}(\begin{array}{c}qx/d,q/dx,q\gamma /\stackrel{~}{d},q/\stackrel{~}{d}\gamma \\ qb/d,qc/d,q^2/ad\end{array};q,q),\hfill \end{array}$$
(3.3)
cf. \[15, (2.8)\], hence $`\varphi _\gamma (x)`$ extends to a meromorphic function in $`x`$ and $`\gamma `$ for generic parameters $`a,b,c`$ and $`d`$, with possible poles at $`x^{\pm 1}=q^{1+k}/d`$ ($`k_+`$) and $`\gamma ^{\pm 1}=q^{1+k}/\stackrel{~}{d}`$ ($`k_+`$). It follows from (3.3) that $`\varphi _{\gamma ^{\pm 1}}(x^{\pm 1})=\varphi _\gamma (x)`$ (all sign combinations possible), and that $`\varphi _\gamma `$ satisfies the duality relation
$$\varphi _\gamma (x;a;b,c;d;q)=\varphi _x(\gamma ;\stackrel{~}{a};\stackrel{~}{b},\stackrel{~}{c};\stackrel{~}{d};q).$$
(3.4)
In the remainder of the paper we use the short-hand notation $`\stackrel{~}{\varphi }_x(\gamma )`$ for the right-hand side of (3.4).
Observe that the meromorphic continuation (3.3) of $`\varphi _\gamma (x)`$ implies that
$$\varphi _{\gamma _n}(x)=\frac{1}{(bc,qa/d,q/ad;q)_{\mathrm{}}}p_n(x),n_+,$$
(3.5)
since the factor $`(\stackrel{~}{a}/\gamma ;q)_{\mathrm{}}`$ in front of the second $`{}_{4}{}^{}\varphi _{3}^{}`$ vanishes when $`\gamma =\gamma _n=\stackrel{~}{a}q^n`$ for $`n_+`$. In particular, the duality (2.3) of the Askey-Wilson polynomials is a special case of the duality (3.4) of $`\varphi _\gamma `$.
###### Definition 1.
The solution $`\varphi _\gamma `$ (see (3.2) and (3.3)) of the eigenvalue equation $`(Lf)(x)=\mu (\gamma )f(x)`$ is called the Askey-Wilson function.
## 4. The $`c`$-function expansion.
We observed in the previous section that the poles of the Askey-Wilson function $`\varphi _\gamma (x)`$ for $`|x|0`$ lie on the $`q`$-line $`\{dq^{1k}\}_{k_+}`$. In this section we consider the Askey-Wilson function $`\varphi _\gamma (x)`$ for $`|x|0`$ and $`x`$ lying on the $`q`$-line $`I=\{dtq^k\}_k`$, where $`tq^{}`$ is an extra degree of freedom.
Observe that the eigenvalue equation (3.1) is asymptotically of the form
$$\stackrel{~}{a}^2\left(f(qx)f(x)\right)+\left(f(q^1x)f(x)\right)=\mu (\gamma )f(x)$$
(4.1)
when $`|x|\mathrm{}`$. For generic $`\gamma `$, the asymptotic eigenvalue equation (4.1) has a basis $`\{\mathrm{\Phi }_\gamma ^{free},\mathrm{\Phi }_{\gamma ^1}^{free}\}`$ of solutions on the $`q`$-line $`I`$, where
$$\mathrm{\Phi }_\gamma ^{free}(dtq^k)=(\stackrel{~}{a}\gamma )^k,k.$$
Furthermore, for generic $`\gamma `$ there exists a unique solution $`\mathrm{\Phi }_\gamma (x)`$ of the eigenvalue equation (3.1) on $`I`$ of the form $`\mathrm{\Phi }_\gamma (x)=\mathrm{\Phi }_\gamma ^{free}(x)g(x)`$, where $`g`$ has a convergent power series expansion around $`|x|=\mathrm{}`$ with constant coefficient equal to one. $`\mathrm{\Phi }_\gamma `$ is the so-called asymptotically free solution of the eigenvalue equation (3.1).
An explicit expression for $`\mathrm{\Phi }_\gamma `$ can be obtained from Ismail’s and Rahman’s \[5, (1.13)\] solution of the three term recurrence relation for the associated Askey-Wilson polynomials. After application of the transformation formula \[4, (III.23)\] for very well poised $`{}_{8}{}^{}\varphi _{7}^{}`$’s, it can be expressed as
$$\begin{array}{cc}\hfill \mathrm{\Phi }_\gamma (x)=& \frac{(qa\gamma /\stackrel{~}{a}x,qb\gamma /\stackrel{~}{a}x,qc\gamma /\stackrel{~}{a}x,q\stackrel{~}{a}\gamma /dx,d/x;q)_{\mathrm{}}}{(q/ax,q/bx,q/cx,q/dx,q^2\gamma ^2/dx;q)_{\mathrm{}}}\hfill \\ & \times {}_{8}{}^{}W_{7}^{}(q\gamma ^2/dx;q\gamma /\stackrel{~}{a},q\gamma /\stackrel{~}{d},\stackrel{~}{b}\gamma ,\stackrel{~}{c}\gamma ,q/dx;q,d/x)\mathrm{\Phi }_\gamma ^{free}(x)\hfill \end{array}$$
(4.2)
for $`xI`$ with $`|x|0`$. We now expand the Askey-Wilson function $`\varphi _\gamma (x)`$ as a linear combination of the asymptotically free solutions $`\mathrm{\Phi }_\gamma (x)`$ and $`\mathrm{\Phi }_{\gamma ^1}(x)`$ for $`xI`$ with $`|x|0`$. The coefficients in this expansion can be expressed in terms of the $`c`$-function $`c(\gamma )=c(\gamma ;a;b,c;d;q;t)`$, which is defined by
$$c(\gamma )=\frac{1}{(ab,ac,bc,qa/d;q)_{\mathrm{}}\theta (qadt)}\frac{(a/\gamma ,b/\gamma ,c/\gamma ;q)_{\mathrm{}}\theta (\gamma /dt)}{(q\gamma /d,1/\gamma ^2;q)_{\mathrm{}}}$$
(4.3)
where $`\theta (x)=(x,q/x;q)_{\mathrm{}}`$ is the renormalized Jacobi theta function. We call $`\stackrel{~}{c}(\gamma )=c(\gamma ;\stackrel{~}{a};\stackrel{~}{b},\stackrel{~}{c};\stackrel{~}{d};q;\stackrel{~}{t})`$ the dual $`c`$-function, with the dual parameter $`\stackrel{~}{t}`$ defined by
$$\stackrel{~}{t}=1/qadt.$$
(4.4)
###### Proposition 1.
Let $`xI`$ with $`|x|0`$. Then we have the $`c`$-function expansion
$$\varphi _\gamma (x)=\stackrel{~}{c}(\gamma )\mathrm{\Phi }_\gamma (x)+\stackrel{~}{c}(\gamma ^1)\mathrm{\Phi }_{\gamma ^1}(x)$$
(4.5)
for generic $`\gamma `$.
###### Proof.
Apply Bailey’s three term recurrence relation \[4, (III.37)\] with its parameters specialized as follows:
$$aq\gamma ^2/dx,bq/dx,cq\gamma /\stackrel{~}{a},dq\gamma /\stackrel{~}{d},e\stackrel{~}{b}\gamma ,f\stackrel{~}{c}\gamma .$$
This gives an expansion of the form (4.5) for explicit coefficients $`\stackrel{~}{c}(\gamma )`$ and $`\stackrel{~}{c}(\gamma ^1)`$, which at a first glance still depend on $`xI`$. Using the functional equation
$$\theta (q^kx)=q^{\frac{1}{2}k(k1)}(x)^k\theta (x),k$$
(4.6)
for theta functions, it is easily seen that the coefficients are independent of $`x`$ and that they coincide with the dual $`c`$-functions $`\stackrel{~}{c}(\gamma ^{\pm 1})`$ as defined above. ∎
## 5. The Askey-Wilson function transform.
In this section we define the Askey-Wilson function transform and we state the main result of this paper. At this stage we need to specify a particular parameter domain for the five parameters $`(a,b,c,d,t)`$ in order to ensure positivity of measures.
###### Definition 2.
Let $`V`$ be the set of parameters $`(a,b,c,d,t)^5`$ satisfying the conditions
$$\begin{array}{cc}& t<0,0<b,ca<d/q,\hfill \\ & bd,cdq,ab,ac<1.\hfill \end{array}$$
Observe that $`b,c<1`$ and $`d>q`$ for all $`(a,b,c,d,t)V`$. The domain $`V`$ is self-dual in the following sense.
###### Lemma 1.
The assignment $`(a,b,c,d,t)(\stackrel{~}{a},\stackrel{~}{b},\stackrel{~}{c},\stackrel{~}{d},\stackrel{~}{t})`$ defined by (2.2) and (4.4), is an involution on $`V`$.
###### Proof.
Direct verification. ∎
We fix parameters $`(a,b,c,d,t)V`$ for the remainder of the paper. For the moment we furthermore assume that $`x1/c(x)c(x^1)`$ only has simple poles. This imposes certain generic conditions on the parameters $`(a,b,c,d,t)`$, which will be removed later by a continuity argument.
It is convenient to renormalize the function $`1/c(x)c(x^1)`$ as follows,
$$W(x)=\frac{1}{c(x)c(x^1)c_0}=\frac{(qx/d,q/dx,x^2,1/x^2;q)_{\mathrm{}}}{(ax,a/x,bx,b/x,cx,c/x;q)_{\mathrm{}}\theta (dtx)\theta (dt/x)},$$
(5.1)
where $`c_0`$ is the positive constant
$$c_0=\frac{(ab,ac,bc,qa/d;q)_{\mathrm{}}^2\theta (adt)^2}{a^2}.$$
(5.2)
It follows from (5.1) that
$$W(x)=\frac{\theta (dx)\theta (d/x)}{\theta (dtx)\theta (dt/x)}\mathrm{\Delta }(x),$$
(5.3)
where $`\mathrm{\Delta }()`$ is the weight function (2.7) for the Askey-Wilson polynomials. By (4.6), the quotient of theta functions in (5.3) is a quasi-constant, i.e. it is invariant under replacement of $`x`$ by $`qx`$. In particular, the weight function $`W(x)`$ differs from $`\mathrm{\Delta }(x)`$ only by a quasi-constant factor.
Let $`S`$ be the discrete subset
$$\begin{array}{cc}\hfill S& =\{x||x|>1,c(x)=0\}=S_+S_{},\hfill \\ \hfill S_+& =\{aq^k|k_+,aq^k>1\},\hfill \\ \hfill S_{}& =\{dtq^k|k,dtq^k<1\}.\hfill \end{array}$$
(5.4)
We denote $`\stackrel{~}{S}`$ and $`\stackrel{~}{S}_\pm `$ for the subsets $`S`$ and $`S_\pm `$ with respect to dual parameters. We define a measure $`d\nu (x)`$ by
$$\begin{array}{cc}\hfill f(x)& d\nu (x)=\frac{K}{4\pi i}_{x𝕋}f(x)W(x)\frac{dx}{x}\hfill \\ & +\frac{K}{2}\underset{xS}{}f(x)\underset{y=x}{\text{Res}}\left(\frac{W(y)}{y}\right)\frac{K}{2}\underset{xS^1}{}f(x)\underset{y=x}{\text{Res}}\left(\frac{W(y)}{y}\right),\hfill \end{array}$$
(5.5)
where the positive constant $`K`$ is given by
$$K=(ab,ac,bc,qa/d,q;q)_{\mathrm{}}\sqrt{\frac{\theta (qt)\theta (adt)\theta (bdt)\theta (cdt)}{qabcdt^2}}.$$
(5.6)
This particular choice of normalization constant for the measure $`\nu `$ is justified in theorem 1.
In view of (5.3), we can relate the discrete masses $`\nu \left(\{x\}\right)(=\nu \left(\{x^1\}\right))`$ for $`xS_+`$ to residues of the weight function $`\mathrm{\Delta }()`$, which were written down explicitly in (see also \[4, (7.5.22)\] in order to avoid a small misprint). Explicitly, we obtain for $`x=aq^kS_+`$ with $`k_+`$ the expression
$$\begin{array}{cc}\hfill \nu \left(\{aq^k\}\right)=& \frac{(qa/d,q/ad,1/a^2;q)_{\mathrm{}}}{(q,ab,b/a,ac,c/a;q)_{\mathrm{}}\theta (adt)\theta (dt/a)}\hfill \\ & \times \frac{(a^2,ab,ac,ad;q)_k}{(q,qa/b,qa/c,qa/d;q)_k}\frac{(1a^2q^{2k})}{(1a^2)}\frac{K}{2\stackrel{~}{a}^{2k}}\hfill \end{array}$$
(5.7)
for the corresponding discrete weight. For fixed $`k_+`$, the right hand side of (5.7) gives the unique continuous extension of the discrete weight $`\nu \left(\{aq^k\}\right)`$ and $`\nu \left(\{a^1q^k\}\right)`$ to all parameters $`(a,b,c,d,t)V`$ satisfying $`aq^k>1`$. Furthermore, the (continuously extended) discrete weight $`\nu \left(\{aq^k\}\right)`$ is strictly positive for these parameter values.
A similar argument can be applied for the discrete weights $`\nu \left(\{x\}\right)(=\nu \left(\{x^1\}\right))`$ with $`xS_{}`$. Explicitly we obtain for $`x=dtq^kS_{}`$ with $`k`$,
$$\begin{array}{cc}\hfill \nu \left(\{dtq^k\}\right)=& \frac{(qt,q/d^2t;q)_{\mathrm{}}}{(q,q,a/dt,b/dt,c/dt,adt,bdt,cdt;q)_{\mathrm{}}}\hfill \\ & \times \frac{(1/t,a/dt,b/dt,c/dt;q)_k}{(q/adt,q/bdt,q/cdt,q/d^2t;q)_k}\left(1\frac{1}{d^2t^2q^{2k}}\right)\frac{K\stackrel{~}{a}^{2k}}{2}.\hfill \end{array}$$
(5.8)
As for $`\nu \left(\{x\}\right)`$ with $`xS_+`$, we use the right hand side of (5.8) to define the strictly positive weight $`\nu \left(\{dtq^k\}\right)(=\nu \left(\{d^1t^1q^k\}\right))`$ for all $`(a,b,c,d,t)V`$ satisfying $`dtq^k<1`$.
We conclude that the definition of the measure $`\nu `$ (see (5.5)) can be extended to arbitrary parameters $`(a,b,c,d,t)V`$ using the continuous extensions of its discrete weights given above. The resulting measure $`\nu `$ is a positive measure for all $`(a,b,c,d,t)V`$.
###### Definition 3.
Let $``$ be the Hilbert space consisting of $`L^2`$-functions $`f`$ with respect to $`\nu `$ which satisfy $`f(x)=f(x^1)`$ $`\nu `$-almost everywhere.
We write $`\stackrel{~}{\nu }`$ for the measure $`\nu `$ with respect to dual parameters $`(\stackrel{~}{a},\stackrel{~}{b},\stackrel{~}{c},\stackrel{~}{d},\stackrel{~}{t})`$, and $`\stackrel{~}{}`$ for the associated Hilbert space $``$.
Let $`𝒟`$ be the dense subspace of functions $`f`$ with compact support, i.e.
$$𝒟=\{f|f(dtq^k)=0,k0\},$$
and define
$$\left(f\right)(\gamma )=f(x)\overline{\varphi _\gamma (x)}𝑑\nu (x),f𝒟$$
(5.9)
for generic $`\gamma \{0\}`$.
###### Remark 1.
Observe that the analytic continuation (3.3) of $`\varphi _\gamma (x)`$ is not defined for parameters $`(a,b,c,d,t)V`$ satisfying $`\theta (ad)=0`$. These apparent poles can be removed in view of the original definition (3.2) for $`\varphi _\gamma (x)=\varphi _{\gamma ^1}(x)`$ in terms of very well poised $`{}_{8}{}^{}\varphi _{7}^{}`$’s (observe that $`q/\stackrel{~}{d}<1`$, so that either $`|q\gamma /\stackrel{~}{d}|<1`$ or $`|q/\stackrel{~}{d}\gamma |<1`$ for $`\gamma \{0\}`$, hence $`\varphi _\gamma `$ can be expressed in terms of the original definition (3.2) for generic $`\gamma `$). In particular, the transform $``$ is well defined for parameters satisfying $`\theta (ad)=0`$.
We write $`\stackrel{~}{𝒟}\stackrel{~}{}`$ (respectively $`\stackrel{~}{}`$) for the dense subspace $`𝒟`$ (respectively the function transform $``$) with respect to dual parameters $`(\stackrel{~}{a},\stackrel{~}{b},\stackrel{~}{c},\stackrel{~}{d},\stackrel{~}{t})`$. The main theorem of this paper can now be stated as follows.
###### Theorem 1.
Let $`(a,b,c,d,t)V`$. The transform $``$ extends to an isometric isomorphism $`:\stackrel{~}{}`$ by continuity. The inverse of $``$ is given by $`\stackrel{~}{}`$.
We will discuss the proof of theorem 1 in detail in the remaining sections.
###### Definition 4.
The isometric isomorphism $`:\stackrel{~}{}`$ is called the Askey-Wilson function transform.
Theorem 1 gives a simultaneous generalization of the $`L^2`$-theory for the Mellin-Fock transform and for the (rank one) Hankel transform, see and . In fact, the Askey-Wilson function transform is on top of an hierarchy of several different generalizations of the Mellin-Fock and Hankel transform, in analogy with the polynomial setting. This point of view is emphasized in . In particular, the Askey-Wilson function generalizes the Jacobi function as well as the Bessel function to the level of very well poised $`{}_{8}{}^{}\varphi _{7}^{}`$ series.
Theorem 1 may be partially reformulated in terms of orthogonality relations for “low degree” Askey-Wilson polynomials with respect to the measures $`\nu `$ ($`t<0`$). We say more about this point of view in section 7, where we prove the orthogonality relations for the Askey-Wilson functions $`\varphi _\gamma `$ ($`\gamma \stackrel{~}{S}`$) with respect to the measure $`\nu `$ as an intermediate step of the proof of theorem 1.
## 6. The Wronskian.
We let $`\chi _k𝒟`$ for $`k`$ with $`k0`$ be the characteristic function which is zero at $`dtq^m,d^1t^1q^m`$ for $`m`$ with $`m<k`$ and which is equal to one otherwise. The sequence $`\{\chi _k\}_k`$ of characteristic functions is an approximation of the unit, in the sense that $`\chi _kff`$ in $``$ as $`k\mathrm{}`$ for all $`f`$.
The starting point for the proof of theorem 1 is the following simple observation.
###### Lemma 2.
Let $`f,g\stackrel{~}{𝒟}`$ and $`k0`$, then
$$\begin{array}{cc}\hfill \chi _k(x)\left(\stackrel{~}{}f\right)(x)& \overline{\left(\stackrel{~}{}g\right)(x)}d\nu (x)=\hfill \\ \hfill =& f(\gamma )\overline{g(\gamma ^{})}\left(\chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)\right)d\left(\stackrel{~}{\nu }\times \stackrel{~}{\nu }\right)(\gamma ,\gamma ^{}).\hfill \end{array}$$
###### Proof.
We substitute the definition of the dual Askey-Wilson function transform $`\stackrel{~}{}`$ into the left hand side of the desired identity, and we use that
$$\overline{\stackrel{~}{\varphi }_x(\gamma )}=\stackrel{~}{\varphi }_x(\gamma )=\varphi _\gamma (x)$$
for $`x\text{supp}(\nu )`$ and $`\gamma \text{supp}(\stackrel{~}{\nu })`$, see (3.4) for the second equality. We arrive at
$$\begin{array}{cc}\hfill \chi _k(x)\left(\stackrel{~}{}f\right)(x)& \overline{\left(\stackrel{~}{}g\right)(x)}d\nu (x)=\hfill \\ \hfill =& \chi _k(x)\left(f(\gamma )\varphi _\gamma (x)𝑑\stackrel{~}{\nu }(\gamma )\right)\left(\overline{g(\gamma ^{})}\varphi _\gamma ^{}(x)𝑑\stackrel{~}{\nu }(\gamma ^{})\right)𝑑\nu (x).\hfill \end{array}$$
Since $`f,g\stackrel{~}{𝒟}`$, all integrations are over compact subsets, so we may use Fubini’s theorem to interchange the order of integration. This yields the desired result. ∎
The proof of theorem 1 now hinges on the explicit evaluation of the weak limit as $`k\mathrm{}`$ of the function
$$(\gamma ,\gamma ^{})\chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)$$
(6.1)
with respect to the product measure $`\stackrel{~}{\nu }\times \stackrel{~}{\nu }`$.
As a first step, we express the integral (6.1) in terms of the Wronskian of the Askey-Wilson functions $`\varphi _\gamma `$ and $`\varphi _\gamma ^{}`$. Here the Wronskian $`[f,g](x)`$ for two functions $`f,g:I`$ is defined by
$$[f,g](x)=2\nu \left(\{x\}\right)\alpha (x)\left(f(x)g(qx)f(qx)g(x)\right),xS_{},$$
(6.2)
where $`\alpha (x)`$ is given by (2.4).
###### Proposition 2.
Let $`(a,b,c,d,t)V`$. For generic spectral parameters $`\gamma ,\gamma ^{}`$ such that $`\mu (\gamma )\mu (\gamma ^{})`$, we have
$$\chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)=\frac{[\varphi _\gamma ,\varphi _\gamma ^{}](dtq^{k1})}{\mu (\gamma )\mu (\gamma ^{})}$$
for $`k`$ sufficiently negative.
###### Proof.
The proof of the desired identity simplifies when we slightly perturb the parameters $`(a,b,c,d,t)`$. The proof for parameters $`(a,b,c,d,t)V`$ can then be derived from the perturbed case using the fact that the left and right hand side of the identity depend continuously on the parameters $`(a,b,c,d,t)`$.
Let us indicate one class of possible perturbations of the parameters which is sufficient for our purposes. Let $`(\alpha ,\beta ,\gamma ,\delta ,t)V`$ with $`\alpha ,\delta ,\alpha \delta ,\delta tq^{}`$. Let $`ϵ>0`$ be sufficiently small, then we set
$$a=\alpha e^{\pi iϵ},b=\beta e^{2\pi iϵ},c=\gamma e^{3\pi iϵ},d=\delta e^{6\pi iϵ},$$
while we keep $`t`$ undisturbed. Observe that $`|a|=\alpha `$, $`|b|=\beta `$ etc., and that $`abcd=\alpha \beta \gamma \delta `$. Note that the weight function $`W(x)`$ (see (5.1)) corresponding to the parameters $`(a,b,c,d,t)`$ only has simple poles. We can now define the (in general complex) measure $`\nu `$ with respect to the parameters $`(a,b,c,d,t)`$ by the same formulas (5.4) and (5.5) as before. We also keep the same notations for the other parameter-dependent objects we have encountered (such as e.g. $`\varphi _\gamma `$, $`\chi _k`$ etc.). Then for spectral parameters $`\gamma ,\gamma ^{}\{0\}`$ with $`\gamma ^{\pm 1},(\gamma ^{})^{\pm 1}q^{1+l}/\stackrel{~}{d}`$ ($`l_+`$) we can write for $`k0`$,
$$\chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)=\frac{K}{4\pi i}_{xC_k}\varphi _\gamma (x)\varphi _\gamma ^{}(x)\left(\frac{W(x)}{x}\right)𝑑x,$$
(6.3)
where $`C_k`$ is a continuous, rectifiable Jordan curve in the complex plane, satisfying the following additional conditions:
1. $`C_k`$ has a parametrization of the form $`r_k(z)e^{2\pi iz}`$ for $`z[0,1]`$ with positive radial function $`r_k:[0,1](0,\mathrm{})`$ (we orientate $`C_k`$ according to this parametrization);
2. $`C_k`$ is invariant under inversion, i.e. $`C_k^1=C_k`$, where $`C_k^1=\{x^1|xC_k\}`$;
3. The sequences $`\{aq^l,bq^l,cq^l,q^{1+l}/d\}_{l_+}`$ are in the interior of $`C_k`$;
4. The intersection of the $`q`$-interval $`I`$ with the interior of $`C_k`$ is given by the sequence $`\{dtq^m|m:mk\}`$.
The existence of $`C_k`$ is easy to prove with our choice of perturbed parameters, and (6.3) is then a direct consequence of Cauchy’s theorem.
Now recall that the Askey-Wilson function $`\varphi _\gamma `$ is an eigenfunction of the second-order $`q`$-difference operator $`L`$ with eigenvalue $`\mu (\gamma )`$, see section 3. This implies that
$$\begin{array}{cc}\hfill \left(\mu (\gamma )\mu (\gamma ^{})\right)& \chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)=\hfill \\ \hfill =& \frac{K}{4\pi i}_{xC_k}\left((L\varphi _\gamma )(x)\varphi _\gamma ^{}(x)\varphi _\gamma (x)\left(L\varphi _\gamma ^{}\right)(x)\right)\left(\frac{W(x)}{x}\right)𝑑x.\hfill \end{array}$$
(6.4)
In view of the explicit expression for $`L`$ (see (2.4)), we can write
$$(L\varphi _\gamma )(x)\varphi _\gamma ^{}(x)\varphi _\gamma (x)\left(L\varphi _\gamma ^{}\right)(x)=\alpha (1/x)\xi _{\gamma ,\gamma ^{}}(x/q)\alpha (x)\xi _{\gamma ,\gamma ^{}}(x)$$
(6.5)
with
$$\xi _{\gamma ,\gamma ^{}}(x)=\varphi _\gamma (x)\varphi _\gamma ^{}(qx)\varphi _\gamma (qx)\varphi _\gamma ^{}(x).$$
Furthermore, by the explicit expression for the weight function $`W(x)`$ (see (5.1)) and for the function $`\alpha (x)`$ (see (2.4)), we have the identity
$$\alpha (1/x)W(x)=\alpha (x/q)W(x/q).$$
(6.6)
Substitution of (6.5) and (6.6) into the right hand side of (6.4) then shows that
$$\begin{array}{cc}\hfill \left(\mu (\gamma )\mu (\gamma ^{})\right)& \chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)=\hfill \\ & =\frac{K}{4\pi i}_{xq^1C_kC_k}\alpha (x)\xi _{\gamma ,\gamma ^{}}(x)\left(\frac{W(x)}{x}\right)𝑑x.\hfill \end{array}$$
(6.7)
Using the explicit form of the contour $`C_k`$, we can shrink $`q^1C_k`$ back to $`C_k`$ in the right hand side of (6.7) while picking up residues of the integrand at the simple poles $`x=dtq^{k1}`$ and $`x=d^1t^1q^k`$. By (6.6) and by the identity $`\xi _{\gamma ,\gamma ^{}}(x)=\xi _{\gamma ,\gamma ^{}}(q^1x^1)`$, we can pull the two residues together, yielding
$$\begin{array}{cc}\hfill \left(\mu (\gamma )\mu (\gamma ^{})\right)& \chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)=\hfill \\ \hfill =& K\underset{y=dtq^{k1}}{\text{Res}}\left(\frac{W(y)}{y}\right)\alpha (dtq^{k1})\xi _{\gamma ,\gamma ^{}}(dtq^{k1})\hfill \\ \hfill =& [\varphi _\gamma ,\varphi _\gamma ^{}](dtq^{k1}),\hfill \end{array}$$
as desired. ∎
## 7. Orthogonality relations.
If we take in (6.1) the values of the spectral parameters $`\gamma ,\gamma ^{}`$ in the discrete support $`\stackrel{~}{S}`$ of the dual measure $`\stackrel{~}{\nu }`$, then the limit of (6.1) as $`k\mathrm{}`$ can be taken point-wise. We evaluate these limits in this section, which lead to explicit orthogonality relations for the Askey-Wilson functions $`\varphi _\gamma `$ ($`\gamma \stackrel{~}{S}`$) with respect to the measure $`\nu `$. This essentially establishes theorem 1 for the “completely discrete part” of the Askey-Wilson function transform.
###### Lemma 3.
The Askey-Wilson functions $`\varphi _\gamma `$ ($`\gamma \stackrel{~}{S}`$) are elements of the Hilbert space $``$. Furthermore,
$$\varphi _\gamma (x)\overline{\varphi _\gamma ^{}(x)}𝑑\nu (x)=0,\gamma ,\gamma ^{}\stackrel{~}{S},\gamma \gamma ^{}.$$
###### Proof.
We first observe that the asymptotic behaviour of the discrete weights $`\nu \left(\{x\}\right)`$ ($`x\stackrel{~}{S}_{}`$) is given by
$$\nu \left(\{dtq^k\}\right)=\frac{M}{2\stackrel{~}{a}^{2k}}\left(1+𝒪(q^k)\right),k\mathrm{},$$
(7.1)
where $`M`$ is the positive constant
$$M=\frac{\theta (qt)K}{(q,q;q)_{\mathrm{}}\theta (adt)\theta (bdt)\theta (cdt)}.$$
(7.2)
This is a direct consequence of the explicit expression (5.8) for $`\nu \left(\{x\}\right)`$ ($`xS_{}`$).
For the moment we assume $`adq^{}`$ and $`\stackrel{~}{a},\stackrel{~}{d}\stackrel{~}{t}\pm q^{\frac{1}{2}}`$ in order to be able to apply the $`c`$-function expansion for $`\varphi _\gamma `$ ($`\gamma \stackrel{~}{S}`$). We fix $`\gamma \stackrel{~}{S}`$, so that $`\varphi _\gamma (x)=\stackrel{~}{c}(\gamma ^1)\mathrm{\Phi }_{\gamma ^1}(x)`$ for $`xS_{}`$ sufficiently negative by proposition 1. Observe that the apparent poles of $`\stackrel{~}{c}(\gamma ^1)`$ as function of the parameters $`(a,b,c,d,t)V`$ are removable. Hence we can extend the definition of the $`c`$-function $`\stackrel{~}{c}(\gamma ^1)`$ to parameter values $`(a,b,c,d,t)V`$ by continuity and we obtain
$$\varphi _\gamma (x)=\stackrel{~}{c}(\gamma ^1)\mathrm{\Phi }_{\gamma ^1}(x),\gamma \stackrel{~}{S}$$
(7.3)
for all $`(a,b,c,d,t)V`$ if $`xI`$ is sufficiently negative. Since
$$\mathrm{\Phi }_{\gamma ^1}(dtq^k)=\left(\stackrel{~}{a}\gamma ^1\right)^k\left(1+𝒪(q^k)\right),k\mathrm{}$$
(7.4)
and $`|\gamma ^1|<1`$ for $`\gamma \stackrel{~}{S}`$, it follows from (7.1) that $`\varphi _\gamma `$ for all $`\gamma \stackrel{~}{S}`$.
Fix now $`\gamma ,\gamma ^{}\stackrel{~}{S}`$ with $`\gamma \gamma ^{}`$. Then $`\mu (\gamma )\mu (\gamma ^{})`$ by (2.5), hence by proposition 2 and by the previous paragraph,
$$\varphi _\gamma (x)\overline{\varphi _\gamma ^{}(x)}𝑑\nu (x)=\frac{1}{(\mu (\gamma )\mu (\gamma ^{}))}\underset{k\mathrm{}}{lim}[\varphi _\gamma ,\varphi _\gamma ^{}](dtq^k).$$
(7.5)
Here we have used the fact that $`\varphi _\gamma ^{}()`$ is real valued on $`\text{supp}(\nu )`$ in order to get rid of the complex conjugate in the integrand.
It remains to prove that the limit in the right hand side of (7.5) tends to zero. Let $`f`$ and $`g`$ be two functions in the Hilbert space $``$. By the asymptotics (7.1), we have $`f(dtq^k)=o(\stackrel{~}{a}^k)`$ as $`k\mathrm{}`$, and similarly for $`g`$. Furthermore, we have the asymptotics
$$\alpha (dtq^k)=\stackrel{~}{a}^2\left(1+𝒪(q^k)\right),k\mathrm{}$$
(7.6)
for the coefficient $`\alpha ()`$ of the Askey-Wilson second order $`q`$-difference operator $`L`$ (see (2.4)). Combined with (7.1), it follows from the definition (6.2) of the Wronskian that
$$\underset{k\mathrm{}}{lim}[f,g](dtq^k)=0,f,g.$$
In particular, the limit in the right hand side of (7.5) tends to zero, hence $`\varphi _\gamma `$ is orthogonal to $`\varphi _\gamma ^{}`$ in $``$. ∎
The quadratic norms of the Askey-Wilson functions $`\varphi _\gamma `$ ($`\gamma \stackrel{~}{S}`$) in $``$ can be evaluated as follows.
###### Lemma 4.
For all $`\gamma \stackrel{~}{S}`$,
$$|\varphi _\gamma (x)|^2𝑑\nu (x)=\frac{1}{2\stackrel{~}{\nu }\left(\{\gamma \}\right)}.$$
###### Proof.
As in the proof of lemma 3, we assume for the moment that $`adq^{}`$ and that $`\stackrel{~}{a},\stackrel{~}{d}\stackrel{~}{t}\pm q^{\frac{1}{2}}`$. In view of (7.3), (7.4) and (7.1), these generic conditions can be removed at the end of the proof by applying the dominated convergence theorem.
We fix $`\gamma \stackrel{~}{S}`$, then by lemma 3 and proposition 2,
$$\begin{array}{cc}\hfill |\varphi _\gamma (x)|^2𝑑\nu (x)=& \underset{k\mathrm{}}{lim}\chi _k(x)\varphi _\gamma (x)^2𝑑\nu (x)\hfill \\ \hfill =& \underset{k\mathrm{}}{lim}\left(\underset{\gamma ^{}\gamma }{lim}\chi _k(x)\varphi _\gamma (x)\varphi _\gamma ^{}(x)𝑑\nu (x)\right)\hfill \\ \hfill =& \underset{k\mathrm{}}{lim}\left(\underset{\gamma ^{}\gamma }{lim}\frac{[\varphi _\gamma ,\varphi _\gamma ^{}](dtq^k)}{(\mu (\gamma )\mu (\gamma ^{}))}\right),\hfill \end{array}$$
(7.7)
where we used that $`\varphi _\gamma ()`$ is real valued on $`\text{supp}(\nu )`$ for the first equality.
It remains to evaluate the limits of the Wronskian in the last equality of (7.7). We use the $`c`$-function expansion for the Askey-Wilson functions (see proposition 1 and (7.3)) to rewrite the Wronskian as
$$\begin{array}{cc}\hfill [\varphi _\gamma ,\varphi _\gamma ^{}](dtq^k)=& \stackrel{~}{c}(\gamma ^1)\stackrel{~}{c}(\gamma ^{})[\mathrm{\Phi }_{\gamma ^1},\mathrm{\Phi }_\gamma ^{}](dtq^k)\hfill \\ & +\stackrel{~}{c}(\gamma ^1)\stackrel{~}{c}(\gamma ^{}{}_{}{}^{1})[\mathrm{\Phi }_{\gamma ^1},\mathrm{\Phi }_{\gamma ^{}^1}](dtq^k)\hfill \end{array}$$
(7.8)
for $`k0`$. Now there exists open neighbourhoods $`U_\pm `$ of $`\gamma ^{\pm 1}`$ in the complex plane such that
$$\mathrm{\Phi }_\gamma ^{}(x)=\mathrm{\Phi }_\gamma ^{}^{free}(x)(1+f_\gamma ^{}(x)),xS_{},x0$$
(7.9)
where $`f_\gamma ^{}(x)`$ admits a convergent power-series expansion around $`x=\mathrm{}`$ with coefficients depending analytically on $`\gamma ^{}U_\pm `$ and with constant coefficient equal to zero. Furthermore, for sufficiently small neighbourhoods $`U_\pm `$, differentiation with respect to $`\gamma ^{}`$ may be interchanged with summation in the power series expansion of $`f_\gamma ^{}(x)`$ around $`x=\mathrm{}`$ when $`|x|N>0`$ for some $`U_\pm `$ independent positive constant $`N`$.
Combined with (7.9), we see that
$$\begin{array}{cc}\hfill [\mathrm{\Phi }_{\gamma ^1},& \mathrm{\Phi }_{\gamma ^{}^{\pm 1}}](dtq^k)=\hfill \\ \hfill =& 2\nu \left(\{dtq^k\}\right)\alpha (dtq^k)\stackrel{~}{a}^{2k1}(\gamma ^1\gamma ^{}{}_{}{}^{\pm 1})^{k1}(\gamma ^1\gamma ^{}{}_{}{}^{\pm 1})(1+𝒪(q^k))\hfill \end{array}$$
(7.10)
as $`k\mathrm{}`$, with $`𝒪(q^k)`$ uniform in $`\gamma ^{}U_\pm `$. Now we substitute (7.9) in (7.8), and we use the fact that $`\stackrel{~}{c}(\gamma ^{})`$ has a simple zero at $`\gamma ^{}=\gamma `$ and that
$$\mu (\gamma )\mu (\gamma ^{})=\frac{\stackrel{~}{a}}{\gamma }(\gamma \gamma ^{})(\gamma \gamma ^{}{}_{}{}^{1}),$$
(7.11)
then using (7.10) and the fact that $`|\gamma |>1`$, we derive that
$$[\varphi _\gamma ,\varphi _\gamma ^{}](dtq^k)=2\nu \left(\{dtq^k\}\right)\alpha (dtq^k)\stackrel{~}{a}^{2k1}\left(s_k(\gamma ,\gamma ^{})+r_k(\gamma ,\gamma ^{})\right)$$
(7.12)
with
$$s_k(\gamma ,\gamma ^{})=\underset{ϵ=\pm 1}{}\stackrel{~}{c}(\gamma ^{}{}_{}{}^{ϵ})\stackrel{~}{c}(\gamma ^1)\gamma ^{}{}_{}{}^{ϵ(k1)}\gamma _{}^{1k}(\gamma ^1\gamma ^{}{}_{}{}^{ϵ})$$
(7.13)
and with a remainder term $`r_k(\gamma ,\gamma ^{})`$ satisfying
$$\underset{\gamma ^{}\gamma }{lim}\frac{r_k(\gamma ,\gamma ^{})}{(\mu (\gamma )\mu (\gamma ^{}))}=𝒪(q^k),k\mathrm{}.$$
(7.14)
Using (7.11), (7.13) and the fact that $`\stackrel{~}{c}(\gamma ^{})`$ has a simple zero at $`\gamma ^{}=\gamma `$, we derive
$$\begin{array}{cc}\hfill \underset{\gamma ^{}\gamma }{lim}\frac{s_k(\gamma ,\gamma ^{})}{(\mu (\gamma )\mu (\gamma ^{}))}=& \frac{\gamma }{\stackrel{~}{a}}\left(\frac{1}{\stackrel{~}{c}(\gamma ^1)}\underset{\gamma ^{}=\gamma }{\text{Res}}\left(\frac{1}{\stackrel{~}{c}(\gamma ^{})}\right)\right)^1+\frac{\gamma }{\stackrel{~}{a}}\frac{\stackrel{~}{c}(\gamma ^1)^2}{(\gamma ^1\gamma )\gamma ^{2k}}\hfill \\ \hfill =& \frac{\stackrel{~}{K}}{\stackrel{~}{a}\stackrel{~}{c}_0}\frac{1}{2\stackrel{~}{\nu }\left(\{\gamma \}\right)}+\frac{\gamma }{\stackrel{~}{a}}\frac{\stackrel{~}{c}(\gamma ^1)^2}{(\gamma ^1\gamma )\gamma ^{2k}}\hfill \end{array}$$
(7.15)
where $`\stackrel{~}{K}`$ and $`\stackrel{~}{c}_0`$ are the constants $`K`$ (see (5.6)) and $`c_0`$ (see (5.2)) respectively with respect to dual parameters. Here we have used (5.1) and (5.5) for the second equality.
We can evaluate now the limits in the last equality of (7.7) by substituting (7.12) for the Wronskian and by using (7.1), (7.6), (7.13) and (7.14) together with the fact that $`|\gamma |>1`$. This gives
$$|\varphi _\gamma (x)|^2𝑑\nu (x)=\frac{\stackrel{~}{K}M}{\stackrel{~}{c}_0}\frac{1}{2\stackrel{~}{\nu }\left(\{\gamma \}\right)}.$$
By a direct computation one verifies that $`\stackrel{~}{K}M/\stackrel{~}{c}_0=1`$, which completes the proof of the lemma. ∎
Let $`𝒟_d𝒟`$ be the sub-space consisting of the functions in $`𝒟`$ which are supported on the discrete support $`S`$ of the measure $`\nu `$, and let $`\stackrel{~}{𝒟}_d`$ be the sub-space $`𝒟_d`$ with respect to dual parameters. In view of lemma 2 and lemma 1, we can now reformulate lemma 3 and lemma 4 in the following way.
###### Proposition 3.
Let $`f,g𝒟_d`$, then $`f,g\stackrel{~}{}`$ and
$$\left(f\right)(\gamma )\overline{\left(g\right)(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )=f(x)\overline{g(x)}𝑑\nu (x).$$
Recall from (3.5) that $`\varphi _\gamma `$ is a multiple of the Askey-Wilson polynomial $`p_k`$ when $`\gamma =\gamma _k=\stackrel{~}{a}q^k\stackrel{~}{S}_+`$ ($`k_+`$). So if $`\stackrel{~}{a}>1`$, then lemma 3 and lemma 4 give explicit orthogonality relations for the “low degree” Askey-Wilson polynomials $`p_k`$ ($`k_+`$, $`kk_0`$) with respect to the one-parameter family of measures $`\nu `$ ($`t<0`$), where $`k_0`$ is the largest positive integer such that $`\stackrel{~}{a}q^{k_0}>1`$. The sub-space spanned by $`\{p_k\}_{k=0}^{k_0}`$ is exactly the sub-space of polynomials in $`x+x^1`$ which are square integrable with respect to $`\nu `$ (this follows easily from (7.1)). Furthermore, the functions $`\varphi _\gamma `$ with $`\gamma \stackrel{~}{S}_{}`$ constitute an explicit family of mutually orthogonal functions with respect to the measure $`\nu `$, which in addition are orthogonal to the finite set of Askey-Wilson polynomials $`p_k`$ ($`k_+`$, $`kk_0`$).
From this point of view, lemma 3 and lemma 4 bear close resemblance with indeterminate moment problems and non-extremal measures. In fact, the measures $`\nu `$ ($`t<0`$) formally reduce to a genuine one-parameter family of non-extremal orthogonality measures for some indeterminate moment problem in certain degenerate cases of the Askey-Wilson function transform. This point of view was emphasized in , where it was for instance shown that a formal limit of the above orthogonality relations leads to a one-parameter family of non-extremal orthogonality measures for the continuous dual $`q^1`$-Hahn polynomials, see also .
###### Remark 2.
Fix $`k`$ sufficiently negative, and let $`\gamma _1(k),\gamma _2(k)`$ be two zeros of the function $`\gamma \varphi _\gamma (dtq^{k1})`$ such that $`\mu (\gamma _1(k))\mu (\gamma _2(k))`$. Then proposition 2 implies that the Askey-Wilson functions $`\varphi _{\gamma _1(k)}`$ and $`\varphi _{\gamma _2(k)}`$ are mutually orthogonal with respect to the measure $`\chi _k(x)d\nu (x)`$ of compact support. These Fourier-Bessel type orthogonality relations for the Askey-Wilson functions were derived and studied by Suslov in and .
## 8. The continuous part of the Askey-Wilson function transform.
In this section we first evaluate the weak limit of (6.1) as $`k\mathrm{}`$ with respect to the product measure $`\stackrel{~}{\nu }|_𝕋\times \stackrel{~}{\nu }|_𝕋`$. We directly formulate the result in terms of suitable unitarity properties of the Askey-Wilson function transform (compare with proposition 3 in the square integrable setting).
Let $`𝒟_c𝒟`$ be the sub-space consisting of the functions in $`𝒟`$ which are supported on $`𝕋`$, and let $`\stackrel{~}{𝒟}_c`$ be the sub-space $`𝒟_c`$ with respect to dual parameters. Observe that $`𝒟=𝒟_c𝒟_d`$ is an orthogonal direct sum decomposition of the dense sub-space $`𝒟`$.
###### Proposition 4.
Let $`f,g𝒟_c`$, then $`f,g\stackrel{~}{}`$ and
$$\left(f\right)(\gamma )\overline{\left(g\right)(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )=f(x)\overline{g(x)}𝑑\nu (x).$$
###### Proof.
The proof is similar to the proof of \[8, prop 7.7\] and of \[7, prop 6.1\], where the analogous statement was derived for the big and the little $`q`$-Jacobi function transforms, respectively. Since some care has to be taken in order to match the constants, we repeat here the proof in some detail.
We prove the proposition with respect to dual parameters (cf. lemma 1). It suffices to prove it for functions $`f,g\stackrel{~}{𝒟}_c`$ which are continuous on $`𝕋`$ and supported within $`𝕋\{\pm 1\}`$.
We use lemma 2, proposition 2 and the explicit form of the measure $`\stackrel{~}{\nu }|_𝕋`$ (see (5.1) and (5.5)), together with the invariance of the Askey-Wilson function $`\varphi _\gamma `$ and the measure $`d\stackrel{~}{\nu }(\gamma )`$ under $`\gamma \gamma ^1`$, to write
$$\begin{array}{cc}\hfill & \chi _{1k}(x)\left(\stackrel{~}{}f\right)(x)\overline{\left(\stackrel{~}{}g\right)(x)}d\nu (x)=\hfill \\ & =\frac{\stackrel{~}{K}^2}{4\pi ^2\stackrel{~}{c}_0^2}_{\theta ^{}=0}^\pi _{\theta =0}^\pi f(e^{i\theta })\overline{g(e^{i\theta ^{}})}\frac{[\varphi _{e^{i\theta }},\varphi _{e^{i\theta ^{}}}](dtq^k)}{\left(\mu (e^{i\theta })\mu (e^{i\theta ^{}})\right)}\frac{d\theta }{|\stackrel{~}{c}(e^{i\theta })|^2}\frac{d\theta ^{}}{|\stackrel{~}{c}(e^{i\theta ^{}})|^2}\hfill \end{array}$$
(8.1)
for $`k`$ sufficiently large. It suffices to prove that the right hand side of (8.1) tends to $`f(\gamma )\overline{g(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )`$ as $`k\mathrm{}`$.
In order to compute the limit $`k\mathrm{}`$ of the right hand side of (8.1), we first observe that the factor $`\mu (e^{i\theta })\mu (e^{i\theta ^{}})`$ occuring in the integrand can be rewritten as
$$\begin{array}{cc}\hfill \mu (e^{i\theta })\mu (e^{i\theta ^{}})& =\stackrel{~}{a}e^{i\theta }\left(e^{i\theta }e^{i\theta ^{}}\right)\left(e^{i\theta }e^{i\theta ^{}}\right)\hfill \\ & =4\stackrel{~}{a}\mathrm{sin}\left(\frac{\theta +\theta ^{}}{2}\right)\mathrm{sin}\left(\frac{\theta \theta ^{}}{2}\right),\hfill \end{array}$$
(8.2)
cf. (7.11). Similarly as in the proof of lemma 4, we derive now from the first equality in (8.2), the $`c`$-function expansion (see proposition 1) and the mean value theorem, that
$$[\varphi _{e^{i\theta }},\varphi _{e^{i\theta ^{}}}](dtq^k)=2\nu \left(\{dtq^k\}\right)\alpha (dtq^k)\stackrel{~}{a}^{2k1}\left(s_k(\theta ,\theta ^{})+r_k(\theta ,\theta ^{})\right)$$
(8.3)
with
$$s_k(\theta ,\theta ^{})=\underset{ϵ,ϵ^{}=\pm 1}{}\stackrel{~}{c}(e^{iϵ\theta })\stackrel{~}{c}(e^{iϵ^{}\theta ^{}})e^{i(k1)ϵ\theta }e^{i(k1)ϵ^{}\theta ^{}}\left(e^{iϵ\theta }e^{iϵ^{}\theta ^{}}\right)$$
(8.4)
and with a remainder term $`r_k(\theta ,\theta ^{})`$ satisfying
$$\underset{\delta \theta \theta ^{}\pi \delta }{\text{sup}}\left|\frac{r_k(\theta ,\theta ^{})}{\left(\mu (e^{i\theta })\mu (e^{i\theta ^{}})\right)}\right|=𝒪(kq^k),k\mathrm{}$$
(8.5)
for all $`0<\delta <\pi /2`$. We substitute the expression (8.3) for the Wronskian in the right hand side of (8.1). In view of the asymptotics (7.1) and (7.6), it then suffices to calculate the limit $`k\mathrm{}`$ of
$$\frac{\stackrel{~}{K}^2M\stackrel{~}{a}}{4\pi ^2\stackrel{~}{c}_0^2}_{\theta ^{}=0}^\pi _{\theta =0}^\pi f(e^{i\theta })\overline{g(e^{i\theta ^{}})}\frac{\left(s_k(\theta ,\theta ^{})+r_k(\theta ,\theta ^{})\right)}{\left(\mu (e^{i\theta })\mu (e^{i\theta ^{}})\right)}\frac{d\theta }{|\stackrel{~}{c}(e^{i\theta })|^2}\frac{d\theta ^{}}{|\stackrel{~}{c}(e^{i\theta ^{}})|^2}.$$
(8.6)
The asymptotic behaviour (8.5) shows that the limit $`k\mathrm{}`$ of the integral (8.6) of the remainder term $`r_k(\theta ,\theta ^{})`$ gives zero by dominated convergence. By the Riemann-Lebesgue lemma, the integral (8.6) of the two terms from $`s_k(\theta ,\theta ^{})`$ corresponding to $`ϵϵ^{}=1`$ in (8.4) tend to zero as $`k\mathrm{}`$. It thus remains to calculate the limit $`k\mathrm{}`$ of (8.6) in which the factor $`\left(s_k(\theta ,\theta ^{})+r_k(\theta ,\theta ^{})\right)`$ of the integrand is replaced by
$$\begin{array}{cc}\hfill t_k(\theta ,\theta ^{})& =\underset{\stackrel{ϵ,ϵ^{}=\pm 1}{ϵϵ^{}=1}}{}\stackrel{~}{c}(e^{iϵ\theta })\stackrel{~}{c}(e^{iϵ^{}\theta ^{}})e^{i(k1)ϵ\theta }e^{i(k1)ϵ^{}\theta ^{}}\left(e^{iϵ\theta }e^{iϵ^{}\theta ^{}}\right)\hfill \\ & =\left(\stackrel{~}{c}(e^{i\theta })\stackrel{~}{c}(e^{i\theta ^{}})\stackrel{~}{c}(e^{i\theta })\stackrel{~}{c}(e^{i\theta ^{}})\right)e^{i(k1)(\theta \theta ^{})}\left(e^{i\theta }e^{i\theta ^{}}\right)\hfill \\ & 4\mathrm{sin}\left(\frac{\theta +\theta ^{}}{2}\right)\mathrm{sin}\left(\frac{(2k1)(\theta \theta ^{})}{2}\right)\stackrel{~}{c}(e^{i\theta })\stackrel{~}{c}(e^{i\theta ^{}}).\hfill \end{array}$$
It follows now by yet another application of the Riemann-Lebesgue lemma that it remains to calculate the limit $`k\mathrm{}`$ of
$$\frac{\stackrel{~}{K}^2M}{4\pi ^2\stackrel{~}{c}_0^2}_{\theta ^{}=0}^\pi _{\theta =0}^\pi f(e^{i\theta })\overline{g(e^{i\theta ^{}})}\left(\frac{\mathrm{sin}\left((2k1)(\theta \theta ^{})/2\right)}{\mathrm{sin}\left((\theta \theta ^{})/2\right)}\right)\frac{d\theta }{\stackrel{~}{c}\left(e^{i\theta }\right)}\frac{d\theta ^{}}{\stackrel{~}{c}\left(e^{i\theta ^{}}\right)}.$$
(8.7)
Now by the well known $`L^2`$-properties of the Dirichlet kernel, the limit $`k\mathrm{}`$ of (8.7) exists, and it equals
$$\frac{\stackrel{~}{K}^2M}{2\pi \stackrel{~}{c}_0^2}_{\theta =0}^\pi f(e^{i\theta })\overline{g(e^{i\theta })}\frac{d\theta }{|\stackrel{~}{c}(e^{i\theta })|^2}=\frac{\stackrel{~}{K}M}{\stackrel{~}{c}_0}f(\gamma )\overline{g(\gamma )}𝑑\stackrel{~}{\nu }(\gamma ).$$
The proposition follows now directly from the fact that $`\stackrel{~}{K}M/\stackrel{~}{c}_0=1`$ (compare with the proof of lemma 4). ∎
In order to completely understand the weak limit of (6.1) as $`k\mathrm{}`$ with respect to the product measure $`\stackrel{~}{\nu }\times \stackrel{~}{\nu }`$, we still have to deal with the mixed continuous-discrete case. This case is covered by the following lemma.
###### Lemma 5.
Let $`f𝒟_c`$ and $`g𝒟_d`$, then
$$\left(f\right)(\gamma )\overline{\left(g\right)(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )=0.$$
###### Proof.
We establish the desired identity with respect to dual parameters. Let $`f\stackrel{~}{𝒟}_c`$ and $`g\stackrel{~}{𝒟}_d`$. By proposition 3 and proposition 4 we have $`\stackrel{~}{}f`$ and $`\stackrel{~}{}g`$ respectively. In particular, we may assume without loss of generality that $`f`$ is continuous on $`𝕋`$, and supported within $`𝕋\{\pm 1\}`$. For $`\gamma \stackrel{~}{S}\stackrel{~}{S}^1`$ we now define $`f_\gamma \stackrel{~}{𝒟}_c`$ by $`f_\gamma (\gamma ^{})=f(\gamma ^{})/\left(\mu (\gamma )\mu (\gamma ^{})\right)`$. Observe that $`f_\gamma `$ is continuous on $`𝕋`$ and supported within $`𝕋\{\pm 1\}`$. By lemma 2 and proposition 2 we can now write
$$\begin{array}{cc}\hfill \left(\stackrel{~}{}f\right)(x)& \overline{\left(\stackrel{~}{}g\right)(x)}d\nu (x)=\hfill \\ & =2\underset{\gamma \stackrel{~}{S}}{}\overline{g(\gamma )}\stackrel{~}{\nu }(\{\gamma \})\underset{k\mathrm{}}{lim}_𝕋f_\gamma (\gamma ^{})[\varphi _\gamma ,\varphi _\gamma ^{}](dtq^k)𝑑\stackrel{~}{\nu }(\gamma ^{}).\hfill \end{array}$$
(8.8)
Now we observe that (7.3), (7.4), proposition 1 and the asymptotic behaviour
$$\underset{\gamma ^{}supp(f_\gamma )}{\text{sup}}\left|\mathrm{\Phi }_{\gamma ^{}^{\pm 1}}(dtq^k)\right|=𝒪\left(\stackrel{~}{a}^k\right),k\mathrm{}$$
imply that
$$\begin{array}{cc}\hfill \underset{\gamma ^{}supp(f_\gamma )}{\text{sup}}\left|[\varphi _\gamma ,\varphi _\gamma ^{}](dtq^k)\right|& =2\nu \left(\{dtq^k\}\right)\alpha (dtq^k)\stackrel{~}{a}^{2k1}|\gamma |^k𝒪(1)\hfill \\ & =𝒪\left(|\gamma |^k\right)\hfill \end{array}$$
as $`k\mathrm{}`$, where the second equality is a consequence of (7.1) and (7.6). But $`|\gamma |^1<1`$ for $`\gamma \stackrel{~}{S}`$, hence we may use Lebesgue’s dominated convergence theorem to interchange limit and integration in the right hand side of (8.8). It follows that the right hand side of (8.8) tends to zero, which completes the proof of the lemma. ∎
## 9. Completion of the proof of theorem 1.
In this section we complete the proof of theorem 1.
First of all, we observe that the results of section 7 and section 8 immediately imply that the Askey-Wilson function transform $``$ extends to an isometry $`:\stackrel{~}{}`$. Indeed, it follows from proposition 3, proposition 4 and lemma 5, together with the orthogonal direct sum decomposition $`𝒟=𝒟_c𝒟_d`$, that
$$\left(f\right)(\gamma )\overline{\left(g\right)(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )=f(x)\overline{g(x)}𝑑\nu (x),f,g𝒟.$$
Since $`𝒟`$ is dense, it follows that the Askey-Wilson function transform $``$ uniquely extends to an isometry $`:\stackrel{~}{}`$ by continuity. In particular, the dual Askey-Wilson function transform $`\stackrel{~}{}:\stackrel{~}{}`$ is an isometry in view of lemma 1.
Fix now arbitrary $`f\stackrel{~}{}`$ and $`g`$, and write $`\stackrel{~}{\chi }_k\stackrel{~}{𝒟}`$ ($`k`$ sufficiently negative) for the characteristic function $`\chi _k`$ with respect to dual parameters (see the beginning of section 6). Since $`:\stackrel{~}{}`$ and $`\stackrel{~}{}:\stackrel{~}{}`$ are continuous, we have
$$\begin{array}{cc}\hfill \left(\stackrel{~}{}f\right)(x)\overline{g(x)}𝑑\nu (x)& =\underset{k,m\mathrm{}}{lim}\left(\stackrel{~}{\chi }_k(\gamma )f(\gamma )\overline{\stackrel{~}{\varphi }_x(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )\right)\overline{\chi _m(x)g(x)}𝑑\nu (x)\hfill \\ & =\underset{k,m\mathrm{}}{lim}\stackrel{~}{\chi }_k(\gamma )f(\gamma )\overline{\left(\chi _m(x)g(x)\overline{\varphi _\gamma (x)}𝑑\nu (x)\right)}𝑑\stackrel{~}{\nu }(\gamma )\hfill \\ & =f(\gamma )\overline{\left(g\right)(\gamma )}𝑑\stackrel{~}{\nu }(\gamma ),\hfill \end{array}$$
where we used Fubini’s theorem, the duality (3.4) of the Askey-Wilson function, and the fact that the Askey-Wilson function $`\varphi _\gamma (x)`$ is real-valued for $`x\text{supp}(\nu )`$ and $`\gamma \text{supp}(\stackrel{~}{\nu })`$. But then we have for all $`f,g`$,
$$\left(\stackrel{~}{}\left(f\right)\right)(x)\overline{g(x)}𝑑\nu (x)=\left(f\right)(\gamma )\overline{\left(g\right)(\gamma )}𝑑\stackrel{~}{\nu }(\gamma )=f(x)\overline{g(x)}𝑑\nu (x)$$
since $`:\stackrel{~}{}`$ is an isometry. It follows that $`\stackrel{~}{}=\text{Id}_{}`$. By a similar argument, we obtain $`\stackrel{~}{}=\text{Id}_\stackrel{~}{}`$ (or simply replace the parameters by the corresponding dual parameters in $`\stackrel{~}{}=\text{Id}_{}`$ and use lemma 1).
We conclude that $`:\stackrel{~}{}`$ and $`\stackrel{~}{}:\stackrel{~}{}`$ are isometric isomorphisms, and that $`\stackrel{~}{}=^1`$. This completes the proof of theorem 1.
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# Contracting the Wigner-Kernel of a Spin to the Wigner-Kernel of a Particle
## 1 Introduction
To represent quantum mechanics in terms of $`c`$-number valued functions has various appealing properties. It becomes possible to situate the quantum mechanical description of a system in a familiar frame, namely the phase space of its classical analog. Similarities and differences of the two descriptions can be visualized particularly well in such an approach. Further, from a structural point of view, to calculate expectation values of operators by means of ‘quasi-probabilities’ in phase space, is strongly analogous to the determination of mean values in classical statistical mechanics . The basic ingredient to set up such a symbolic calculus is a one-to-one correspondence between (self-adjoint) operators $`\widehat{A}`$ (acting on a Hilbert space $``$) on the one hand, and (real) functions $`W_A`$ defined on the phase-space $`\mathrm{\Gamma }`$ of the classical system on the other.
The quantum mechanics of spin and particle systems can be represented faithfully in terms of functions defined on the surface of a sphere with radius $`s`$, and on a plane, respectively. Intuitively, one expects these phase space-formulations to approach each other for increasing values of the spin quantum number since the surface of a sphere is then approximated by a plane with increasing accuracy. Therefore, appropriate Wigner functions of a spin, say, should go over smoothly into particle Wigner-functions in the limit of large $`s`$. It will be shown how this transition can be transformed in a rigorous and general way. The derivation is based on the group theoretical technique of contraction. The group $`SU(2)`$ of quantum mechanical rotations is contracted to the Heisenberg-Weyl group $`HW_1`$ associated with the particle. In this procedure, rotations go over into translations. Subsequently, the operator kernel which defines the spin Wigner-formalism in a condensed manner will be shown to contract to the operator kernel for a particle in the limit of infinite $`s`$.
## 2 Wigner-kernel for a particle
Consider a particle on the real line $`IR^1`$, with position and momentum operators satisfying $`[\widehat{q},\widehat{p}]=i\mathrm{}`$. The Stratonovich-Weyl correspondence, associating operators with functions in phase space, can be characterized elegantly by means of a kernel ,
$$\widehat{\mathrm{\Delta }}(\alpha )=2\widehat{T}(\alpha )\widehat{\mathrm{\Pi }}\widehat{T}^{}(\alpha ),\alpha =\frac{1}{\sqrt{2}}(q+ip)\mathrm{\Gamma }Cl,$$
(1)
which has an interpretation as a parity operator displaced by $`\alpha `$. The unitary
$$\widehat{T}(\alpha )=\mathrm{exp}[\alpha a^+\alpha ^{}a],$$
(2)
effects translations in phase space $`\mathrm{\Gamma }`$,
$$a\widehat{T}(\alpha )a\widehat{T}^{}(\alpha )=a\alpha ,$$
(3)
where $`a^{}a=(\widehat{q}i\widehat{p})/\sqrt{2}`$ and $`a^+=a^{}`$ are the standard annihilation and creation operators ($`\mathrm{}=1`$). At the origin $`\alpha =0`$, the kernel equals (two times) the unitary, involutive parity operator $`\widehat{\mathrm{\Pi }}`$,
$$\widehat{\mathrm{\Pi }}a\widehat{\mathrm{\Pi }}^{}=a,$$
(4)
corresponding to a reflection at the origin of $`\mathrm{\Gamma }`$. Using the number operator $`\widehat{N}=a^+a`$ and its eigenstates,
$$\widehat{N}|n=n|n,n=0,1,2,\mathrm{},$$
(5)
parity can be given a simple form which will be useful later,
$$\widehat{\mathrm{\Pi }}=\mathrm{exp}[i\pi \widehat{N}]=\underset{n=0}{\overset{\mathrm{}}{}}()^n|nn|.$$
(6)
The kernel $`\widehat{\mathrm{\Delta }}(\alpha )`$ can be derived from the Stratonovich-Weyl postulates which are natural conditions on a quantum mechanical phase-space representation. The correspondence between a (self-adjoint) operator $`\widehat{A}`$ and a (real) function is defined by
$$W_A(\alpha )=\text{ Tr }\left[\widehat{\mathrm{\Delta }}(\alpha )\widehat{A}\right],$$
(7)
while its inverse reads
$$\widehat{A}=_\mathrm{\Gamma }𝑑\alpha W_A(\alpha )\widehat{\mathrm{\Delta }}(\alpha ).$$
(8)
If $`\widehat{A}`$ is the density operator of a pure state, $`\widehat{\rho }=|\psi \psi |`$, the symbol defined in (7) is the Wigner function of the state $`|\psi `$,
$$W_\psi (p,q)=\frac{2}{h}_\mathrm{\Gamma }𝑑x\psi ^{}(q+x)\psi (qx)\mathrm{exp}[2ipx/\mathrm{}].$$
(9)
It is important to note that the kernel $`\widehat{\mathrm{\Delta }}(\alpha )`$ is entirely defined in terms of the operators $`a^\pm `$ and $`\widehat{N}`$, forming a closed algebra under commutation if the identity is included:
$$[a,a^+]=1,[\widehat{N},a^\pm ]=\pm a^\pm .$$
(10)
This algebra generates the Heisenberg-Weyl group $`HW_1`$, and the kernel $`\widehat{\mathrm{\Delta }}(\alpha )`$ is an element of it (apart from the factor of two).
## 3 Wigner-kernel for a spin
For a quantum spin, the symbol associated with an operator is a continuous function defined on the sphere $`𝒮^2`$, being the phase space of the classical spin. When setting up a phase-space formalism, rotations take over the role of translations. The group $`SU(2)`$ is generated by the components of the spin operator $`\widehat{𝐒}`$. The three operators $`\widehat{S}^\pm =(\widehat{S}^x\pm i\widehat{S}^y)`$ and $`\widehat{S}^z`$, satisfy the commutation relations
$$[\widehat{S}^+,\widehat{S}^{}]=2\widehat{S}^z,[\widehat{S}^z,\widehat{S}^\pm ]=\widehat{S}^\pm .$$
(11)
The standard basis
$$𝐧_z\widehat{𝐒}|s,m=m|s,m,m=s,\mathrm{},s,$$
(12)
is given by the eigenstates of the $`z`$ component $`\widehat{S}^z`$ of the spin.
For a quantum spin, it is natural to expect that the elements of the Wigner kernel will be labeled by points of the sphere $`𝒮^2`$, corresponding to unit vectors $`𝐧`$ $`=(\mathrm{sin}\vartheta \mathrm{cos}\phi ,`$$`\mathrm{sin}\vartheta \mathrm{sin}\phi ,\mathrm{cos}\vartheta )`$, parametrized by standard spherical coordinates. Replacing intuitively translations in (1) by rotations leads to the expression
$$\widehat{\mathrm{\Delta }}(𝐧)=\widehat{U}(𝐧)\widehat{\mathrm{\Pi }}_s\widehat{U}^{}(𝐧),$$
(13)
where
$$\widehat{U}(𝐧)=\mathrm{exp}[i\vartheta 𝐤\widehat{𝐒}]$$
(14)
with a unit vector $`𝐤=(\mathrm{sin}\phi ,\mathrm{cos}\phi ,0)`$ in the $`xy`$ plane. Thus, $`\widehat{U}(𝐧)`$ represents a finite rotation which maps the operator $`\widehat{S}^z=𝐧_z\widehat{𝐒}`$ into $`𝐧\widehat{𝐒}`$, i.e. $`𝐧_z𝐧`$. What are natural choices for the operator $`\widehat{\mathrm{\Pi }}_s`$?
Two possibilities come to one’s mind. First, try to transfer the concept of reflection about some point in phase space. Introduce canonical coordinates $`(q,p)=(\phi ,\mathrm{cos}\vartheta )`$ on the sphere. Then, ‘parity’ would correspond to the map $`(\phi ,\mathrm{cos}\vartheta )(\phi ,\mathrm{cos}\vartheta )`$, or $`(\phi ,\vartheta )(2\pi \phi ,\pi \vartheta )`$. This is just a rotation by $`\pi `$ about the $`x`$ axis. Since all points of the sphere are equivalent, one could also chose a rotation by $`\pi `$ about the $`z`$ axis as candidate for parity. Second, $`\widehat{\mathrm{\Pi }}_s`$ might be considered to generate reflections about the center of the sphere, $`𝐧𝐧`$, that is, $`(\phi ,\vartheta )(\phi +\pi ,\pi \vartheta )`$. It can be shown that both possibilities do not give rise to a symbolic calculus on the sphere , violating bijectivity between operators and phase-space functions, for example.
Nevetheless, acceptable operator kernels $`\widehat{\mathrm{\Delta }}_\epsilon (𝐧)`$ do exist as shown by Stratonovich , Várilly and Gracia-Bondía , and by Amiet and Cibils . For example, the condition that the kernel should satisfy appropriate Stratonovich-Weyl postulates implies that
$$\widehat{\mathrm{\Delta }}_\epsilon (𝐧)=\underset{m,m^{}=s}{\overset{s}{}}Z_{mm^{}}^\epsilon (𝐧)|s,ms,m^{}|.$$
(15)
The coefficients,
$$Z_{mm^{}}^\epsilon (𝐧)=\frac{\sqrt{4\pi }}{2s+1}\underset{l=0}{\overset{2s}{}}\epsilon _l\sqrt{2l+1}\begin{array}{ccc}s& l& s\\ m& m^{}m& m^{}\end{array}Y_{l,m^{}m}(𝐧),$$
(16)
where $`\epsilon _0=1`$ and $`\epsilon _l=\pm 1,l=1,\mathrm{},2s`$, are linear combinations of Clebsch-Gordan coefficients multiplied by spherical harmonics $`Y_{l,m}(𝐧),l=0,1,\mathrm{},2s,`$ $`m=l,\mathrm{},l`$. Note that there is no unique kernel but, due to the factors $`\epsilon _l`$, one can define $`2^{2s}`$ different Stratonovich-Weyl correspondence rules.
Unfortunatley, the expression (15) does not admit a simple interpretation of the operator in analogy to (1). It follows from an independent derivation of $`\widehat{\mathrm{\Delta }}(𝐧)`$ that (15) can be written in the form (13) where
$$\widehat{\mathrm{\Pi }}_s=\widehat{\mathrm{\Delta }}_\epsilon (𝐧_z)=\underset{m=s}{\overset{s}{}}\mathrm{\Delta }_\epsilon (m)|s,ms,m|,$$
(17)
with coefficients
$$\mathrm{\Delta }_\epsilon (m)=\underset{l=0}{\overset{2s}{}}\epsilon _l\frac{2l+1}{2s+1}\begin{array}{ccc}s& l& s\\ m& 0& m\end{array}.$$
(18)
Still, the operator $`\widehat{\mathrm{\Pi }}_s`$ does not have an obvious interpretation but a new strategy to justify its form emerges. Consider a plane tangent to the sphere at its north pole. For increasing radius, the sphere is approximated locally better and better by the plane. Therefore, one might expect that for $`s\mathrm{}`$ objects defined on the sphere turn into objects defined on the plane. It has been conjectured in that in this limit the Wigner kernel of a spin goes over into the kernel for a particle. It is the purpose of this paper to show that
$$\underset{s\mathrm{}}{lim}\widehat{U}(𝐧)\widehat{\mathrm{\Delta }}(𝐧_z)\widehat{U}^{}(𝐧)=\widehat{\mathrm{\Delta }}(\alpha ),$$
(19)
is indeed true for the kernel $`\widehat{\mathrm{\Delta }}_\epsilon (𝐧_z)`$ with parameters $`\epsilon _1=\epsilon _2=\mathrm{}=\epsilon _{2s}=1`$, denoted by $`\widehat{\mathrm{\Delta }}(𝐧_z)`$ for short. Thus, while the rotations $`\widehat{U}(𝐧)`$ should go over into translations, the operator $`\widehat{\mathrm{\Delta }}(𝐧_z)`$ corresponds, in one way or another, to parity for a spin. A convenient framewok to prove (19) is the contraction of groups as shown in the next section.
## 4 Contracting $`SU(2)`$ to $`HW_1`$
Introduce three operators $`\widehat{A}^\pm `$ and $`\widehat{A}^z`$ defined as linear combinations of the generators of the algebra $`su(2)`$ in polar form,
$$\widehat{A}^\pm =c\widehat{S}^{},\widehat{A}^z=\widehat{S}^z+\frac{1_s}{2c^2},$$
(20)
plus the identity $`1_s`$. This transformation is invertible for each value of the parameter $`c>0`$. The non-zero commutators of the new generators are given by
$$[\widehat{A}^{},\widehat{A}^+]=1_s2c^2\widehat{A}^z,[\widehat{A}^\pm ,\widehat{A}^z]=\widehat{A}^\pm .$$
(21)
These relations have a well defined limit if $`c0`$, nonwithstanding that the transformation (20) is not invertible for $`c=0`$. In fact, they reproduce the commutation relations (10) of the Heisenberg-Weyl algebra after identifying
$$\underset{c0}{lim}\widehat{A}^\pm =a^\pm ,\underset{c0}{lim}\widehat{A}^z=\widehat{N},\underset{c0}{lim}1_s=1.$$
(22)
How do rotations behave in this limit? Any finite rotation $`\widehat{U}(𝐧)SU(2)`$ in (14) can be written in the form
$$\widehat{U}(𝐧)=\mathrm{exp}\left[\xi _{}\widehat{S}^{}\xi _+\widehat{S}^+\right],\xi _{}=\frac{\vartheta }{2}e^{i\phi },\xi _+=\xi _{}^{},$$
(23)
or, expressed in terms of the operators (20),
$$\widehat{U}(𝐧)=\mathrm{exp}\left[c(\xi _{}\widehat{A}^+\xi _+\widehat{A}^{})\right].$$
(24)
Consequently, if the coefficients $`\xi _\pm `$ shrink with the parameter $`c`$ according to
$$\underset{c0}{lim}\frac{\xi _{}}{c}=\underset{c0}{lim}\frac{\vartheta e^{i\phi }}{2c}=\alpha ,\underset{c0}{lim}\frac{\xi _+}{c}=\underset{c0}{lim}\frac{\vartheta e^{i\phi }}{2c}=\alpha ^{},$$
(25)
a rotation $`\widehat{U}(𝐧)`$ tends to a well-defined element of the Heisenberg-Weyl group, Eq. (2):
$$\underset{c0}{lim}\widehat{U}(𝐧)=\widehat{T}(\alpha ).$$
(26)
For consistency, the limit $`c0`$ must correctly reproduce the eigenvalues of the operator $`\widehat{N}`$, given by the non-negative integers. Let us look at the fate of the eigenvalue equation (12) for $`m=s`$, which is expected to give $`\widehat{N}|n=0`$. One has
$$\underset{c0}{lim}\left(\widehat{A}^z|s,s\right)\underset{c0}{lim}\left[\left(\widehat{S}^z+\frac{1_s}{2c^2}\right)|s,s\right]=\underset{c0}{lim}\left(s+\frac{1}{2c^2}\right)\underset{c0}{lim}|s,s=0$$
(27)
implying $`2c^2s=1`$ for $`lim_{c0}|s,s=|n=0`$. Consequently, the radius of the sphere, $`s`$, increases with decreasing values of $`c`$. The state $`|s,s`$ turns indeed into the ground state associated with the operator $`\widehat{N}`$ since one has in general
$$\underset{c0}{lim}|s,m=\underset{c0}{lim}|s,sn=|n,n=smIN_0,$$
(28)
as follows from
$$\widehat{N}|n=\underset{c0}{lim}\left[\left(\widehat{S}^z+\frac{1_s}{2c^2}\right)|s,sn\right]=\underset{c0}{lim}\left((sm)+(\frac{1}{2c^2}s)\right)|n=n|n.$$
(29)
Now it is obvious why one needs to associate the creation operator $`\widehat{S}^+`$ with the annihilation operator $`a`$ (cf. (20)): the eigenstates with maximal $`s`$ are linked to the oscillator ground state with minimal $`n=0`$. In , a different convention has been used. Nevertheless, it remains true that not only spin eigenstates are mapped into number eigenstates but many other expressions related to the group $`U(2)`$ turn into an equivalent expression for the group $`HW_1`$.
This is good news for the present purpose to establish a relation between the Moyal formalism of a particle and a spin. Consider the limit of the kernel (13) under contraction using (26),
$$\underset{c0}{lim}\widehat{\mathrm{\Delta }}(𝐧)=\widehat{T}(\alpha )\left(\underset{c0}{lim}\widehat{\mathrm{\Pi }}_s\right)\widehat{T}^{}(\alpha ).$$
(30)
The middle term can be written as
$$\underset{c0}{lim}\widehat{\mathrm{\Pi }}_s=\underset{c0}{lim}\underset{m=s}{\overset{s}{}}\mathrm{\Delta }_\epsilon (m)|s,ms,m|=\underset{n=0}{\overset{\mathrm{}}{}}\left(\underset{c0}{lim}\mathrm{\Delta }_\epsilon (sn)\right)|nn|.$$
(31)
Upon comparison with (6), the Wigner kernel of a spin is seen to turn into the Wigner kernel of the particle if
$$\underset{s\mathrm{}}{lim}\underset{l=0}{\overset{2s}{}}\epsilon _l\left(\frac{2l+1}{2s+1}\right)^{1/2}\begin{array}{ccc}s& s& l\\ sn& ns& 0\end{array}=2$$
(32)
holds for all non-negative integers $`n`$. In the next section, this will be shown to be true for the choice $`\epsilon _l=+1`$, $`l=1,\mathrm{}2s`$.
## 5 Summing the series
Evaluating the sum (32) in the limit $`s\mathrm{}`$ proceeds in two steps. First, the asymptotic form of the terms
$$\mathrm{\Delta }_{l,n}^s=\left(\frac{2l+1}{2s+1}\right)^{1/2}\begin{array}{ccc}s& s& l\\ sn& ns& 0\end{array}$$
(33)
to be summed is determined with the help of a recurrence formula for Clebsch-Gordan coefficients. Then, the sums are transformed into integrals which can be evaluated. All approximations drop terms of the order $`1/s`$ at least, hence the result is exact in the limit of infinite $`s`$.
Clebsch-Gordan coefficients satisfy the following recursion relation :
$`[l(l+1)2s(s+1)+2m^2]\begin{array}{ccc}s& s& l\\ m& m& 0\end{array}`$ (36)
$`=`$ $`[s(s+1)m(m+1)]\begin{array}{ccc}s& s& l\\ m+1& (m+1)& 0\end{array}`$ (42)
$`+[s(s+1)m(m1)]\begin{array}{ccc}s& s& l\\ m1& (m1)& 0\end{array},`$
implying that
$`(n+1)\left(1{\displaystyle \frac{n+1}{2s+1}}\right)\mathrm{\Delta }_{l,n+1}^s+\left(2n+1{\displaystyle \frac{2n^2+2n+1}{2s+1}}\right)\mathrm{\Delta }_{l,n}^s+`$
$`+n\left(1{\displaystyle \frac{n}{2s+1}}\right)\mathrm{\Delta }_{l,n1}^s={\displaystyle \frac{l(l+1)}{2s+1}}\mathrm{\Delta }_{l,n}^s.`$ (43)
For any finite $`n`$ the terms subtracted on the left-hand-side become less and less important if $`s\mathrm{}`$. Assume now that one can write the terms with large values of $`n`$ in the form
$$\mathrm{\Delta }_{l,n}^s(x_l)=\mathrm{\Lambda }_n(x_l)\mathrm{\Delta }_{l,0}^s,\mathrm{\Lambda }_0(x_l)=1,x_l=\frac{l(l+1)}{2s+1}.$$
(44)
The polynomial $`\mathrm{\Lambda }_n(x_l)`$ of order $`n`$ in $`x_l`$ satisfies a three-term recursion relation,
$$(n+1)\mathrm{\Lambda }_{n+1}(x_l)+(2n+1)\mathrm{\Lambda }_n(x_l)+n\mathrm{\Lambda }_{n1}(x_l)=x_l\mathrm{\Lambda }_n(x_l),$$
(45)
where terms of order $`1/s`$ have been dropped in (45). Its solutions are proportional to the Laguerre polynomials, and the ‘normalization’ condition $`\mathrm{\Lambda }_0(x_l)=1`$ implies that
$$\mathrm{\Lambda }_n(x_l)=()^nL_n(x_l)=()^n\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)\frac{(x_l)^k}{k!},n=0,1,2,\mathrm{}$$
(46)
The term $`\mathrm{\Delta }_{l,0}^s`$ in (44) can be determined in the following way. If $`s`$ is large, one writes for each finite $`k`$
$$\left(1\frac{k}{2s+1}\right)^{2s+1}\mathrm{exp}[k],$$
(47)
which leads to the approximation
$`\mathrm{\Delta }_{l,0}^s`$ $`=`$ $`\left({\displaystyle \frac{2l+1}{2s+1}}\right)^{1/2}\begin{array}{ccc}s& s& l\\ s& s& 0\end{array}={\displaystyle \frac{2l+1}{2s+1}}\left({\displaystyle \frac{(2s)!}{(2sl)!}}{\displaystyle \frac{(2s)!}{(2s+l+1)!}}\right)^{1/2}`$ (50)
$`=`$ $`{\displaystyle \frac{2l+1}{2s+1}}\left({\displaystyle \frac{\mathrm{\Pi }_{k=0}^l(1k/(2s+1))}{\mathrm{\Pi }_{k=0}^l(1+k/(2s+1))}}\right)^{1/2}{\displaystyle \frac{2l+1}{2s+1}}\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{l(l+1)}{2s+1}}\right].`$ (51)
Collecting the results, one has
$$\underset{s\mathrm{}}{lim}\underset{l=0}{\overset{2s}{}}\mathrm{\Delta }_{l,n}^s()^n\underset{s\mathrm{}}{lim}\underset{l=0}{\overset{2s}{}}\mathrm{\Delta }x_lL_n(x_l)e^{x_l/2},$$
(52)
where $`\mathrm{\Delta }x_l=(x_{l+1}x_l)=(2l+1)/(2s+1)+𝒪(1/s)`$. Transforming now the Riemann sum into an integral, one obtains the final result
$$\underset{s\mathrm{}}{lim}\underset{l=0}{\overset{2s}{}}\mathrm{\Delta }_{l,n}^s=()^n_0^{\mathrm{}}𝑑xL_n(x)e^{x/2}=2,$$
(53)
using the formula
$$_0^{\mathrm{}}𝑑xL_n(x)e^{x/t}=t(1t)^n,$$
(54)
for $`t=2`$. This identity is proven easily by means of the expansion in (46).
## 6 Discussion
Starting from a new form of the kernel defining the familiar Wigner formalism for a spin, its limit for infinite values of $`s`$ has been shown to be the Wigner kernel of a particle. As the kernel defines entirely a phase-space representation, this result guarantees that the Moyal formalism for a particle is reproduced automatically and in toto, if the limit $`s\mathrm{}`$ of the spin Moyal formalism is taken.
In fact, slightly more has been shown. The result removes an ambiguity of the Moyal formalism for a spin: the Stratonovich-Weyl postulates are compatible with a discrete family of $`2^{2s}`$ distinct kernels $`\widehat{\mathrm{\Delta }}_\epsilon (𝐧)`$. However, only one of these kernels turns into the particle kernel. This kernel had been singled out before for other reasons . In summary, the group theoretical contraction shows that the phase-space representations à la Wigner for spin and particle systems are structurally equivalent.
### Acknowledgements
St. W. acknowledges financial support by the Schweizerische Nationalfonds.
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# Linear response conductance and magneto-resistance of ferromagnetic single-electron transistors
## I Introduction
The spin of the electron can provide new functionality in electronic devices and spin-polarized transport is therefore an active research field. Spin-transport was pioneered by Tedrow and Meservey and the discovery of the giant magnetoresistance effect (GMR) in metallic multilayers has motivated a substantial re-newed interest in this field. Recently spin-polarized current injection in all-semiconductor devices and carbon nanotubes have been realized. The fabrication of tunnel junctions made of two ferromagnetic leads separated by an insulating layer has become well under control. These technological advancements can make it possible to design new generation of electronic devices such as magnetic RAM’s, sensors and ultimately perhaps quantum computers.
In this article, we will focus on an elementary structure, the ferromagnetic single-electron transistor and discuss the role of the ferromagnetic electrodes, spin-accumulation and Coulomb charging effects on the transport properties. A schematic picture of the system under consideration is shown in Fig. 1. A normal (ferromagnetic) metal island is coupled to two ferromagnetic reservoirs by tunnel junctions. There is a capacative coupling $`C`$ between the metal island and the metal reservoirs. The island is so small so that the charging energy $`E_c=e^2/2C`$ associated with the addition of a single electron to the island can be larger than both the temperature and the applied source-drain bias. The electrostatic potential on the island can be controlled by a gate voltage $`V_g`$ capacitively coupled via a capacitance $`C_g`$ to the island. We consider an island weakly coupled to the reservoirs by many channels with transmission probabilities much less than unity, but the total tunnel conductance $`G`$ can be smaller or larger than the quantum conductance $`G_K`$. In the simplest case, the single-electron transistor (SET) comprising of tunnel junctions with junction resistances much larger than the quantum resistance $`R_K=h/e^2`$ exhibits Coulomb blockade of the single electron tunneling at low temperatures when the gate voltage vanishes, and a conductance peak when $`C_gV_g=e/2`$. For spin-polarized charge transfer the different transport properties of the spin-up and spin-down electrons have to be taken into account. Therefore the behavior of the SET is determined by the interplay between spin-dependent transport and Coulomb charging effects.
All-ferromagnetic (F-F-F) SET’s have recently been investigated experimentally and theoretically. The tunneling magneto-resistance (TMR) in F-F-F systems is defined as the relative conductance difference by changing the magnetization of the ferromagnetic island from being parallel to anti-parallel to the magnetizations of the ferromagnetic reservoirs (which are kept parallel). Ferromagnetic islands have a short spin-flip relaxation time and several works have computed the transport properties of F-F-F systems disregarding spin-accumulation. An enhanced TMR in the Coulomb blockade regime was found experimentally, and it was sub-sequentially shown that the TMR could be doubled due to co-tunneling and that higher order tunneling processes could give even larger enhancements of the TMR. In what follows, we will use a non-perturbative approach to calculate the tunneling magneto-resistance (TMR) of the F-F-F SET showing that even larger enhancements of the TMR than what has been reported in Ref. is possible. We will also discuss the TMR in the ’on’-state when the Coulomb blockade effect is lifted. It will be demonstrated that the TMR in this regime is a non-monotonic function of the temperature.
The physics of the TMR in F-N-F systems is completely different from the mechanism of the TMR in F-F-F systems. F-F-F systems exhibit a finite TMR even in the absence of any spin-accumulation on the ferromagnetic island. The transport properties can be modeled by an all-normal metal device with magnetization configuration dependent tunnel conductances. In contrast, the TMR of F-N-F systems is uniquely related to the spin-accumulation on the island. Spin-accumulation governs the transport properties and leads to a measurable TMR when the spin-flip relaxation time $`\tau _{\text{sf}}`$ is longer than the dwell time so that
$$\tau _{\text{sf}}/(\chi h)>G_K/G,$$
(1)
where $`\chi `$ is the spin-susceptibility, $`G`$ is the tunnel conductance and $`G_K=e^2/h`$ is the quantum conductance. In order to realize spin-accumulation a small spin-susceptibility is required (1) which implies a small normal metal island with a large Coulomb charging energy $`E_c`$ and consequently Coulomb blockade effects should be taken into account. When the tunnel conductance $`G`$ is much smaller than the quantum conductance $`G_K`$ and the temperature is not too low, transport can be described by sequential tunneling processes and this was performed in Refs. . However, Eq. (1) shows that spin-accumulation is easier to realize when the tunnel conductance is large, $`G>G_K`$. In this strong tunneling limit the quantum fluctuations of the charge on the island as well as the spin-acccumulation should be taken into account with a theoretical description that goes beyond the low order perturbative treatment presented in Refs. . The tunneling magneto-resistance (TMR) in F-N-F systems is defined in a different way than for (F-F-F) systems: The TMR is the relative conductance difference on going from a configuration where the magnetizations in the ferromagnetic reservoirs are parallel to a configuration where the magnetizations are anti-parallel $`(g_\text{P}g_{\text{AP}})/g_\text{P}`$. Changing the magnetization configuration in F-N-F systems only changes the spin-dependence of the tunnel conductances but conserves the total tunnel conductances. Hence, a much weaker coupling to the charge degrees of freedom of the island and a reduced influence of Coulomb charging should be expected. In the Coulomb blockade regime at low temperatures the total conductance through the system is suppressed. Consequently, also the transport of spins into the island decreases and at sufficiently low temperatures the dwell time is longer than the spin-flip relaxation time so that (1) is no longer satisfied. The spin-accumulation and thus the TMR must therefore be reduced in the Coulomb blockade regime which is contrary to the behavior of the TMR in F-F-F systems. We will show that the linear response TMR in F-N-F systems in the Coulomb blockade regime equals or is smaller than the classical high-temperature TMR, which is in contrast to the results for F-F-F systems.
We present in this paper a unified analytical expression for the linear response conductance of F-F-F and F-N-F SET’s with dirty metallic islands. Our formula is valid to all orders of the ratio between the tunnel conductances and the quantum conductance, for arbitrary spin-flip relaxation times and takes into account general band-structures for the ferromagnetic reservoirs and the normal/ferromagnetic metal island. Our results are in this sense general. The analytical expression for the linear response conductance is used to discuss the magneto-resistance of F-F-F and F-N-F SET’s.
The paper is organized in the following way. The system is described and the model Hamiltonian introduced in Section II. In Section III the linear response conductance is calculated, including the multi-band effect of the electrodes and the spin-accumulation on the central island. The TMR is evaluated in Section IV in the case of ferromagnet-normal metal-ferromagnet SET’s and in Section V in the case of all-ferromagnet SET’s. Section VI concludes the paper.
## II Model
We consider a ferromagnetic single-electron transistor comprising of a normal/ferromagnetic metal island connected to ferromagnetic leads by tunnel junctions. The central island is capacitively coupled to a gate voltage $`V_g`$ via a capacitance $`C_g`$. The Hamiltonian of such a device is
$$\widehat{H}=\left(\widehat{H}_l+\widehat{H}_i+\widehat{H}_r\right)+\widehat{H}_c+\left(\widehat{H}_{li}+\widehat{H}_{ri}+\text{h.c.}\right).$$
(2)
The quasi-particles in the left reservoir are described by
$$\widehat{H}_l=\underset{𝐥\sigma nm}{}W_{𝐥\sigma }^{nm}\widehat{c}_{𝐥\sigma }^n\widehat{c}_{𝐥\sigma }^m,$$
(3)
where $`n`$ and $`m`$ are the band-indices, $`\sigma `$ is the electron spin, $`𝐥`$ is the momentum in the left reservoir, $`W_{𝐥\sigma }^{nm}`$ determine the quasi-particle energy bands and the hybridization between the bands for a given spin, and $`\widehat{c}_𝐥^n`$ creates an electron with spin $`\sigma `$ and momentum $`𝐥`$ in band $`n`$ in the left reservoir. The hat ($`\widehat{}`$) denotes an operator. There are similar Hamiltonians for the quasi-particles on the island ($`li`$) and the quasi-particles in the right reservoir ($`lr`$). The Coulomb charging effects are included via
$$\widehat{H}_c=\frac{e^2}{2C}\left(\underset{\sigma }{}\widehat{N}_\sigma ^in_{\text{ex}}\right)^2,$$
(4)
where the induced electron number is $`n_{\text{ex}}=C_gV_g/e`$ and the spin-dependent number of excess electrons on the island is $`\widehat{N}_\sigma ^i=_{𝐢n}\widehat{c}_{𝐢\sigma }^n\widehat{c}_{𝐢\sigma }^n`$. We will below also make use of the spin-dependent excess number of electrons in the left (right) reservoir, $`\widehat{N}_\sigma ^l=_{𝐥n}\widehat{c}_{𝐥\sigma }^n\widehat{c}_{𝐥\sigma }^n`$ ($`\widehat{N}_\sigma ^r=_{𝐫n}\widehat{c}_{𝐫\sigma }^n\widehat{c}_{𝐫\sigma }^n`$). The tunneling between the reservoirs and the island is taken into account via the tunneling Hamiltonians $`\widehat{H}_{pi}=_\sigma \widehat{H}_{pi\sigma }`$ ($`p=l`$ denotes left and $`p=r`$ denotes right):
$$\widehat{H}_{pi\sigma }=\underset{\mathrm{𝐩𝐢}nm}{}t_{\mathrm{𝐩𝐢}\sigma }^{nm}\widehat{c}_{𝐩\sigma }^n\widehat{c}_{𝐢\sigma }^m.$$
(5)
The tunneling matrix elements $`t_{\mathrm{𝐩𝐢}\sigma }^{nm}`$ allow tunneling between different bands with the same spin. Spin-flip scatterings during the tunneling processes are disregarded, but can be included in the formalism leading to a renormalization of the spin-asymmetries of the junction resistances.
We consider the situation when the temperature is much larger than the level spacing and the level spacing is much smaller than the Coulomb charging energy. The level spacing is consequently taken to be continous. Co-tunneling and higher order tunneling processes have both elastic and inelastic contributions, e.g. in the case of co-tunneling the elastic contribution means that the same electron tunnels through both of the two junctions whereas the inelastic contribution corresponds to two different electrons that tunnel in the two junctions: One tunnels into the island above its Fermi level, and another jumps out of the island from below the Fermi level. The elastic contribution is sensitive to the phase of the electron on the island and consequently becomes vanishingly small for dirty systems when the characteristic time of tunneling through the macroscopic barrier caused by the Coulomb energy $`\mathrm{}/E_c`$ ($`E_c=e^2/(2C)`$) is larger than the classical time $`L^2/D`$ for diffusion through the island ($`L`$ is the size of the island and $`D`$ is the diffusion coefficient). We consider the inelastic transport regime exclusively and disregard the elastic contributions to the tunneling processes, i.e. for an estimate of a diffusion constant on the island of $`D=10cm^2/s`$ and a charging energy of $`E_c=1K`$ we assume that the size of the island $`L<0.1\mu m`$ (in reality there is also scattering at the boundary of the island that further relax this assumption). The occupation of the energy levels on the island can then be described by an energy $`ϵ`$ and spin $`\sigma `$ dependent non-equilibrium distribution $`f_\sigma ^i(ϵ)`$. The non-equilibrium distribution $`f_\sigma ^i(ϵ)`$ is spin-dependent to allow a non-equilibrium spin-accumulation and only equals the equilibrium Fermi-Dirac distribution $`f(ϵ\mu )`$ when the bias voltage vanishes ($`\mu `$ is the equlibrium chemical potential). Furthermore, the leads and the island are metallic so that the tunnel conductances (see Eq. (17) below) are energy independent on the scale of the Coulomb charging energy which is much smaller than the Fermi energy. We consider the linear response regime where the bias voltage is the smallest energy scale in the system. Under these assumptions, only the total occupation of spin-up and spin-down electrons on the island, $`𝑑ϵ[f_\sigma ^if(ϵ\mu )]`$, determine the macroscopic Coulomb barrier and the total effective tunneling rates. Consequently to all orders in perturbation theory of the tunneling Hamiltonian, we can describe the electron occupation on the island by a local chemical potential $`\mu _\sigma ^i\mu +𝑑ϵ[f_\sigma ^i(ϵ)f(ϵ\mu )]`$ since the energy-dependence of the distribution $`f_\sigma ^i(ϵ)`$ does not influence the transport properties. It is thus e.g. irrelevant if there is substantial energy-relaxation on the island so that the non-equilibrium distribution $`f_\sigma ^i(ϵ)`$ approaches the Fermi-Dirac distribution $`f(ϵ\mu _\sigma ^i)`$ or if there is no energy-relaxation so that the non-equilibrium distribution $`f_\sigma ^i(ϵ)`$ is a linear combination of the Fermi-Dirac distributions in the left and right reservoirs with different local chemical potentials.
The electrons in the right reservoir and in the left reservoir are in local equilibrium with chemical potentials $`\mu ^l(t)`$ and $`\mu ^r(t)`$, respectively and as seen above the electrons on the island can be described as in local equilibrium with a spin-dependent chemical potential $`\mu _\sigma ^i(t)`$. The temperature is the same in all subsystems. The thermal average of the spin-dependent current from the left ($`p=l`$) or right ($`p=r`$) reservoir to the island can then be written as
$$I_\sigma ^p(t)=ie\text{Tr}\left[\left(\widehat{H}_{pi\sigma }(t)\widehat{H}_{pi\sigma }^{}(t)\right)\mathrm{exp}(\beta K(t))\right]$$
(6)
with the non-equilibrium time-dependent grand canonical potential
$$\widehat{K}(t)=\widehat{H}\underset{\sigma }{}[\mu ^l(t)\widehat{N}_\sigma ^l+\mu ^r(t)\widehat{N}_\sigma ^r+\mu _\sigma ^i(t)\widehat{N}_\sigma ^i].$$
(7)
The spin-dependent chemical potentials on the island $`\mu _{}^i`$ and $`\mu _{}^i`$ should be determined from the flux of particles and spins into the island. The spin-flux into the island is
$$(I_{}^l+I_{}^r)(I_{}^l+I_{}^r)=\frac{es}{\tau _{\text{sf}}}.$$
(8)
Current conservation through the system requires
$$(I_{}^l+I_{}^l)+(I_{}^r+I_{}^r)=0.$$
(9)
The excess number of spins on the island $`s`$ is related to the non-equilibrium difference in the chemical potentials by the spin-susceptibility $`\chi _s`$ via $`s=\chi _s(\mu _{}^i\mu _{}^i)`$. For non-interacting electrons $`\chi _s=D`$, where $`D`$ is the density of states. The expression for the current (6) and the conservation laws for spins (8) and particles (9) uniquely determine the spin-dependent local chemical potentials on the island, $`\mu _i`$ and $`\mu _i`$, and consequently the current through the device. We will focus on the linear response regime in what follows.
## III Linear response conductance
In the linear response regime a perturbation expansion in terms of the small differences in the chemical potentials on the left reservoir, on the island and on the right reservoir can be performed: $`\mu ^l(t)=\mu +\delta \mu ^l(t)`$, $`\mu ^r(t)=\mu +\delta \mu ^r(t)`$, and $`\mu _\sigma ^i(t)=\mu +\delta \mu _\sigma ^i(t)`$. The non-equilibrium grand canonical potential is $`\widehat{K}(t)=\widehat{K}_0+\delta \widehat{K}(t)`$, where the time-independent unperturbed grand canonical potential is $`\widehat{K}_0=\widehat{H}\mu _\sigma (\widehat{N}_\sigma ^l+\widehat{N}_\sigma ^r+\widehat{N}_\sigma ^i)`$ and the time-dependent perturbation is $`\delta \widehat{K}=_\sigma (\delta \mu ^l\widehat{N}_\sigma ^l+\delta \mu ^r\widehat{N}_\sigma ^r+\delta \mu _\sigma ^i\widehat{N}_\sigma ^i)`$. We perform a time-dependent unitary transformation with the unitary matrix $`\widehat{U}(t)=\mathrm{exp}[i_{\mathrm{}}^t𝑑t_1\delta \widehat{K}(t_1)]`$. The transformed grand canonical potential becomes $`\widehat{K}_t(t)=\widehat{K}_0+\delta \widehat{H}(t)`$, where to the lowest order in the non-equilibrium chemical potential differences
$$\delta \widehat{H}(t)=\frac{1}{e}\underset{qs}{}\widehat{I}_s^q(t)_{\mathrm{}}^t𝑑t_1\left[\delta \mu ^q(t_1)\delta \mu _s^i(t_1)\right]$$
(10)
and the sum over the parameter $`q`$ corresponds to the left reservoir $`q=l`$ and the right reservoir $`q=r`$. The spin-currents through the left and the right barriers can be found from a standard linear response calculation:
$$I_\sigma ^p(\omega )=\frac{i}{e\omega }\underset{qs}{}\mathrm{\Pi }_{\sigma s}^{pq}(\omega )\left[\delta \mu ^q(\omega )\delta \mu _s^i(\omega )\right],$$
(11)
where the retarded current-current correlation function is
$$\mathrm{\Pi }_{\sigma s}^{pq}(t)=i[\widehat{I}_\sigma ^p(t),\widehat{I}_s^q(0)]_{\text{eq}},$$
(12)
This correlation function will be evaluated using path integrals. The computation is similar to the case of ultrasmall tunnel junctions or metallic SET that has been studied extensively. The main steps in the calculation of the correlation function is as follows. We use the equilibrium finite-temperature Matsubara formalism to find the retarded current-current correlation function. The imaginary time generating functional is defined as
$$Z[\left\{\eta \right\}]=D[\{c^{},c\}]D[\phi ]e^{S_g},$$
(13)
with the generalized action
$`S_g`$ $`=`$ $`{\displaystyle _0^\beta }𝑑\tau \text{Tr}_{𝐩n\sigma }c^{}(\tau )_\tau c(\tau )`$ (15)
$`+{\displaystyle _0^\beta }𝑑\tau K_0(\tau ){\displaystyle _0^\beta }𝑑\tau {\displaystyle \underset{p\sigma }{}}I_\sigma ^p(\tau )\eta _\sigma ^p(\tau ).`$
A Hubbard-Stratonovich transformation is performed to reduce the quadratic term in the charging energy to a linear coupling of the particle number operator with an auxiliary field. After integrating out the electronic degrees of freedom and including the principal contribution of the tunneling processes the generalized effective action is derived. The current-current correlation function can then be found by performing the functional derivative $`\mathrm{\Pi }_{\sigma s}^{pq}(\tau ,\tau ^{})=Z^1^2Z/(\eta _\sigma ^p(\tau )\eta _s^q(\tau ^{}))|_{\eta 0}`$. The resulting Fourier transform of the current-current correlation function is
$$\frac{1}{\omega }\mathrm{\Pi }_{\sigma s}^{pq}=\frac{e^2}{h}\left[\delta _{p\sigma ,qs}g_\sigma ^pg_0+2g_\sigma ^pg_s^qg_1\right].$$
(16)
The dimensionless tunnel conductance ($`p=l,r`$) is
$$g_\sigma ^p=\underset{jj^{}}{}D_{p\sigma }^jD_{i\sigma }^j^{}\left|\stackrel{~}{t}_{p\sigma }^{jj^{}}\right|^2,$$
(17)
Here $`D_{l\sigma }^j`$, $`D_{r\sigma }^j`$ and $`D_{i\sigma }^j`$ are the spin and band dependent densities of states and the renormalized transmission coefficient at the Fermi energy is $`\stackrel{~}{t}_{p\sigma }=\left(U_{p\sigma }\right)^{}t_{p\sigma }U_{i\sigma }`$, where $`U_{p\sigma }`$ ($`U_{i\sigma }`$) is the material-specific matrix that diagonalizes the band Hamiltonian in reservoir $`p`$ (island) for electrons with spin $`\sigma `$. In a single-band model the conductance is simply proportional to the density of states in the reservoir and the island, $`g_\sigma ^p=D_{p\sigma }D_{i\sigma }|t_\sigma |^2`$. The relation is in general more complicated as shown in (17). The correlation function $`g_0(i\omega _n)`$ is
$$g_0=\frac{4\pi }{i\omega _n}𝑑\tau e^{i\omega _n\tau }\chi (\tau )\mathrm{\Gamma }(\tau )$$
(18)
with the phase-phase correlation function
$`\mathrm{\Gamma }(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}e^{2\pi ikn_{\text{ex}}}\times `$ (20)
$`{\displaystyle _{b_k}}D[\phi ]e^{S[\phi ]}\mathrm{cos}\left[\phi (\tau )\phi (0)\right].`$
The effective action is
$`S[\phi ]`$ $`=`$ $`{\displaystyle 𝑑\tau \frac{\left[\dot{\phi }(\tau )\right]^2}{4E_c}}`$ (22)
$`g^{\text{cl}}{\displaystyle 𝑑\tau 𝑑\tau ^{}\chi \left(\tau \tau ^{}\right)\mathrm{cos}\left[\phi (\tau )\phi (\tau ^{})\right]}.`$
with the boundary condition $`b_k\phi (\beta )=\phi (0)+2\pi k`$. Since the leads and island are (ferro-)metallic with continuous spectra, the damping kernel $`\chi (\tau )`$ is an even function of the imaginary time $`\tau `$ with a period $`\beta `$ and the Matsubara-Fourier components $`\chi (\omega _l)=|\omega _l|/4\pi `$. The coupling of the island to the reservoirs is determined by the algebraic total dimensionless classical conductance $`g^{\text{cl}}=_{p\sigma }g_\sigma ^p`$. The partition function is $`Z=_k_{b_k}D[\phi ]\mathrm{exp}(S[\phi ])`$. Another correlation function in (16) is
$`g_1`$ $`=`$ $`{\displaystyle \frac{4\pi }{i\omega _n}}{\displaystyle }d\tau e^{i\omega _n\tau }{\displaystyle \frac{1}{Z}}\times `$ (24)
$`{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}e^{2\pi ikn_{\text{ex}}}{\displaystyle _{b_k}}D[\phi ]e^{S[\phi ]}I_t(\phi ,\tau )I_t(\phi ,0)`$
where
$$I_t(\phi ,\sigma )=𝑑\tau \chi \left(\sigma \tau \right)\mathrm{sin}\left[\phi (\sigma )\phi (\tau )\right].$$
(25)
Due to the current conservation (9) the correlation function $`g_1`$ (24) does not appear in the final form of the stationary current.
Knowing the current-current correlation function, we can find the non-equilibrium chemical potential on the island by using the conservation laws for spins (8) and particles (9). We find the spin-accumulation
$$\frac{\mu _{}^i\mu _{}^i}{\mu ^l\mu ^r}=\frac{(g_{}^lg_{}^rg_{}^lg_{}^r)/(g_{}g_{})}{1+g_{\text{sf}}\left[(g_0g_{})^1+(g_0g_{})^1\right]}$$
(26)
and the average chemical potential
$$\frac{\mu _{}^i+\mu _{}^i}{2}=\frac{g^l\mu ^l+g^r\mu ^r}{g^l+g^r}\frac{g_{}g_{}}{g_{}+g_{}}\frac{\mu _{}^i\mu _{}^i}{2},$$
(27)
where the spin-flip conductance $`g_{\text{sf}}`$ is related to the spin-susceptibility $`\chi _s`$ and the spin-flip relaxation time $`\tau _{\text{sf}}`$ via $`g_{\text{sf}}=\pi \chi _s/\tau _{\text{sf}}`$. From the non-equilibrium chemical potentials on the island we obtain the current through the systems by using the relation for the current (11) and finally we calculate the conductance $`g=eI/(\delta \mu ^l\delta \mu ^r)`$: The linear conductance in units of the quantum conductance $`e^2/h`$ of the FSET is
$$g=\frac{g^rg^l}{g^r+g^l}g_0\left[1\frac{(g_{}^lg_{}^rg_{}^lg_{}^r)^2/(g^rg^lg_{}g_{})}{1+g_{\text{sf}}\left[(g_0g_{})^1+(g_0g_{})^1\right]}\right],$$
(28)
where $`g_{}^l`$ and $`g_{}^l`$ ($`g_{}^r`$ and $`g_{}^r`$) are the spin-up and spin-down conductances of the left (right) junction in units of the quantum conductance $`e^2/h`$, and we have also introduced the combined conductances $`g^l=g_{}^l+g_{}^l`$, $`g^r=g_{}^r+g_{}^r`$, $`g_{}=g_{}^l+g_{}^r`$, and $`g_{}=g_{}^l+g_{}^r`$. The retarded correlation function $`g_0`$ (18) is the same as for an all-normal metal single-electron transistor. It renormalizes the conductance due to the Coulomb charging effects and depends on the induced electron number $`C_gV_g/e`$, the ratio of the charging energy to the temperature $`\beta E_c`$, and the coupling between the island and the reservoirs through the total conductance $`g^r+g^l`$. Through Eq. (28) we have thus reduced the problem of finding the conductance of ferromagnetic single-electron transistors to a well-studied problem of the renormalized conductance in all-normal metallic single-electron transistors. The second term in the bracket of (28) is a signature of the spin-accumulation and vanishes when $`\tau _{\text{sf}}0`$ ($`g_{\text{sf}}\mathrm{}`$). The spin-accumulation causes a reduction of the conductance. The resulting conductance (28) can be understood in terms of the equivalent circuit shown in Fig. (2), which is identical to the high-temperature classical circuit for two spins, but now generalized to include Coulomb charging effects by renormalizing the conductnaces so that $`g_{}^lg_0g_{}^l`$, $`g_{}^lg_0g_{}^l`$, $`g_{}^rg_0g_{}^r`$ and $`g_{}^rg_0g_{}^r`$. The spin-flip conductance $`g_{\text{sf}}`$ is not renormalized, since spin-flip processes do not change the number of electrons on the island and consequently not the charging energy. We will now use (28) to find the TMR ratio in F-N-F and F-F-F SET’s.
## IV F-N-F magneto-resistance
We first focus on F-N-F systems. The quantity of interest is the TMR $`(g_\text{P}g_{\text{AP}})/g_\text{P}`$. When the magnetizations in the left and right leads are parallel, the tunnel conductances can be characterized as $`g_s^l=g^l(1+sP_l)/2`$ and $`g_s^r=g^r(1+sP_r)/2`$, where $`s=+1`$ ($`s=1`$) denotes spin-up (spin-down) and $`P_l=(g_{}^lg_{}^l)/(g_{}^l+g_{}^l)`$ and $`P_r=(g_{}^rg_{}^r)/(g_{}^r+g_{}^r)`$ denote the polarizations of the left and right tunnel junctions, respectively. In the anti-parallel configuration the tunnel conductances are $`g_s^l=g^l(1+sP_l)/2`$ and $`g_s^r=g^r(1sP_r)/2`$. The total conductances of each junction, $`g^l=g_{}^l+g_{}^l`$ and $`g^r=g_{}^r+g_{}^r`$, are thus independent of the magnetization configuration. Hence, the correlation function $`g_0`$ and consequently also the prefactor in (28) remain invariant on going from the parallel to the anti-parallel configuration. The magneto-resistive effect is thus solely due to the spin-accumulation which is contained in the second term of Eq. (28). The TMR of F-N-F systems is
$$\frac{g_\text{P}g_{\text{AP}}}{g_\text{P}}=\frac{P_lP_r}{1+\alpha _g^2\beta _g^2}(1\gamma _g^2),$$
(29)
where the parameter $`\alpha _g^2=4g_{\text{sf}}/[g_0(g^l+g^r)]`$ ($`\alpha _g^20`$) reduces the spin-accumulation and hence the TMR due to spin-flip relaxation, $`\gamma _g=(g^lg^r)/(g^l+g^r)`$ ($`0\gamma _g^21`$) reduces the TMR due to the asymmetries in the conductances of the left and right junctions, and $`\beta _g^2=(P_l^2(1+\gamma _g)^2+P_r^2(1\gamma _g)^22P_lP_r(1\gamma _g^2))/4`$ ($`0\beta _g^21`$). The TMR is proportional to the polarization of the left and the right tunnel conductance. This result (29) reduced to the high temperature or large bias result in Ref. when $`g_01`$.
The spin-flip relaxation time is governed by spin-orbit coupling and magnetic impurities at low temperatures. We assume that $`\tau _{\text{sf}}`$ is temperature-independent below the Coulomb charging energy (at very low temperatures this assumption can become invalid, but that is a regime beyond the investigation in this article). Firstly, since the correlation function $`g_0`$ can only decrease due to Coulomb charging effects the factor $`\alpha _g`$ appearing in (29) can only increase at low temperatures and the TMR in the Coulomb blockade regime is equal to or smaller than the classical high-temperature TMR. Secondly, when the spin-flip relaxation time is much larger than the dwell time, the TMR is independent of the charging energy, temperature and tunnel resistances for practically all temperatures since $`\alpha _g0`$ except at extremely low temperatures. Thirdly, in this regime we expect the nonlinear response TMR to have a re-entrant behavior: The linear response TMR is identical to the TMR when the bias is much larger than the Coulomb charging energy which equals the classical high-temperature TMR. In the intermediate bias regime effects of the Coulomb charging energy appear in the sequential tunneling limit and an enhancement of the TMR has been demonstrated in Ref. using the co-tunneling formalism. Finally (28) and (29) can be generalized to finite frequencies by letting $`g_0g_0(\omega )`$ and $`1/\tau _{\text{sf}}1/\tau _{\text{sf}}+i\omega `$ so that the AC response can be used to detect the spin-accumulation. Our general result (29) is consistent with the calculations performed in the sequential tunneling regime.
## V F-F-F magneto-resistance
The spin-flip relaxation time in a ferromagnet is typically much shorter than the transport dwell time. In the limit $`\tau _{\text{sf}}0`$ the conductance of F-F-F systems reduces to
$$g=\frac{g^lg^r}{g^l+g^r}g_0(\beta E_c,g^l+g^r,n_{\text{ex}}).$$
(30)
The TMR of F-F-F systems is defined as
$$\gamma =\frac{R_{\text{AP}}R_\text{P}}{R_\text{P}}=\frac{g_\text{P}g_{\text{AP}}}{g_{\text{AP}}},$$
(31)
where $`R_\text{P}=R_K/g_\text{P}`$ ($`R_{\text{AP}}=R_K/g_{\text{AP}}`$) is the resistance of the FSET and $`g_\text{P}`$ ($`g_{\text{AP}}`$) is the dimensionless conductance of the FSET when the magnetizations in the leads and the central island are parallel (anti-parallel). (Note that we use a slightly different definition for the relative tunneling magnetoresistance in F-N-F systems than in F-F-F systems in order to make our discussions coherent with previous work). For F-F-F systems it is assumed that the magnetizations in the leads remain parallel, and that only the magnetization direction of the central island changes by applying an external magnetic field.
The TMR is caused by the magnetization configuration dependence of the total junction conductances $`g^l=g_{}^l+g_{}^l`$ and $`g^r=g_{}^r+g_{}^r`$ and therefore also by the magnetization configuration dependence of the renormalization factor $`g_0`$. The latter dependence gives rise to the enhancement of the TMR in the Coulomb blockade regime. The first term of the effective action (22) governing the renormalization of the conductance $`g_0`$ (18) only depends on the Coulomb charging energy of the device, and thus in the weak tunneling limit $`g_01`$ and the dimensionless conductances and the TMR of the FSET reduce to the classical values. The second term in the action (22) is proportional to the sum of the dimensionless conductances of the tunnel junctions and consequently is magnetization configuration dependent.
In order to evaluate the TMR the conductances of the FSET in different alignments must be calculated. We characterize the polarizations of the junctions by the relative polarizations $`P`$ and $`P^{}`$ associated with the reservoirs and the island, respectively. In a simple two-band model with spin-independent tunneling probabilities the polarizations $`P`$ and $`P^{}`$ are directly related to the density of states in the reservoirs and the island. In general there is not such a one-to-one correspondence between the polarizations of the tunnel barriers and the polarizations of the density of states, but we can still characterize the spin-dependent asymmetries in the tunnel conductances by the parameters $`P`$ and $`P^{}`$. We denote the dimensionless conductance of the left (right) junction averaged over the parallel and anti-parallel alignments by $`\overline{g}^{l/r}`$. In the parallel configuration the spin-dependent conductances are
$$g_{\text{P},s}^{l/r}=\frac{1}{2}\overline{g}^{l/r}(1+sP)(1+sP^{}),$$
(32)
where $`s=1`$ ($`s=1`$) for spin-up (spin-down) electrons. Hence, the total classical dimensionless junction conductance is
$$g_\text{P}^{\text{cl}}=g_\text{P}^l+g_\text{P}^r=\overline{g}(1+PP^{}),$$
(33)
where $`\overline{g}=\overline{g}^l+\overline{g}^r`$. Similarly, in the anti-parallel configuration the conductance is
$$g_{\text{AP},s}^{l/r}=\frac{1}{2}\overline{g}^{l/r}(1+sP)(1sP^{}).$$
(34)
The total classical dimensionless junction conductance is thus
$$g_{\text{AP}}^{\text{cl}}=g_{\text{AP}}^l+g_{\text{AP}}^r=\overline{g}(1PP^{}).$$
(35)
In the high temperature limit, one can substitute the forms of $`g_\text{P}^{\text{cl}}`$ and $`g_{\text{AP}}^{\text{cl}}`$ described above into the definition of the TMR (31), and obtain a simple expression $`\gamma ^{\text{cl}}=2PP^{}/(1PP^{})`$ for the classical TMR.
If the thermal energy is larger than the Coulomb charging energy, the path integrals can be evaluated semiclassically, and the lowest order quantum corrections to the TMR can be obtained. In this regime the conductance of the FSET is
$`g_0(\omega =0)`$ $`1{\displaystyle \frac{\beta E_c}{3}}+(0.0667+0.0185g_\eta ^{\text{cl}})(\beta E_c)^2`$ (37)
$`{\displaystyle \frac{8\pi ^2}{\beta E_c}}e^{\pi ^2/\beta E_cg_\eta ^{\text{cl}}/2}\mathrm{cos}(2\pi n_{\text{ex}})`$
with $`\eta =P(AP)`$ for the parallel (anti-parallel) alignment. Substituting this expression into the definition of the TMR (31), we obtain
$`\gamma ^{\text{se}}`$ $`=`$ $`\gamma ^{\text{cl}}+\mu ^{\text{se}}\left(1+\gamma ^{\text{cl}}\right)\left(\beta E_c\right)^2+\delta \gamma (n_{\text{ex}}),`$ (38)
where $`\mu ^{\text{se}}=0.02g_\text{P}^{\text{cl}}`$, and the leading term of the gate voltage dependent part of the semiclassical TMR is
$`\delta \gamma (n_{\text{ex}})`$ $`=`$ $`{\displaystyle \frac{8\pi ^2}{\beta E_c}}(1+\gamma ^{\text{cl}})e^{\pi ^2/\beta E_c}\times `$ (40)
$`[e^{g_{\text{AP}}^{\text{cl}}/2}e^{g_\text{P}^{\text{cl}}/2}]\mathrm{cos}(2\pi n_{\text{ex}}).`$
In Ref. , $`n_{\text{ex}}=0`$ was considered and consequently the gate voltage dependent contribution (40) was disregarded. For large junction conductances or high temperatures, this term is exponentially suppressed and the influence of the gate voltage on the TMR is negligible. The gate voltage dependence part of the TMR is important at low temperatures. Although Eq. (40) is not accurate at very low temperatures, it can be used to estimate when $`\delta \gamma (n_{\text{ex}})`$ becomes relevant. For FSET’s with positive polarizations $`PP^{}>0`$ the conductance in the anti-parallel configuration is smaller than the conductance in the parallel configuration $`g_{\text{AP}}^{\text{cl}}<g_\text{P}^{\text{cl}}`$, and the maximum TMR occurs when $`n_{\text{ex}}=0`$. The TMR attains its minimum at $`n_{\text{ex}}=1/2`$, and returns to the maximum at $`n_{\text{ex}}=1`$ since the TMR is a periodic function of $`n_{\text{ex}}`$ with period 1.
At low temperatures, an analytical expression of the TMR is not available. We will study the maximum and minimum value of the TMR at low temperatures via Monte-Carlo simulations. We use Ni leads with $`P=0.23`$, and Co island with $`P^{}=0.35`$ in all numerical calculations below.
### A TMR in the ’off’-state
In the absence of induced charges, the geometric phase factor $`2\pi ikn_{\text{ex}}`$ in the phase-phase correlation function (20) vanishes. The density matrix for each winding number $`k`$ is then positive and the phase-phase correlation function can be computed directly using standard Monte Carlo simulation techniques. In our earlier work a large enhancement of the TMR, roughly 4 times larger than the classical TMR for $`\overline{g}=12`$ and $`E_c/k_BT=40`$, was found. In this section we will demonstrate that the TMR is further enhanced at even lower temperatures. We will also address whether even larger junction conductances increase or decrease the TMR ratio.
The enhancement of the TMR ratio can be clearly demonstrated in the co-tunneling regime, where the conductance can be calculated analytically. Since co-tunneling is a second order tunneling process and the first order sequential tunneling vanishes in the Coulomb blockade regime, the conductance of the FSET is proportional to the square of the classical conductance with a temperature dependent prefactor. Consequently the TMR (31) is enhanced to a value slightly larger than twice the classical TMR. A naive extension of this result is that beyond the sequential and co-tunneling regimes, the TMR is enhanced by a factor of $`n`$ when the $`n`$-th order tunneling process dominates. However this picture is too simplified since contributions from different orders of tunneling have different temperature-dependent prefactors and one cannot in general specify the tunneling process that dominates for given junction resistances and temperature. In the strong tunneling regime perturbation theory breaks down and a number of higher order tunneling processes have to be taken into account simultaneously.
When the thermal energy is not much smaller than the charging energy, the semiclassical formula (38) can be modified to fit the TMR values by replacing the Coulomb charging energy by a renormalized charging energy. This substitution agrees well with the results of earlier investigations on the single-electron box. The renormalization scheme also predicts that the effective charging energy at sufficiently low temperature decreases for larger junction conductances due to the coupling to the leads. This means that for a given temperature, the TMR in the strong tunneling regime is smaller for larger junction conductances. The role of strong tunneling is therefore dual. The constructive role is to increase the TMR by the inclusion of higher order tunneling processes at low temperature. The destructive role is caused by the renormalized charging energy to a smaller value $`E_c^{}E_c`$ so that the TMR is closer to the classical value at a given temperature.
When the thermal energy is much smaller than the Coulomb charging energy strong tunneling is manifested as an enhancement of the TMR by lowering the temperature. We compute the TMR values for two values of the classical total dimensionless conductance $`\overline{g}=12`$ and $`\overline{g}=24`$. For a given conductance $`\overline{g}`$ and temperature the phase-phase correlation function (20) is computed via Monte-Carlo simulations for a sufficiently large number of discrete points in the imaginary time interval $`[0,\beta ]`$. The phase-phase correlation functions after two succeeding $`N`$ samplings are compared. If the curves are smooth and sufficiently close to each other, they predict roughly the same conductance after the analytical continuation of the Matsubara-Fourier components. The small difference determines the uncertainty of the final result, which is shown as the size of the data point in the figures below. When the deviation between the two phase-phase correlation functions is larger than the requested accuracy they are averaged to give the phase-phase correlation function for the case of $`2N`$ samplings, and a further simulation is performed for the next case of $`2N`$ samplings. This procedure is followed until the final result converges. Using the computed phase-phase correlation functions, the conductances of the TMR can be calculated via the Kubo’s formula Eq. (18). The analytical continuation from the Matsubara frequency to the real frequency is performed via the Páde approximate.
The resulting TMR ratio in the absence of the gate voltage is shown in Fig. (3) for $`\overline{g}=12`$ and 24. In both cases, the TMR values at low temperature are larger than the classical values, which shows that the TMR is enhanced by higher order tunneling processes in the Coulomb blockade regime. For $`\overline{g}=12`$ at $`k_BT/E_c=0.01`$, the TMR is enhanced by a factor of 8 compared to the classical value, which is almost twice as large as the value that we reported for $`k_BT/E_c=0.025`$ in our previous work. For a given temperature, the TMR for $`\overline{g}=12`$ is larger than the one for $`\overline{g}=24`$ in the temperature region shown in Fig. (3), which means that the smearing of the Coulomb charging energy via higher order tunneling processes is stronger for larger junction conductances. Therefore for experimentally easily accessible temperatures, only moderate enhancement of the TMR for the FSET formed by tunnel junctions with extremely large conductances can be obtained. In order to experimentally observe large TMR ratios in SET’s tunnel conductances of the order of 10 times the quantum conductance are recommended.
### B TMR in the ’on’-state
In the ’on’-state when $`n_{\text{ex}}=C_gV_g/e=1/2`$ module $`1`$, there is no Coulomb blockade effect even when the thermal energy is much smaller than the charging energy. In the sequential tunneling regime, the conductance of the FSET approaches half of the classical value when the temperature is lowered. Beyond the sequential tunneling regime, the reduction of the conductance of the FSET by lowering the temperature is slow at relatively high temperatures, but at sufficiently low temperature the conductance attains values smaller than the conductance in the sequential tunneling regime. The conductance in the ’on’-state in the strong tunneling limit and its implication for the TMR will be investigated in this section.
In order to calculate the ’on’-state conductance in the strong tunneling regime by Monte-Carlo simulations, we reformulate the phase-phase correlation function as
$$\mathrm{\Gamma }^{\text{on}}(\tau )=\mathrm{\Gamma }^{\text{re}}(\tau )/Z^{\text{re}},$$
(41)
where the function $`\mathrm{\Gamma }^{\text{re}}(\tau )`$ is the mean value of the phase factor
$`\mathrm{\Gamma }^{\text{re}}(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{Z^0}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _{b_k}}D\phi e^{S[\phi ]}\mathrm{cos}(\pi k)\times `$ (43)
$`\mathrm{cos}[\phi (\tau )\phi (0)],`$
with respect to the partition function $`Z^0`$ which is a direct summation of the density matrices for all winding numbers $`k`$, $`Z^0=_{k=\mathrm{}}^{\mathrm{}}_{b_k}D\phi e^{S[\phi ]}`$ and the partition function at resonance is renormalized in a similar way,
$$Z^{\text{re}}=\frac{1}{Z^0}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}_{b_k}D\phi e^{S[\phi ]}\mathrm{cos}(\pi k).$$
(44)
In the above formulas, both the numerator and the denominator of the phase-phase correlation function $`\mathrm{\Gamma }^{\text{on}}(\tau )`$ can be calculated with respect to the density matrices of the FSET with a geometric factor which equals one. The computation can therefore be performed in the same way as in the case of vanishing gate voltage. The complication is that we now have to calculate both $`\mathrm{\Gamma }^{\text{re}}(\tau )`$ and $`Z^{\text{re}}`$ instead of a single correlation function.
When the thermal energy is not much smaller than the charging energy, quantum fluctuations are small and the simulations are rather fast. Consequently the stable phase-phase correlation functions can be readily obtained using the approach described above. A typical result is shown in Fig. (4). At lower temperatures, the imaginary time $`\beta `$ is longer, and more lattice points are needed to guarantee the convergence of the Trotter product. Moreover, the computed correlation functions for the same amount of sampling numbers turn out to fluctuate more strongly at lower temperatures as shown in Fig. (5). For $`\overline{g}=12`$ at $`\beta E_c60`$ and $`\overline{g}=24`$ at $`\beta E_c100`$, the correlation functions computed via the approach of doubling the sampling number still fluctuate considerably within a realistic computing time. Therefore we do not attempt to perform simulations for temperatures lower than $`\beta E_c=100`$. At the lowest temperature the correlation function at the points with large fluctuations is replaced by a best fit smooth function so that the analytical continuations in the Kubo’s formula lead to a converging conductance of the FSET. The accuracy of this smoothing approach needs to be determined by long time Monte Carlo simulations in the future, since we cannot guarantee that the single data point at the lowest temperature is accurate within its size.
The TMR ratio as a function of the inverse dimensionless temperature is shown in Fig. (6). For $`\overline{g}=12`$, the TMR first increases by lowering the temperature. The values are almost the same as the case of the vanishing gate voltage. Around $`\beta E_c=10`$, the gate voltage dependent contributions to the TMR becomes important shifting the TMR value for $`C_gV_g=e/2`$ farther away from the value for $`C_gV_g=0`$ at lower temperatures. Therefore while the TMR in the Coulomb blockade regime is enhanced further at lower temperatures, the TMR at resonance decreases rapidly to the values smaller than the classical TMR for $`\beta E_c30`$. Negative TMR values could possibly occur at even lower temperatures, which needs to be determined quantitatively by further investigations which at present cannot be performed due to the computing time required. For $`\overline{g}=24`$, in contrast to the case of smaller junction conductances $`\overline{g}=12`$, the TMR increases constantly and slowly by lowering the temperature down to a value of $`0.01E_c/k_B`$. The curve more or less follows the tendency of the TMR in the case of vanishing gate voltage. This is understandable because the gate voltage dependent contributions to the TMR is proportional to the factor $`\mathrm{exp}(g_\eta ^{\text{cl}}/2)`$ as expressed in (40), which is smaller for larger $`\overline{g}`$, therefore the difference between the maximal TMR at $`C_gV_g=0`$ and the minimal TMR at $`C_gV_g=e/2`$ determined by the gate voltage dependent contribution is also smaller, and the variation of the minimal TMR will change the direction at lower temperatures. In both cases, the difference between the maximal TMR and the minimal TMR is found to be larger at lower temperatures, therefore a wider variation range of the TMR can be achieved by tuning the gate voltage.
## VI Conclusion
We have derived a general formula for the linear conductance of single-electron transistors containing ferromagnetic elements. In F-N-F structures, the TMR is almost independent of the Coulomb charging energy and is only reduced at sufficiently low temperatures when the effective transport dwell time becomes shorter than the spin-flip relaxation time of the normal-metallic island. In F-F-F systems, we have calculated the magneto-resistance as a function of gate voltages for a wide range of temperatures in the strong tunneling regime. In the ’off’-state, the magneto-resistance can be enhanced at low temperatures by higher order tunneling processes. However, higher order tunneling processes also reduce the effective charging energy, and slow down the rate of the enhancement of the magneto-resistance when the junction conductances are much larger than the quantum conductance. Consequently the largest magneto-resistance for a given temperature is obtained when the junction conductances are larger, but not very much larger than the quantum conductance. The magneto-resistance of the device in the ’on’ state has a more complicated temperature dependence. Due to the contributions from the induced charges, the magneto-resistance increases slowly for very large junction conductances. In contrast to this, the magneto-resistance is reduced well below the classical magneto-resistance for not very large junction conductances at sufficiently low temperatures.
We are grateful to G. E. W. Bauer, K.-A. Chao, B. I. Halperin, J. Inoue, Yu. V. Nazarov, and A. A. Odintsov for stimulating discussions. A. B. would like to thank the Norwegian Research Council for financial support and X. H. W. acknowledges support from the Swedish Research Council.
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# Broad iron lines in Active Galactic Nuclei
## 1 Introduction
In recent years, observations with both ground-based and space-based instruments have led to realization that most, if not all, nucleated galaxies harbor a massive black hole at their center (Kormendy & Richstone 1995; Magorrian et al. 1998). While many of these black holes appear to be relatively isolated, some fraction accrete significant amounts of material from the surrounding galaxy. The angular momentum of the incoming material leads to the formation of a flattened rotating disk — the accretion disk. The gravitational potential energy of material flowing through the accretion disk is converted into radiative (i.e. electromagnetic) and kinetic energy. These powerful and compact energy sources, observed in approximately 1–10% of galaxies are termed active galactic nuclei (AGN). AGN are also observed to be copious X-ray emitters. These X-rays are thought to originate from the innermost regions of an accretion disk around a central supermassive black hole. Since the accretion disk itself is expected to be an optical/UV emitter, the most likely mechanism producing the X-rays is inverse Compton scattering of these soft photons in a hot and tenuous corona that sandwiches the accretion disk. Thus, in principle, the study of these X-rays should allow the immediate environment of the accreting black hole as well as the exotic physics, including strong-field general relativity, that operates in this environment to be probed.
This review discusses how, in the past decade, X-ray astronomy has begun to fulfill that promise. Guided by observations with the Ginga, ASCA, RXTE and BeppoSAX satellites, there is a broad concensus that X-ray irradiation of the surface layers of the accretion disk in a class of AGN known as Seyfert 1 galaxies gives rise to fluorescent K$`\alpha `$ emission line of cold iron via the process of “X-ray reflection”. Since this line is intrinsically narrow in frequency, the observed energy profile of the line is shaped by both special relativistic (i.e. Doppler shifting) and general relativistic (i.e. gravitational redshifting and light bending) effects into a characteristic skewed profile predicted over a decade ago (Fabian et al 1989) and first clearly seen in ASCA data (Tanaka et al 1995). Since these lines are typically broadened to a full-width half maximum of $`5\times 10^4\mathrm{km}\mathrm{s}^1`$ or more, they are often referred to as “broad iron lines”. After discussing the physical processes responsible for the production of these spectral signatures, we will summarize the current observational status of broad iron line studies. We will show how current observations are already addressing the nature of the accretion disk within a few gravitational radii of the black hole. Observations of the broad iron line also provide valuable insights into the physical differences behind AGN of differing luminosities and type. Finally, we discuss and attempt to predict the results that will emerge from high throughput X-ray spectroscopy with XMM–Newton, Constellation-X and XEUS. We argue that these future data will provide unprecedented constraints on the spacetime geometry near the black hole (thereby measuring the spin of the black hole), well as the physical nature of the accretion disk.
## 2 The basics of the broad iron line
### 2.1 Line production
A substantial amount of the power in AGN is thought to be emitted as X-rays from the accretion disk corona in active or flaring regions. Thermal Comptonization (i.e. multiple inverse Compton scattering by hot thermal electrons; Zdziarski et al. 1994) of soft optical/UV disk photons by the corona naturally gives rise to a power-law X-ray spectrum. The flares irradiate the accretion disk which is relatively cold resulting in the formation of a “reflection” component within the X-ray spectrum. A similar component is produced in the Solar spectrum by flares on the solar photosphere (Bai & Ramaty 1978), in X-ray binaries by irradiation of the stellar companion (Basko 1978) and in accreting white dwarfs.
The basic physics of X-ray reflection and iron line fluorescence can be understood by considering a hard X-ray (power-law) continuum illuminating a semi-infinite slab of cold gas. When a hard X-ray photon enters the slab, it is subject to a number of possible interactions: Compton scattering by free or bound electrons<sup>1</sup><sup>1</sup>1Whether the electrons are bound or free is of little consequence for X-rays above 1 keV incident on gas mostly composed of hydrogen (Vainshtein et al 1983)., photoelectric absorption followed by fluorescent line emission, or photoelectric absorption followed by Auger de-excitation. A given incident photon is either destroyed by Auger de-excitation, scattered out of the slab, or reprocessed into a fluorescent line photon which escapes the slab.
Figure 1 shows the results of a Monte Carlo calculation which includes all of the above processes (Reynolds 1996; based on similar calculations by George & Fabian 1991). Due to the energy dependence of photoelectric absorption, incident soft X-rays are mostly absorbed, whereas hard photons are rarely absorbed and tend to Compton scatter back out of the slab. The reflected continuum is therefore a factor of about $`\sigma _\mathrm{T}/\sigma _{\mathrm{pe}}`$ below the incident one. Above energies of several tens of kilovolts, Compton recoil reduces the backscattered photon flux. These effects give the reflection spectrum a broad hump-like shape. In addition, there is an emission line spectrum resulting primarily from fluorescent K$`\alpha `$ lines of the most abundant metals. The iron K$`\alpha `$ line at $`6.4\mathrm{keV}`$ is the strongest of these lines. For most geometries relevant to this discussion, the observer will see this reflection component superposed on the direct (power-law) primary continuum. Under such circumstances, the main observables of the reflection are a flattening of the spectrum above approximately 10 keV (as the reflection hump starts to emerge) and an iron line at $`6.4\mathrm{keV}`$.
The fluorescent iron line is produced when one of the 2 K-shell (i.e. $`n=1`$) electrons of an iron atom (or ion) is ejected following photoelectric absorption of an X-ray. The threshold for the absorption by neutral iron is 7.1 keV. Following the photoelectric event, the resulting excited state can decay in one of two ways. An L-shell ($`n=2`$) electron can then drop into the K-shell releasing 6.4 keV of energy either as an emission line photon (34 per cent probability) or an Auger electron (66 per cent probability). (This latter case is equivalent to the photon produced by the $`n=2n=1`$ transition being internally absorbed by another electron which is consequently ejected from the ion.) In detail there are two components to the K$`\alpha `$ line, K$`\alpha _1`$ at 6.404 and K$`\alpha _2`$ at 6.391 keV, which are not separately distinguished in our discussion here. There is also a K$`\beta `$ line at 7.06 keV and a nickel K$`\alpha `$ line at 7.5 keV is expected.
For ionized iron, the outer electrons are less effective at screening the inner K-shell from the nuclear charge and the energy of both the photoelectric threshold and the K$`\alpha `$ line are increased. (The line energy is only significantly above 6.4 keV when the M-shell is lost, i.e. FeXVII and higher states.) The fluorescent yield (i.e. the probability that a photoelectric absorption event is followed by fluorescent line emission rather than the Auger effect) is also a weak function of the ionization state from neutral iron (FeI) upto FeXXIII. For Lithium-like iron (FeXXIV) through to Hydrogen-like iron (FeXXVI), the lack of at least 2 electrons in the L-shell means that the Auger effect cannot occur. For He- and H-like iron ions the line is produced by the capture of free electrons, i.e. recombination. The equivalent fluorescent yield is high and depends on the conditions (see Matt, Fabian & Reynolds 1997).
The fluorescent yield for neutral matter varies as the fourth power of atomic number $`Z^4`$, for example being less than one half per cent for oxygen. Predicted equivalent widths for low $`Z`$ lines are given in Matt et al (1997). Fluorescent X-ray spectroscopy is a well-known, non-invasive way to determine the surface composition of materials in the laboratory, or even of a planetary surface.
For cosmic abundances the optical depth to bound-free iron absorption is higher than, but close to, the Thomson depth, The iron line production in an X-ray irradiated surface therefore takes place in the outer Thomson depth. This is only a small fraction of the thickness (say 1 to 0.1 per cent) of a typical accretion disk and it is the ionization state of this thin skin which determines the nature of the iron line.
The strength of the iron line is usually measured in terms of its equivalent width with respect to the direct emission. (The equivalent width is the width of the continuum in, say eV, at the position of the line which contains the same flux as the line. Its determination is not entirely straightforward when the line is very broad.) It is a function of the geometry of the accretion disk (primarily the solid angle subtended by the “reflecting” matter as seen by the X-ray source), the elemental abundances of the reflecting matter, the inclination angle at which the reflecting surface is viewed, and the ionization state of the surface layers of the disk. We will address the last three of these dependences in turn.
#### 2.1.1 Elemental abundance
Elemental abundances effect the equivalent width of the iron line both through the amount of iron that is present to fluoresce, and the absorption of the line photons by L-shell photoelectric absorption of iron and K-shell photoelectric absorption of lower-$`Z`$ elements. These competing effects, together with the fact that the edge is saturated (i.e. most incident photons just above the photoelectric edge are absorbed by iron ions), leads to a roughly logarithmical dependence on abundance. For example, using the cosmic abundance values from Anders & Grevesse (1982), the equivalent width as a function of the iron abundance $`A_{\mathrm{Fe}}`$ is given by
$`W(A_{\mathrm{Fe}})=W(A_{\mathrm{Fe}}=1)(A_{\mathrm{Fe}})^\beta `$ $`(0.1<A_{\mathrm{Fe}}<1)`$ (1)
$`W(A_{\mathrm{Fe}})=W(A_{\mathrm{Fe}}=1)[1+b\mathrm{log}(A_{\mathrm{Fe}})]`$ $`(1<A_{\mathrm{Fe}}<20)`$ (2)
where
$`(\beta ,b)=(0.85,0.95)`$ $`\mathrm{edge}\mathrm{on}`$ (3)
$`(\beta ,b)=(0.75,0.48)`$ $`\mathrm{face}\mathrm{on}`$ (4)
$`(\beta ,b)=(0.78,0.58)`$ $`\mathrm{angle}\mathrm{averaged}`$ (5)
where $`A_{\mathrm{Fe}}=1`$ refers to cosmic abundances (Matt, Fabian & Reynolds 1997).
#### 2.1.2 Inclination angle
As the inclination angle at which the disk is viewed is increased, the observed equivalent width is depressed due to the extra absorption and scattering suffered by the iron line photon as it leaves the disk surface at an oblique angle. Ghisellini, Haardt & Matt (1994) find that
$$I(\mu )=\frac{I(\mu =1)}{\mathrm{ln}2}\mu \mathrm{log}(1+\frac{1}{\mu }),$$
where $`\mu =\mathrm{cos}i`$, $`i`$ being the angle between the line-of-sight and the normal to the reflecting surface.
#### 2.1.3 Ionization of the disk surface
X-ray irradiation can photoionize the surface layers of a disk (Ross & Fabian 1993; Ross, Fabian & Young 1999). As discussed above, the fluorescent line that the illuminated matter produces depends upon its ionization state. A useful quantity in this discussion is the ionization parameter
$$\xi (r)=4\pi F_\mathrm{x}(r)/n(r),$$
where $`F_\mathrm{x}(r)`$ is the X-ray flux received per unit area of the disc at a radius $`r`$, and $`n(r)`$ is the comoving electron number density: it measures the ratio of the photoionization rate (which is proportional to $`n`$) to the recombination rate (proportional to $`n^2`$). The iron line emission for various ionization parameters has been investigated by Matt et al (1993, 1996). They found that the behaviour split into four regimes depending on the value of $`\xi `$ (also see Fig. 2).
1. $`\xi <100`$ ergs cm s<sup>-1</sup>: The material is weakly ionized. X-ray reflection from the accretion disk produces a cold iron line at 6.4 keV. Since the total photoelectric opacity of the material is large even below the iron edge, the Compton backscattered continuum only weakly contributes to the observed spectrum at 6 keV, and the observed iron K-shell absorption edge is small. This regime is termed ‘cold’ reflection, since the reflection spectrum around the energy of the iron-K features resembles that from cold, neutral gas.
2. 100 ergs cm s$`{}_{}{}^{1}<\xi <500`$ ergs cm s<sup>-1</sup>: In this intermediate regime, the iron is in the form of FeXVII-FeXXIII and there is a vacancy is the L-shell ($`n=2`$) of the ion. Thus, these ions can resonantly absorb the corresponding K$`\alpha `$ line photons. Successive fluorescent emission followed by resonant absorption effectively traps the photon in the surafec layers of the disk until it is terminated by the Auger effect. Only a few line photons can escape the disk leading to a very weak iron line. The reduced opacity below the iron edge due to ionization of the lower-$`Z`$ elements leads to a moderate iron absorption edge.
3. 500 ergs cm s$`{}_{}{}^{1}<\xi <5000`$ ergs cm s<sup>-1</sup>: In this regime, the ions are too highly ionized to permit the Auger effect. While the line photons are still subject to resonant scattering, the lack of a destruction mechanism ensurs that they can escape the disk and produces a ‘hot’ iron line at $`6.8\mathrm{keV}`$. There is a large absorption edge.
4. $`\xi >5000`$ ergs cm s<sup>-1</sup>: When the disk is highly ionized, it does not produce an iron line because the iron is completely ionized. There is no absorption edge.
Note that ionization of the reflector paradoxically causes the observed iron edge to strengthen at moderate values of $`\xi `$. This is because the edge is saturated in reflection from a cold absorber, as is absorption at lower energies where oxygen and iron-L are the main absorbers. Ionization of oxygen and iron leads to the iron-K edge being revealed, and so apparently becoming stronger, as the reflected flux below the edge increases.
The Matt et al. (1993, 1996) calculations assume a fixed density structure in the atmosphere of the accretion disk. Nayakshin, Kallman & Kazanas (1999) have relaxed this assumption and included the effect of thermal instability in the irradiated disk atmosphere. In their solutions, the cold dense disk that produces the X-ray reflection features is blanketed with an overlying low-density, highly ionized, region. For weak irradiation, the ionized blanket is thin and does not affect the observed spectrum. However, for strong irradiation, the ionized blanket scatters and smears the ionized reflection features. In their models, it can be difficult to produce highly ionized iron lines in reflection — the effect of increasing ionization is to dilute the ‘cold’ reflection signature. The extent of this effect will depend on the Compton temperature of the radiation field.
### 2.2 The profile of the broad iron line
The iron K$`\alpha `$ line is intrinsically a rather narrow line. Hence, we can use broadening of the line to study the dynamics of the accretion disk. The line profiles is shaped by the effects of Doppler shifts and gravitational redshifting. Figure 3 demonstrates these effects at work in a schematic way. In a non-relativistic disk, each radius of the disk produces a symmetric double-horned line profile corresponding to emission from material on both the approaching (blue-shifted) and receding (red-shifted) side. The inenr regions of the disk, where the material is moving the fastest, produce the broadest parts of the line. Near a black hole, where the orbital velocities of the disk are mildly relativistic, special relativistic beaming enhances the blue peak of the line from each radius (second panel of Fig. 3). Finally, the comparable influences of the transverse Doppler effect (i.e. “moving clocks run slowly”) and gravitational redshifting (i.e. “clocks near black holes run slowly”) shifts the contribution from each radius to a lower energy. Summing the line emission from all radii of the relativistic disk gives a skewed and highly broadened line profile. It has been suggested by Pariev & Bromley (1998) that turbulence in the accretion disk may also significantly broaden the line. While a detailed assessment of this possibility must await future magneto-hydrodynamic disk simulations, it seems unlikely that the turbulent velocity field in a thin accretion disk will be large enough to broaden the line.
Some fully relativistic model line profiles are plotted in Figs. 4 and 5. In Fig. 4, we show the line profile from an accretion disk in orbit around a non-rotating black hole (described by the Schwarzschild metric). The line is assumed to be emitted from an annulus of the disk extending between 6$`r_\mathrm{g}`$ and 30$`r_\mathrm{g}`$ from the black hole, where $`r_\mathrm{g}=GM/c^2`$ is the standard gravitational radius. It is seen that the high energy “bluewards” extent of the line is a strong function of the inclination of the disk. In fact, the blue extent of the line is almost entirely a function of the inclination, thereby providing a robust way to measure the inclination of the disk. On the other hand, the redward extent of the line is a sensitive funcion of the inner radius of the line emitting annulus. In Fig. 5, we show model iron lines from a Schwarzschild black hole and a rapidly rotating black hole (described by a near extremal kerr metric). In this figure, we have made the assumption that the line emission extends down to the innermost stable orbit of the acretion disk. For these purposes, the principal difference between these two space-time geometries is the location of the innermost stable orbit (and hence the inner edge of the line emission) — this critical radius is at $`6r_g`$ in the Schwarzschild case, and $`r_g`$ in the extremal kerr case.
Model line profiles are given in the Schwarzschild case by Fabian et al (1989) and for the maximal Kerr (spinning black hole) case by Laor (1991). Iron lines in extreme kerr metrics are also computed by Bromley, Miller & Pariev (1998), and Martocchia, Karas & Matt (2000). These last two sets of authors have presented diagnostics that can be used by observers who wish to avoid full spectral fitting of complex relativistic models. The formalism of computing relativistic line profiles is also discussed by Fanton et al. (1997).
### 2.3 Observations of broad iron lines
The X-ray reflection spectrum was first clearly seen by the Japanese X-ray observatory Ginga (Pounds et al 1990; Matsuoka et al 1990). The CCD detectors on board the Advanced Satellite for Cosmology and Astrophysics (ASCA) were the first X-ray spectrometers to provide sufficient spectral resolution and sensitivity for investigating the profile of the iron line in detail. The first clear example of a broad skewed iron line came from a long ASCA observation of the Seyfert 1 galaxy MCG–6-30-15 (Fig. 6; Tanaka et al 1995). The sharp drop seen at about 6.5 keV both demonstrates the good spectral resolution of the CCD detector and, as discussed above, constrains the inclination of the disk to be about 30 degrees. If the inclination were greater then this blue edge to the line moves to higher energies (as seen in the broad iron line of the Seyfert 2 galaxy IRAS 18325–5926, Iwasawa et al 1996a). The redward extent of the line constraints the inner radius of the emission to be 7 r<sub>g</sub>and the overall shape means that most of the line emission is peaked within 20 r<sub>g</sub>.
Nandra et al. (1997a) and Reynolds (1997) used ASCA data to study the iron line in over 20 Seyfert 1 galaxies and found that most are significantly broader than the instrumental resolution. In a typical ASCA observation of an AGN, the signal-to-noise ratio of the detected iron line is insufficient to study the line beyond simply measuring that it is broad. To combat this problem, Nandra et al. (1997a) have summed together the data from many AGN to produce an average iron line profile. They find that the average line has clear extension to low energies. To the extent that individual sources can be studied, the inferred inclinations of the accretion disks are clustered around 30 deg, indicating some bias to the selected galaxies. Such a bias is expected within the context of the “unified model” of Seyfert galaxies (Antonucci 1993). In brief, the unified model states that Seyfert galaxies possess an obscuring torus on scales larger than the accretion disk. When one views the central regions of the AGN along a line of sight that is not blocked by the torus, one seens a type-1 Seyfert galaxy. Otherwise, one would see a type-2 Seyfert galaxy. If the accretion disk and tori are co-aligned (as might be expected on the basis of dynamical models; Krolik & Begelman 1988), and the tori have an average opening angle of 30–40 degrees, then we would naturally expect a bias in the measured disk inclinations in a sample of Seyfert 1 galaxies.
It should be noted that, in some cases, results on the inclination of the disk implied from the broad iron line and orientations of the systems inferred from the other techniques, e.g., ionization cone, radio jet, broad line clouds etc., differ from each other (e.g., Nishiura, Murayama & Taniguchi 1998). The prime example is NGC4151, for which the iron line suggests the inner accretion disk to be almost face-on (Yaqoob et al 1996) whilst the other observations point to an edge-on system. The recent suggestion by Wang et al. (1999) that a significant proportion of the iron line in NGC4151 is scattered into our line-of-sight by an electron scattering disk atmosphere may explain this discrepancy. Also, such a difference in geometry depending on scales of interest may be expected due to warping of the accretion disk or a multiple merger (e.g., NGC1068, Bland-Hawthorn & Begelman 1997). Iron lines from warped accretion disks have been studied theoretically in some detail by Hartnoll & Blackman (2000). In other objects the inclination inferred from the iron line agrees with the classification of the object (e.g. MCG–5-23-16; Weaver, Krolik & Pier 1998).
The observed iron line profiles in active galaxies are not necessarily solely from the accretion disk. Absorption and extra emission may alter the iron line emitted by the accretion disk. Nandra et al (1999) reported detection of an absorption feature at about 5.7 keV imposed on the broad iron line profile obtained from a long observation of NGC3516 (see Fig. 7). They suggest that this feature is due to K$`\alpha `$ resonant line scattering by highly ionized iron (with an intrinsic energy of 6.9 keV). The redshift of the absorption feature has been interpreted as evidence for matter infalling onto a black hole. However, gravitational redshifting of resonant absorption which could occur when the line photons are passing through the hot corona above the disk can also account for the observed feature if it occurs close to the black hole (Ruszkowski & Fabian 2000).
The iron line profiles observed in some Seyfert galaxies, especially Compton-thin Seyfert 2s (or those classified as Seyfert 1.8 or 1.9 in optical), have significant contribution of a narrow line component originating from matter far away from a the central black hole, e.g., a molecular torus (Weaver & Reynolds 1998). Such a narrow component would become clear if a primary X-ray source had faded away as has been seen in NGC2992 (Weaver et al 1996) and NGC4051 (Guainazzi et al 1998), since light travel times puts a fundamental limit on how rapidly line emission from the torus can respond to the central source. The iron line variability observed in NGC7314 (Yaqoob et al 1996) is exactly what is expected from a line consisting of a broad component originating from the accretion disk and a torus line. In this source, the broad iron line responds to changes in the primary power-law X-ray flux, while a narrow line component is found to be constant. It is interesting to note that this composite nature is not always found in Seyfert 1 galaxies. Any contribution of a stable, narrow line component in MCG–6-30-15 has been found to be very small (Iwasawa et al 1996; 1999). This might suggest that, in comparison to Seyfert 2 galaxies, Seyfert 1 galaxies have tori that have either smaller optical depths or smaller geometric covering factors.
Not all bright Seyfert 1 galaxies have iron lines so broad that the disk is required to extend down to the marginally stable orbit (i.e. $`<6r_\mathrm{g}`$). IC4329a is found to have a relatively narrow line (Done et al 2000) which can be modelled with a disk of inner radius of about $`50r_\mathrm{g}`$. The disk within this radius may either be missing or highly ionized.
Rapid X-ray variability in active galaxies predicts the iron line also to vary in response to the continuum with a small time lag. The light crossing time for 10 r<sub>g</sub>in a black hole with a mass of $`10^7M_7\mathrm{M}_{}`$ is about 500$`r_1M_7`$ s, which is much shorter than an integration time required for ASCA to collect enough line photons to perform a meaningful measurement in X-ray bright AGN ($`>10^{11}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$). It means that no reverberation effects in the line can be detected with ASCA. We will return to reverberation effects in Section 5.
Despite the inability of current detectors to measure reverberation, significant and complex variations of the iron line in profile as well as intensity have been observed in the Seyfert 1 galaxy MCG–6-30-15 (see Fig. 8; Iwasawa et al 1996b). During the first long observation of this galaxy in 1994, a line profile with unusually strong blue peak was found in a time interval of a bright flare whilst the line showed a very-broad, red-wing dominated profile during a deep minimum period. In this very-broad state, line emission from within 6 r<sub>g</sub> is required to explain the line profile and width. Possible theoretical interpretations of this are discussed in Section 4. A succession of large flares on the approaching side of the disk could produce the blue-peak dominated line profile, although it can also be explained if the line is produced predominantly at large radii ($`100`$ r<sub>g</sub>). It is worth noting that this bright flare showed a continuum spectral evolution similar to that seen in a shot in a Galactic black hole candidate, e.g., Cyg X-1 (Negoro et. al 1995).
Another peculiar line shape seen in a brief period ($`1`$ hr) of a flare during the 1997 long observation of MCG–6-30-15 also requires a large redshift within 6 r<sub>g</sub>. The blue wing of this line is, this time, shifted well below 6 keV and no significant line was detected around 6.4 keV. A possible explanation is that the line production occurred either in a thin annulus at 4 r<sub>g</sub> or a small patch at $``$2.5 r<sub>g</sub> on the approaching side of the disk (Iwasawa et al 1999).
Detailed studies of the RXTE data on the iron line variability in MCG–6-30-15 have presented a puzzling problem; most of the line flux appears to be constant in spite of strong continuum changes (Lee et al 1999; Reynolds 2000: see also studies of NGC5548 by Chiang et al 1999 and of the Galactic Black Hole Candidate Cygnus X-1 by Revnivtsev, Gilfanov & Churazov 1999). There should therefore be some self-regulating mechanism to produce a constant line flux, which has yet to be understood. However, a separate investigation of the broad red-wing and narrow core of the iron line for the 1994 observation (Iwasawa et al 1996) has revealed interesting behaviours of each component. The narrow core remains constant on time scales shorter than $`10^3`$ s but follows the continuum variations on longer time scales ($`>10^4`$ s). In contrast, the broad red-wing appears to follow the continuum on the short time scales. This is consistent with a line produced from a relativistic disk (see also Blackman 1999). Also puzzling is an anti-correlation between the reflected fraction and the equivalent width of the iron line measured in MCG–6-30-15 (Lee et al 2000) and NGC5548 (Chiang et al 2000).
## 3 Alternative models for a broad line
The claim that iron line studies are probing the region within a few gravitational radii of the black hole is a bold one, and should be tested against other models at every opportunity. Furthermore, the internal consistencies of the accretion disk hypothesis must be critically examined. Given the quality of data, the July-1994 MCG$``$6-30-15 line profile has become a testbed for such comparisons.
Fabian et al. (1995) examined many alternative models including lines from mildly relativistic outflows, the effect of absorption edges on the observed spectrum, and broadening of the line via comptonization. Fabian et al. found that none of these models were viable alternatives for the MCG$``$6-30-15 line profile. The idea of producing the broad line via Comptonization has been revived recently by Misra & Kembhavi (1997) and Misra & Sutaria (1999). They argue that the spectrum initially consists of a narrow iron line superposed on a power-law continuum and that Comptonization in a surrounding cloud with optical depth $`\tau 4`$ produces the broad line. The Comptonizing cloud must be both cold ($`kT<0.5\mathrm{keV}`$ in order to predominately downscatter rather than upscatter the line photons), and fully-ionized (since no strong iron absorption edges are seen in the continuum spectrum). The cloud is kept fully ionized and yet cool by postulating that the continuum source has a very luminous optical/UV component.
There are strong arguments against such a model. Since the power-law continuum emission also passes through any such Comptonizing cloud, one would observe a break in the continuum spectrum at $`E_{\mathrm{br}}m_\mathrm{e}c^2/\tau ^230\mathrm{keV}`$. Such a break is not observed in the BeppoSAX (Guainazzi et al. 1999) or RXTE data (Lee et al. 1999) for MCG–6-30-15 (see Misra 1999). Also, both continuum variability (which is seen on timescales as short as 100 s) and ionization arguments limit the size of the Comptonizing cloud in MCG–6-30-15 to $`R<10^{12}\mathrm{cm}`$. The essence of this ionization argument is that the ionization parameter at the outer edge of the cloud (which, for a fixed cloud optical depth, scales with cloud size as $`\xi 1/R`$) must be sufficiently high that all abundant metals, including iron, are fully ionized (Fabian et al. 1995; Reynolds & Wilms 2000). In the case of MCG–6-30-15, these constraints on the cloud size turn out to so tight that the postulated optical/UV component required to Compton cool the cloud would violate the black body limit (Reynolds & Wilms 2000). Comptonization moreover provides a poor fit (Ruszkowski & Fabian 2000). Hence, we consider the Comptonization model for the broad iron line not to be viable.
In another alternative model, Skibo (1997) has proposed that energetic protons transform iron in the surface of the disk into chromium and lower $`Z`$ metals via spallation which then enhances their fluorescent emission (see Fig. 1). With limited spectral resolution, such a line blend might appear as a broad skewed iron line. This model suffers both theoretical and observational difficulties. On the theoretical side, high-energy protons have to be produced and slam into the inner accretion disk with a very high efficiency (Skibo assumes $`\eta =0.1`$ for this process alone). On the observational side, it should be noted that the broad line in MCG–6-30-15 (Tanaka et al 1995) is well resolved by the ASCA SIS (the instrumental resolution is about 150 eV at these energies) and it would be obvious if it were due to several separate and well-spaced lines spread over 2 keV. There can of course be doppler-blurring of all the lines, as suggested by Skibo (1997), but it will still be considerable and require that the redward tails be at least 1 keV long.
Finally it is worth noting that the line profile indicates that most of the Doppler shifts are due to matter orbiting at about 30 degrees to the line of sight. The lack of any large blue shifted component rules out most models in which the broad line results from iron line emission from bipolar outflows or jets. What we cannot determine at present is the geometry in more detail. For example, we cannot rule out a ‘blobby’ disk (Nandra & George 1994). We do, however, require that any corona be either optically-thin or localized, in order that passage of the reflection component back through the corona does not smear out the sharp features. (Note though that an optically-thick corona over the inner regions of a disc would explain the lack of an iron line from that region.)
## 4 What happens inside 6 r<sub>g</sub>: Kerr black holes
Circular particle orbits around a black hole are only stable outside the radius of marginal stability, $`r_{\mathrm{ms}}`$. Within this radius, it is normally assumed that the material plunges ballistically into the black hole. The location of $`r_{\mathrm{ms}}`$ depends on the black hole spin, decreasing from 6 r<sub>g</sub> for a static Schwarzschild black hole ($`a/m=0`$) to 1.235 r<sub>g</sub> for prograde orbits around a maximally spinning Kerr black hole ($`a/m=0.998`$). Thus the inner edge of a Keplerian accretion disk will be at $`r_{\mathrm{ms}}`$. In the standard model the X-ray emission from the accretion disk takes place in a (possibly patchy) disk-hugging corona above the (almost) Keplerian part of the disk. The material inflowing within $`r_{\mathrm{ms}}`$ cannot support a corona and is assumed to receive an insignificant fraction of the X-ray illumination. In such models fluorescent line emission is expected to extend in only as far as $`r_{\mathrm{ms}}`$.
As discussed previously the degree of redshifting seen in the iron line is an indication of how close the line emitting region extends to the black hole. During a deep minimum in the X-ray light curve of the 1994 *ASCA* observation of MCG–6-30-15 the iron line was seen to broaden and shift to lower energies (see Fig. 8). The only way to have such significantly redshifted line flux is for the source of emission to move within the innermost stable orbit for a static Schwarzschild black hole (i.e. 6 r<sub>g</sub>). The line is well fit by the profile of a maximally spinning Kerr black hole leading to the tentative conclusion that the line was the first spectroscopic evidence for a Kerr black hole (Iwasawa et al 1996). (It was tentative because of the difficulty in measuring the continuum precisely at that time due to an increase in the strength of the warm absorber; Otani et al 1996). Later work by Dabrowski et al (1997) quantified the spin of the black hole required to produce such highly redshifted emission to exceed 95 per cent of its maximal value. The equivalent width of the line also increased dramatically (by a factor of 3–4) during this deep-minimum.
Of course, the black hole spin cannot change on such short timescales. It seems, instead, that the pattern of X-ray illumination across the disk must change in the sense that it becomes more concentrated that average during the deep minimum. Since such a dramatic change in illumination pattern is unlikely to occur just by chance, some structural change in the X-ray emitting corona is inferred. A plausible correspondence of the accretion disk thermal timescale, and the timescale on which the source enters the deep-minimum, may suggest that such structural changes are mediated by thermal instabilities in the accretion disk or accretion disk corona. An unsolved problem of this model is that one would expect the continuum level to increase significantly (rather than suffer the observed decrease) when the emission is originating from the energetically-dominant, innermost regions of the accretion disk.
Reynolds & Begelman (1997) have pointed out that although the inner edge of the accretion disk around a Schwarzschild (i.e. non-rotating) black hole is at 6 r<sub>g</sub> the accretion flow does not immediately become optically thin at smaller radii, and may remain optically thick almost all the way down to the horizon. If this it is illuminated in the right manner, ionized iron fluorescence from this region can give rise to an extremely broad and highly redshifted iron line very similar to that observed during the deep minimum (see Fig. 9). In this model, the illuminating source is assumed to be a point source located on the symmetry axis of the accretion disk. The change in line profile, line equivalent width, and continuum level are explained by assuming that the X-ray source changes its height above the accretion disk with approximately constant luminosity. However, the modelling of Reynolds & Begelman (1997) did not account for the Compton reflected continuum which would be expected to accompany the line. Young, Ross & Fabian (1998) have computed both the line and continuum spectrum expected from such a model. Within $`6m`$ the density of the infalling material drops rapidly and it becomes photoionized by the X-ray illumination, generating a large iron absorption edge (see Fig. 10). Preliminary indications suggest that such an edge may not be present in the data, although detailed spectral fitting has yet to be performed in order to address this issue. While the situation with current data remains ambiguous, searching for these type of spectral features in future XMM–Newton data may allow us to unambiguously determine the spin of the black hole in MCG–6-30-15 and other objects.
The detailed nature of the accretion flow within 6 r<sub>g</sub> of a Schwarzschild black hole is likely to be a lot more complicated than we have assumed in the above discussion. Magnetic fields within the innermost stable orbit may be rapidly amplified until their energy density is comparable to the rest-mass energy density of the accretion flow (Krolik 1999), and hence will be dynamically significant. Such strong field enhancement within the radius of marginal stability may lead to the creation of an “inner” X-ray emitting corona (note that Lee et al. 2000 have argued for an inner, highly variable, X-ray corona on the basis of the observed RXTE variability). The presence of magnetic fields may also exert a torque on the inner edge of the disk (usually assumed to have a zero-stress boundary condition) thereby increasing the outward flow of angular momentum and the efficiency of the disk (Agol & Krolik 2000; although also see Paczyński 2000 for an opposing view). If there is appreciable line fluorescence from within the radius of marginal stability, reverberation mapping (discussed below) may allow this region to be mapped out in detail.
## 5 Reverberation
The rapid X-ray variability of many Seyfert galaxies leads us to believe that the primary X-rays are emitted during dramatic flare-like events in the accretion disk corona. When a new flare becomes active, the hard X-rays from the flare will propagate down to the cold disk and excite iron fluorescence. Due to the finite speed of light, the illumination from the flare sweeps across the disk, and the reflected X-rays act as an ‘echo’ of this flare. Such flaring will cause temporal changes in the iron line profile and strength due to the changing illumination pattern of the disk and, more interestingly, time delays between the observed flare and the its fluorescent echo. This latter effect is known as reverberation.
In principal, reverberation provides powerful diagnostics of the space-time geometry and the geometry of the X-ray source. When attempting to understand reverberation, the basic unit to consider is the point-source transfer function, which gives the response of the observed iron line to an X-ray flash at a given location. As a starting point, one could imagine studying the brightest flares in real AGN and comparing the line variability to these point-source transfer functions in an attempt to measure the black hole mass, spin and the location of the X-ray flare. By studying such transfer functions, it is found that rapidly rotating black holes possess a characteristic reverberation signature. In the case where the fluorescing part of the accretion disk extends down to the radius of marginal stability of a near-extremal kerr black hole, the instantaneous iron line profile displays a low-energy bump which moves to lower-energies on a timescale of $`GM/c^3`$. This feature corresponds to highly redshifted and delayed line emission that originates from an inwardly moving ring of illumination/fluorescence that asymptotically freezes at the horizon (see Reynolds et al. 1999 and Young & Reynolds 2000 for a detailed discussion of this feature).
The primary observational difficulty in characterizing iron line reverberation will be to obtain the required signal-to-noise. One must be able to measure an iron line profile on a timescale of $`t_{\mathrm{reverb}}GM/c^3500M_8\mathrm{s}`$, where we have normalized to a central black hole mass of $`10^8\mathrm{M}_{}`$. This requires an instrument with at least the collecting area of XMM–Newton, and probably Constellation-X. Figure 11 shows that Constellation-X can indeed detect reverberation from a bright AGN with a mass of $`10^8\mathrm{M}_{}`$. Furthermore, the signatures of black hole spin may well be within reach of Constellation-X (Young & Reynolds 2000). Although these simulations make the somewhat artificial assumption that the X-ray flare is instantaneous and located on the axis of the system, it provides encouragement that reverberation signatures may be observable in the foreseeable future.
Of course, the occurrence of multiple, overlapping flares will also hamper the interpretation of iron line reverberation. In the immediate future, the most promising (although observationally expensive) approach will be to observe a bright source while it undergoes a very large flare. (Although the current data show little evidence in general for correlated behaviour of the continuum and iron line, as mentioned in Section 2.3, changes in the line are seen during some of the largest flares and dips in MCG–6-30-15.) Assuming that the single flare temporarily dominates the flux from the source, we might expect the point source transfer function to be an adequate description of the reverberation. Ultimately, reconstruction techniques for extracting reverberation signatures from overlapping flares should be explored.
## 6 Other types of black hole systems
As outlined in the preceding discussion, most of the progress in the field of broad iron line observations has been for type-1 Seyfert galaxies. In this section, we briefly discuss X-ray reflection diagnostics of other types of black hole systems.
### 6.1 Quasars
There is no evidence of an iron line or any reflection features in the spectra of most luminous quasars (Nandra et al 1995; 1997). The equivalent width of the iron line is observed to decrease with increasing luminosity as one moves from the Seyfert 1 regime ($`L_\mathrm{X}<10^{44}\mathrm{erg}\mathrm{s}^1`$) to the quasar regime, a phenomenon termed the ‘X-ray Baldwin effect’ (Iwasawa & Taniguchi 1993; Nandra et al 1995). It has been suggested that the accretion disk becomes increasingly ionized. This might be due to the most luminous objects possessing accretion rates closer to the Eddington limit. One puzzling aspect of this explanation is that the disk must jump from being ‘cold’ to being completely ionized otherwise we would observe instances of intermediate ionization in which the iron in the surface layers of the disk is H or He-like and the equivalent width of the line is even larger (Fig. 2 with $`\xi =1\times 10^3`$). There should also be a large absorption edge that is not seen. As mentioned in Section 2.1, a thermal instability in the surface layers of the disk (Nayakshin et al 1999) may lead to the formation of a highly ionized blanket and circumvent this problem.
The absence of reflection features can also be explained if the fluorescing accretion disk subtends less than $`2\pi \mathrm{sr}`$ as seen from the illuminating X-ray source, e.g. if the fluorescing accretion disk truncates at a few tens of Schwarzschild radii. A transition to an advection dominated accretion flow (ADAF), which is hot and optically-thin, would produce such a geometry. However, such structures can only exist at small acretion rates. Together with the fact that most of the energy in such a structure is advected into the black hole, ADAFs are expected to be much less luminous than most quasars. A more likely possibility is that, as an object approaches its Eddington limit, the region of the disk that is radiation pressure dominated moves further out, causing the surface layers of the disk to become more tenuous and highly ionized. The formation of a super-Eddington accretion disk in which much of the radiative energy is trapped in the accretion flow may also be relevant to quasars.
Curiously, one of the most luminous low redshift quasars, PDS 456, shows significant features at the iron-K energies (Reeves et al 2000). The features appear as a deep ionized edge and a possible broad line, and are modelled as either an ionized reflector (disk), or less likely as a strong highly ionized warm absorber.
### 6.2 Low luminosity AGN
Most massive black holes in the local universe are accreting at rates that are much lower than those found in AGN. If the accretion is so small as to render the black hole undetectable, the galaxy is termed a quiescent galaxy. Slightly higher accretion rates will lead to the classification as a low luminosity AGN (LLAGN).
The nature of black hole accretion when the accretion rates are very low is a topic of active research. It was realized by several authors that when the accretion rate is low (relative to the Eddington rate), an accretion disk may switch into a hot, radiatively-inefficient mode (Ichimaru 1977; Rees 1982; Narayan & Yi 1994; Narayan & Yi 1995). In essence, the plasma becomes so tenuous that the timescale for energy transfer from the protons to the electrons (via Coulomb interactions) becomes longer than the inflow timescale. The energy remains as thermal energy in the protons (which are very poor radiators) and gets advected through the event horizon of the black hole. These are the so-called Advection Dominated Accretion Flows (ADAFs). ADAFs are to be contrasted with ‘standard’ radiatively-efficient accretion disks in which the disk remains cool and geometrically thin all of the way down to the black hole (Shakura & Sunyaev 1973; Novikov & Thorne 1974). Broad iron line studies of LLAGN provide a potentially important probe of the physics of accretion when the accretion rate is low — the iron line traces only the radiatively-efficient portions of the disk since ADAFs are far too hot to produce fluorescent iron line emission.
Observationally, LLAGN have proven difficult to study due to the fact that they are X-ray faint. In addition, their X-ray spectra are typically complex with non-nuclear spectral components (such as starburst regions and/or thermal emission from hot gas) rivalling the nuclear component (e.g. see Ptak 1997). One of the best studied LLAGN resides in the nearby galaxy NGC 4258 (M 106). A short ASCA observation of this galaxy hinted at the presence of an iron line (Makishima et al. 1994). However, it took a deep ASCA observation to unambiguously detect the line and allow a detailed study (Reynolds, Nowak & Maloney 2000). It was found that the line in NGC 4258 is fairly weak (with an equivalent width of about 100 eV) and narrow (with a FWHM of less than $`22000\mathrm{km}\mathrm{s}^1`$). Reynolds et al. (2000) argue that this line does indeed originate from the accretion disk, implying that the X-ray emitting corona has a size greater than $`100GM/c^2`$. The contrast between the iron lines found in NGC 4258 and its higher luminosity Seyfert cousins is consistent with an ADAF scenerio for LLAGNs. However, the observational results are not yet conclusive. If the iron line seen in NGC 4258 comes from material not associated with the accretion disk (such as a distant torus that is misaligned with the almost edge-on accretion disk so as not to obscure the central engine from our view), then the data are consistent with the presence of a “Seyfert-like” broad iron line. See Reynolds et al. (2000) for further details.
While it is significantly more luminous than NGC 4258, the well studied active nucleus in the galaxy NGC 4051 is also often classified as a LLAGN. This object display a classic relativistic iron line indicating the presence of a radiatively-efficient accretion disk in this object (Guainazzi et al. 1996). Wang et al. (1999) have recently discovered interesting temporal variability in this iron line which displays opposite trends to the variability found in MCG–60-30-15 — both the equivalent width and energy width of the line positively correlate with the source flux.
### 6.3 Radio-loud AGN
It has long been known that most AGN can be readily characterized as being radio-loud (i.e. with relativistic radio jets) or radio-quiet (i.e. with no well-defined radio jets). One of the greatest mysteries in the field of AGN research is the physical mechanism underlying this division. A first step in solving this puzzle is to compare and contrast the central engine structures of radio-quiet and radio-loud AGN. Since they originate from deep within the central engine, X-rays are a good tool for probing any such differences.
Radio-loud AGN are rarer, and hence typically fainter, that their radio-quiet counterparts. Furthermore, many of the best candidates for study are found in clusters of galaxies and it can be difficult to observationally distinguish AGN emission from thermal cluster emission. For these reasons, the quality of te observational constraints is rather poorer in the case of radio-loud AGN as compared with radio-quiet sources. Having stated those caveats, there does appear to be a difference between the X-ray properties of radio-loud nuclei and radio-quiet nuclei. Broad iron lines, and the associated Compton reflection continua, are generally weak or absent in the radio-loud counterparts (Eracleous, Halpern & Livio 1996; Woźniak et al. 1997; Reynolds et al. 1997; Sambruna, Eracleous & Mushotzky 1999; Grandi et al 1999; Eracleous et al 2000). This effect might be due to the swamping of a normal ‘Seyfert-like’ X-ray spectrum by a beamed jet component (similar to the swamping of optical emission lines in a blazar spectrum). Alternatively, the inner disk might be in a physical state incapable of producing reflection signatures (such as an ADAF or some similarly hot state). Future observations with XMM–Newton should be able to distinguish these possibilities by searching for very weak broad components to the iron line.
### 6.4 Galactic Black hole Candidates
Smeared edges with little evidence of line emission have been observed in the spectra of Galactic Black Hole Candidates (GBHC) (Ebisawa et al 1996). These observations can be explained if the surface of the disk is moderately ionized with a mean ionization parameter of a few hundred (Fig. 2 with $`\xi =100300`$; Ross, Fabian & Brandt 1996). As discussed in Section 2.1, for this relatively narrow range of ionization parameters line photons are resonantly trapped and eventually lost as Auger electrons. Hence, the X-ray reflection produces very little line emission but an appreciable absorption edge. Sharp features will be smeared as a result of the Doppler and relativistic effects, and this blurring has possibly been detected in the spectrum of Cygnus X–1 (Done & Zycki 1999). Similar spectra but with broader emission and absorption features are produced for higher values of the ionization parameter $`\xi >3\times 10^3`$, which appear to match those seen in GBHC. The smearing in this instance is due to the line photons being generated a few Thomson scattering depths into the disk (the very outermost layers are completely ionized) and being Compton scattered on leaving it.
An interesting correlation has been claimed by Zdziarski, Lubiǹski & Smith (1999) between the reflection fraction seen in accreting black holes (AGN or GBHC) and the spectra index. Sources with flat spectra tend to have a low reflection fraction. Models involving large central holes in the disc, or ionized discs, may explain the correlation, as may mild relativistic motion, thus beaming, of the continuum radiation (Beloborodov 1999; see also Reynolds & Fabian 1997).
It is currently difficult to discriminate between models in which a cold disk truncates at a few tens of Schwarzschild radii (e.g. Gierliński et al 1997) and models in which an ionized disk extend in to the innermost stable orbit (Young et al 2000) since both provide good fits to present data. Future observations will hopefully resolve these issues.
## 7 Summary
In this section we summarize some of the key points of this review.
* *X-ray continuum.* The hard X-ray continuum in AGN and GBHCs is thought to be produced in active or flaring regions in a corona above the accretion disk. Thermal electrons multiply inverse Compton scatter optical and UV photons from the disk to X-ray energies. The hard X-ray power law that results irradiates the accretion disk and produces a “reflection” component in the spectrum.
* *Reflection component.* The reflection component causes the observed spectrum to flatten above 10 keV as Compton recoil reduces the backscattered flux and also results in a strong iron fluorescence line at approximately 6.4 keV. The precise energy and strength of the line depend on a number of factors such as the iron abundance, the inclination of the disk and its ionization state. An ionized disk may also produce a strong iron absorption edge.
* *Iron line profile.* The line profile is determined by Doppler shifts and relativistic boosting due to the motion of the disk and the gravitational redshifting of the black hole. This produces a broad, skewed line profile. Since the line originates from the innermost regions of the accretion disk, these effects are very pronounced. From observations of the line profile the black hole spin and the inclination angle of the accretion disk may be determined. In most Seyfert 1 galaxies (i.e. AGN in which we can view the accreting black hole directly) the accretion disk is inclined at about 30 to the observer. This is consistent with the standard model in which the Seyfert nucleus is surrounded by an optically thick torus with an opening angle of 30–40.
* *Observations.* The iron line was first clearly detected by *Ginga* and a line profile subsequently resolved by *ASCA* confirming the broad and skewed shape expected from an accretion disk around a Schwarzschild black hole. *RXTE* has been able to study the line and continuum variability on much shorter timescales, although with reduced energy resolution.
* *Black hole spin.* The radius of the smallest stable circular orbit around the black hole decreases with the spin of the black hole. Since the line profile is sensitive to the innermost radius of fluorescent emission this may be used (with some assumptions about the astrophysics of this region) to estimate the spin of the black hole. With present time-averaged observations, however, such measurements may be ambiguous as alternative models with very different values for the black hole spin may produce almost identical line profiles.
* *Alternative models for the production of the broad iron line.* Models for the broad iron line that do not require a black hole accretion disk appear to fail. In particular the line width cannot be entirely due to Comptonization. Hybrid models in which both Comptonization and Doppler/gravitational effects produce the line profile are heavily constrained.
* *Variability.* Rapid X-ray continuum variability is observed in most AGN and the iron line is expected to vary in response to this with a short time lag. Whilst these timescales are too short to be probed with present instruments, significant and complex iron line variability has been observed. Curiously, the line flux is seen to remain constant whilst the continuum changes, and there appears to be an anti-correlation between the reflected fraction and the equivalent width of the line. In another study the reflected fraction and the photon index of the power law are correlated, both for an individual object, and between different objects (including both AGN and GBHC). Such observations need to be explained, especially since they appear contrary to our simple model of reflection. Flux-correlated changes in the ionization state of the disk may explain some of these facts.
* *Reverberation mapping.* The rapid X-ray variability is associated with the activation of new flares in the corona above the accretion disk. X-ray reverberation mapping is the technique of using observations of the iron line response, or “echo”, to sudden changes in the continuum to study the accretion disk and black hole. In principle this may be used to determine the geometry of the X-ray emission and the black hole spin and mass. Such observations will be within the capabilities of the next generation of X-ray observatories.
* *Other classes of object.* Iron lines are observed in other classes of object in addition to Seyfert galaxies. In quasars the strength of the iron line decreases with increasing luminosity. This may be because the more luminous sources accreting closer to the Eddington limit and more highly ionized. The observation of iron lines in LLAGN may determine whether the accretion with low rates is an ADAF or a thin disk. Weak iron lines have also been seen in LLAGN suggesting their low accretion rate flows are thin disks as opposed to geometrically thick ADAFs. In radio loud AGN, broad iron lines and reflection humps are weak or absent, perhaps because the reflection signature is swamped by a beamed continuum. All of these require further detailed observations. In GBHC the accretion disk is ionized and the reflection spectra show smeared absorption and emission feature, and there is debate as to the precise nature of the accretion flow within a few tens of Schwarzschild radii of the black hole.
Over the past decade observations of the broad iron line have provided an unprecedented probe of the region within a few tens of Schwarzschild radii of the black hole event horizon. The next generation of X-ray observatories, beginning with *XMM–Newton*, will address many of the puzzling questions we have, and significantly enhance our understanding of these enigmatic objects.
## Acknowledgements
We thank Mateusz Ruszkowski for comments on the manuscript. ACF thanks the Royal Society for support. CSR appreciates support from Hubble Fellowship grant HF-01113.01-98A. This grant was awarded by the Space Telescope Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5-26555.
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# Dynamics of deviations from the Gaussian state in a freely cooling homogeneous system of smooth inelastic particles
## 1 Introduction
Freely cooling systems of smooth inelastically colliding spheres or discs (in the following denoted by inelastic hard spheres systems (IHS)) have been investigated by means of kinetic theory and computer simulations by several groups (see e.g. goldhirsch ; goldhirsch2 ; deltour ; esipoe ; mcnamara ; noijeprl ; Orza ). Most of the studies focus on latest times where interesting phenomena like formation of vortex patterns mcnamara ; noijeprl and clustering goldhirsch ; deltour ; Cahn can be observed. For short times or not too high inelasticities, however, the system remains homogeneous with a decreasing temperature, or equivalently, a decreasing average velocity. It is the so called Homogeneous Cooling State (HCS) the regime that will be studied in this paper. The HCS admits a scaling solution, i.e. if one scales all velocities with the average velocity and assumes that the shape of the scaled velocity distribution function remains constant in time, the entire time dependence is given by the time dependence of the average velocity only. This scaling solution is the starting point for a hydrodynamic analysis. Although many of the existing theories use a Gaussian velocity distribution function (which may be valid for small inelasticity) in general the shape is not Gaussian. First evidence, at very late evolution stages, was obtained by Goldhirsch et al. goldhirsch2 by measuring the fourth moment of the velocity distribution function. Later Goldshtein and Shapiro goldshtein proposed a solution based on an expansion in Sonine polynomials. Van Noije and Ernst noije correctly calculated the first term to linear order and Brilliantov and Pöschel poeschel included nonlinear corrections. Numerical solutions of the Boltzmann equation were calculated using Direct Simulation Monte Carlo method (DSMC) brey . Finally, extensions to viscoelastic particles have been recently presented in brilliantov .
Further confirmation of non Gaussian behavior is given by Esipov and Pöschel esipoe by studying the high-velocity tails of the velocity distribution function. They find that the tails are of an exponential type instead of a Gaussian, results confirmed by DSMC by Brey et al. brey2 and also in experiments performed by Losert and coworkers losert . In fact, an asymptotic exponential tail of the form $`\mathrm{exp}(v^\beta )`$ with $`\beta <2`$ seems to be a more fundamental behavior, as it is also present in driven or vibrated granular materials noije ; losert . Recently, a detailed study by DSMC in driven systems with three different types of forcing has been presented Andres . Only for the so-called Non-Gaussian thermostat (where there is a balance between energy input and dissipation) $`\beta =2`$, while for the other forcings, $`\beta =3/2`$ and 1. Surprisingly, no molecular dynamics results have been presented so far for calculating moments and high energy tails for the freely evolving case.
Our starting point is similar to goldshtein , i.e. expand the scaled velocity distribution function in a series of Generalized or Associated Laguerre polynomials around the Gaussian distribution with coefficients denoted with $`a_l`$. However, we assume that the coefficients are time dependent poeschel . With these ideas we try to achieve two goals. Firstly, investigate the influence of higher coefficients $`a_3,\mathrm{},a_6`$, in the expansion of the distribution function. Secondly study their time evolution. We find that for not too high inelasticities the above expansion seems to be convergent. Furthermore, the cooling proceeds in two stages: (1) A fast decay (in the order of few collisions per particle) of all coefficients $`a_l`$ to their asymptotically constant values. (2) An algebraically slow decay of the kinetic energy $`T`$ determined by Haff’s law $`\frac{d}{dt}T=CT^{3/2}`$ and time independent coefficient $`C`$ depending on the asymptotic values of $`a_l`$. Two dimensional event driven molecular dynamics simulations are performed in order to test the theory with good agreement for moderate inelasticity. For higher inelasticity the perturbations expansions seems to fail and in the simulations we were able to observe a transition to an exponential high-energy tail.
The paper is organized as follows: In Sec. 2 we propose the expansion of the velocity distribution in Laguerre polynomials with time dependent coefficients. In Sec. 3 we determine formally the full time dependence of the HCS expressed by the time evolution of the coefficients of the expansion. We obtain an infinitely large system of ordinary differential equations, which can only be investigated approximately. This is done in Sec. 4, where a truncation scheme is proposed and analyzed to different orders. In Sec. 5 we compare the analytical theory to results from event-driven simulations and the validity of the perturbation expansion is discussed. Results for the exponential high-energy tail are presented here. We summarize the results in Sec. 6.
## 2 The system under consideration
We consider a system of $`N`$ smooth, inelastically colliding spheres with diameter $`\sigma `$ confined to a $`d`$-dimensional volume $`V`$, so that the homogeneous density is given by $`n:=\frac{N}{V}`$. The positions of each sphere are denoted by $`\stackrel{}{r}_i`$ and each particle has a velocity $`\stackrel{}{v}_i`$. The particles interact via a hard-core potential and in each collision (i.e. if $`r_{ij}:=|\stackrel{}{r}_{ij}|:=|\stackrel{}{r}_i\stackrel{}{r}_j|=\sigma `$) the velocities are instantaneously changed by the following collision rules determined by a constant coefficient $`e_n[0,1]`$ of restitution
$`\stackrel{}{v}_i^{}`$ $`=\stackrel{}{v}_i{\displaystyle \frac{1+e_n}{2}}(\stackrel{}{v}_{ij}\widehat{\stackrel{}{r}}_{ij})\widehat{\stackrel{}{r}}_{ij},`$
$`\stackrel{}{v}_j^{}`$ $`=\stackrel{}{v}_j+{\displaystyle \frac{1+e_n}{2}}(\stackrel{}{v}_{ij}\widehat{\stackrel{}{r}}_{ij})\widehat{\stackrel{}{r}}_{ij},`$ (1)
where $`\stackrel{}{v}_{ij}:=\stackrel{}{v}_i\stackrel{}{v}_j`$ and $`\widehat{\stackrel{}{r}}_{ij}=\stackrel{}{r}_{ij}/r_{ij}`$. Velocities after collision are primed quantities given by velocities before collision (unprimed quantities).
The IHS is described statistically by the single particle distribution function $`\rho (\stackrel{}{r},\stackrel{}{v},t)d\stackrel{}{r}d\stackrel{}{v}`$, the (average) number of particles at positions between $`\stackrel{}{r}`$ and $`\stackrel{}{r}+d\stackrel{}{r}`$ and with velocities between $`\stackrel{}{v}`$ and $`\stackrel{}{v}+d\stackrel{}{v}`$ at time $`t`$. As proposed in goldshtein , for a homogeneous inelastic system the distribution function can be expressed by a scaling function as:
$$\rho (\stackrel{}{v},t):=n\frac{1}{(v_0(t)\sqrt{\pi })^d}\stackrel{~}{\rho }(\stackrel{}{c},t),$$
(2)
where $`\stackrel{}{c}=\stackrel{}{v}/v_0(t)`$ and $`v_0(t)`$ is the thermal velocity defined as the square root of the second moment of the distribution function:
$$𝑑\stackrel{}{v}\rho (\stackrel{}{v},t)\stackrel{}{v}^2=n\frac{d}{2}v_0^2(t).$$
(3)
The temperature is then defined as $`T(t):=\frac{m}{2}v_0^2(t)`$. For elastic systems the distribution function is Gaussian and it is expected that it will remain close to a Gaussian for small inelasticity. Therefore, we expand the scaled distribution function in a series of Generalized or Associated Laguerre polynomials oberhettinger around the Gaussian distribution function. The expansion is carried out in the scaled velocity variable $`\stackrel{}{c}`$ and with time dependent coefficients $`a_l(t)`$ poeschel . The general ansatz for the single particle distributions function for the homogeneous cooling then reads
$$\stackrel{~}{\rho }(\stackrel{}{c},t):=\mathrm{exp}\left(\stackrel{}{c}^2\right)\underset{l=0}{\overset{\mathrm{}}{}}a_l(t)L_l^\alpha \left(\stackrel{}{c}^2\right),$$
(4)
where $`\alpha =d/21`$ in $`d`$ dimensions. In the context of kinetic theory Laguerre polynomials are called Sonine polynomials chapcowl .
The normalization condition for $`\rho `$, $`𝑑\stackrel{}{v}\rho =n`$, leads to $`a_0=1`$. We express $`\stackrel{}{v}^2`$ by the first and second Laguerre polynomial
$$\stackrel{}{v}^2=L_1^\alpha (\stackrel{}{v}^2)+(\alpha +1)L_0^\alpha (\stackrel{}{v}^2),$$
(5)
and using the orthogonality relations for the Laguerre polynomials we find
$$𝑑\stackrel{}{v}\rho (\stackrel{}{v},t)\stackrel{}{v}^2=nv_0^2\left(\frac{d}{2}\left(\genfrac{}{}{0pt}{}{1+\alpha }{1}\right)a_1\right),$$
(6)
which implies together with Eq. (3) that $`a_1=0`$ for all times goldshtein ; noije ; poeschel . We denoted the binomial coefficients by $`\left(\genfrac{}{}{0pt}{}{a}{b}\right)`$.
Finally, as Laguerre polynomials are orthogonal, the coefficients $`a_l`$ are given by
$$na_l(t)=\frac{1}{\left(\genfrac{}{}{0pt}{}{l+\alpha }{l}\right)}𝑑\stackrel{}{v}\rho (\stackrel{}{v},t)L_l^\alpha \left(\left(\frac{\stackrel{}{v}}{v_0(t)}\right)^2\right).$$
(7)
## 3 The Homogeneous Cooling State
### The Boltzmann Equation
We assume that dynamics of the one particle distribution function $`\rho `$ is given by the Enskog Boltzmann equation, which can be written in $`d`$ dimensions without external forces as
$$_t\rho (\stackrel{}{r},\stackrel{}{v}_1,t)+(\stackrel{}{v}_1\stackrel{}{}_\stackrel{}{r})\rho (\stackrel{}{r},\stackrel{}{v}_1,t)=J[\rho ,\rho ],$$
(8)
with collision integral
$$\begin{array}{c}J[\rho ,\rho ]=\sigma ^{d1}\chi 𝑑\stackrel{}{v}_2𝑑\widehat{\stackrel{}{\sigma }}\mathrm{\Theta }(\stackrel{}{v}_{12}\widehat{\stackrel{}{\sigma }})(\stackrel{}{v}_{12}\widehat{\stackrel{}{\sigma }})\hfill \\ \hfill \left(\frac{b^1}{e_{n}^{}{}_{}{}^{2}}1\right)\left(\rho (\stackrel{}{r},\stackrel{}{v}_1,t)\rho (\stackrel{}{r},\stackrel{}{v}_2,t)\right).\end{array}$$
(9)
$`\widehat{\stackrel{}{\sigma }}`$ is the unit vector pointing from particle 2 to particle 1, $`\chi `$ the pair correlation function at contact, and $`b^1`$ describes ‘restituting collisions’ by changing velocities in $`\rho `$, i.e. $`b^1\rho (\stackrel{}{r},\stackrel{}{v}^{\prime \prime },t)=\rho (\stackrel{}{r},\stackrel{}{v},t)`$ in a way that $`\stackrel{}{v}^{\prime \prime }`$ are the velocities before collision leading to $`\stackrel{}{v}`$ after collision. The operator $`b`$ describes ‘direct collisions’ given in Eqs. (2). The inverse operator of $`b`$, i.e. $`b^1`$, is simply given by substituting $`e_n`$ by $`1/e_n`$ in Eqs. (2).
By multiplying the Boltzmann equation Eq. (8) with some function $`\psi (\stackrel{}{v}_1)`$ and integrating over $`\stackrel{}{v}_1`$ one gets
$$\begin{array}{c}_t𝑑\stackrel{}{v}_1\psi (\stackrel{}{v}_1)\rho (\stackrel{}{r},\stackrel{}{v}_1,t)+\stackrel{}{}𝑑\stackrel{}{v}_1\stackrel{}{v}_1\psi (\stackrel{}{v}_1)\rho (\stackrel{}{r},\stackrel{}{v}_1,t)=\hfill \\ \hfill 𝑑\stackrel{}{v}_1\psi (\stackrel{}{v}_1)J[\rho ,\rho ],\end{array}$$
(10)
which can be rewritten in the form of a balance equation
$$_t\overline{\psi }+\stackrel{}{}\stackrel{}{j}_\psi =_t^{\mathrm{coll}}\overline{\psi },$$
(11)
describing that the time change of an averaged quantity $`\overline{\psi }`$ is due to flux $`\stackrel{}{j}_\psi `$ or due to change through collisions. The right hand side of Eq. (10) can be written as chapcowl
$$\begin{array}{c}𝑑\stackrel{}{v}_1\psi (\stackrel{}{v}_1)J[f,f]=\sigma ^{d1}\chi 𝑑\stackrel{}{v}_2𝑑\stackrel{}{v}_1𝑑\widehat{\stackrel{}{\sigma }}\hfill \\ \hfill \mathrm{\Theta }(\stackrel{}{v}_{12}\widehat{\stackrel{}{\sigma }})(\stackrel{}{v}_{12}\widehat{\stackrel{}{\sigma }})\rho (\stackrel{}{v}_1)\rho (\stackrel{}{v}_2)\mathrm{\Delta }\psi ,\end{array}$$
(12)
and $`\mathrm{\Delta }\psi `$ is the change of $`\psi `$ in a direct collision for both particles $`\mathrm{\Delta }\psi =\frac{1}{2}(\psi (\stackrel{}{v}_1^{})+\psi (\stackrel{}{v}_2^{})\psi (\stackrel{}{v}_1)\psi (\stackrel{}{v}_2))`$.
### Dynamics of Moments
Using Eqs. (3) and (12) the time dependence of $`T(t)=\frac{m}{2}v_0^2`$ in the homogeneous case is given by
$$\frac{d}{dt}T=\frac{d}{dt}\frac{m}{2}v_0^2=2\gamma \omega _0T,$$
(13)
where $`\gamma `$ is defined as
$$\begin{array}{c}\gamma :=\frac{\sqrt{2\pi }}{dS_d}\frac{1}{\pi ^d}d\stackrel{}{c}_1d\stackrel{}{c}_2d\widehat{\stackrel{}{\sigma }}\mathrm{\Theta }(\stackrel{}{c}_{12}\widehat{\stackrel{}{\sigma }})(\stackrel{}{c}_{12}\widehat{\stackrel{}{\sigma }})\times \hfill \\ \hfill \stackrel{~}{\rho }(\stackrel{}{c}_1)\stackrel{~}{\rho }(\stackrel{}{c}_2)(b1)\frac{1}{2}(\stackrel{}{c}_1^2+\stackrel{}{c}_2^2),\end{array}$$
(14)
and $`\stackrel{~}{\rho }`$ as in Eq. (4). $`S_d`$ is the surface of a unit sphere in $`d`$ dimensions and $`\omega _0`$ the Enskog collision frequency for a classical gas of hard spheres with temperature $`T`$, given by:
$$S_d=\frac{2\pi ^{d/2}}{\mathrm{\Gamma }(d/2)}\text{and}\omega _0=\frac{S_d}{\sqrt{2\pi }}\chi n(2\sigma )^{d1}v_0.$$
(15)
If the velocity distribution function $`\stackrel{~}{\rho }`$ in Eq. (14) is a Maxwellian, $`\gamma `$ takes the value of $`\gamma _0:=(1e_n^2)/(2d)`$. This is the Gaussian value of the energy decay rate obtained by Haff haff . The fact that this distribution function is not a Gaussian modifies the cooling rate, as it will be calculated later on.
In order to obtain the time evolution of $`a_l`$, we take the time derivative of Eq. (7), and it has to be considered the time dependence of $`\rho (\stackrel{}{v},t)`$ as well as the time dependence of $`L_l^\alpha ((\frac{\stackrel{}{v}}{v_0(t)})^2)`$ via $`v_0(t)`$. The time dependence of $`\rho (\stackrel{}{v},t)`$ is given by Boltzmann equation, and the time dependence of $`v_0(t)`$ is given by equation (13). After a straight forward calculation using differential formulas for the Laguerre polynomials, we get
$$\frac{d}{dt}a_l=\omega _0\gamma _l+l2\gamma \omega _0(a_la_{l1}),$$
(16)
and
$$\begin{array}{c}\gamma _l=\frac{\sqrt{2\pi }}{S_d}\frac{1}{\left(\genfrac{}{}{0pt}{}{l+\alpha }{l}\right)}\frac{1}{\pi ^d}d\stackrel{}{c}_1d\stackrel{}{c}_2d\widehat{\stackrel{}{\sigma }}\mathrm{\Theta }(\stackrel{}{c}_{12}\widehat{\stackrel{}{\sigma }})(\stackrel{}{c}_{12}\widehat{\stackrel{}{\sigma }})\times \hfill \\ \hfill \stackrel{~}{\rho }(\stackrel{}{c}_1)\stackrel{~}{\rho }(\stackrel{}{c}_2)(b1)\frac{1}{2}(L_l^\alpha (\stackrel{}{c}_1^2)+L_l^\alpha (\stackrel{}{c}_2^2)).\end{array}$$
(17)
All collision integrals $`\gamma `$ and $`\gamma _l`$ depend on $`a_l`$ for all $`l`$ via $`\stackrel{~}{\rho }`$. We mention here that our approach is equivalent to the dynamics proposed in poeschel , but has the advantage to give immediately the explicit time dependence of all coefficients, at least formally.
### Collision frequency
The set of equations (13) and (16) for $`T`$ and $`a_l`$ with $`\gamma _l`$ and $`\gamma `$ given by (17) and (14) are the main results of this paper. However, before analyzing them in detail in the next section and comparing them with computer simulations in Sec. 5, it is instructive to study the collision frequency $`\omega `$. In order to do so, we introduce the average number of collisions, $`\tau `$, that a particle has suffered in a time $`t`$. Then, the collision frequency is defined as $`\omega =\frac{d}{dt}\tau (t)`$. In elastic fluids $`\omega `$ is a constant number depending only on the density and temperature, so that $`\tau `$ and $`t`$ are proportional quantities. In granular fluids, however, $`\omega `$ depends on time, as the temperature (and more precisely also the shape of the distribution function) of the system changes with time. Therefore, it is more natural from a physical point of view to express the time evolution equations in terms of the variable $`\tau `$ rather than $`t`$. Moreover, the hydrodynamic matrix become time independent when the hydrodynamic equations are expressed in the variable $`\tau `$ Cahn .
To determine $`\omega =\frac{d}{dt}\tau (t)`$ we use Eq. (12) and the fact that in each collision the number of collisions that each particle has suffered increases by one and we obtain chapcowl
$$\frac{d}{dt}\tau =\omega _0\gamma _\tau ,\mathrm{and}$$
(18)
$$\gamma _\tau =\frac{\sqrt{2\pi }}{S_d}\frac{1}{\pi ^d}𝑑\stackrel{}{c}_1𝑑\stackrel{}{c}_2𝑑\widehat{\stackrel{}{\sigma }}\mathrm{\Theta }(\stackrel{}{c}_{12}\widehat{\stackrel{}{\sigma }})(\stackrel{}{c}_{12}\widehat{\stackrel{}{\sigma }})\stackrel{~}{\rho }(\stackrel{}{c}_1)\stackrel{~}{\rho }(\stackrel{}{c}_2).$$
(19)
$`\gamma _\tau `$ depends on all $`a_l`$ and for the case that all $`a_l=0`$ for $`l>1`$ we would get $`\gamma _\tau =1`$ and thus the Enskog value $`\omega _0`$. We define a time $`\stackrel{~}{\tau }`$ by
$$d\stackrel{~}{\tau }=\omega _0dt.$$
(20)
Note that $`\stackrel{~}{\tau }`$ is only an approximation of $`\tau `$ defined in Eq. (18), so it does not really measure time in collisions, but we will show later that the deviations of $`\stackrel{~}{\tau }`$ from $`\tau `$ remain small for not too high inelasticities. In other words we hope that the collision frequency is approximately determined by the Enskog value and corrections due to deviations from the Gaussian affecting the collision frequency are small.
### Cooling rate
How can the dynamics be described in a state where all coefficients have already reached their asymptotic values? Note that the quantities $`\gamma `$ and $`\gamma _\tau `$ are entirely given by the values of $`a_l`$. Assuming that all $`a_l`$ have reached their asymptotic values for some time $`t>t^{}`$ or equivalently $`\tau >\tau ^{}`$ the quantities $`\gamma `$ and $`\gamma _\tau `$ also remain constant and we denote their asymptotic values by $`\gamma ^{}`$ and $`\gamma _\tau ^{}`$. Then we consider Eq. (13) and (18):
$$\frac{d}{dt}T=2\gamma ^{}\omega _0(T)T,$$
(21)
$$\frac{d}{dt}\tau =\omega _0(T)\gamma _\tau ^{},$$
(22)
which is solved analytically
$$\begin{array}{c}T=\frac{T(t^{})}{\left[1+\gamma ^{}\omega _0\left(T(t^{})\right)\left(tt^{}\right)\right]^2}=\hfill \\ \hfill T(\tau ^{})\mathrm{exp}\left(2\gamma ^{}/\gamma _\tau ^{}(\tau \tau ^{})\right),\end{array}$$
(23)
so that
$$\tau (t)\tau ^{}=\frac{\gamma _\tau ^{}}{\gamma ^{}}\mathrm{ln}\left[1+\gamma ^{}\omega _0\left(T(t^{})\right)\left(tt^{}\right)\right].$$
(24)
Eq. (24) provides a relation between collisional time and real time. Eq. (23), i.e. the algebraic decay of the temperature in time or the exponential behavior in $`\tau `$ is called Haff’s law. Furthermore, this equation makes explicit the fact that the shape of the velocity distribution function modifies the energy decay rate.
## 4 Analytical Results
### Truncation scheme
Up to now we have determined the full time dependence of the HCS in terms of the time dependence of all its moments in Eqs. (13) and (16). This infinitely large system of differential equations can only be solved by truncation. The approximate solution found by truncation only makes sense if all neglected terms are small as compared with the remaining ones. On the other hand, for small inelasticities the velocity distribution function is close to a Gaussian, so that $`a_21=a_0`$. We generalize this unequality and assume that $`a_{l+1}a_l`$ for all $`l`$, i.e. contributions from higher order coefficients get smaller the higher the index. Therefore, we propose a truncation scheme in which we assume that $`a_l`$ is of order $`\lambda ^l`$, where $`\lambda `$ is a small parameter. If we now make “an approximation of order $`𝒪(\lambda ^l)`$” we neglect all terms in all considered equations higher than $`\lambda ^l`$. This truncation scheme produces a finite set of differential equations that can be solved.
We will concentrate on two aspects: (i) To investigate the dynamics we integrate the full set of differential equations (up to a certain order $`\lambda `$). The asymptotic values of the coefficients can then be obtained by taking the long-time limit if they become stationary in time. (ii) To discuss the stationary state we set the left hand side of Eq. (16) equal to zero. This set of coupled and, as the case may be, non-linear equations can be solved with the numerical tool provided by the computer algebra program. Note that not all of these stationary values are necessarily dynamically stable solutions of the corresponding differential equation.
### Results to order 2
In a first step we only take into account $`a_2`$ to linear order. Then the functional form of the equation for $`a_2`$ Eq. (16) is given by
$$\frac{d}{dt}a_2=\omega _0(\gamma _2+4\gamma a_2)\frac{d}{d\stackrel{~}{\tau }}a_2=\gamma _2+4\gamma a_2.$$
(25)
Again, the use of $`\stackrel{~}{\tau }`$ simplifies the form of the equations, as it eliminates the time dependent factor $`\omega _0`$. We recall here that both $`\gamma _2`$ and $`\gamma `$ depend on all $`a_l`$, but for the approximation treated here, they only depend on $`a_2`$ in a linear manner. Therefore, we can express Eq. (25) as
$$\frac{d}{d\stackrel{~}{\tau }}a_2=A+Ba_2+𝒪(\lambda ^3),$$
(26)
where A and B are constants given by the collision integrals in $`\gamma `$ and $`\gamma _2`$ with the explicit expressions in two dimensions: $`A=\frac{1}{8}(e_n^21)(2e_n^21)`$ and $`B=\frac{1}{128}(30e_n^45e_n^232e_n57)`$.
(i) Dynamics– The evolution equation (26) for $`a_2`$ where time is expressed in collisions per particle is linear in $`a_2`$, so it can be easily integrated to give, when $`a_2(0)=0`$,
$$a_2(\stackrel{~}{\tau })=\frac{A}{B}(1\mathrm{exp}(B\stackrel{~}{\tau })),$$
(27)
so that the asymptotic value of $`a_2`$ is reached exponentially fast on a time scale of the order of $`\stackrel{~}{\tau }_0:=B^1`$. The decay time $`\stackrel{~}{\tau }_0`$ ranges between 1.7 and 2.25. As an important consequence the asymptotic solution is quickly reached on a kinetic time scale of few collisions per particle.
(ii) Stationary state– For times larger than $`\stackrel{~}{\tau }_0`$ $`a_2`$ reaches the stationary value of $`A/B`$ which coincides with the values calculated in noije .
### Results to order 3
To keep the discussion simple and to compare results from order to order, we first take into account only $`a_2`$ and $`a_3`$ i.e. up to $`𝒪(\lambda ^3)`$. We then still have to deal with equations which are linear in the coefficients ($`a_2^2`$ is already of order $`𝒪(\lambda ^4)`$). In the next section we will discuss the non-linear case up to order $`𝒪(\lambda ^6)`$. The equations read
$`{\displaystyle \frac{d}{dt}}T`$ $`=`$ $`2\gamma \omega _0T,`$
$`{\displaystyle \frac{d}{dt}}a_2`$ $`=`$ $`w_0\gamma _2+4\gamma \omega _0(a_20),`$
$`{\displaystyle \frac{d}{dt}}a_3`$ $`=`$ $`w_0\gamma _3+6\gamma \omega _0(a_3a_2),`$
$`{\displaystyle \frac{d}{dt}}\tau `$ $`=`$ $`\omega _0\gamma _\tau ,`$ (28)
$`{\displaystyle \frac{d}{dt}}\stackrel{~}{\tau }`$ $`=`$ $`\omega _0\text{ neglecting corrections of }a_2\text{ and }a_3\text{.}`$
We use computer algebraic programs to calculate the collisions integrals $`\gamma `$, $`\gamma _l`$ and $`\gamma _\tau `$ up to order $`𝒪(\lambda ^3)`$. The analytical solutions are rather lengthy and we will only show here results for a system with $`e_n=0.9`$ in Figs. 1–3.
(i) Dynamics– We have solved the simultaneous Eqs. (28) numerically<sup>1</sup><sup>1</sup>1We have used the built-in numerical procedure dsolve of MAPLE to integrate the differential equation. for the case $`e_n=0.9`$ and in the following we always plot time in units of $`1/\omega _0(T(0))`$ and temperature in units of $`T(0)`$. We have chosen $`a_2(0)=a_3(0)=0`$ as initial condition. In a first step we proof that the approximation to use $`\stackrel{~}{\tau }`$ instead of $`\tau `$ can be justified (at least to this order). In Fig. 1 a) we show the relative deviations of the true number of collisions to the approximation given by the Enskog Boltzmann value, i.e $`(\tau \stackrel{~}{\tau })/\tau `$ as function of time in a semilogarithmic plot.
We see that the relative deviations remain smaller than 0.2 %. This allows us, at least in the homogeneous cooling state, to use $`\stackrel{~}{\tau }`$ instead of $`\tau `$ in Eqs. (23) and (24).
In the asymptotic state we get $`\gamma ^{}=0.04723`$ for $`e_n=0.9`$ and values from the numerical integration of Eqs. (28) coincide with Eq. (23) and (24) within the graphical accuracy, so we plotted here only the numerical solution. Fig. 1 b) shows $`\tau (t)`$ which has the same form as predicted in Eq. (24).
In Fig. 2 a), we show $`T`$ as a function of time in a double logarithmic plot. We see the well-known asymptotic time dependence $`Tt^2`$. In Fig. 2 b), we show $`T`$ as a function of $`\tau `$ in a semi logarithmic plot resulting in a straight line with slope $`2\gamma ^{}`$ as predicted by Eq. (23).
In Fig. 3 we show the time dependence of $`a_2`$ and $`a_3`$ as a function of time a) and as a function of $`\tau `$ b).
We see that $`a_2`$ and $`a_3`$ reach their asymptotic value on a very short time scale which is of the order of few $`\tau `$’s. Therefore, few collisions per particle are necessary to reach the asymptotic state for $`a_2`$ and $`a_3`$.
(ii) Stationary state– As mentioned above we calculate the stationary values by setting the l.h.s of Eq. (4) equal to zero. In Fig. 4 we show the results for the stationary values of $`a_2`$ and $`a_3`$ to $`𝒪(\lambda ^3)`$ as well as $`a_2`$ to $`𝒪(\lambda ^2)`$.
As long as $`e_n>0.6`$, $`a_2`$ to $`𝒪(\lambda ^2)`$ does not differ significantly from $`a_2`$ to $`𝒪(\lambda ^3)`$ and $`a_3`$ remains small. We see stronger differences for smaller $`e_n`$ and $`a_3`$ becomes as important as $`a_2`$ which indicates stronger deviations from the Gaussian state. We also cannot assume anymore that corrections of higher orders remain small since we do not have any indication that the series is converging in the sense that the $`|a_l|`$ are small and decreasing.
Since to $`𝒪(\lambda ^3)`$ we have to deal with a set of linear equations we only find one unique solution. Considering higher orders one will find many solutions whose validity must be investigated. We will discuss this problem in the next section.
### Results to order 6
In this section we go to $`𝒪(\lambda ^6)`$, which is the highest order we were able to calculate with the computer algebra program.
(i) Dynamics– In Fig. 5 we show for $`e_n=0.8`$ the dynamics for the 5 non-vanishing coefficients $`a_2,\mathrm{},a_6`$ as a function of time. We have chosen the initial condition $`a_2(0)=\mathrm{}=a_6(0)=0`$. We see again a very fast decay to their asymptotic values. We observe that $`|a_l|>|a_{l+1}|`$ and in this sense the perturbation expansion seems to converge.
(ii) Stationary state– We calculate the stationary values by setting the l.h.s of Eq. (16) equal to zero. In Fig. 6 we show the results of the stationary values as a function of $`e_n>0.3`$.
For $`e_n>0.7`$ the coefficients remain small and the expansion seems to converge in the sense that $`|a_l|>|a_{l+1}|`$ for all $`l`$. For $`e_n<0.7`$ the absolute values of the coefficients start to grow and seem to diverge with $`e_n`$ approaching $`0.3`$.
To discuss the validity of these results we compare in Fig. 7 the asymptotic values of $`a_2`$ of order $`𝒪(\lambda ^2)`$ up to order $`𝒪(\lambda ^6)`$. As long as $`e_n>0.6`$, we do not find significant differences between the two orders, only the first order differ slightly from the other ones. This is a further hint that for these values of $`e_n`$ the perturbation method works. Moreover, the ratios of $`a_{l+1}/a_l`$ are small and of the same order: for instance, for $`e_n=0.85`$, these ratios are: $`a_3/a_2=0.39,a_4/a_3=0.29,a_5/a_4=0.24`$ and $`a_6/a_5=0.16`$.
For $`e_n<0.6`$ the results differ drastically from order to order and the proposed truncation scheme for Eq. (16) fails. We conjecture that around $`e_n0.6`$ an essential change in the distribution function occurs. Then the distribution function is no more described by small deviations around a Gaussian and might be better expressed by an expansion around an exponential as suggested in esipoe and confirmed by DSMC simulations of brey2 . We will go back to this point at the end of Sec. 5.
In view of Fig. 7, another possible explanation for the failure of the convergence is that the expansion in $`\lambda `$ is of asymptotic type, in such a way that including higher orders the expansion would break down in the whole range $`0<e_n<1`$. Unfortunately, we cannot decide which is the correct option, as we can only calculate up to order $`𝒪(\lambda ^6)`$.
### Further unstable solutions
Since we have to deal with non-linear equations, the solution is not unique and e.g. for $`e_n=0.8`$ two further stationary solutions can be found, similar as in poeschel . We list the values of the other coefficients for these stationary solutions:
| | Solution 1 | Solution 2 |
| --- | --- | --- |
| $`a_2`$ | 8.95 | -24.62 |
| $`a_3`$ | -14.39 | -4.50 |
| $`a_4`$ | 59.11 | 39.53 |
| $`a_5`$ | -109.17 | 178.6 |
| $`a_6`$ | 127.8 | -197.7 |
Both solutions are dynamically unstable which we have shown by numerical integration of the corresponding differential equations (13) and (16). In addition we observe that the higher coefficients are not at all negligible so that our assumptions, which should allow us to truncate the system of differential equations, are severely violated. Hence for these cases we can assume that we have not even found an approximate solution of the homogeneous Boltzmann equation.
## 5 Computer simulation results
In the literature only Direct Simulation Monte Carlo methods (DSMC) have been used for measuring the values of $`a_2`$ and $`a_3`$ brey and $`a_2`$ agrees very well with the value calculated in noije . However, DSMC lacks some features of the real IHS fluid, as correlations among the particles.
We present in this paper for the first time results for $`a_2`$ and for high-energy tails obtained from Molecular Dynamics (MD) simulations of the IHS system in 2 dimensions. Our code closely follows the event driven molecular dynamics code presented in Allen+Tildesley , adapted to the collision rules described in Eqs. (2) and accelerated by techniques described in lub . Typical simulations are performed with $`N=50000`$ particles in a square box of size $`L`$, being its area fraction $`\varphi =\frac{\pi \sigma ^2}{4}n=\frac{\pi N\sigma ^2}{4L^2}`$. The initial configuration is that of an elastic fluid at equilibrium (Maxwellian distribution for velocities and equilibrium correlations for the positions due to excluded volume effects), prepared by running the system with $`e_n=1`$ (elastic interactions) for not less than 50 collisions per particle. Therefore, in the initial state $`a_l=0`$ for $`l1`$.
At the beginning of the inelastic evolution the system remains for some time in the HCS, where the assumptions made in Sec. 2 are fully applicable and where computer simulations will serve to test those predictions. Later, vortices and clusters start to develop through the system and homogeneity is lost goldhirsch ; noijeprl . The higher the density $`\varphi `$ and the inelasticity the sooner these structures appear and important deviations from the theoretical values of $`a_l`$ are expected. We will come back to this point later. Furthermore, the analytical results are independent of the density<sup>2</sup><sup>2</sup>2To eliminate the dependency of real time on the density, time can be scaled by the collision frequency $`\omega (T(0))`$ at time $`t=0`$ when expressed in $`\tau `$, but only depend on the inelasticity $`e_n`$. Hence we have performed our simulations at low density of $`\varphi =0.05`$, although simulations at higher and lower densities have also been carried out.
### Results for moments
The typical time evolution of the 4th cumulant is shown in Fig. 8, where we have plotted the value of $`a_2`$ versus the number of collisions per particle $`\tau `$ for a low density case $`\varphi =0.05`$ and low inelasticity $`e_n=0.92`$. The dotted line is the result of Eq. (27) while the solid line is the result of the numerical simulation averaged over two realizations to slightly improve the accuracy. In this plot we observe the typical features of the IHS evolution described in former sections. Initially $`a_2`$ is equal to zero, as the system starts from a Maxwellian distribution, with $`a_l=0`$ for $`l>1`$. Then, within a very short time of a few collisions per particle, deviations from a Maxwellian build up in the system and the asymptotic values of $`a_l`$ are reached. This is a very fast process on a hydrodynamic time scale, as it involves only a few collisions per particle and, therefore, a few mean free times. Then the moments stay constant (within the accuracy of the computer simulations) as long as the system remains in the HCS.
A best fit of the simulation data to an expression like Eq. (27) gives that $`\stackrel{~}{\tau }_0=1.6\pm 0.5`$ and $`a_2(\mathrm{})=0.026\pm 0.004`$, while the theoretical values developed in Sec. 4 are $`\stackrel{~}{\tau }_0=1.85`$ and $`a_2(\mathrm{})=0.0246`$. We observe excellent agreement with the theory. Unfortunately, the accuracy of our computer simulations is not high enough to distinguish between the lowest order and the order 6. The DSMC method also finds good agreement with the value of $`a_2`$ brey .
At later times the system is no longer in the HCS and the assumptions used in the theoretical sections break down. This is best illustrated in the Fig. 9, where a simulation at low density $`\varphi =0.03`$ but at very high inelasticity $`e_n=0.4`$ is presented. We observe the same features described in Fig. 8 with values from simulations $`\stackrel{~}{\tau }_0=1.5\pm 0.5`$ and $`a_2(\mathrm{})=0.125\pm 0.007`$, a best fit for $`\tau <10`$, compared to the theoretical values $`\stackrel{~}{\tau }_0=1.83`$ and $`a_2(\mathrm{})=0.130`$. Again, the agreement is excellent. However, after a short time $`\tau \stackrel{>}{}12`$ the values of the moments start to deviate from the theoretical predictions. For $`\tau \stackrel{>}{}12`$ the homogeneity assumption breaks down. This can be checked, e.g. by plotting the curve of energy vs time goldhirsch2 ; mcnamara : deviations from Haff’s law imply lack of homogeneity Orza ; B+E . Visual inspection of the system (not presented here) show that, at this high inelasticity, the system immediately develops currents and dense clusters where particles move almost parallel inside them. The description in terms of $`\stackrel{~}{\rho }(c)`$ is wrong, as it does not take into account the local macroscopic currents. The values of $`a_2`$ can, at this late evolution stages, grow up to values of $`a_2=2`$ goldhirsch2 . This regime is outside the scope of this article.
Concerning the duration of HCS, a hydrodynamic analysis Orza shows that currents develop through the system on a time scale of $`\tau `$ of the order of $`\gamma _0^1`$. In Fig. 8 the deviation with respect to $`a_2(\mathrm{})`$ appears at $`\tau 35=1.3\gamma _0^1`$, while in Fig. 9 it appears at $`\tau 12=2.5\gamma _0^1`$.
The results of the MD simulations for $`a_2`$ in the HCS compared with the asymptotic solution of Eq. (26) are given in Fig. 10. The agreement is excellent even down to high inelasticities as $`e_n=0.4`$, as shown in Fig. 9. However our simulations do not have precision enough to test if higher order corrections are important. Concerning the values of $`\stackrel{~}{\tau }_0`$, MD results are very close, but always smaller than the theoretical values quoted below Eq. (27) and they are affected by large errors.
To conclude this section, we have also measured the 6th moment, $`c^6=v^6/v^2^3`$, and from this quantity we have obtained $`a_3`$ given by $`a_3=\frac{1}{6}c^6+(1+\frac{d}{4})c^4+(\frac{d^3}{48}\frac{7d}{12}1)`$. The results also agree with the solutions of Eq. (28), but, as expected, the discrepancies are now larger because the absolute value of $`a_3`$ is very small.
### Cooling rate
Another way to test the analytical results of Sec. 4 is to measure the dissipation rate given by the decay of the temperature $`T`$ or energy $`E`$ versus $`\tau `$. Haff’s calculations predicts an exponential law $`\mathrm{exp}(2\gamma _0\tau )`$ with $`\gamma _0=(1e_n^2)/2d`$ while higher order corrections modify it as shown in Eq. (23). These corrections respect to $`\gamma _0`$ are very small of only few parts in a thousand. If the IHS system is too large or inelastic, these deviations cannot be measured, as the curve of energy vs $`\tau `$ bends apart of the exponential law mcnamara ; B+E . This is due to the appearance of currents and vortices as quantitatively explained by B+E . Moreover, temperature and energy are no longer proportional, and, in contrast to energy, temperature is difficult to measure in a numerical experiment.
However, if the system is small enough no shear or clustering instability is excited (see e.g. Refs.mcnamara ; Orza for detailed explanations) and the system is forced to remain in the HCS for all times, where Eq. (23) is valid. This is called ‘kinetic regime’ in Ref.mcnamara . The drawback of this method is that, for a given density, it sets a maximum number of particles, that decreases with increasing inelasticity. As described in Ref. Orza a lower bound for $`k_{min}=\frac{2\pi }{L}<k_{}^{}`$ has to be satisfied, in the notation of Orza . We keep $`k_{min}=2k_{}^{}`$. This condition, together with our chosen low density of $`\varphi =0.05`$, restricts the number of particles to $`N=80`$ at $`e_n=0.70`$, and the results are no longer reliable. Even for $`e_n=0.95`$ the maximum number of particles is only $`N300`$. Hence we restrict ourselves to $`e_n>0.7`$.
Following this method we have performed simulations and have measured the energy decay rate as a function of $`\tau `$. We have verified that it is indeed exponential for all times, and, in order to improve the statistics, results are averaged over 1000 realizations. The results are presented in Fig. 11, where we plot the difference between the stationary energy decay rate $`\gamma ^{}/\gamma _\tau ^{}`$ from Eq. (23) and $`\gamma _0`$. Circles are data obtained by MD simulations with their errors bars, while dotted and dashed lines are the results from our theory to order $`𝒪(\lambda ^2)`$ and $`𝒪(\lambda ^3)`$ respectively. The small number of particles does not allow to obtain any significant result below $`e_n<0.7`$. For $`e_n>0.75`$ we observe a reasonable agreement, although we find significant deviations with respect to the theoretical results, MD data are always smaller that theoretical values. Unfortunately, no direct measurements with DSMC of the deviations respect to $`\gamma _0`$ have been reported so far. They would allow us to compare our results and elucidate the nature of these deviations.
It is important to note here that due to the small size of the deviations of the cooling rate with respect to $`\gamma _0`$ it is necessary to use in the theory the real number of collisions $`\tau `$ instead of $`\stackrel{~}{\tau }`$. The approximation $`\stackrel{~}{\tau }\tau `$ is correct within a few parts in a thousand as shown in Sec. 4, which is of the same order of magnitude as the correction $`\gamma ^{}/\gamma _\tau ^{}`$ with respect to $`\gamma _0`$. This is not the case, however, for calculations of the dynamics of $`a_l`$, for example given in Eq. (27), where we are not interested in such small deviations but want to give a first estimate of the time scales. Hence the use of $`\stackrel{~}{\tau }`$ and the approximation made in Eq. (27) is fully justified.
Finally, there is still the open question if deviations from the theoretical $`a_2`$ and the cooling rate are due to the existence of correlation in the HCS and the breakdown of the molecular chaos hypothesis in Eq. (8) reported in TC ; Jago . These effects cannot be tested in DSMC either, as this method is based on the factorization of the two particle distribution function.
### High-energy tails
In Refs. esipoe ; noije it has been shown that strong deviations with respect to the Maxwellian are present in the tails of the distribution, where particles have large energies. More precisely, they have found if $`c=v/v_0(t)(1e_n^2)^1`$ the velocity distribution function is no longer Maxwellian, but a simple exponential $`\stackrel{~}{\rho }𝒜\mathrm{exp}(Ac)`$ instead. This exponential distribution has been verified by numerical solutions of the Boltzmann equation using the DSMC method brey2 . The shape of the exponential is very different to the Maxwellian and therefore it is understandable that an expansion of the type given in Eq. (4) might be non-convergent as suggested by our theoretical analysis. Another reasonable possibility is that the series is indeed convergent but we have not gone high enough in the truncation scheme or there is a better choice of truncation, because $`a_{l+1}`$ is as large as $`a_l`$, as shown in Fig. 6.
In order to investigate the velocity distribution function and make the exponential range accessible, we have performed MD simulations with extreme inelasticities of $`e_n=0.1,0.2,0.4`$ to compare to moderate inelasticities $`0.6`$ and $`0.8`$. Moreover, as we need high accuracy in the tails, where populations are small, we have simulated systems with 250 000 particles at $`\varphi =0.05`$. However, even with this large number of particles we are not able to obtain the accuracy that can be achieved by the DSMC method brey2 . We will only be able to give evidences of the exponential tail. We measured the distribution function at times, where $`a_2`$ has already reached its asymptotic value, but the homogeneity assumption is still valid. In the example of Fig. 9 this would be at times $`5\stackrel{<}{}\tau \stackrel{<}{}10`$.
To estimate the distribution function from data, we use the kernel estimator technique described e. g. in eubank . In general, given a set of outcomes $`x_i`$, $`i=1,\mathrm{},M`$ of a random experiment the distribution function $`\rho (x)`$ can be estimated by
$$\rho (x)\frac{1}{M}\underset{i=1}{\overset{M}{}}\frac{1}{\sqrt{2\pi }\delta }\mathrm{exp}\left(\frac{(x_ix)^2}{2\delta ^2}\right).$$
(29)
The idea behind this method is that each data point also gives some information about its surrounding, which can be justified for smooth distribution functions. It is not necessary to use a Gaussian as kernel in Eq. (29), any normalized and more or less sharply peaked function can be used. The value $`\delta `$ is a free parameter which was chosen such that the measured distribution function for the elastic gas (i.e. the initial condition) is fitted best (to the eye) to the Maxwellian. We choose $`\delta =0.05`$. Since the distribution function $`\stackrel{~}{\rho }(c^2)`$ –with $`c`$ the modulo of the velocity– is not continuous at $`c=0`$, i.e. $`\stackrel{~}{\rho }=0`$ for $`c<0`$ and $`\stackrel{~}{\rho }(0)1/\pi `$ this technique gives bad results around $`c=0`$, but better results than the histogram method for the interesting high-velocity limit, where only few data points are given.
The simulation results for the velocity distribution function $`\stackrel{~}{\rho }`$ as a function of $`c`$ as well as the Maxwellian are shown in Fig. 12 in a semi-logarithmic plot. For $`c\stackrel{>}{}3.5`$ the statistical accuracy is poor, so results are only plotted up to $`c=3.5`$.
We observe that the deviations with respect to the Maxwellian are larger for lower values of $`e_n`$. On the contrary, as $`e_n`$ increases, the measured distribution approaches the Maxwellian. Closer inspection shows that for $`e_n=0.1`$ the distribution gets possibly close to an exponential (straight line in the semilogarithmic plot of Fig. 12) for $`2\stackrel{<}{}c\stackrel{<}{}3.5`$, while for $`e_n=0.6`$ the range where $`\mathrm{log}\stackrel{~}{\rho }`$ seems to be linear shrinks to $`3\stackrel{<}{}c\stackrel{<}{}3.5`$. We will show below that in this case ($`e_n=0.6`$ and $`c\stackrel{<}{}3.5`$) the distribution function can be reasonably well described by the results of the perturbation expansion around the Gaussian given in Eq. (4). If we perform a linear fit to an exponential in these ranges, we obtain values of the coefficient $`A`$ quite close to those reported by noije , tested in brey2 by DSMC method. For instance, for $`e_n=0.1`$, $`A3.2`$, that increases to $`A3.8`$ at $`e_n=0.4`$ and further to $`A4.7`$ at $`e_n=0.6`$.
If we go beyond $`e_n>0.6`$ the perturbation expansions of Sec. 4 seems to converge and it make sense to compare the analytical results with the simulation data. In Fig. 13 we show for $`e_n=0.6`$ and $`e_n=0.8`$ results of the simulations and of the analytical theory to orders $`\lambda ^2`$ and $`\lambda ^6`$.
In the whole range of $`c`$ which is plotted, the analytical results to order $`\lambda ^6`$ coincides fairly well with the measured data. However, results to order $`\lambda ^2`$ agree at small velocity, $`c\stackrel{<}{}3`$ for $`e_n=0.8`$, but fails for higher velocities. It seems reasonable that higher orders are needed to describe higher energy tails, where deviations respect the Gaussian are larger.
Therefore, for $`e_n>0.6`$ there is no need to describe data by an exponential tail for high velocities, although it cannot be assured, that the theory does not fail for even higher velocities. Note that the distribution function at $`c3.5`$ is already smaller than $`10^5`$ so that for 250 000 particles only 2 or 3 particles might not be correctly described.
## 6 Conclusion
In this article we investigated by means of analytical theory and simulations the dynamics of a freely cooling system of smooth granular particles as long as it remains homogeneous. Starting from a pure Gaussian state the system develops on a fast time scale to a state where the deviations from the Gaussian (described by cumulants) are stationary in time and the dynamics is entirely described by a decreasing kinetic energy.
More technically, we determined formally the full dynamics of the homogeneous cooling state in terms of the dynamics of the temperature $`T(t)`$ and the time dependent coefficients $`a_l(t)`$ of an expansion of the velocity distribution function in Generalized or Associated Laguerre polynomials around the Gaussian state. We obtained an infinitely large system of non linear ordinary differential equations, which can be solved numerically under the assumption that higher coefficients do not contribute.
Analytically, we found two main results. i) As far as dynamics is concerned, the HCS is characterized by the fact that only a few collisions per particle are necessary to reach a state where the coefficients are stationary in time. Then the entire time dependence is given by a slow algebraically decay of the temperature obeying $`\frac{d}{dt}T=2\omega _0\gamma T`$, with $`\gamma `$ depending on all the asymptotic values of the coefficients $`a_l`$. ii) As far as the asymptotic values of the coefficients are concerned, the expansion seem to converge in the sense that $`|a_{l+1}|<|a_l|`$ for $`e_n>0.6`$ and for this range of $`e_n`$ we do not find significant differences between orders. There exist further stationary but dynamically instable solutions of the considered differential equation, which are far away from the assumption of absolutely decreasing and small coefficients, so we cannot make any prediction for that cases. For $`e_n<0.6`$ the perturbation procedure seems to fail, the assumption we made to truncate the system of differential equations are severely violated and we find a strong dependency on the order of approximation. We have no answer if going to higher order or choosing a more suitable truncation scheme would show that the perturbation procedure nevertheless works, or, on the other hand, if the expansion presented here is only of an asymptotic type. A reasonable conclusion is that the system develops to a state which is very far from a Gaussian and might be better described by an expansion around an exponential as discussed in poeschel and brey .
Although much numerical work has been done on the HCS and clustering regimes (see, e.g. goldhirsch , mcnamara , Orza , TC and Stefan ), for the first time event-driven simulations are used in the present work to investigate deviations from the Gaussian distribution in the HCS. Mainly, three aspects were considered: 1) The dynamics and asymptotic values of the coefficients, 2) the influence on the decay rate of the temperature, 3) the shape of the velocity distribution function.
1. As long as the system remains in the homogeneous state the dynamical behavior as well as the static value of $`a_2`$ could be confirmed. Nevertheless, the statistics was too poor to distinguish if higher order analytical results give better values.
2. Similarly, the decay of the temperature was measured showing deviations with respect to the analytical results and opening the possibility to study other effects as correlations. Here we had to assure that the system remains in the homogeneous state, restricting the number of particles and therefore the quality of the statistics.
3. Measuring the full velocity distribution function we found that the distribution function can be described very well by the expansion around the Gaussian as long as $`e_n>0.6`$. For smaller $`e_n`$ the high-energy tails show an exponential shape confirming previous results found by analytical theory esipoe ; noije and DSMC simulations brey2 .
A possible extension of our work are systems of rough spheres with constant coefficient of restitution or Coulomb friction. Strong deviations from the Gaussian are observed in the angular velocity distribution function, surprisingly for the cases where the particles are almost smooth huthmann ; herbst . It would be interesting to perform a similar dynamical analysis along these lines.
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# Topics in Finite Temperature Field Theory
## 1 Introduction
Studies of physical systems at finite temperature have led, in the past, to many interesting properties such as phase transitions, blackbody radiation etc. However, the study of complicated quantum mechanical systems at finite temperature has had a systematic development only in the past few decades. There are now well developed and well understood formalisms to describe finite temperature field theories, as they are called. In fact, as we know now, there are three distinct, but equivalent formalisms \[1-3\] to describe such theories and each has its advantages and disadvantages. But, the important point to note is that we now have a systematic method of calculating thermal averages perturbatively in any quantum field theory.
This, of course, has led to a renewed interest in the study of finite temperature field theories for a variety of reasons. We can now study questions such as phase transitions involving symmetry restoration in theories with spontaneously broken symmetry . We can study the evolution of the universe at early times which clearly is a system at high temperature. More recently, even questions such as the chiral symmetry breaking phase transition or the confinement-deconfinement phase transition in QCD \[5-6\] have drawn a lot of attention in view of the planned experiments involving heavy ion collisions. This would help us understand properties of the quark-gluon plasma better.
The goal of this article is to share, with the readers, some of the developments in finite temperature field theories in the recent past and the plan of the article is as follows. In the next section, we will describe some basic ideas behind describing a quantum mechanical theory in terms of path integrals . This is the approach which generalizes readily to the study of finite temperature field theory. In section 3, we will discuss one of the formalisms, in fact, the oldest one, of describing finite temperature field theory. This goes under the name of the imaginary time formalism or the Matsubara formalism \[1, 5, 8-10\]. In this description, the dynamical time is traded in for the temperature. In contrast, the real time formalisms of finite temperature field theory contain both time and temperature. In section 4, we discuss one of the real time formalisms known as thermo field dynamics \[3, 10-11\]. This is an ideal description to understand operator related issues involving finite temperature field theories although it has a path integral representation which is quite nice for calculations as well. The other real time formalism, which is much older and is known as the closed time path formalism , is described in section 5. This formalism is very nice because it describes both equilibrium and non-equilibrium phenomena, at finite temperature, with equal ease. Temperature leads to many subtle features in field theories. In section 6, we discuss one such subtlety, namely, how one needs a generalization of the Feynman combination formula to perform calculations at finite temperature . In section 7, the issue of large gauge invariance is discussed within the context of a simple quantum mechanical model \[14-15\]. In section 8, we discuss in some detail how temperature can lead to breaking of some symmetries like supersymmetry (Temperature normally has the effect of restoring symmetries). Finally, we present a brief conclusion in section 9. The subject of finite temperature field theories is quite technical and to keep the contents simple, we have chosen, wherever possible, simple, quantum mechanical models to bring out the relevant ideas. Finally, we would like to note that there are many works in the literature and the references, at the end, are only representative and are not meant to be exhaustive in any way.
## 2 Path Integrals at Zero Temperature
In studying a quantum mechanical system or a system described by a quantum field theory, we are basically interested in determining the time evolution operator. In the standard framework of quantum mechanics, one solves the Schrödinger equation to determine the energy eigenvalues and eigenstates simply because the time evolution operator is related to the Hamiltonian. There is an alternate method for evaluating the matrix elements of the time evolution operator which is useful in studying extremely complicated physical systems. This goes under the name of path integral formalism \[7, 17-18\].
In stead of trying to develop the ideas of the path integral formalism here, let us simply note that, for a bosonic system described by a time independent quantum mechanical Hamiltonian, the transition amplitude can be represented as (The subscript $`H`$ denotes the Heisenberg picture.)
$${}_{H}{}^{}x_f,t_f|x_i,t_i_{H}^{}=x_f|e^{\frac{i}{\mathrm{}}H(t_ft_i)}|x_i=𝒟xe^{\frac{i}{\mathrm{}}S[x]}$$
(1)
There are several comments in order. First, the transition amplitude is nothing other than the matrix element of the time evolution operator in the coordinate basis. Second, the integral on the right hand side is known as a path integral. It is an integral over all possible paths connecting the initial coordinate $`x_i`$ and the final coordinate $`x_f`$ which are held fixed. The simplest way to evaluate such an integral is to divide the time interval of the path between $`x_i`$ and $`x_f`$ into $`N`$ intervals of equal length. Integrating over all possible values of the coordinates of the intermediate points (which are ordinary integrals) and taking $`N\mathrm{}`$ such that the time interval is held fixed is equivalent to integrating over all possible paths. Finally, the action $`S[x]`$ in the exponent of the integrand is nothing other than the classical action for the bosonic system under study. This is true for most conventional physical systems where the Hamiltonian depends quadratically on the momentum. If this is not the case (and there are some cases where it is not), the right hand side of (1) needs to be modified. However, for most systems that we will discuss, we do not have to worry about this fine point.
The advantage of the path integral is that while the left hand side involves quantum mechanical operators, the right hand side is described only in terms of classical variables and, therefore, the manipulations become quite trivial. Furthermore, the transition amplitude defined in eq. (1) can be generalized easily to incorporate sources and this allows us to derive various Greens functions of the theory in a very simple and straightforward manner. As an example, let us simply note here that for a harmonic oscillator, the action is quadratic in the dynamical variables, namely,
$$S[x]=_{t_i}^{t_f}𝑑t\left[\frac{1}{2}m\dot{x}^2\frac{1}{2}m\omega ^2x^2\right]$$
and, in this case, the path integral can be exactly evaluated and has the form
$`x_f|e^{\frac{i}{\mathrm{}}HT}|x_i`$ $`=`$ $`{\displaystyle 𝒟xe^{\frac{i}{\mathrm{}}S[x]}}`$ (2)
$`=`$ $`\left({\displaystyle \frac{m\omega }{2\pi i\mathrm{}\mathrm{sin}\omega T}}\right)^{\frac{1}{2}}e^{\frac{i}{\mathrm{}}S[x_{cl}]}`$
Here, we have defined $`T=t_ft_i`$. $`S[x_{cl}]`$ represents the action associated with the classical trajectory (satisfying the Euler-Lagrange equation) and has the form
$$S[x_{cl}]=\frac{m\omega }{2\mathrm{sin}\omega T}\left[(x_i^2+x_f^2)\mathrm{cos}\omega T2x_ix_f\right]$$
(3)
The path integrals can also be extended to quantum mechanical systems describing fermionic particles. However, one immediately recognizes that there are no classical variables which are fermionic. Therefore, in order to have a path integral description of such systems in terms of classical variables, we must supplement our usual notions of classical variables with anti-commuting Grassmann variables . With this, for example, we can write a classical action for the fermionic oscillator as
$$S[\psi ,\overline{\psi }]=_{t_i}^{t_f}𝑑t(i\overline{\psi }\dot{\psi }\omega \overline{\psi }\psi )$$
(4)
Here $`\psi `$ and $`\overline{\psi }`$ are anti-commuting Grassmann variables and in the quantum theory, as operators, can be identified with the fermionic annihilation and creation operators respectively. The action in eq. (4) is also quadratic in the variables much like the bosonic oscillator and the path integral for the fermionic oscillator can also be exactly evaluated giving
$`\psi _f,\overline{\psi }_f|e^{\frac{i}{\mathrm{}}HT}|\psi _i,\overline{\psi }_i`$ $`=`$ $`{\displaystyle 𝒟\overline{\psi }𝒟\psi e^{\frac{i}{\mathrm{}}S[\psi ,\overline{\psi }]}}`$ (5)
$`=`$ $`e^{\frac{i\omega T}{2}}e^{\left(e^{i\omega T}\overline{\psi }_f\psi _i\overline{\psi }_f\psi _f\right)}`$
In a quantum field theory, we are often interested in evaluating time ordered correlation functions in the vacuum because the S-matrix elements can be obtained from such Greens functions. These can be derived in a natural manner from what is known as the vacuum to vacuum transition functional which can be obtained from the transition amplitude in eq. (1) in a simple manner and also has a path integral representation of the form
$$\underset{T\mathrm{}}{lim}0|e^{\frac{i}{\mathrm{}}HT}|0=𝒟xe^{\frac{i}{\mathrm{}}S[x]}$$
(6)
where
$$S[x]=_{\mathrm{}}^{\mathrm{}}𝑑tL(x,\dot{x})$$
(7)
Furthermore, the path integral in eq. (6) has no end-point restriction unlike in eq. (1). This vacuum to vacuum transition amplitude is also commonly denoted by $`0|0`$ with the limiting process understood. We note here that an analogous formula also holds for fermionic systems.
The vacuum to vacuum amplitude in the presence of a source has the form
$$Z[J]=0|0_J=𝒟xe^{\frac{i}{\mathrm{}}S[x,J]}$$
(8)
where
$$S[x,J]=S[x]+_{\mathrm{}}^{\mathrm{}}𝑑tJ(t)x(t)$$
(9)
Here $`J(t)`$ is a classical source and it can be easily checked that, in the limit of vanishing source, the functional derivatives of $`Z[J]`$ give rise to time ordered Greens functions in the vacuum.
With this very brief review of the path integral description for zero temperature quantum mechanical theories, we are now ready to describe the different formalisms available to study quantum mechanical systems at finite temperature.
## 3 Imaginary Time Formalism
The properties of a quantum mechanical system, at finite temperature, can also be given a path integral description. There are various, but equivalent ways of doing this. Of the different formalisms available to study a quantum mechanical system at finite temperature, the imaginary time formalism is the oldest . To appreciate this, let us recall some of the features of a statistical ensemble. A statistical ensemble in equilibrium at a finite temperature $`\frac{1}{\beta }`$ (in units of Boltzmann constant) is described in terms of a partition function
$$Z(\beta )=\mathrm{Tr}\rho (\beta )=\mathrm{Tr}e^\beta $$
(10)
Here $`\rho (\beta )`$ is known as the density matrix (operator) and $``$ can be thought of as the generalized Hamiltonian of the system. If
$$=H$$
where $`H`$ is the Hamiltonian of the system, we say that the ensemble is a canonical ensemble where the particle number is fixed and the system is allowed to exchange only energy with a heat bath. On the other hand, if
$$=H\mu N$$
where $`N`$ is the number operator, then, the ensemble is known as a grand canonical ensemble where the system can exchange not only energy with a heat bath, but can also exchange particles with a reservoir. The constant $`\mu `$ is known as the chemical potential. In a statistical ensemble, of course, the important observables are the ensemble averages and, for any observable $`𝒪`$, they are defined as
$$𝒪_\beta =\frac{1}{Z(\beta )}\mathrm{Tr}\rho (\beta )𝒪$$
(11)
Let us also note here that since the partition function involves a trace, it leads to an interesting identity following from the cyclicity of the trace, namely, (we will assume from now on, unless otherwise specified, that $`\mathrm{}=1`$)
$`𝒪_1(t)𝒪_2(t^{})_\beta `$ $`=`$ $`{\displaystyle \frac{1}{Z(\beta )}}\mathrm{Tr}e^\beta 𝒪_1(t)𝒪_2(t^{})`$ (12)
$`=`$ $`{\displaystyle \frac{1}{Z(\beta )}}\mathrm{Tr}e^\beta 𝒪_2(t^{})e^\beta 𝒪_1(t)e^\beta `$
$`=`$ $`{\displaystyle \frac{1}{Z(\beta )}}\mathrm{Tr}e^\beta 𝒪_2(t^{})𝒪_1(t+i\beta )`$
$`=`$ $`𝒪_2(t^{})𝒪_1(t+i\beta )_\beta `$
Such a relation is known as the KMS (Kubo-Martin-Schwinger) relation which generalizes to all statistical ensemble averages and plays a crucial role in the study of finite temperature field theories.
It was observed quite early by Bloch that the operator $`e^\beta `$ in the definition of the partition function is like the time evolution operator in the imaginary time axis. This is really at the heart of the imaginary time formalism. In fact, let us note that the canonical partition function can be written as (with the trace taken in the coordinate basis)
$$Z(\beta )=𝑑xx|e^{\beta H}|x$$
(13)
It is clear now that if we identify $`T=i\beta `$ in eq. (1), then, we can give the partition function a path integral representation as ($`\mathrm{}=1`$)
$$Z(\beta )=𝒟xe^{S_E[x]}$$
(14)
where $`S_E[x]`$ is the Euclidean (imaginary time) action for the system defined over a finite time interval as
$$S_E[x]=_0^\beta 𝑑tL_E(x,\dot{x})$$
(15)
Furthermore, it is clear from eq. (13) that the variable $`x`$ must satisfy the periodic boundary condition
$$x(\beta )=x(0)$$
(16)
for eq. (14) to represent a trace (namely, the initial and the final states must be the same) and that the end point is being integrated over in the path integral in eq. (14) unlike in eq. (1). (It is important to note that the original work of Matsubara is an operator description of the imaginary time, but we will not discuss it in the present article.)
In fact, as an example, let us evaluate the canonical partition function for the bosonic oscillator using this formalism . The transition amplitude is already given for zero temperature in eq. (2). Now making the identifications
$$T=i\beta ,x_i=x_f=x$$
(17)
we obtain from eqs. (2) and (13)
$`Z(\beta )`$ $`=`$ $`{\displaystyle 𝑑x\left(\frac{m\omega }{2\pi \mathrm{sinh}\beta \omega }\right)^{\frac{1}{2}}e^{(m\omega \mathrm{tanh}\frac{\beta \omega }{2})x^2}}`$ (18)
$`=`$ $`\left({\displaystyle \frac{m\omega }{2\pi \mathrm{sinh}\beta \omega }}\right)^{\frac{1}{2}}\left({\displaystyle \frac{\pi }{m\omega \mathrm{tanh}\frac{\beta \omega }{2}}}\right)^{\frac{1}{2}}`$
$`=`$ $`{\displaystyle \frac{e^{\frac{\beta \omega }{2}}}{e^{\beta \omega }1}}`$
This is, indeed, the partition function for the bosonic oscillator as can be directly verified.
The partition function, for a fermionic system, can also be similarly given a path integral representation. However, the anti-commuting nature of the fermion variables introduces one crucial difference, namely, for a fermion theory, we have
$$Z(\beta )=𝒟\overline{\psi }𝒟\psi e^{S_E[\psi ,\overline{\psi }]}$$
(19)
with anti-periodic boundary conditions
$$\psi (\beta )=\psi (0),\overline{\psi }(\beta )=\overline{\psi }(0)$$
(20)
The Euclidean (imaginary time) action is again defined over a finite time interval as in eq. (15). In fact, let us calculate the canonical partition function for a fermionic oscillator, as an example, from the result in eq. (5) as well as the identifications in (20) . Using
$$\psi _f=\psi _i=\psi ,\overline{\psi }_f=\overline{\psi }_i=\overline{\psi }$$
we obtain (remember $`T=i\beta `$)
$`Z(\beta )`$ $`=`$ $`{\displaystyle 𝑑\overline{\psi }𝑑\psi e^{\frac{\beta \omega }{2}}e^{(1+e^{\beta \omega })\overline{\psi }\psi }}`$ (21)
$`=`$ $`e^{\frac{\beta \omega }{2}}(1+e^{\beta \omega })=2\mathrm{cosh}{\displaystyle \frac{\beta \omega }{2}}`$
In evaluating this, we have made use of the Berezin rules of integration for Grassmann variables and we note that eq. (21), indeed, gives the correct partition function for a fermionic oscillator as can be directly calculated.
Although our discussion so far has been within the context of simple quantum mechanical systems, everything we have said can be carried over to a quantum field theory. The partition function for a quantum field theory can again be written as a path integral involving a Euclidean action as
$$Z(\beta )=𝒟\overline{\psi }𝒟\psi 𝒟\varphi e^{S_E[\varphi ,\psi ,\overline{\psi }]}$$
(22)
where the Euclidean action is defined over a finite time interval and the fields satisfy the periodicity (anti-periodicity) conditions
$$\varphi (\beta ,\stackrel{}{x})=\varphi (0,\stackrel{}{x}),\psi (\beta ,\stackrel{}{x})=\psi (0,\stackrel{}{x})$$
(23)
and so on. The discussion is slightly more involved for gauge theories and to keep things simple, we will not discuss gauge theories.
This formulation of a field theory at finite temperature is known as the imaginary time formalism or the Matsubara formalism and is the oldest formalism. There are several distinguishing features of this formalism. For example, since the time interval is finite, Fourier transformation of the time variable would involve discrete energies. In other words, the Fourier transform of the propagator, say for example, at finite temperature in the imaginary time formalism, would take the general form
$$𝒢_\beta (\tau ,\stackrel{}{x})=\frac{1}{\beta }\underset{n}{}e^{i\omega _n\tau }𝒢_\beta (\omega _n,\stackrel{}{x})$$
(24)
where $`\omega _n=\frac{n\pi }{\beta }`$ with $`n=0,\pm 1,\pm 2,\mathrm{}`$. However, from the definition of the time ordered product
$$T_\tau (\varphi (\tau )\varphi ^{}(\tau ^{}))=\theta (\tau \tau ^{})\varphi (\tau )\varphi ^{}(\tau ^{})\pm \theta (\tau ^{}\tau )\varphi ^{}(\tau ^{})\varphi (\tau )$$
(25)
where we have allowed for both bosonic and fermionic fields and the KMS condition in eq. (12), it follows that, for $`\tau <0`$,
$$𝒢_\beta (\tau ,\stackrel{}{x})=\pm 𝒢_\beta (\tau +\beta ,\stackrel{}{x}))$$
(26)
It is important to recognize that the periodicity (anti-periodicity) of the propagator arises from the definition of the time ordered product for the bosonic (fermionic) fields and the KMS condition and is not directly connected with the periodicity (anti-periodicity) of the corresponding field variables which we have discussed earlier. This periodicity (anti-periodicity) of the propagator, on the other hand, leads to the restriction that eq. (24) holds with
$$\omega _n=\{\begin{array}{cc}\frac{2n\pi }{\beta }& \mathrm{for}\mathrm{bosons}\hfill \\ \frac{(2n+1)\pi }{\beta }& \mathrm{for}\mathrm{fermions}\hfill \end{array}$$
(27)
where $`n=0,\pm 1,\mathrm{}`$. These are conventionally known as the Matsubara frequencies .
Given this, one can now calculate the propagators for bosonic and fermionic field theories in the Matsubara formalism and they take the forms (in the momentum space)
$`𝒢_\beta (\omega _n,\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{1}{\omega _n^2+\stackrel{}{k}^2+m^2}}={\displaystyle \frac{1}{(\frac{2n\pi }{\beta })^2+\stackrel{}{k}^2+m^2}}`$ (28)
$`𝒮_\beta (\omega _n,\stackrel{}{k})`$ $`=`$ $`{\displaystyle \frac{\gamma ^0\omega _n+\stackrel{}{\gamma }\stackrel{}{k}+m}{\omega _n^2+\stackrel{}{k}^2+m^2}}={\displaystyle \frac{\gamma ^0(\frac{(2n+1)\pi }{\beta })+\stackrel{}{\gamma }\stackrel{}{k}+m}{(\frac{(2n+1)\pi }{\beta })^2+\stackrel{}{k}^2+m^2}}`$ (29)
Perturbative calculations can now be developed quite analogously to the zero temperature field theory. For example, given a field theory, we can read out the vertices from the Euclidean form of the action and use the propagators of eq. (28, 29) to carry out a diagrammatic calculation which would lead to the ensemble average for a given observable. It is clear that, because the time interval is finite in this formalism, the coordinate space calculation of any diagram is cumbersome. However, much like at zero temperature, the momentum space calculation is much simpler. However, one should keep the difference in mind, namely, that, at finite temperature, the external and the internal energies are discrete as in eq. (27). Consequently, the integration over internal energies (of zero temperature) is replaced by a sum over the internal energies. More specifically, we must use
$$\frac{d^4k}{(2\pi )^4}\frac{1}{\beta }\underset{n}{}\frac{d^3k}{(2\pi )^3}$$
(30)
As an example, let us consider the self-interacting scalar theory described by
$$=\frac{1}{2}_\mu \varphi ^\mu \varphi \frac{m^2}{2}\varphi ^2\frac{\lambda }{4!}\varphi ^4\lambda >0$$
(31)
We note that the only one loop correction in this theory is the mass correction. Rotating to Euclidean space and using the propagator for a scalar theory as given in eq. (28) as well as (30), we obtain the one loop mass correction to be
$`\mathrm{\Delta }m^2`$ $`=`$ $`{\displaystyle \frac{\lambda }{2\beta }}{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{(\frac{2n\pi }{\beta })^2+\stackrel{}{k}^2+m^2}}`$ (32)
$`=`$ $`{\displaystyle \frac{\lambda }{2\beta }}\left({\displaystyle \frac{\beta }{2\pi }}\right)^2{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{n^2+(\frac{\beta \omega _k}{2\pi })^2}}`$
Here, we have introduced the notation,
$$\omega _k=(\stackrel{}{k}^2+m^2)^{\frac{1}{2}}$$
(33)
The sum, in eq. (32), can be easily evaluated using the method of residues leading to
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{n^2+y^2}=\frac{\pi }{y}\mathrm{coth}\pi y\mathrm{for}y>0$$
(34)
Using this, the one loop mass correction can be determined to be
$`\mathrm{\Delta }m^2`$ $`=`$ $`{\displaystyle \frac{\lambda }{4}}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{\omega _k}\mathrm{coth}\left(\frac{\beta \omega _k}{2}\right)}`$ (35)
$`=`$ $`{\displaystyle \frac{\lambda }{4}}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{\omega _k}}+{\displaystyle \frac{\lambda }{2}}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{\omega _k}\frac{1}{e^{\beta \omega _k}1}}`$
$`=`$ $`\mathrm{\Delta }m_0^2+\mathrm{\Delta }m_\beta ^2`$
There are several things to note from this calculation. First, the mass correction separates into two parts – one independent of temperature and the other genuinely a finite temperature correction. The temperature independent part (zero temperature part) is divergent as is expected at zero temperature and the divergence has to be handled by the usual process of renormalization. However, the finite temperature part is completely free from ultraviolet divergence. This is a general feature of finite temperature field theories that temperature does not introduce any new ultraviolet divergence. We will return to this question later within the context of real time formalisms for finite temperature field theories. Let us also note that (see (35)) the finite temperature integrals are, in general hard to evaluate and cannot be evaluated in a closed form. However, we can always make a high temperature expansion (small $`\beta `$) which would give the temperature dependent correction to the mass as
$$\mathrm{\Delta }m_\beta ^2\frac{\lambda }{24\beta ^2}=\frac{\lambda T^2}{24}$$
(36)
This shows that temperature induces a mass correction which is positive. Intuitively, it is clear that this is the behavior we would expect from a particle moving in a medium and, furthermore, the positivity of this correction is crucial in the study of symmetry restoration in field theories with spontaneous symmetry breaking.
This gives a flavor of calculations at finite temperature, particularly, in the imaginary time (Matsubara) formalism. It is worth noting here that, by construction, the imaginary time formalism would describe physical systems in equilibrium quite well. Since we have traded the time variable for temperature, it is well suited to calculate static, equilibrium quantities. Slow temperature dependence can, however, be brought in by analytically rotating the final result to Minkowski time . This rotation is, on the other hand, nontrivial since we only have information about quantities at discrete energy values in the Euclidean space. The imaginary time formalism is not suitable to discuss non-equilibrium phenomena.
## 4 Thermo Field Dynamics
As we have seen, in the imaginary time formalism, the time variable is traded for the temperature. However, in studying various processes, it is desirable to have the time coordinate in addition to the temperature. Formalisms where this can be achieved are known as the real time formalisms and there are two distinct, but equivalent such formalisms. In this section, we will discuss the formalism of thermo field dynamics returning to the alternate formalism in the next section.
Let us recall from (11) that the ensemble average of any observable is given by
$`𝒪_\beta `$ $`=`$ $`{\displaystyle \frac{1}{Z(\beta )}}\mathrm{Tr}e^\beta 𝒪`$ (37)
$`=`$ $`{\displaystyle \frac{1}{Z(\beta )}}{\displaystyle \underset{n}{}}e^{\beta E_n}n|𝒪|n`$
Here, we have assumed that the eigenvalues of $``$ are discrete, for simplicity, and that
$`|n`$ $`=`$ $`E_n|n`$
$`m|n`$ $`=`$ $`\delta _{mn}`$
$`{\displaystyle \underset{n}{}}|nn|`$ $`=`$ $`I`$ (38)
At zero temperature, we know that the Feynman diagrams correspond to vacuum expectation values of time ordered products. Thus, intuitively, it is clear that if we can express the ensemble averages as expectation values in some vacuum (say, a thermal vacuum), then, we can take over all the diagrammatic machinery of the zero temperature field theory. The question, therefore, is whether we can define a vacuum, say $`|0,\beta `$, such that we can write any ensemble average as
$$𝒪_\beta =0,\beta |𝒪|0,\beta =\frac{1}{Z(\beta )}\underset{n}{}e^{\beta E_n}n|𝒪|n$$
(39)
Let us suppose that we can define such a thermal vacuum state as a linear superposition of the states in our physical Hilbert space, namely,
$$|0,\beta =\underset{n}{}|nn|0,\beta =\underset{n}{}f_n(\beta )|n$$
(40)
This would lead to
$$0,\beta |𝒪|0,\beta =\underset{n,m}{}f_n^{}(\beta )f_m(\beta )n|𝒪|m$$
(41)
Consequently, this would coincide with eq. (39) only if
$$f_n^{}(\beta )f_m(\beta )=\frac{1}{Z(\beta )}e^{\beta E_n}\delta _{mn}$$
(42)
Since $`f_n`$’s are ordinary numbers and eq. (42) is more like an orthonormality condition, it is clear that we cannot satisfy this condition (and, therefore, define a thermal vacuum with the right properties) if we restrict ourselves to the original Hilbert space.
On the other hand, it is also clear from this analysis that if $`f_n`$’s, somehow, behave like a state vector, then, the condition in eq. (42) can be easily satisfied. In fact, let us introduce a fictitious system identical to our original system and denote it by a tilde system. The states in the combined Hilbert space of this doubled system would have the form
$$|n,\stackrel{~}{m}=|n|\stackrel{~}{m}$$
Let us assume that the thermal vacuum can be written as a linear superposition of states in this doubled Hilbert space of the form
$$|0,\beta =\underset{n}{}f_n(\beta )|n,\stackrel{~}{n}=\underset{n}{}f_n(\beta )|n|\stackrel{~}{n}$$
(43)
This would lead to
$`0,\beta |𝒪|0,\beta `$ $`=`$ $`{\displaystyle \underset{n,m}{}}f_n^{}(\beta )f_m(\beta )n,\stackrel{~}{n}|𝒪|m,\stackrel{~}{m}`$ (44)
$`=`$ $`{\displaystyle \underset{n,m}{}}f_n^{}(\beta )f_m(\beta )n|𝒪|m\delta _{n,m}`$
$`=`$ $`{\displaystyle \underset{n}{}}f_n^{}(\beta )f_n(\beta )n|𝒪|n`$
In deriving this result, we have used the fact that an operator of the original system does not act on states of the tilde system and vice versa. The result in eq. (44) is quite interesting because it says that if we choose
$$f_n^{}(\beta )f_n(\beta )=\frac{e^{\beta E_n}}{Z(\beta )}\mathrm{or},f_n(\beta )=f_n^{}(\beta )=\frac{e^{\beta E_n/2}}{Z^{1/2}(\beta )}$$
(45)
then, eq. (44) would, indeed, coincide with the ensemble average in eq. (39).
This analysis shows that it is possible to introduce a thermal vacuum such that the ensemble average of any operator can be written as the expectation value of the operator in the thermal vacuum. The price one has to pay is that the Hilbert space needs to be doubled. The advantage, on the other hand, lies in the fact that the description would now involve both time and temperature (since we have not traded time for temperature) and all the diagrammatic methods of zero temperature field theory can now be taken over directly.
### Fermionic Oscillator
To get a flavor for things in this formalism, let us analyze in some detail the simple quantum mechanical system of the fermionic oscillator. The Hamiltonian for the system is given by ($`\mathrm{}=1`$)
$$H=\omega a^{}a$$
(46)
Here, the fermionic creation and annihilation operators satisfy the canonical anti-commutation relations
$`[a,a^{}]_+`$ $`=`$ $`1`$
$`[a,a]_+`$ $`=`$ $`[a^{},a^{}]_+=0`$ (47)
In this case, the spectrum of the Hamiltonian is quite simple and the Hilbert space is two dimensional with the basis states given by $`|0`$ and $`|1=a^{}|0`$.
According to the general philosophy of thermo field dynamics, we are supposed to introduce a fictitious tilde system which is identical to our original system. Thus, we define
$$\stackrel{~}{H}=\omega \stackrel{~}{a}^{}\stackrel{~}{a}$$
(48)
with the anti-commutation relations
$`[\stackrel{~}{a},\stackrel{~}{a}^{}]_+`$ $`=`$ $`1`$
$`[\stackrel{~}{a},\stackrel{~}{a}]_+`$ $`=`$ $`[\stackrel{~}{a}^{},\stackrel{~}{a}^{}]_+=0`$ (49)
Furthermore, we assume that the creation and the annihilation operators for the tilde and the non-tilde systems anti-commute.
The Hilbert space for the combined space is now four dimensional and, following our earlier discussion, we choose the thermal vacuum to be
$$|0,\beta =f_0(\beta )|0|\stackrel{~}{0}+f_1(\beta )|1|\stackrel{~}{1}$$
(50)
The normalization of the thermal vacuum gives
$$0,\beta |0,\beta =|f_0(\beta )|^2+|f_1(\beta )|^2=1$$
(51)
while the expectation value of the number operator gives
$$0,\beta |N|0,\beta =0,\beta |a^{}a|0,\beta =|f_1(\beta )|^2=\frac{1}{e^{\beta \omega }+1}$$
(52)
From these, we can obtain
$$f_0(\beta )=\frac{1}{\sqrt{1+e^{\beta \omega }}},f_1(\beta )=\frac{e^{\beta \omega /2}}{\sqrt{1+e^{\beta \omega }}}$$
(53)
so that we can write
$$|0,\beta =\frac{1}{\sqrt{1+e^{\beta \omega }}}\left(|0,\stackrel{~}{0}+e^{\beta \omega /2}|1,\stackrel{~}{1}\right)$$
(54)
To further understand the properties of this system, let us note that we can define a Hermitian operator in this doubled space
$$G(\theta )=i\theta (\beta )\left(\stackrel{~}{a}aa^{}\stackrel{~}{a}^{}\right)$$
(55)
This would, in turn, lead to a formally unitary operator
$$U(\beta )=e^{iG(\theta )}$$
(56)
which would connect the thermal vacuum to the vacuum of the doubled space, namely,
$$U(\beta )|0,\stackrel{~}{0}=\mathrm{cos}\theta (\beta )|0,\stackrel{~}{0}+\mathrm{sin}\theta (\beta )|1,\stackrel{~}{1}=|0,\beta $$
(57)
provided
$$\mathrm{cos}\theta (\beta )=f_0(\beta )=\frac{1}{\sqrt{1+e^{\beta \omega }}},\mathrm{sin}\theta (\beta )=f_1(\beta )=\frac{e^{\beta \omega /2}}{\sqrt{1+e^{\beta \omega }}}$$
(58)
The unitary operator would also induce a transformation on the operators of the form
$$𝒪(\beta )=U(\beta )𝒪U^{}(\beta )$$
(59)
In particular, this would give
$`a(\beta )`$ $`=`$ $`\mathrm{cos}\theta (\beta )a\mathrm{sin}\theta (\beta )\stackrel{~}{a}^{}`$
$`\stackrel{~}{a}(\beta )`$ $`=`$ $`\mathrm{cos}\theta (\beta )\stackrel{~}{a}+\mathrm{sin}\theta (\beta )a^{}`$ (60)
as well as their Hermitian conjugates. These operators would satisfy the same anti-commutation relations as the original ones and we can think of them as the thermal creation and annihilation operators. Consequently, we can build up the thermal Hilbert space starting from $`|0,\beta `$ and the thermal creation operators.
In particular, it is trivial to check, using (57) that the thermal vacuum satisfies
$`a(\beta )|0,\beta `$ $`=`$ $`(\mathrm{cos}\theta (\beta )a\mathrm{sin}\theta (\beta )\stackrel{~}{a}^{})|0,\beta =0`$
$`\stackrel{~}{a}(\beta )|0,\beta `$ $`=`$ $`(\mathrm{cos}\theta (\beta )\stackrel{~}{a}+\mathrm{sin}\theta (\beta )a^{})|0,\beta =0`$ (61)
This is quite interesting for it says that annihilating a particle in the thermal vacuum is equivalent to creating a tilde particle and vice versa. Consequently, we can intuitively think of the tilde particles as kind of hole states of the particles or particle states of the heat bath. This gives a nice intuitive meaning to the doubling of the degrees of freedom in thermo field dynamics. Namely, an isolated system in thermal equilibrium really consists of two components – the original system and the heat bath.
We also note here that although the operator connecting the thermal vacuum to the vacuum in the doubled space is formally unitary, it is more like a Bogoliubov transformation. In more complicated models with an infinite number of degrees of freedom (namely, in field theories) such an operator takes us to a unitarily inequivalent Hilbert space. Let us also note here, for future use, the simple formula following from eq. (60) that
$$\left(\begin{array}{c}a(\beta )\\ \stackrel{~}{a}^{}(\beta )\end{array}\right)=\stackrel{~}{U}(\beta )\left(\begin{array}{c}a\\ \stackrel{~}{a}^{}\end{array}\right)$$
(62)
where
$$\stackrel{~}{U}(\beta )=\left(\begin{array}{cc}\hfill \mathrm{cos}\theta (\beta )& \hfill \mathrm{sin}\theta (\beta )\\ \hfill \mathrm{sin}\theta (\beta )& \hfill \mathrm{cos}\theta (\beta )\end{array}\right)$$
(63)
Finally, let us conclude the discussion of this example by noting that the states in the thermal Hilbert space are eigenstates of neither $`H`$ nor $`\stackrel{~}{H}`$. Rather they are the eigenstates of the operator
$$\widehat{H}=H\stackrel{~}{H}$$
(64)
Furthermore, this combination of the Hamiltonians is also invariant under the unitary transformation of (57). This is, indeed, the Hamiltonian that governs the dynamics of the combined system.
### Bosonic Oscillator
The analysis for the case of the bosonic oscillator is quite analogous to the discussion of the fermionic oscillator. Therefore, without going into too much detail, let us summarize the results. First, the Hamiltonian for the system is given by
$$H=\omega a^{}a$$
(65)
much like the fermionic oscillator. However, the creation and annihilation operators satisfy canonical commutation relations of the form
$`[a,a^{}]`$ $`=`$ $`1`$
$`[a,a]`$ $`=`$ $`[a^{},a^{}]=0`$ (66)
The Hilbert space for the bosonic oscillator is infinite dimensional with the energy eigenstates given by
$$H|n=n\omega |n,n=0,1,2,\mathrm{}$$
(67)
According to the general discussions of thermo field dynamics, we introduce an identical, but fictitious tilde system with the Hamiltonian
$$\stackrel{~}{H}=\omega \stackrel{~}{a}^{}\stackrel{~}{a}$$
(68)
The tilde creation and annihilation operators are expected to satisfy commutation relations analogous to (66). Furthermore, the tilde operators are supposed to commute with the original operators of the theory.
Following the discussion of the earlier section, we can determine the thermal vacuum state in this case to be
$$|0,\beta =(1e^{\beta \omega })^{1/2}\underset{n=0}{\overset{\mathrm{}}{}}e^{n\beta \omega /2}|n,\stackrel{~}{n}$$
(69)
As in the fermionic oscillator, we can introduce the Hermitian operator
$$G(\theta )=i\theta (\beta )\left(\stackrel{~}{a}aa^{}\stackrel{~}{a}^{}\right)$$
(70)
and the unitary operator
$$U(\beta )=e^{iG(\theta )}$$
(71)
Then, it is straightforward to check and see that the unitary operator connects the thermal vacuum to the vacuum of the doubled space provided
$$\mathrm{cosh}\theta (\beta )=\frac{1}{\sqrt{1e^{\beta \omega }}}\mathrm{sinh}\theta (\beta )=\frac{e^{\beta \omega /2}}{\sqrt{1e^{\beta \omega }}}$$
(72)
The unitary operator induces a transformation of the operators of the form
$$𝒪(\beta )=U(\beta )𝒪U^{}(\beta )$$
(73)
leading to
$`a(\beta )`$ $`=`$ $`\mathrm{cosh}\theta (\beta )a\mathrm{sinh}\theta (\beta )\stackrel{~}{a}^{}`$
$`\stackrel{~}{a}(\beta )`$ $`=`$ $`\mathrm{cosh}\theta (\beta )\stackrel{~}{a}\mathrm{sinh}\theta (\beta )a^{}`$
and similarly for the Hermitian conjugates. As we have seen in the last section, these can be thought of as the creation and annihilation operators for the thermal Hilbert space. In particular, the thermal vacuum is easily seen to satisfy
$`a(\beta )|0,\beta `$ $`=`$ $`(\mathrm{cosh}\theta (\beta )a\mathrm{sinh}\theta (\beta )\stackrel{~}{a}^{})|0,\beta =0`$
$`\stackrel{~}{a}(\beta )|0,\beta `$ $`=`$ $`(\mathrm{cosh}\theta (\beta )\stackrel{~}{a}\mathrm{sinh}\theta (\beta )a^{})|0,\beta =0`$ (74)
This, again, reinforces the intuitive picture of doubling in thermo field dynamics. Let us also note here, for future use, the simple formula following from (73)
$$\left(\begin{array}{c}a(\beta )\\ \stackrel{~}{a}^{}(\beta )\end{array}\right)=\stackrel{~}{U}(\beta )\left(\begin{array}{c}a\\ \stackrel{~}{a}^{}\end{array}\right)$$
(75)
where
$$\stackrel{~}{U}(\beta )=\left(\begin{array}{cc}\hfill \mathrm{cosh}\theta (\beta )& \hfill \mathrm{sinh}\theta (\beta )\\ \hfill \mathrm{sinh}\theta (\beta )& \hfill \mathrm{cosh}\theta (\beta )\end{array}\right)$$
(76)
### Field Theory
The extension of these results to a field theory is quite straightforward once we keep in mind that, at the free level, a quantum field theory is simply an infinite collection of oscillators with frequencies dependent on the momentum of the mode. Consequently, the thermal vacuum, in this case, would be connected to the vacuum of the doubled space as
$$|0,\beta =U(\beta )|0,\stackrel{~}{0}=e^{iG(\theta )}|0,\stackrel{~}{0}$$
(77)
where
$$G(\theta )=i\underset{\stackrel{}{k}}{}\theta _\stackrel{}{k}(\beta )\left(\stackrel{~}{a}_\stackrel{}{k}a_\stackrel{}{k}a_\stackrel{}{k}^{}\stackrel{~}{a}_\stackrel{}{k}^{}\right)$$
(78)
with, say, for bosons,
$$\mathrm{cosh}\theta _\stackrel{}{k}(\beta )=\frac{1}{\sqrt{1e^{\beta \omega _k}}}\mathrm{sinh}\theta _\stackrel{}{k}(\beta )=\frac{e^{\beta \omega _k/2}}{\sqrt{1e^{\beta \omega _k}}}$$
(79)
Here, for a relativistic theory, we have
$$\omega _k=\sqrt{\stackrel{}{k}^2+m^2}$$
Let us next note that, at zero temperature, the original fields are decoupled from the tilde fields. Thus, if we were to define a doublet of fields (real scalar field) as in eq. (75)
$$\mathrm{\Phi }=\left(\begin{array}{c}\varphi \\ \stackrel{~}{\varphi }\end{array}\right)$$
(80)
then, at zero temperature the propagator is defined to be (This is not to be confused with the generator of the Bogoliubov transformations in eqs. (55), (70) and (78))
$$iG(xy)=0,\stackrel{~}{0}|T(\mathrm{\Phi }(x)\mathrm{\Phi }(y))|0,\stackrel{~}{0}$$
which has the momentum space representation
$$G(k)=\left(\begin{array}{cc}\frac{1}{k^2m^2+iϵ}& 0\\ 0& \frac{1}{k^2m^2iϵ}\end{array}\right)$$
(81)
Given this, the finite temperature propagator can be determined to be
$`iG_\beta (xy)`$ $`=`$ $`0,\beta |T(\mathrm{\Phi }(x)\mathrm{\Phi }(y))|0,\beta `$ (82)
$`=`$ $`0,\stackrel{~}{0}|U^{}(\beta )T(\mathrm{\Phi }(x)\mathrm{\Phi }(y))U(\beta )|0,\stackrel{~}{0}`$
Using now the generalization of eqs. (73,75,76), the momentum representation for the propagator can be determined to be
$`G_\beta (k)`$ $`=`$ $`\stackrel{~}{U}(\theta _\stackrel{}{k})G(k)\stackrel{~}{U}^T(\theta _\stackrel{}{k})`$ (87)
$`=`$ $`\left(\begin{array}{cc}\frac{1}{k^2m^2+iϵ}& 0\\ 0& \frac{1}{k^2m^2iϵ}\end{array}\right)2i\pi n_B(|k^0|)\delta (k^2m^2)\left(\begin{array}{cc}\hfill 1& \hfill e^{\beta |k^0|/2}\\ \hfill e^{\beta |k^0|/2}& \hfill 1\end{array}\right)`$
There are several things to note from the structure of the propagator which are quite general for a real time formalism. First, the propagator is a $`2\times 2`$ matrix, a consequence of the doubling of the degrees of freedom. Second, the propagator is a sum of two parts – one representing the zero temperature part and the other representing the true temperature dependent corrections. The propagator is still the Greens function for the free operator of the theory, but corresponding to different boundary conditions (remember the KMS condition in eq. (12)). While the zero temperature part of the propagator corresponds, as usual, to the exchange of a virtual particle, the temperature dependent part represents an on-shell contribution (because of the delta function). In fact, the intuitive meaning of the temperature dependent correction is quite clear. In a hot medium, there is a distribution of real particles and the temperature dependent part merely represents the possibility that a particle, in addition to having virtual exchanges, can also emit or absorb a real particle of the medium.
Since the temperature dependent part of the propagator is on-shell, it is clear that there can be no new ultraviolet divergence generated at finite temperature. All the counter terms needed to renormalize the theory at zero temperature would be sufficient for studies at finite temperature as well. (Of course, the infrared behavior is another story. Infrared divergence, in a field theory, becomes much more severe at finite temperature, a topic that I will not get into.) There is an alternate way to visualize this. At finite temperature the distribution of the real particles is Boltzmann suppressed as we go up in energy and, consequently, thermal corrections corresponding to infinite energy cannot arise.
Once, we have the propagator, we can venture to do a diagrammatic calculation in this formalism. The only things missing are the interaction vertices of the theory. There is a well defined procedure (called the tilde conjugation rule) to construct the complete Lagrangian from which to construct the vertices. Very simply, it corresponds to what we have noted earlier, namely, the dynamical Hamiltonian (and, similarly, the Lagrangian) is as given in eq. (64). It is simply the difference between the original and the tilde Hamiltonians. Thus, we see that the complete theory would contain two kinds of vertices – one for the original fields while the second for the tilde fields. The vertices for the tilde fields will have a relative negative sign corresponding to the original vertices. Given the vertices and the propagator, it is now straightforward to carry out any diagrammatic calculation to any order. Let me emphasize here that although, at the tree level, there is no vertex containing both the original and the tilde fields, such vertices would be generated at higher loops because of the nontrivial matrix structure of the propagator.
Thermo Field dynamics is a real time formalism. But, more than that, it is really an operator formalism and hence very well suited to study various operator questions such as the structure of the thermal vacuum, the theorems on symmetry breaking etc. It can also be given a path integral representation and corresponds to choosing a specific time contour in the complex $`t`$ plane (remember that the imaginary time formalism also corresponds to choosing a specific time contour, namely, along the imaginary time axis) and I will come back to this question in the next section. However, once again from the philosophy of thermo field dynamics, it is clear that, it is a natural formalism to describe equilibrium phenomena where quantities depend on both time and temperature. While there are several attempts to generalize this to include non-equilibrium phenomena, there does not yet exist a complete description.
## 5 Closed Time Path Formalism
The closed time path formalism is also a real time formalism which was formulated much earlier than thermo field dynamics within the context of non-equilibrium phenomena . The two formalisms are, in some sense, complementary to each other although the closed time path formalism can describe both equilibrium and non-equilibrium phenomena with equal ease.
The basic idea behind the closed time path formalism is the fact that when a quantum mechanical system is in a mixed state, as is the case in the presence of a heat bath, the system can be naturally described in terms of a density matrix defined, in the Schrödinger picture, as
$$\rho (t)=\underset{n}{}p_n|\psi _n(t)\psi _n(t)|$$
(88)
Here, $`p_n`$ represents the probability for finding the quantum mechanical system in the state $`|\psi _n(t)`$ and, for simplicity, we have assumed the quantum mechanical states to form a discrete set. It is $`p_n`$ which contains information regarding the surrounding which is hard to determine, but, being a probability, it satisfies
$$\underset{n}{}p_n=1$$
Given the density matrix, the ensemble average of any operator can be calculated in the Schrödinger picture as
$$𝒪(t)=\underset{n}{}p_n\psi _n(t)|𝒪|\psi _n(t)=\mathrm{Tr}\rho (t)𝒪$$
(89)
The ensemble average, in this case, naturally develops a time dependence from the time dependence of the density matrix. In this formalism, we can naturally define an entropy as
$$S=\underset{n}{}p_n\mathrm{ln}p_n$$
which is by definition positive semi-definite and measures the order (or lack of it) in an ensemble.
The state vectors satisfy the Schrödinger equation ($`\mathrm{}=1`$)
$$i\frac{|\psi (t)}{t}=H|\psi (t)$$
From this, we can determine the time evolution of the density matrix which turns out to be the Liouville equation
$$i\frac{\rho (t)}{t}=[H,\rho (t)]$$
(90)
In deriving this, we have assumed that the probabilities do not change with time (appreciably) implying that entropy remains constant during such an evolution. The reason for this assumption is our lack of knowledge about the time evolution of the surrounding such as the heat bath. On the other hand, adiabatic evolutions do arise frequently in physical systems and, consequently, we would continue with this assumption.
Let us note that eq. (90) has a simple solution of the form
$$\rho (t)=U(t,0)\rho (0)U^{}(t,0)=U(t,0)\rho (0)U(0,t)$$
(91)
where the time evolution operator has the general form
$$U(t,t^{})=T\left(e^{i_t^{}^t𝑑t^{\prime \prime }H(t^{\prime \prime })}\right)$$
(92)
Furthermore, it satisfies the semi-group properties
$`U(t_1,t_2)U(t_2,t_1)`$ $`=`$ $`1`$
$`U(t_1,t_2)U(t_2,t_3)`$ $`=`$ $`U(t_1,t_3)\mathrm{for}t_1>t_2>t_3`$ (93)
In particular, let us note that if the Hamiltonian is time independent, eq. (91) takes the simple form
$$\rho (t)=e^{iHt}\rho (0)e^{iHt}$$
and, furthermore, if the Hamiltonian commutes with $`\rho (0)`$, the density matrix would be time independent, describing a system in equilibrium. This would be true, for example, if the states in eq. (88) are stationary states. This is also true if the probabilities have a Boltzmann distribution in which case, we refer to the system as being in thermal equilibrium. However, we will not restrict to any such special case allowing for the formalism to accommodate both equilibrium and non-equilibrium phenomena.
Keeping in mind the fact that we are ultimately interested in a thermal ensemble, let us choose
$$\rho (0)=\frac{e^{\beta H_i}}{\mathrm{Tr}e^{\beta H_i}}$$
(94)
for some $`H_i`$. Since the density matrix is a positive Hermitian matrix with unit trace, mathematically, this is allowed. But, more important is the physical reason behind such a choice. Namely, we can think of the dynamical Hamiltonian of our system as
$$H(t)=\{\begin{array}{ccc}H_i& \mathrm{for}\hfill & \mathrm{Re}t0\hfill \\ (t)& \mathrm{for}\hfill & \mathrm{Re}t0\hfill \end{array}$$
(95)
This would correspond to the fact that we prepare our system in a equilibrium state at temperature $`\frac{1}{\beta }`$ for negative times and let the system evolve, for positive times, with the true Hamiltonian $``$ which may be time dependent. If $`(t)=H_i`$, then, the system will evolve in equilibrium and not otherwise.
With eq. (95) in mind, we note that we can write
$$\rho (0)=\frac{U(Ti\beta ,T)}{\mathrm{Tr}U(Ti\beta ,T)}$$
(96)
where $`T`$ is assumed to be a large negative time (and not the temperature) and $`T\mathrm{}`$ at the end. Using the semi-group properties of the time evolution operator, it is easy to see that the ensemble average of any operator can now be represented as
$`𝒪_\beta `$ $`=`$ $`\mathrm{Tr}\rho (t)𝒪`$ (97)
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}U(t,0)U(Ti\beta ,T)U(0,t)𝒪}{\mathrm{Tr}U(Ti\beta ,T)}}`$
$`=`$ $`{\displaystyle \frac{\mathrm{Tr}U(Ti\beta ,T)U(T,T^{})U(T^{},t)𝒪U(t,T)}{\mathrm{Tr}U(Ti\beta ,T)U(T,T^{})U(T^{},T)}}`$
where we have introduced a large positive time $`T^{}`$ and assume that $`T^{}\mathrm{}`$ at the end. This gives a nice representation to the ensemble average of any operator. Namely, we let the system evolve from a large negative time $`T`$ to $`t`$ where the appropriate operator $`𝒪`$ is inserted. The system then, evolves from $`t`$ to a large positive time $`T^{}`$ and back from $`T^{}`$ to $`T`$ and then, continues evolving along the imaginary branch from $`T`$ to $`Ti\beta `$. Since the matrix elements of the time evolution operator can be given a path integral representation, it is clear that the ensemble average of any operator can also be given a path integral representation in this formalism corresponding to the specific contour in the complex time plane as described above. Although the specific contour has three branches – one along the real axis increasing with time, the second also along the real axis decreasing with time and the third along the negative imaginary axis – in the limit $`T\mathrm{}`$ and $`T^{}\mathrm{}`$, it can be shown that the third branch gets decoupled from the other two (the factors in the propagators connecting such branches are asymptotically damped). Consequently, in this limit, we are effectively dealing with two branches leading to the name “closed time path formalism” . In this contour, then, the time integration has to be thought of as
$$_c𝑑t=_{\mathrm{}}^{\mathrm{}}𝑑t_+_{\mathrm{}}^{\mathrm{}}𝑑t_{}$$
(98)
where the relative negative sign arises because time is decreasing in the second branch of the time contour.
The doubling of the degrees of freedom, in this formalism, is now clear. To have a path integral description, we must specify the fields on both the branches of the contour. Or, equivalently, we can use just the positive branch and double the field degrees of freedom. Namely, corresponding to every original field, say $`\varphi _+`$, we must introduce a second field $`\varphi _{}`$ and remember that the action for the $`\varphi _{}`$ fields must have a relative negative sign arising from eq. (98), namely, that time is decreasing along the second branch.
### Scalar Field Theory
Just as an example, let us study next the self-interacting scalar field theory in some detail. The Lagrangian density is the same as in eq. (30), but following the earlier discussion, we should take the complete Lagrangian density for the system to be
$$=(\varphi _+)(\varphi _{})$$
(99)
where
$$(\varphi )=\frac{1}{2}_\mu \varphi ^\mu \varphi \frac{m^2}{2}\varphi ^2\frac{\lambda }{4!}\varphi ^4\lambda >0$$
(100)
The Feynman propagator can again be determined for this theory and would have a $`2\times 2`$ matrix structure because of the doubling of the field degrees of freedom. It can be determined subject to compatibility with the KMS conditions and has the form in the momentum space
$$G(k)=\left(\begin{array}{cc}G_{++}(k)& G_+(k)\\ G_+(k)& G_{}(k)\end{array}\right)$$
(101)
with
$`G_{++}(k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2m^2+iϵ}}2i\pi n_B(|k^0|)\delta (k^2m^2)`$
$`G_+(k)`$ $`=`$ $`2i\pi \left(\theta (k^0)+n_B(|k^0|)\right)\delta (k^2m^2)`$
$`G_+(k)`$ $`=`$ $`2i\pi \left(\theta (k^0)+n_B(|k^0|)\right)\delta (k^2m^2)`$
$`G_{}(k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2m^2iϵ}}2i\pi n_B(|k^0|)\delta (k^2m^2)`$ (102)
There are several things to note from the structure of this propagator. First, as in the case of the propagator in thermo field dynamics, here, too, we see that the propagator naturally is a sum of two parts – the temperature independent part and the temperature dependent part. But, more interestingly, here the propagator has the simplification that the temperature dependent part of every component is the same which leads to various simplifications in actual studies of thermal quantities. Furthermore, not all the components of the propagator are independent. In fact, it is easily seen that (this can be traced back to their definition)
$$G_{++}(k)+G_{}(k)=G_+(k)+G_+(k)$$
These are known as the causal propagators of the theory and are useful in diagrammatic evaluation. There is, of course, another kind of propagator, conventionally known as the physical propagators and is defined as
$$\widehat{G}(k)=\left(\begin{array}{cc}0& G_A(k)\\ G_R(k)& G_C(k)\end{array}\right)$$
(103)
where $`G_A`$, $`G_R`$ and $`G_C`$ are known as the advanced, retarded and the correlated Greens functions. These are quite useful in the study of various phenomena such as the linear response theory. The important thing to observe is that the causal and the physical propagators are connected through a unitary transformation
$$\widehat{G}(k)=QGQ^{}$$
(104)
where
$$Q=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)$$
(105)
It can be determined from this that, at the tree level,
$`G_A(k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2m^2iϵk^0}}`$
$`G_R(k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2m^2+iϵk^0}}`$
$`G_C(k)`$ $`=`$ $`2i\pi \left(1+2n_B(|k^0|)\right)\delta (k^2m^2)`$ (106)
as they should be.
The diagrammatic calculations can now be easily understood in this formalism. The vertices can be read out from the Lagrangian density in eq. (99). There are two kinds of vertices, one for the original fields, $`\varphi _+`$, and the other for the doubled fields, $`\varphi _{}`$. The vertices for the $`\varphi _{}`$ fields are the same as those for the $`\varphi _+`$ fields except for a relative sign. With the vertices and the causal propagators, one can now carry out the calculation of any observable to any order in perturbation theory. As before, we note that, although there is no coupling between the $`\varphi _+`$ and $`\varphi _{}`$ fields at the tree level, higher order corrections would, in general, couple them.
As an example, let us calculate the one loop mass correction in this theory. There will be two such diagrams to calculate – one for the $`\varphi _+`$ field and the other for the $`\varphi _{}`$ field. The mass correction for the $`\varphi _+`$ field is readily seen to be
$`i\mathrm{\Delta }m_+^2`$ $`=`$ $`{\displaystyle \frac{(i\lambda )}{2}}{\displaystyle \frac{d^4k}{(2\pi )^4}iG_{++}(k)}`$ (107)
$`=`$ $`{\displaystyle \frac{\lambda }{2}}{\displaystyle \frac{d^4k}{(2\pi )^4}\left(\frac{1}{k^2m^2+iϵ}2i\pi n_B(|k^0|)\delta (k^2m^2)\right)}`$
$`=`$ $`i(\mathrm{\Delta }m_0^2+\mathrm{\Delta }m_\beta ^2)`$
where it is easily seen that the temperature independent part has the form
$$\mathrm{\Delta }m_0^2=\frac{\lambda }{4}\frac{d^3k}{(2\pi )^3}\frac{1}{\omega _k}$$
(108)
while the temperature dependent part is given by
$$\mathrm{\Delta }m_\beta ^2=\frac{\lambda }{2}\frac{d^3k}{(2\pi )^3}\frac{n_B(\omega _k)}{\omega _k}=\frac{\lambda }{2}\frac{d^3k}{(2\pi )^3}\frac{1}{\omega _k}\frac{1}{e^{\beta \omega _k}1}$$
(109)
These can be compared with the corresponding terms in eq. (35). We can also calculate the mass correction for the $`\varphi _{}`$ field. With a little bit of analysis, it is seen that
$$\mathrm{\Delta }m_{}^2=\mathrm{\Delta }m_+^2$$
## 6 Feynman Parameterization
So far, we have described the various formalisms that can be used to do calculations at finite temperature. However, actual calculations lead to many subtle, but interesting features of theories at finite temperature. One immediate and obvious feature, of course, is that finite temperature effects break Lorentz invariance. Namely, in studying a system at finite temperature, one has to go to a specific frame where the heat bath is at rest and, consequently, Lorentz invariance will no longer be manifest. This is, of course, already manifest at the level of propagators. For example, the structure of the propagators in eqs. (87) or (102) clearly displays a Lorentz non-invariant structure. The consequence of this is that an amplitude calculated at finite temperature, say for example, the self-energy $`\mathrm{\Pi }(p^0,\stackrel{}{p})`$ depends on the external energy and momentum independently. In fact, the self-energy becomes a non-analytic function of these two variables at the origin and two different ways of approaching the origin in this space leads to distinct plasmon and screening masses . Thus, such non-analyticities are quite physical and their origin can be traced back to the fact that, at finite temperature, there are new channels of reactions possible leading to new branch cuts which give rise to such discontinuities . (To be absolutely fair, it is worth noting that statistical mechanics can be formulated in a covariant way. In such a case, one finds that there is a larger number of Lorentz invariant variables that can be defined on which amplitudes can depend. The non-analyticity in $`p^0`$ and $`\stackrel{}{p}`$ can then be translated to a non-analyticity in these new, Lorentz invariant variables .)
There are, of course, some other kinds of subtlety that arise which influence the calculations directly at finite temperature. We will discuss one such subtlety in this section. Let us note that a particularly useful formula in the evaluation of amplitudes at zero temperature is the Feynman combination formula given by
$$\frac{1}{A+iϵ}\frac{1}{B+iϵ}=_0^1\frac{dx}{[x(A+iϵ)+(1x)(B+iϵ)]^2}$$
(110)
This can be directly checked by evaluating the $`x`$ integral on the right hand side.
This formula is extremely useful and works at zero temperature mainly because the Feynman propagators have the same analytic structure, namely, they have the same “$`iϵ`$” dependence. In contrast, we note that the finite temperature propagators contain delta functions (see eqs. (87) and (102)) and recalling that
$$\delta (x)=\underset{ϵ0^+}{lim}\frac{1}{2i\pi }\left(\frac{1}{xiϵ}\frac{1}{x+iϵ}\right)$$
we recognize that, at finite temperature, the propagators no longer have the same “$`iϵ`$” dependence. Consequently, in evaluating Feynman amplitudes at finite temperature, we have to combine denominators which do not necessarily have the same “$`iϵ`$” dependence. Keeping this in mind, let us examine the combination of two different denominators with arbitrary analytic dependence.
Without loss of generality, let us choose $`\alpha ,\beta =\pm 1`$ and note that
$`{\displaystyle _0^1}{\displaystyle \frac{dx}{\left[x(A+i\alpha ϵ)+(1x)(B+i\beta ϵ)\right]^2}}`$ $`=`$ $`{\displaystyle \frac{1}{(AB)+i(\alpha \beta )ϵ}}{\displaystyle \frac{1}{x(A+i\alpha ϵ)+(1x)(B+i\beta ϵ)}}|_0^1`$ (111)
$`=`$ $`{\displaystyle \frac{1}{A+i\alpha ϵ}}{\displaystyle \frac{1}{B+i\beta ϵ}}`$
This is, of course, the standard Feynman combination formula. However, let us note that this will not hold if $`0<x_0<1`$ such that
$$x_0=\frac{\beta }{\beta \alpha },\beta A=\alpha B$$
(112)
because, in such a case, the integrand will have a singularity on the real $`x`$-axis inside the interval of integration. In this case, we have
$`{\displaystyle _0^1}{\displaystyle \frac{dx}{\left[x(A+i\alpha ϵ)+(1x)(B+i\beta ϵ)\right]^2}}`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \frac{dx}{\left[x(AB+i(\alpha \beta )ϵ)+B+i\beta ϵ\right]^2}}`$ (113)
$`=`$ $`\underset{\eta 0}{lim}({\displaystyle _0^{\frac{\beta }{\beta \alpha }\eta }}+{\displaystyle _{\frac{\beta }{\beta \alpha }+\eta }^1}){\displaystyle \frac{dx}{\left[x(AB+i(\alpha \beta )ϵ)+B+i\beta ϵ\right]^2}}`$
$`=`$ $`{\displaystyle \frac{1}{A+i\alpha ϵ}}{\displaystyle \frac{1}{B+i\beta ϵ}}2i\pi {\displaystyle \frac{(\alpha \beta )\delta (\beta A\alpha B)}{AB+i(\alpha \beta )ϵ}}`$
In other words, when the parameters of the integrand satisfy eq. (112), the Feynman combination formula of eq. (110) will modify and the general formula follows from eq. (113) to be
$$\frac{1}{A+i\alpha ϵ}\frac{1}{B+i\beta ϵ}=_0^1\frac{dx}{\left[x(A+i\alpha ϵ)+(1x)(B+i\beta ϵ)\right]^2}+2i\pi \frac{(\alpha \beta )\delta (\beta A\alpha B)}{AB+i(\alpha \beta )ϵ}$$
(114)
Note that, condition (112) can only be satisfied (with $`0<x_0<1`$) if $`\alpha `$ and $`\beta `$ are of opposite sign. Indeed, let us note from eq. (114) that the second term vanishes when $`\alpha =\beta =1`$ as is the case at zero temperature. Namely, when propagators with identical “$`iϵ`$” dependence are combined, the standard combination formula of eq. (110) holds. However, if denominators with opposite “$`iϵ`$” dependence are combined, the correct combination formula involves a second term. This is quite crucial at finite temperature and without this second term, one ends up with a wrong result as was discovered in finite temperature calculations the hard way .
## 7 Large Gauge Invariance
Gauge theories are beautiful theories which describe physical forces in a natural manner and because of their rich structure, the study of gauge theories at finite temperature is quite interesting in itself. However, to avoid getting into technicalities, we will not discuss the intricacies of such theories either at zero temperature or at finite temperature. Rather, we will discuss a simple quantum mechanical model, in this section, to bring out some of the new features that temperature brings into such theories – features which are very different from what we expect at zero temperature.
To motivate, let us note that gauge invariance is realized as an internal symmetry in quantum mechanical systems. Consequently, we do not expect a macroscopic external surrounding such as a heat bath to modify gauge invariance. This is more or less what is also found by explicit computations at finite temperature, namely, that gauge invariance and Ward identities continue to hold even at finite temperature . This is certainly the case when one is talking about small gauge transformations for which the parameters of transformation vanish at infinity.
However, there is a second class of gauge invariance, commonly known as large gauge invariance where the parameters do not vanish at infinity and this brings in some new topological character to physical theories. For example, let us consider a $`2+1`$ dimensional Chern-Simons theory of the form
$``$ $`=`$ $`M_{\mathrm{CS}}+_{\mathrm{fermion}}`$ (115)
$`=`$ $`Mϵ^{\mu \nu \lambda }\mathrm{tr}A_\mu (_\nu A_\lambda {\displaystyle \frac{2}{3}}A_\nu A_\lambda )+\overline{\psi }(\gamma ^\mu (i_\mu gA_\mu )m)\psi `$
where $`M`$ is a mass parameter, $`A_\mu `$ a matrix valued non-Abelian gauge field and “tr” stands for the matrix trace. The first term, on the right hand side, is known as the Chern-Simons term which exists only in odd space-time dimensions. We can, of course, also add a Maxwell like term to the Lagrangian and, in that case, the Chern-Simons term behaves like a mass term for the gauge field. Consequently, such a term is also known as a topological mass term (topological because it does not involve the metric). For simplicity of discussion, however, we will not include a Maxwell like term to the Lagrangian.
Under a gauge transformation of the form
$`\psi `$ $``$ $`U^1\psi `$
$`A_\mu `$ $``$ $`U^1A_\mu U{\displaystyle \frac{i}{g}}U^1_\mu U`$ (116)
it is straightforward to check that the action in eq. (115) is not invariant, rather it changes as
$$S=d^3xS+\frac{4\pi M}{g^2}\mathrm{\hspace{0.17em}2}i\pi W$$
(117)
where
$$W=\frac{1}{24\pi ^2}d^3xϵ^{\mu \nu \lambda }\mathrm{tr}_\mu UU^1_\nu UU^1_\lambda UU^1$$
(118)
is known as the winding number. It is a topological quantity which is an integer (Basically, the fermion Lagrangian density is invariant under the gauge transformations, but the Chern-Simons term changes by a total divergence which does not vanish if the gauge transformations do not vanish at infinity. Consequently, the winding number counts the number of times the gauge transformations wrap around the sphere.). For small gauge transformations, the winding number vanishes since the gauge transformations vanish at infinity.
Let us note from eq. (117) that even though the action is not invariant under a large gauge transformation, if $`M`$ is quantized in units of $`\frac{g^2}{4\pi ^2}`$, the change in the action would be a multiple of $`2i\pi `$ and, consequently, the path integral would be invariant under a large gauge transformation. Thus, we have the constraint coming from the consistency of the theory that the coefficient of the Chern-Simons term must be quantized. We have derived this conclusion from an analysis of the tree level behavior of the theory and we have to worry if the quantum corrections can change the behavior of the theory. At zero temperature, an analysis of the quantum corrections shows that the theory continues to be well defined with the tree level quantization of the Chern-Simons coefficient provided the number of fermion flavors is even. The even number of fermion flavors is also necessary for a global anomaly of the theory to vanish and so, everything is well understood at zero temperature.
At finite temperature, however, the situation appears to change drastically. Namely, the fermions induce a temperature dependent Chern-Simons term effectively making
$$MM\frac{g^2}{4\pi }\frac{mN_f}{2|m|}\mathrm{tanh}\frac{\beta |m|}{2}$$
(119)
Here, $`N_f`$ is the number of fermion flavors and this shows that, at zero temperature ($`\beta \mathrm{}`$), $`M`$ changes by an integer (in units of $`g^2/4\pi `$) for an even number of flavors. However, at finite temperature, this becomes a continuous function of temperature and, consequently, it is clear that it can no longer be an integer for arbitrary values of the temperature. It seem, therefore, that temperature would lead to a breaking of large gauge invariance in such a system. This is, on the other hand, completely counter intuitive considering that temperature should have no direct influence on gauge invariance of the theory.
### C-S Theory in $`0+1`$ Dimension
As we have noted, Chern-Simons terms can exist in odd space-time dimensions. Consequently, let us try to understand this puzzle of large gauge invariance in a simple quantum mechanical theory. Let us consider a simple theory of an interacting massive fermion with a Chern-Simons term in $`0+1`$ dimension described by
$$L=\overline{\psi }_j(i_tAm)\psi _j\kappa A$$
(120)
Here, $`j=1,2,\mathrm{},N_f`$ labels the fermion flavors. There are several things to note from this. First, we are considering an Abelian gauge field for simplicity. Second, in this simple model, the gauge field has no dynamics (in $`0+1`$ dimension the field strength is zero) and, therefore, we do not have to get into the intricacies of gauge theories. There is no Dirac matrix in $`0+1`$ dimension as well making the fermion part of the theory quite simple as well. And, finally, the Chern-Simons term, in this case, is a linear field so that we can, in fact, think of the gauge field as an auxiliary field.
In spite of the simplicity of this theory, it displays a rich structure including all the properties of the $`2+1`$ dimensional theory that we have discussed earlier. For example, let us note that under a gauge transformation
$$\psi _je^{i\lambda (t)}\psi _j,AA+_t\lambda (t)$$
(121)
the fermion part of the Lagrangian is invariant, but the Chern-Simons term changes by a total derivative giving
$$S=𝑑tLS2\pi \kappa N$$
(122)
where
$$N=\frac{1}{2\pi }𝑑t_t\lambda (t)$$
(123)
is the winding number and is an integer which vanishes for small gauge transformations. Let us note that a large gauge transformation can have a parametric form of the form, say,
$$\lambda (t)=iN\mathrm{log}\left(\frac{1+it}{1it}\right)$$
(124)
The fact that $`N`$ has to be an integer can be easily seen to arise from the requirement of single-valuedness for the fermion field. Once again, in light of our earlier discussion, it is clear from eq. (122) that the theory is meaningful only if $`\kappa `$, the coefficient of the Chern-Simons term, is an integer.
Let us assume, for simplicity, that $`m>0`$ and compute the correction to the photon one-point function arising from the fermion loop at zero temperature.
$$iI_1=(i)N_f\frac{dk}{2\pi }\frac{i(k+m)}{k^2m^2+iϵ}=\frac{iN_f}{2}$$
(125)
This shows that, as a result of the quantum correction, the coefficient of the Chern-Simons term would change as
$$\kappa \kappa \frac{N_f}{2}$$
As in $`2+1`$ dimensions, it is clear that the coefficient of the Chern-Simons term would continue to be quantized and large gauge invariance would hold if the number of fermion flavors is even. At zero temperature, we can also calculate the higher point functions due to the fermions in the theory and they all vanish. This has a simple explanation following from the small gauge invariance of the theory. Namely, suppose we had a nonzero two point function, then, it would imply a quadratic term in the effective action of the form
$$\mathrm{\Gamma }_2=\frac{1}{2}𝑑t_1𝑑t_2A(t_1)F(t_1t_2)A(t_2)$$
(126)
Furthermore, invariance under a small gauge transformation would imply
$$\delta \mathrm{\Gamma }_2=𝑑t_1𝑑t_2\lambda (t_1)_{t_1}F(t_1t_2)A(t_2)=0$$
(127)
The solution to this equation is that $`F=0`$ so that there cannot be a quadratic term in the effective action which would be local and yet be invariant under small gauge transformations. A similar analysis would show that small gauge invariance does not allow any higher point function to exist at zero temperature.
Let us also note that eq. (127) has another solution, namely,
$$F(t_1t_2)=\mathrm{constant}$$
In such a case, however, the quadratic action becomes non-extensive, namely, it is the square of an action. We do not expect such terms to arise at zero temperature and hence the constant has to vanish for vanishing temperature. As we will see next, the constant does not have to vanish at finite temperature and we can have non-vanishing higher point functions implying a non-extensive structure of the effective action.
The fermion propagator at finite temperature (in the real time formalism) has the form
$`S(p)`$ $`=`$ $`(p+m)\left({\displaystyle \frac{i}{p^2m^2+iϵ}}2\pi n_F(|p|)\delta (p^2m^2)\right)`$ (128)
$`=`$ $`{\displaystyle \frac{i}{pm+iϵ}}2\pi n_F(m)\delta (pm)`$
and the structure of the effective action can be studied in the momentum space in a straightforward manner. However, in this simple model, it is much easier to analyze the amplitudes in the coordinate space. Let us note that the coordinate space structure of the fermion propagator is quite simple, namely,
$$S(t)=\frac{dp}{2\pi }e^{ipt}\left(\frac{i}{pm+iϵ}2\pi n_F(m)\delta (pm)\right)=(\theta (t)n_F(m))e^{imt}$$
(129)
In fact, the calculation of the one point function is trivial now
$$iI_1=(i)N_fS(0)=\frac{iN_f}{2}\mathrm{tanh}\frac{\beta m}{2}$$
(130)
This shows that the behavior of this theory is completely parallel to the $`2+1`$ dimensional theory in that, it would suggest
$$\kappa \kappa \frac{N_f}{2}\mathrm{tanh}\frac{\beta m}{2}$$
and it would appear that large gauge invariance would not hold at finite temperature.
Let us next calculate the two point function at finite temperature.
$`iI_2`$ $`=`$ $`(i)^2{\displaystyle \frac{N_f}{2!}}S(t_1t_2)S(t_2t_1)`$ (131)
$`=`$ $`{\displaystyle \frac{N_f}{2}}n_F(m)(1n_F(m))`$
$`=`$ $`{\displaystyle \frac{N_f}{8}}\mathrm{sech}^2{\displaystyle \frac{\beta m}{2}}={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2!}}{\displaystyle \frac{i}{\beta }}{\displaystyle \frac{(iI_1)}{m}}`$
This shows that the two point function is a constant as we had noted earlier implying that the quadratic term in the effective action would be non-extensive.
Similarly, we can also calculate the three point function trivially and it has the form
$$iI_3=\frac{iN_f}{24}\mathrm{tanh}\frac{\beta m}{2}\mathrm{sech}^2\frac{\beta m}{2}=\frac{1}{2}\frac{1}{3!}\left(\frac{i}{\beta }\right)^2\frac{^2(iI_1)}{m^2}$$
(132)
In fact, all the higher point functions can be worked out in a systematic manner. But, let us observe a simple method of computation for these. We note that because of the gauge invariance (Ward identity), the amplitudes cannot depend on the external time coordinates as is clear from the calculations of the lower point functions. Therefore, we can always simplify the calculation by choosing a particular time ordering convenient to us. Second, since we are evaluating a loop diagram (a fermion loop) the initial and the final time coordinates are the same and, consequently, the phase factors in the propagator (129) drop out. Therefore, let us define a simplified propagator without the phase factor as
$$\stackrel{~}{S}(t)=\theta (t)n_F(m)$$
(133)
so that we have
$$\stackrel{~}{S}(t>0)=1n_F(m),S(t<0)=n_F(m)$$
(134)
Then, it is clear that with the choice of the time ordering, $`t_1>t_2`$, we can write
$`{\displaystyle \frac{\stackrel{~}{S}(t_1t_2)}{m}}`$ $`=`$ $`\beta \stackrel{~}{S}(t_1t_3)\stackrel{~}{S}(t_3t_2)t_1>t_2>t_3`$
$`{\displaystyle \frac{\stackrel{~}{S}(t_2t_1)}{m}}`$ $`=`$ $`\beta \stackrel{~}{S}(t_2t_3)\stackrel{~}{S}(t_3t_1)t_1>t_2>t_3`$ (135)
In other words, this shows that differentiation of a fermionic propagator with respect to the mass of the fermion is equivalent to introducing an external photon vertex (and, therefore, another fermion propagator as well) up to constants. This is the analogue of the Ward identity in QED in four dimensions except that it is much simpler. From this relation, it is clear that if we take a $`n`$-point function and differentiate this with respect to the fermion mass, then, that is equivalent to adding another external photon vertex in all possible positions. Namely, it should give us the $`(n+1)`$-point function up to constants. Working out the details, we have,
$$\frac{I_n}{m}=i\beta (n+1)I_{n+1}$$
(136)
Therefore, the $`(n+1)`$-point function is related to the $`n`$-point function recursively and, consequently, all the amplitudes are related to the one point function which we have already calculated. (Incidentally, this is already reflected in eqs. (131,132)).
With this, we can now determine the full effective action of the theory at finite temperature to be
$`\mathrm{\Gamma }`$ $`=`$ $`i{\displaystyle \underset{n}{}}a^n(iI_n)`$ (137)
$`=`$ $`{\displaystyle \frac{i\beta N_f}{2}}{\displaystyle \underset{n}{}}{\displaystyle \frac{(ia/\beta )^n}{n!}}\left({\displaystyle \frac{}{m}}\right)^{n1}\mathrm{tanh}{\displaystyle \frac{\beta m}{2}}`$
$`=`$ $`iN_f\mathrm{log}\left(\mathrm{cos}{\displaystyle \frac{a}{2}}+i\mathrm{tanh}{\displaystyle \frac{\beta m}{2}}\mathrm{sin}{\displaystyle \frac{a}{2}}\right)`$
where we have defined
$$a=𝑑tA(t)$$
(138)
There are several things to note from this result. First of all, the higher point functions are no longer vanishing at finite temperature and give rise to a non-extensive structure of the effective action. More importantly, when we include all the higher point functions, the complete effective action is invariant under large gauge transformations, namely, under
$$aa+2\pi N$$
(139)
the effective action changes as
$$\mathrm{\Gamma }\mathrm{\Gamma }+NN_f\pi $$
(140)
which leaves the path integral invariant for an even number of fermion flavors. This clarifies the puzzle of large gauge invariance at finite temperature in this model. Namely, when we are talking about large changes (large gauge transformations), we cannot ignore higher order terms if they exist. This may provide a resolution to the large gauge invariance puzzle in the $`2+1`$ dimensional theory as well. However, in spite of several nice analysis , this puzzle has not yet been settled in all its generality in the $`2+1`$ dimensional case.
### Exact Result
In the earlier section, we discussed a perturbative method of calculating the effective action at finite temperature which clarified the puzzle of large gauge invariance. However, this quantum mechanical model is simple enough that we can also evaluate the effective action directly and, therefore, it is worth asking how the perturbative calculations compare with the exact result.
The exact evaluation of the effective action can be done easily using the imaginary time formalism. But, first, let us note that the fermionic part of the Lagrangian in eq. (120) has the form
$$L_f=\overline{\psi }(i_tAm)\psi $$
(141)
where we have suppressed the fermion flavor index for simplicity. Let us note that if we make a field redefinition of the form
$$\psi (t)=e^{i_0^t𝑑t^{}A(t^{})}\stackrel{~}{\psi }(t)$$
(142)
then, the fermionic part of the Lagrangian becomes free, namely,
$$L_f=\overline{\stackrel{~}{\psi }}(i_tm)\stackrel{~}{\psi }$$
(143)
This is a free theory and, therefore, the path integral can be easily evaluated. However, we have to remember that the field redefinition in (142) changes the periodicity condition for the fermion fields. Since the original fermion field was expected to satisfy anti-periodicity
$$\psi (\beta )=\psi (0)$$
it follows now that the new fields must satisfy
$$\stackrel{~}{\psi }(\beta )=e^{ia}\stackrel{~}{\psi }(0)$$
(144)
Consequently, the path integral for the free theory (143) has to be evaluated subject to the periodicity condition of (144).
Although the periodicity condition (144) appears to be complicated, it is well known that the effect can be absorbed by introducing a chemical potential , in the present case, of the form
$$\mu =\frac{ia}{\beta }$$
(145)
With the addition of this chemical potential, the path integral can be evaluated subject to the usual anti-periodicity condition. The effective action can now be easily determined
$`\mathrm{\Gamma }`$ $`=`$ $`i\mathrm{log}\left({\displaystyle \frac{det(i_tm+\frac{ia}{\beta })}{(i_tm)}}\right)^{N_f}`$ (146)
$`=`$ $`iN_f\mathrm{log}\left({\displaystyle \frac{\mathrm{cosh}\frac{\beta }{2}(m\frac{ia}{\beta })}{\mathrm{cosh}\frac{\beta m}{2}}}\right)`$
$`=`$ $`iN_f\mathrm{log}\left(\mathrm{cos}{\displaystyle \frac{a}{2}}+i\mathrm{tanh}{\displaystyle \frac{\beta m}{2}}\mathrm{sin}{\displaystyle \frac{a}{2}}\right)`$
which coincides with the perturbative result of eq. (138).
## 8 Supersymmetry Breaking
One of the reasons for studying finite temperature field theory is to understand questions such as phase transitions in such systems. It is by now well understood that most field theoretic models of spontaneous symmetry breaking display a phase structure much like what one sees in a magnet, namely, above a certain critical temperature, the system is in a symmetric phase while below the critical temperature, the system is in a broken symmetry phase. Thus, temperature has the almost universal effect that if a symmetry is spontaneously broken at low temperature, it is restored at temperatures above a certain critical value. Qualitatively, it can be understood as follows. Temperature, particularly high temperature, provides a lot of thermal energy to a physical system to wash out any structure in the zero temperature potential which may be responsible for symmetry breaking. There is, however, one class of symmetries where temperature has the inverse effect, namely, in a supersymmetric theory, a symmetric phase at low temperature goes to a broken phase at high temperature. (Of course, if supersymmetry is broken at low temperature, it continues to be broken even at high temperature.) We will discuss this phenomenon with a simple quantum mechanical model in this section.
### Supersymmetric Oscillator at $`T=0`$
Let us note that supersymmetry is an ultimate form of symmetry that one can dream of, namely, it transforms bosons into fermions and vice versa \[34-35\]. To introduce supersymmetry, let us consider a simple quantum mechanical model, commonly known as the supersymmetric oscillator . It consists of a bosonic and a fermionic oscillator of the same frequency. Therefore, we can write the Hamiltonian, for the system as
$$H=H_B+H_F=\omega \left(a_B^{}a_B+a_F^{}a_F\right)$$
(147)
where $`a_B`$ and $`a_F`$ describe, respectively, the bosonic and the fermionic annihilation operators.
The immediate thing to note from the structure of the Hamiltonian in eq. (147) is that there is no zero point energy. We will see this shortly as a general feature of supersymmetric theories. Let us also define two fermionic operators of the form
$$Q=a_B^{}a_F,\overline{Q}=a_F^{}a_B$$
(148)
With the usual canonical commutation relations for the bosonic operators (see eq. (66)) and anti-commutation relations for the fermionic operators (see eq. (47)), it is easy to check that
$$[Q,H]=0=[\overline{Q},H]$$
Namely, these fermionic operators are conserved. In fact, together with the Hamiltonian, they satisfy the algebra (it is straightforward to check this)
$`[Q,H]`$ $`=`$ $`0=[\overline{Q},H]`$
$`[Q,Q]_+`$ $`=`$ $`0=[\overline{Q},\overline{Q}]_+`$
$`[Q,\overline{Q}]_+`$ $`=`$ $`{\displaystyle \frac{1}{\omega }}H`$ (149)
Such an algebra, where both commutators and anti-commutators are involved (or alternately, where there is a grading of the multiplication rule of the algebra), is known as a graded Lie algebra and supersymmetric theories are realizations of graded Lie algebras.
As we know from the study of symmetries, conserved quantities generate infinitesimal symmetries of the theory. Since both $`Q`$ and $`\overline{Q}`$ are conserved, it is worth asking what kind of symmetry transformations of the theory they generate. In fact, let us keep in mind that they are fermionic operators and hence the symmetry they will generate cannot be conventional. Explicitly, we can check that
$`[Q,a_B^{}]`$ $`=`$ $`0=[Q,a_F]_+`$
$`[Q,a_B]`$ $`=`$ $`a_F`$
$`[Q,a_F^{}]_+`$ $`=`$ $`a_B^{}`$
$`[\overline{Q},a_B]`$ $`=`$ $`0=[\overline{Q},a_F^{}]_+`$
$`[\overline{Q},a_B^{}]`$ $`=`$ $`a_F^{}`$
$`[\overline{Q},a_F]_+`$ $`=`$ $`a_B`$ (150)
Namely, $`Q`$ and $`\overline{Q}`$ take bosonic operators to fermionic ones and vice versa which is the bench mark of supersymmetry. Thus, our Hamiltonian in eq. (147) is invariant under supersymmetric transformations of the form (150).
There are several things to note from the structure of the supersymmetry algebra in eq. (149). First, the energy eigenvalues of our supersymmetric theory have to be positive semi-definite since the operator on the left hand side of the last relation in (149) is. Furthermore, if the ground state is supersymmetric satisfying
$$Q|0=0=\overline{Q}|0$$
(151)
then, the ground state will have vanishing energy, as we had pointed out earlier as the case for our system. Both these results are, in fact, quite general for any supersymmetric theory. We also note from the structure of the algebra that the spectrum of the Hamiltonian will be doubly degenerate except for the ground state. Namely, if $`|\psi `$ is an eigenstate of the Hamiltonian, then, $`Q|\psi `$ (or, $`\overline{Q}|\psi `$ – only one of them would be nontrivial depending on the form of $`|\psi `$) would be degenerate in energy.
Let us, in fact, examine some of these general results explicitly. The spectrum of the Hamiltonian in eq. (147) is, in fact, quite straightforward. The Hilbert space is a product space containing bosonic and fermionic oscillator states and a general state has the structure
$$|n_B,n_F=|n_B|n_F=\frac{(a_B^{})^{n_B}(a_F^{})^{n_F}}{\sqrt{n_B!}}|0,0$$
(152)
with energy eigenvalues
$$E_{n_B,n_F}=\omega (n_B+n_F),n_F=0,1;n_B=0,1,2,\mathrm{}$$
(153)
where the ground state is expected to satisfy
$$a_B|0=0=a_F|0$$
(154)
We note that an immediate consequence of (154) is that
$$Q|0=0=\overline{Q}|0$$
and, consequently, the ground state is supersymmetric and that the ground state energy is seen from (153) to vanish. All the higher states have positive energy. Furthermore, we note that all the states (except the ground state) of the form $`|n_B,0`$ and $`|n_B1,1`$ are degenerate in energy. Let us also note the effect of $`Q`$ and $`\overline{Q}`$ acting on the states of the Hilbert space, namely,
$`Q|n_B,n_F`$ $`=`$ $`\{\begin{array}{cc}\sqrt{n_B+1}|n_B+1,n_F1& \mathrm{if}n_F0\hfill \\ 0& \mathrm{if}n_F=0\hfill \end{array}`$ (157)
$`\overline{Q}|n_B,n_F`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{\sqrt{n_B}}|n_B1,n_F+1& \mathrm{if}n_B0\mathrm{or},n_F1\hfill \\ 0& \mathrm{if}n_B=0\mathrm{or},n_F=1\hfill \end{array}`$ (160)
### Supersymmetric Oscillator at $`T0`$
Let us next analyze the supersymmetric oscillator at finite temperature in the formalism of thermo field dynamics. As we had noted earlier, this is the ideal setting to discuss questions such as symmetry breaking. Let us note, even before carrying out the calculations, that we expect supersymmetry to be broken at finite temperature. Intuitively, this is quite clear. Namely, supersymmetry takes bosons to fermions and vice versa and, consequently, any boundary condition that distinguishes between the two would lead to a breaking of this symmetry. Temperature, in fact, introduces such a condition, namely, bosons and fermions behave differently at finite temperature (they obey distinctly different statistics). However, what is not clear a priori is whether such a breaking would be explicit or spontaneous.
To study the system at finite temperature within the framework of thermo field dynamics, let us look at the complete system, including the tilde oscillators, described by
$$\widehat{H}=H\stackrel{~}{H}=\omega (a_B^{}a_B+a_F^{}a_F)\omega (\stackrel{~}{a}_B^{}\stackrel{~}{a}_B+\stackrel{~}{a}_F^{}\stackrel{~}{a}_F)$$
(161)
The Hilbert space of the doubled system has the structure
$$|n_B,n_F;\stackrel{~}{n}_B,\stackrel{~}{n}_F=|n_B,n_F|\stackrel{~}{n}_B,\stackrel{~}{n}_F$$
(162)
The thermal vacuum can now be defined (as discussed in section 3). Let us define
$$G(\theta _B,\theta _F)=i\theta _B(\beta )(\stackrel{~}{a}_Ba_Ba_B^{}\stackrel{~}{a}_B^{})i\theta _F(\beta )(\stackrel{~}{a}_Fa_Fa_F^{}\stackrel{~}{a}_F^{})$$
with (see eqs. (58)and (72))
$$\mathrm{tan}\theta _F(\beta )=e^{\beta \omega /2}=\mathrm{tanh}\theta _B(\beta )$$
(163)
Then, the thermal vacuum can be defined as
$$|0,\beta =e^{iG(\theta _B,\theta _F)}|0$$
(164)
This also allows us to calculate the thermal operators in a straightforward manner.
Let us note next that the expectation value of the Hamiltonian in the thermal vacuum is given by
$`E_0(\beta )`$ $`=`$ $`0,\beta |H|0,\beta =0,\beta |\omega (a_B^{}a_B+a_F^{}a_F)|0,\beta `$ (165)
$`=`$ $`\omega (\mathrm{sinh}^2\theta _B(\beta )+\mathrm{sin}^2\theta _F(\beta ))={\displaystyle \frac{2\omega e^{\beta \omega }}{(1e^{2\beta \omega })}}`$
This shows that the energy of the thermal vacuum is nonzero for any finite temperature signaling that supersymmetry is broken. Furthermore, let us note that
$`Q|0,\beta `$ $`=`$ $`a_B^{}a_F|0,\beta ={\displaystyle \frac{e^{\beta \omega /2}}{\sqrt{1e^{2\beta \omega }}}}|n_B(\beta )=1,n_F(\beta )=0;\stackrel{~}{n}_B(\beta )=0,\stackrel{~}{n}_F(\beta )=1`$
$`\overline{Q}|0,\beta `$ $`=`$ $`a_F^{}a_B|0,\beta ={\displaystyle \frac{e^{\beta \omega /2}}{\sqrt{1e^{2\beta \omega }}}}|n_B(\beta )=0,n_F(\beta )=1;\stackrel{~}{n}_B(\beta )=1,\stackrel{~}{n}_F(\beta )=0`$ (166)
This, in fact, shows that supersymmetry breaking is spontaneous at finite temperature and the new states on the right hand side of (166) would correspond to the appropriate quasi particle Goldstino states associated with such a symmetry breaking.
There are various other order parameters for the breaking of supersymmetry and all of them lead to the same conclusion that supersymmetry is spontaneously broken at finite temperature .
## 9 Conclusion
In this article, we have tried to describe some of the interesting features of finite temperature field theories. There are, of course, many more topics that we have not been able to discuss. However, it is our hope that the topics discussed, in this article, would raise the curiosity of the readers to pursue various other questions in this field.
This work was supported in part by the U.S. Dept. of Energy Grant DE-FG 02-91ER40685.
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# References
Introduction
The evaluation of cross sections for hadron collisions typically involves a large number of subprocesses. The reason for this is the quark and gluon content of the initial hadrons and the Cabbibo-Kobayashi-Maskava (CKM) mixing of the down quarks. At Tevatron and LHC energies in many cases 5 flavors, $`u`$, $`d`$, $`c`$, $`s`$ and $`b`$, give a sizable contribution through the corresponding parton densities. Note that subprocesses sometimes contribute only because of the non-diagonality of the CKM matrix. An example is the subprocess $`u\overline{d}s\overline{d}W^+`$ which contributes to the $`W+2jets`$ production. Furthermore, the number of subprocesses increases with the number of quarks in the final state, in particular, because the final state jets produced by various subprocesses are indistinguishable <sup>1</sup><sup>1</sup>1An important exclusion is a $`b`$-quark jet which can be distinguished from light quark jets with a micro-vertex detector.. The process of $`W`$-boson and jet production (an important background to various Standard Model and ”new physics” processes) exemplifies the problem: 180 subprocesses contribute to $`W+2jets`$ and 292 subprocesses to $`W+3jets`$ production if only the quarks of first two generations and only QCD diagrams are taken into account.
The huge multiplication of channels stands as a real computation problem despite the fact that the matrix elements of some subprocesses may have similar analytical structures. Indeed, one should separately organize Monte Carlo integration and/or event generation for each subprocess separately because the convolution with the parton distribution functions (PDF) is flavour dependent. This is a present-day problem for automatic calculations of collision processes (see e.g. and references therein for review of this new computation approach).
In this letter we propose a new method which simplifies the flavour combinatorics and reduces the multiplication of channels due to the mixture of quark states. Implementation of the proposed algorithm in Monte Carlo codes is straightforward. It has been incorporated into the CompHEP program and checked for many examples. In the method the quarks of the first two generations are taken to be massless and do not mix with the third generation which is obviously a good approximation for many processes in the Tevatron and LHC energy range. The method is based on a rotation of down quarks, thus, transporting the mixing matrix elements from vertices of subprocess Feynman diagrams to the parton distribution functions. The complete set of subprocess diagrams are divided into gauge invariant classes, and depending on the topology of the class, new computation rules are formulated. The method drastically simplifies the calculation of matrix elements but leads to a modification of the PDF convolution procedure. We demonstrate the power of the proposed technique with the example of $`W+2jets`$ production at LHC. Here only 21 subprocesses have to be evaluated with the new technique, which should be compared with 180 subprocesses in the standard approach.
The paper is organized as follows. In Section 1 we discuss CKM diagonalization at the level of Feynman amplitudes and introduce basic notations. Here the topologies of the gauge invariant classes of diagrams are defined, for which in Section 2 new PDF convolution rules are formulated. The $`W+2jets`$ production process at LHC is discussed in Section 3. In the Appendix we briefly comment on the implementation of the new rules in the CompHEP program.
1. Feynman amplitudes
Let us consider a parton subprocess with a quark in the initial state. There are two possible topologies of any Feynman diagram with the corresponding quark line: 1) the quark line goes through the diagram from the initial to final state – ”scattering topology”, and 2) the quark line connects both initial state partons – ”annihilation topology”. In Fig.1 and Fig.2a one can see examples of scattering topology, and diagrams of annihilation topology are shown in Fig.2b,c. Here the interior of the circle represents the diagram, and the solid line represents a quark line (quark current). An important point is that according to the theorem proved in , the two sets of diagrams with the scattering and annihilation topologies form gauge invariant classes with respect to the SM gauge group. Thus, one can change the computational rules for each class independently. Note, that a specific parton subprocess could contain only scattering topology diagrams, or only annihilation topology diagrams, or both. For the process $`u\overline{d}d\overline{d}W^+`$ shown in Fig.3 both annihilation and scattering diagrams contribute. For our further consideration it is necessary to determine whether the quark current is charged (CC) or neutral (NC). This is done by counting the quark vertices with $`W`$-bosons. The charged current involves an odd and the neutral current an even number of $`W`$ vertices. Both cases might include an arbitrary number of vertices with gluons, photons and Z-bosons.
Two approximations are used in our method. We neglect: 1)the masses of quarks from the two first generations, and 2) the mixing with the third generation. Thus, the $`3\times 3`$ CKM mixing matrix reduces to the $`2\times 2`$ Cabbibo matrix:
$$V_{CKM}\left(\begin{array}{cc}V& 0\\ 0& 1\end{array}\right),V=\left(\begin{array}{cc}\mathrm{cos}\vartheta _c& \mathrm{sin}\vartheta _c\\ \mathrm{sin}\vartheta _c& \mathrm{cos}\vartheta _c\end{array}\right)$$
where $`\vartheta _c`$ is the Cabbibo angle.
The zero quark mass limit is a good approximation in many applications, in particular, when one evaluates matrix elements of hard subprocesses. Indeed, the energy scale, e.g. the partonic collision energy, for Tevatron-LHC hard subprocesses is $`𝒪(100)`$ GeV and higher. Then, minimal energy-like cuts on jets (momentum transfer, invariant masses of jet pairs etc.) are of order $`𝒪(10)`$ GeV or more. This means that the momenta in the propagators are of the same order and mass effects can be neglected in matrix elements. At the same time, the flavour dependence of parton distributions is sizable and can substantially affect observables.
The second approximation (no mixing with the third generation) also works well in many applications. The $`V_{tb}`$ mixing matrix element is very close to unity and the non-diagonal CKM matrix elements for bottom and top quarks are very small. Of course, these elements are responsible for important physical phenomena, such as B-meson oscillations and rare decays of heavy mesons. In these cases the proposed method can, obviously, not be applied. Another reason why the second approximation is reasonable is that $`t`$ and $`b`$ quarks produce final state objects in a detector which are very different from light quark jets: the $`b`$ quark jet has a secondary vertex, and the $`t`$-quark appears as a heavy narrow resonance.
The starting point of our consideration is the diagonalization of the quark mixing matrix in vertices. This means that rotated down quark states are used in Feynman rules, $`d^{}=d\mathrm{cos}\vartheta _c+s\mathrm{sin}\vartheta _c`$ and $`s^{}=d\mathrm{sin}\vartheta _c+s\mathrm{cos}\vartheta _c`$ rather than $`d`$ and $`s`$ states being eigenstates of the mass matrix. It is worth recalling that all electroweak vertices are diagonal over the isodoublets $`\left(\begin{array}{c}u\\ d^{}\end{array}\right)_L`$ and $`\left(\begin{array}{c}c\\ s^{}\end{array}\right)_L`$ . For example, the charged current (CC) electroweak vertex is diagonal in this rotated basis,
$$W_\mu ^+J_{CC}^\mu =W_\mu ^+\overline{u}_{i}^{}{}_{}{}^{L}\gamma ^\mu V_{ij}d_{j}^{}{}_{}{}^{L}=W_\mu ^+\overline{u}_{i}^{}{}_{}{}^{L}\gamma ^\mu d_{j}^{}{}_{}{}^{L}.$$
Here $`d_i^{}=V_{ij}d_j`$ (where $`i,j=1,2`$ are the generation indices, so $`d=d_1`$ and $`s=d_2`$). As a result, elements of the CKM mixing matrix do not enter in the matrix element if one calculates in terms of these rotated down quarks.
Let us consider an electroweak model with the only one generation of quarks (referred to as the $`EW_{ud}`$ model), and denote generalized up and down quarks by $`q_u`$ and $`q_d`$ respectively. Then, Feynman amplitudes in the EW model with two generations and Cabbibo mixing can be evaluated in the $`EW_{ud}`$ model with a single quark generation. Indeed, in the $`EW_{ud}`$ model, the CC vertex represents the Lagrangian term $`W_\mu ^+\overline{q_u}^L\gamma ^\mu q_{d}^{}{}_{}{}^{L}`$, and its contribution to the amplitude is analytically the same as the contribution of the standard CC vertex (remember that we have neglected quark masses). The only difference is a multiplication of some amplitudes by mixing matrix elements. There are only two generic variants for the case of scattering topology diagrams illustrated in Fig.1: CC in Fig.1a and NC in Fig.1b. One can easily see that the rotation of down quarks results in only one factor $`V_{ij}`$ (or $`V_{ij}^1`$ in conjugated cases) in front of the each CC line, while the mixing matrix elements cancel each other inside the NC line due to the unitarity of the mixing matrix $`_kV_{ik}^1V_{kj}=\delta _{ij}`$. One should stress that the summation over internal flavours in propagators is performed when the amplitude is evaluated. In the Fig.1 only one $`W`$ vertex is shown in the CC case and two $`W`$ vertices in the NC case. However the statements are correct for the general CC and NC cases. Obviously the same conclusions are valid in the case of CC and NC diagrams of the annihilation topology as shown in Fig.2b,c.
An important difference between the annihilation and scattering topologies appears at the level of squared diagrams when the convolution with the parton distribution functions is performed.
2. Squared diagrams and PDF convolution
The next step is to square the matrix element and convolute with the PDF’s. Here we derive four basic rules for the technique under discussion.
Scattering topology. 1st Rule
Let us consider Feynman diagrams with the scattering topology. The second in-state can be a quark, an anti-quark or a gluon. We denote the sum of this type of Feynman diagrams as $`𝒟_{sc}`$. An example is presented in Fig. 2a for the case of a CC upper quark line and an NC lower quark line. After squaring of diagrams of this type and convoluting with PDF’s one can write the contribution to the cross section in the form (for the example in the Fig. 2a):
$`|𝒟_{sc}|^2`$ $``$ $`{\displaystyle \underset{ijkn}{}}{\displaystyle 𝑑x_1𝑑x_2f_{d_i}(x_1)f_{u_j}(x_2)|M_{ijkn}|^2}=`$
$`=`$ $`{\displaystyle 𝑑x_1𝑑x_2\left[\underset{ik}{}f_{d_i}(x_1)V_{ik}V_{ki}^1\right]\left[\underset{jn}{}f_{u_j}(x_2)\delta _{jn}\delta _{nj}\right]||^2},`$
where by $`M`$ and $``$ we denote matrix elements evaluated in the standard EW theory and in the $`EW_{ud}`$ model respectively. Due to the unitarity of the mixing matrix one has $`V_{ik}V_{ki}^1=\delta _{ii}`$, and the first rule can be written as
$$|𝒟_{sc}|^2𝑑x_1𝑑x_2[f_d(x_1)+f_s(x_1)][f_{\overline{u}}(x_2)+f_{\overline{c}}(x_2)]||^2.$$
(1)
This means that one can evaluate squared diagrams of the scattering topology class with only one quark generation, but should convolute the corresponding gauge invariant squared matrix element with modified structure function(s) – a sum of down (or up) PDF’s.
Annihilation CC topology. 2nd Rule
In the case of the annihilation topology the CC and NC cases lead to different rules for convolution with the PDF. We start from the CC annihilation case which is generically shown in Fig. 2b. Recall that this variant occurs only if the quark line has an odd number of $`W`$ vertices. We denote the sum of this class of Feynman diagrams as $`𝒟_a^{CC}`$. When one squares this sum of diagrams and convolutes with the PDF, one gets:
$`|𝒟_a^{CC}|^2`$ $``$ $`{\displaystyle \underset{ij}{}}{\displaystyle 𝑑x_1𝑑x_2f_{d_i}(x_1)f_{\overline{u}_j}(x_2)|M_{ij}|^2}`$
$`=`$ $`{\displaystyle 𝑑x_1𝑑x_2\left[\underset{ij}{}f_{d_i}(x_1)f_{\overline{u}_j}(x_2)V_{ij}V_{ji}^1\right]||^2}.`$
In contrast to the scattering case here one can not use the unitarity condition to cancel two elements of the mixing matrix because the summations over the indices, $`i`$ and $`j`$, also include the structure functions $`f_{d_i}`$ and $`f_{\overline{u}_j}`$. Thus, the second rule can be written in the following form (see the example in Fig. 2b):
$`|𝒟_a^{CC}|^2{\displaystyle 𝑑x_1𝑑x_2}`$ $`[`$ $`f_d(x_1)f_{\overline{u}}(x_2)\mathrm{cos}^2\vartheta _c+f_s(x_1)f_{\overline{c}}(x_2)\mathrm{cos}^2\vartheta _c+`$
$`+`$ $`f_d(x_1)f_{\overline{c}}(x_2)\mathrm{sin}^2\vartheta _c+f_s(x_1)f_{\overline{u}}(x_2)\mathrm{sin}^2\vartheta _c]||^2.`$
where we have explicitly expressed the mixing matrix elements in terms of the Cabbibo mixing angle.
One can see that the annihilation-type contribution to the squared matrix element is convoluted with non-factorizable products of PDF’s.
Annihilation NC topology. 3d Rule
We denote the sum of all diagrams with the annihilation NC quark line as $`𝒟_a^{NC}`$. The generic example is shown in Fig. 2c. In this case the quark line has an even number of $`W`$-boson vertices. In the same way as above one can derive the following formula:
$`|𝒟_a^{NC}|^2`$ $``$ $`{\displaystyle \underset{ij}{}}{\displaystyle 𝑑x_1𝑑x_2f_{d_i}(x_1)f_{\overline{d}_j}(x_2)|M_{ij}|^2}=`$
$`=`$ $`{\displaystyle 𝑑x_1𝑑x_2\left[\underset{ij}{}f_{d_i}(x_1)f_{\overline{d}_j}(x_2)\delta _{ij}\delta _{ji}\right]||^2}.`$
Therefore, the third rule can be written in the form (if the in-quarks are down):
$$|𝒟_a^{NC}|^2𝑑x_1𝑑x_2[f_d(x_1)f_{\overline{d}}(x_2)+f_s(x_1)f_{\overline{s}}(x_2)]||^2.$$
(3)
with the obvious generalization for up quarks.
Interference of $`𝒟_{sc}`$ and $`𝒟_a`$ topologies
In the general case diagrams of both, scattering and annihilation, topologies could contribute. The interference between gauge invariant diagram classes with these topologies is also gauge invariant and, therefore, can be independently convoluted with the PDF. As in the above derivations the PDF convolution for the interference of $`𝒟_{sc}`$ diagrams with $`𝒟_a^{CC}`$ diagrams is given by the same formula as for $`|𝒟_a^{CC}|^2`$ (2nd Rule), and for the interference of $`𝒟_{sc}`$ with $`𝒟_a^{NC}`$ by the corresponding formula for $`|𝒟_a^{NC}|^2`$ (3rd Rule).
The final state quark-antiquark line. 4th Rule
Finally let us consider Feynman diagrams where a quark line connects two out-states, as in Fig. 2d. Here the summation over generation indices does not involve parton distribution functions and can, therefore, be performed explicitly. The result is that the contribution of the corresponding squared diagrams, evaluated in the $`EW_{ud}`$ model, should be multiplied by two for both the CC and NC cases. Indeed, according to Fig. 2d each of these summations give in the squared diagram the factors: $`_{ij}\delta _{ij}\delta _{ji}=2`$ in the NC and $`_{ij}V_{ij}V_{ji}^1=_{ij}\delta _{ij}=2`$ in the CC cases.
Note that the 4th Rule is valid not only in cases where the quark loop in squared diagrams connects out-state(s) and never passes through in-state(s), but also for each quark loop in next-to-leading corrections.
Of course, the 4th Rule is valid only under the assumption that the fragmentation of the four light quarks and antiquarks leads to indistinguishable jets. If one includes nontrivial fragmentation functions, e.g. for a $`c`$-quark, all the above rules have to be modified. We do not present here the corresponding formulas which however could be easily derived.
3. Test: $`W+2jets`$ production at LHC
In this section we illustrate the proposed technique with the example of $`W^++2jets`$ production at LHC. Here a total of 180 subprocesses contribute in the standard technique when all four light quarks contribute separately and the CKM matrix is present in $`W`$-boson vertices. With the new technique only 21 subprocesses need to be evaluated.
We shall not discuss this example in full detail but present the results of a numerical test for the $`u\overline{d}d\overline{d}W^+`$ subprocess with permutations of quarks within pairs ($`u`$,$`c`$) and ($`d`$,$`s`$). Note that in this example three rules are used: 1st, 2nd and 4th. In the standard technique 12 subprocesses are involved and the corresponding contributions to the cross section are collected in Tab. 1 for two values of the kinematical cut on the transverse momenta of the final partons and the $`W`$ boson, $`p_T^{jet}>p_T^0`$ and $`p_T^W>p_T^0`$ with $`p_T^0=20`$ or 200 GeV. The cross sections were calculated in the Standard Model with the averaged values for CKM matrix elements and quark masses taken from the Particle Data Group . For calculations the CompHEP program has been used, and the accumulated MC error in all cases was less than 0.6%. We neglect the contributions of subleading diagrams with electroweak boson propagators, calculating the cross section to leading $`\alpha \alpha _s^2`$ order.
Using the new method one should evaluates only one subprocess in the one quark doublet $`EW_{ud}`$ model, $`q_u\overline{q_d}q_d\overline{q_d}W^+`$. The corresponding Feynman diagrams are shown in Fig. 3. This subprocess is of the mixed type where two gauge invariant classes of diagrams: the annihilation CC topology (Fig. 3a) and the scattering topology (Fig. 3b), contribute. The 1st Rule is used for the squared scattering topology contribution. The 2nd Rule is used for the squared annihilation CC topology contribution with a multiplication by the factor 2 according to the 4th Rule. Finally, the 2nd Rule is used to evaluate the interference between the two gauge invariant classes of Feynman diagrams. We have calculated all three contributions to the cross section of the subprocess $`q_u\overline{q_d}q_d\overline{q_d}W^+`$ using CompHEP code in which the $`EW_{ud}`$ model and Rules 1-3 have been implemented (see the Appendix for details). The results (see Tab. 2) for the total rate, $`\sigma (p_T>20\text{GeV})=112.95`$ pb and $`\sigma (p_T>200\text{GeV})=0.30213`$ pb, are in an agreement with the ”standard” calculations of Tab. 1 within the statistical error of less than 0.6%.
Conclusions
We have shown that hard collision processes at hadron colliders can be evaluated in an economical way, greatly reducing the number of contributing subprocesses. The proposed computational technique can be applied only if the quark masses of the first two generations and mixing with the 3rd generation can be neglected. These assumptions are valid for most applications at the Tevatron and LHC. In the proposed technique the Standard Model with a single generation of up and down quarks is used and a squared matrix element is evaluated without involving elements of the mixing matrix. The resulting matrix squared matrix element is convoluted with modified parton distribution functions according to formulas given above as Rules 1-4. Each of these Rules corresponds to a gauge invariant class of squared diagrams.
In Rule 1 the squared matrix element is convoluted with $`f_u(x)+f_c(x)`$ or $`f_d(x)+f_s(x)`$ for an in-state of up-type or down-type quarks respectively However, in the cases of Rules 2 and 3 the squared matrix element is convoluted over Bjorken variables $`x_1`$ and $`x_2`$ with a non-factorizable function. When applying the 2nd Rule the function
$$[f_d(x_1)f_{\overline{u}}(x_2)+f_s(x_1)f_{\overline{c}}(x_2)]\mathrm{cos}^2\vartheta _c+[f_d(x_1)f_{\overline{c}}(x_2)+f_s(x_1)f_{\overline{u}}(x_2)]\mathrm{sin}^2\vartheta _c,$$
or a similar function where the substitution $`d\overline{d}`$ and $`\overline{u}u`$ are made, is used. In the case of the 3rd Rule the convolution is with either
$$f_u(x_1)f_{\overline{u}}(x_2)+f_c(x_1)f_{\overline{c}}(x_2)or,$$
$$f_d(x_1)f_{\overline{d}}(x_2)+f_s(x_1)f_{\overline{s}}(x_2),$$
depending on which quark-antiquark pair occurs in the in-state.
The authors are indebted to A.E. Pukhov, Th. Ohl and B. Straub for useful discussions. This work was partially supported by the CERN-INTAS 99-377, RFBR-DFG 99-02-04011 and RFBR 00-01-00704 grants, by the Russian Ministry of Science and Technology, by St.Petersburg Grant Center and by the program ”Universities of Russia” (grant 990588).
Appendix
The Rules derived in this letter have been implemented in the CompHEP code v.33. New model, referred to above as $`EW_{ud}`$, was created on the base of the SM where only one quark generation, up and down quarks denoted as $`q_u`$ and $`q_d`$ was kept without any CKM matrix elements. The masses of these generalized quarks were set to zero. Then, the option for numerical convolution of squared diagrams with parton distribution functions was modified in accordance with Rules 1-3. The code of this version of CompHEP is available from the Web .
In the most general case the user has to subdivide the whole set of squared diagrams into two parts: 1) $`|𝒟_{sc}|^2`$, and 2) $`|𝒟_a|^2`$ plus the interference diagrams $`2Re(𝒟_{sc}𝒟_a^{})`$. Each of these parts should be calculated separately. In particular, for each part the user has to set the variable ”PDFfactor” in the menu option ”User menu” as follows: $`PDFfactor=1`$ for $`|𝒟_{sc}|^2`$ (referred to in CompHEP as the ”t-channel” case), and $`PDFfactor=0`$ for $`|𝒟_a|^2+2Re(𝒟_{sc}𝒟_a^{})`$ (referred to in CompHEP as the ”s-channel” case). The program automatically recognizes which Rule, 2nd or 3rd, should be used in the latter case.
Rule 4 (multiplication by a factor 2) has to be applied by hand if the corresponding quark line is presented in a squared diagram.
One should note that this is not a completely automatic realization of Rules 1-4. This is planned for future CompHEP development.
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# Post-Newtonian Gravitational Radiation
## I Introduction
### A On approximation methods in general relativity
Let us declare that the most important devoir of any physical theory is to draw firm predictions for the outcome of laboratory experiments and astronomical observations. Unfortunately, the devoir is quite difficult to fulfill in the case of general relativity, essentially because of the complexity of the Einstein field equations, to which only few exact solutions are known. For instance, it is impossible to settle the exact prediction of this theory when there are no symmetry in the problem (as is the case in the problem of the gravitational dynamics of separated bodies). Therefore, one is often obliged, in general relativity, to resort to approximation methods.
It is beyond question that approximation methods do work in general relativity. Some of the great successes of this theory were in fact obtained using approximation methods. We have particularly in mind the test by Taylor and collaborators regarding the orbital decay of the binary pulsar PSR 1913+16, which is in agreement to within 0.35% with the general-relativistic post-Newtonian prediction. However, a generic problem with approximation methods (especially in general relativity) is that it is non trivial to define a clear framework within which the approximation method is mathematically well-defined, and such that the results of successive approximations could be considered as theorems following some precise (physical and/or technical) assumptions. Even more difficult is the problem of the relation between the approximation method and the exact theory. In this context one can ask: What is the mathematical nature of the approximation series (convergent, asymptotic, $`\mathrm{}`$)? What its “reliability” is (i.e., does the approximation series come from the Taylor expansion of a family of exact solutions)? Does the approximate solution satisfy some “exact” boundary conditions (for instance the no-incoming radiation condition)?
Since the problem of theoretical prediction in general relativity is complex, let us distinguish several approaches (and ways of thinking) to it, and illustrate them with the example of the prediction for the binary pulsar. First we may consider what could be called the “physical” approach, in which one analyses the relative importance of each physical phenomena at work by using crude numerical estimates, and where one uses only the lowest-order approximation, relating if necessary the local physical quantities to observables by means of balance equations (perhaps not well defined in terms of basic theoretical concepts). The physical approach to the problem of the binary pulsar is well illustrated by Thorne in his beautiful Les Houches review (see also the round table discussion moderated by Ashtekar ): one derives the loss of energy by gravitational radiation from the (Newtonian) quadrupole formula applied formally to point-particles, assumed to be test-masses though they are really self-gravitating, and one argues “physically” that the effect comes from the variation of the Newtonian binding energy in the center-of-mass frame – indeed, on physical grounds, what else could this be (since we expect the rest masses won’t vary)? The physical approach yields the correct result for the rate of decrease of the period of the binary pulsar. Of course, thinking physically is extremely useful, and indispensable in a preliminary stage, but certainly it should be completed by a solid study of the connection to the mathematical structure of the theory. Such a study would a posteriori demote the physical approach to the status of “heuristic” approach. On the other hand, the physical approach may fall short in some situations requiring a sophisticated mathematical modelling (like in the problem of the dynamics of singularities), where one is often obliged to follow one’s mathematical rather than physical insight.
A second approach, that we shall qualify as “rigorous”, has been advocated mainly by Jürgen Ehlers (see, e.g., ). It consists of looking for a high level of mathematical rigor, within the exact theory if possible, and otherwise using an approximation scheme that we shall be able to relate to the exact theory. This does not mean that we will be so much wrapped up by mathematical rigor as to forget about physics. Simply, in the rigorous approach, the prediction for the outcome of an experiment should follow mathematically from first theoretical principles. Clearly this approach is the one we should ideally adhere to. As an example, within the rigorous approach, one was not permitted, by the end of the seventies, to apply the standard quadrupole formula to the binary pulsar. Indeed, as pointed out by Ehlers et al , it was not clear that gravitational radiation reaction on a self-gravitating system implies the standard quadrupole formula for the energy flux, notably because computing the radiation reaction demands a priori three non-linear iterations of the field equations , which were not fully available at that time. Ehlers and collaborators remarked also that the exact results concerning the structure of the field at infinity (notably the asymptotic shear of null geodesics whose variation determines the flux of radiation) were not connected to the actual dynamics of the binary.
Maybe the most notable result of the rigorous approach concerns the relation between the exact theory and the approximation methods. In the case of the post-Newtonian approximation (limit $`c\mathrm{}`$), Jürgen Ehlers has provided with his frame theory a conceptual framework in which the post-Newtonian approximation can be clearly formulated (among other purposes). This theory unifies the theories of Newton and Einstein into a single generally covariant theory, with a parameter $`1/c`$ taking the value zero in the case of Newton and being the inverse of the speed of light in the case of Einstein. Within the frame theory not only does one understand the limit relation of Einstein’s theory to Newton’s, but one explains why it is legitimate when describing the predictions of general relativity to use the common-sense language of Newton (for instance thinking that the trajectories of particles in an appropriately defined coordinate system take place in some Euclidean space, and viewing the coordinate velocities as being defined with respect to absolute time). It was shown by Lottermoser that the constraint equations of the (Hamiltonian formulation of the) Ehlers frame theory admit solutions with a well-defined post-Newtonian limit. Further in the spirit of the rigorous approach, we quote the work of Rendall on the definition of the post-Newtonian approximation, and the link to the post-Newtonian equations used in practical computations. (See also for an attempt at showing, using restrictive assumptions, that the post-Newtonian series is asymptotic.)
The important remarks of Jürgen Ehlers et al on the applicability of the quadrupole formula to the binary pulsar stimulated research to settle down this question with (al least) acceptable mathematical rigor. The question was finally answered positively by Damour and collaborators , who obtained in algebraically closed form the general-relativistic equations of motion of two compact objects, up to the requisite 5/2 post-Newtonian order (2.5PN order or $`1/c^5`$) where the gravitational radiation reaction force appears. This extended to 2.5PN order the work at 1PN of Lorentz and Droste , and Einstein, Infeld and Hoffmann . The net result is that the dynamics of the binary pulsar as predicted by (post-Newtonian) general relativity is in full agreement both with the prediction of the quadrupole formula, as derived earlier within the “physical” approach, and with the observations by Taylor et al (see for discussion).
Motivated by the success of the theoretical prediction in the case of the binary pulsar , we shall try to follow in this article the spirit of the “rigorous” approach of Jürgen Ehlers, notably in the way it emphasizes the mathematical proof, but we shall also differ from it by a systematic use of approximation methods. This slightly different approach recognizes from the start that in certain difficult problems, it is impossible to derive a physical result all the way through the exact theory without any gap, so that one must proceed with approximations. But, in this approach, one implements a mathematically well-defined framework for the approximation method, and within this framework one proves theorems that (ideally) guarantee the correctness of the theoretical prediction to be compared with experiments. Because the comparison with experiments is the only thing which matters in fine for a pragmatist, we qualify this third approach as “pragmatic”.
In this article we describe the pragmatic approach to the problem of gravitational radiation emitted by a general isolated source, based on the rigorous post-Minkowskian iteration of the field outside the source , and on the general connection of the exterior field to the field inside a slowly-moving source . Note that for this particular problem the pragmatic approach is akin to the rigorous one in that it permits to establish some results on the connection between approximate and exact methods. For instance it was proved by Damour and Schmidt (see also ) that the post-Minkowskian algorithm generates an asymptotic approximation to exact solutions, and it was shown that the solution satisfies to any order in the post-Minkowskian expansion a rigorous definition of asymptotic flatness at future null infinity. However it remains a challenge to analyse in the manner of the rigorous approach the relation to exact theory of the whole formalism of .
By combining the latter post-Minkowskian approximation and a post-Newtonian expansion inside the system, it was proved (within this framework of approximations) that the quadrupole formula for slowly-moving, weakly-stressed and self-gravitating systems is correct, even including post-Newtonian corrections ; and idem for the radiation reaction forces acting locally inside the system, and for the associated balance equations . These results answered positively Ehlers’ remarks in the case of slowly-moving extended (fluid) systems. However we are also interested in this article to the application to binary systems of compact objects modelled by point-masses. Indeed the latter sources of radiation are likely to be detected by future gravitational-wave experiments, and thus concern the pragmatist. We shall see how one can address the problem in this case. (When specialized to point-mass binaries, the results on radiation reaction are in agreement with separate work of Iyer and Will .) For other articles on the problem of gravitational radiation from general and binary point-mass sources, see .
### B Field equations and the no-incoming radiation condition
The problem is to find the solutions, in the form of analytic approximations, of the Einstein field equations in $`IR^4`$,
$$R^{\mu \nu }\frac{1}{2}g^{\mu \nu }R=\frac{8\pi G}{c^4}T^{\mu \nu },$$
(1)
and thus also of their consequence, the equations of motion of the matter source, $`_\nu T^{\mu \nu }=0`$. Throughout this work we assume the existence and unicity of a global harmonic (or de Donder) coordinate system. This means that we can choose the gauge condition
$$_\nu h^{\mu \nu }=0;h^{\mu \nu }\sqrt{g}g^{\mu \nu }\eta ^{\mu \nu },$$
(2)
where $`g`$ and $`g^{\mu \nu }`$ denote the determinant and inverse of the covariant metric $`g_{\mu \nu }`$, and where $`\eta ^{\mu \nu }`$ is an auxiliary flat metric \[i.e. $`\eta ^{\mu \nu }=\mathrm{diag}(1,1,1,1)=\eta _{\mu \nu }`$\]. The Einstein field equations (1) can then be replaced by the relaxed equations
$$h^{\mu \nu }=\frac{16\pi G}{c^4}\tau ^{\mu \nu },$$
(3)
where the box operator is the flat d’Alembertian, $`_\eta =\eta ^{\mu \nu }_\mu _\nu `$, and where the source term is the sum of a matter part and a gravitational part,
$$\tau ^{\mu \nu }|g|T^{\mu \nu }+\frac{c^4}{16\pi G}\mathrm{\Lambda }^{\mu \nu }.$$
(4)
In harmonic coordinates the field equations take the form of simple wave equations, but whose source term is actually a complicated functional of the gravitational field $`h^{\mu \nu }`$; notably the gravitational part depends on $`h^{\mu \nu }`$ and its first and second space-time derivatives:
$`\mathrm{\Lambda }^{\mu \nu }=`$ $``$ $`h^{\rho \sigma }_{\rho \sigma }^2h^{\mu \nu }+_\rho h^{\mu \sigma }_\sigma h^{\nu \rho }+{\displaystyle \frac{1}{2}}g^{\mu \nu }g_{\rho \sigma }_\lambda h^{\rho \tau }_\tau h^{\sigma \lambda }`$ (5)
$``$ $`g^{\mu \rho }g_{\sigma \tau }_\lambda h^{\nu \tau }_\rho h^{\sigma \lambda }g^{\nu \rho }g_{\sigma \tau }_\lambda h^{\mu \tau }_\rho h^{\sigma \lambda }+g_{\rho \sigma }g^{\lambda \tau }_\lambda h^{\mu \rho }_\tau h^{\nu \sigma }`$ (6)
$`+`$ $`{\displaystyle \frac{1}{8}}(2g^{\mu \rho }g^{\nu \sigma }g^{\mu \nu }g^{\rho \sigma })(2g_{\lambda \tau }g_{ϵ\pi }g_{\tau ϵ}g_{\lambda \pi })_\rho h^{\lambda \pi }_\sigma h^{\tau ϵ}.`$ (7)
The point is that $`\mathrm{\Lambda }^{\mu \nu }`$ is at least quadratic in $`h`$, so the relaxed field equations (3) are very naturally amenable to a perturbative non-linear expansion. As an immediate consequence of the gauge condition (2), the right side of the relaxed equations is conserved in the usual sense, and this is equivalent to the equations of motion of matter:
$$_\nu \tau ^{\mu \nu }=0_\nu T^{\mu \nu }=0.$$
(8)
We refer to $`\tau ^{\mu \nu }`$ as the total stress-energy pseudo-tensor of the matter and gravitational fields in harmonic coordinates. Since the harmonic coordinate condition is Lorentz covariant, $`\tau ^{\mu \nu }`$ is a tensor with respect to Lorentz transformations (but of course not with respect to general diffeomorphisms).
In order to select the physically sensible solution of the field equations in the case of a bounded system, one must choose some boundary conditions at infinity, i.e. the famous no-incoming radiation condition, which ensures that the system is truly isolated (no radiating sources located at infinity). In principle the no-incoming radiation condition is to be formulated at past null infinity $`𝒥^{}`$. Here, we shall simplify the formulation by taking advantage of the presence of the Minkowski background $`\eta _{\mu \nu }`$ to define the no-incoming radiation condition with respect to the Minkowskian past null infinity $`𝒥_\mathrm{M}^{}`$. Of course, this does not make sense in the exact theory where only exists the metric $`g_{\mu \nu }`$ and where the metric $`\eta _{\mu \nu }`$ is fictituous, but within approximate (post-Minkowskian) methods it is legitimate to view the gravitational field as propagating on the flat background $`\eta _{\mu \nu }`$, since $`\eta _{\mu \nu }`$ does exist at any finite order of approximation.
We formulate the no-incoming radiation condition in such a way that it suppresses any homogeneous, regular in $`IR^4`$, solution of the d’Alembertian equation $`h=0`$. We have at our disposal the Kirchhoff formula which expresses $`h(𝐱^{},t^{})`$ in terms of values of $`h(𝐱,t)`$ and its derivatives on a sphere centered on $`𝐱^{}`$ with radius $`\rho |𝐱^{}𝐱|`$ and at retarded time $`tt^{}\rho /c`$:
$$h(𝐱^{},t^{})=\frac{d\mathrm{\Omega }}{4\pi }\left[\frac{}{\rho }(\rho h)+\frac{1}{c}\frac{}{t}(\rho h)\right](𝐱,t)$$
(9)
where $`d\mathrm{\Omega }`$ is the solid angle spanned by the unit direction $`(𝐱𝐱^{})/\rho `$. From the Kirchhoff formula we obtain the no-incoming radiation condition as a limit at $`𝒥_\mathrm{M}^{}`$, that is $`r+\mathrm{}`$ with $`t+r/c=`$const (where $`r=|𝐱|`$). In fact we obtain two conditions: the main one,
$$\underset{\genfrac{}{}{0pt}{}{r+\mathrm{}}{t+r/c=\mathrm{const}}}{lim}\left[\frac{}{r}(rh^{\mu \nu })+\frac{1}{c}\frac{}{t}(rh^{\mu \nu })\right](𝐱,t)=0,$$
(10)
and an auxiliary condition, that $`r_\lambda h^{\mu \nu }`$ should be bounded at $`𝒥_\mathrm{M}^{}`$, coming from the fact that $`\rho `$ in the Kirchhoff formula (7) differs from $`r`$ \[we have $`\rho =r𝐱^{}.𝐧+O(1/r)`$ where $`𝐧=𝐱/r`$\].
In fact, we adopt in this article a much more restrictive condition of no-incoming radiation, namely that the field is stationary before some finite instant $`𝒯`$ in the past:
$$t𝒯\frac{}{t}[h^{\mu \nu }(𝐱,t)]=0.$$
(11)
In addition we assume that before $`𝒯`$ the field $`h^{\mu \nu }(𝐱)`$ is of order $`O(1/r)`$ when $`r+\mathrm{}`$. These restrictive conditions are imposed for technical reasons following , since they allow constructing rigorously (and proving theorems about) the metric outside some time-like world tube $`r|𝐱|>`$. We shall assume that the region $`r>`$ represents the exterior of an actual compact-support system with constant radius $`d<`$ \[i.e. $`d`$ is the maximal radius of the adherence of the compact support of $`T^{\mu \nu }(𝐱,t)`$, for any time $`t`$\].
Now if $`h^{\mu \nu }`$ satisfies for instance (9), so does the pseudo-tensor $`\tau ^{\mu \nu }`$ built on it, and then it is clear that the retarded integral of $`\tau ^{\mu \nu }`$ satisfies itself the same condition. Therefore one infers that the unique solution of the Einstein equation (3) satisfying the condition (9) is
$$h^{\mu \nu }=\frac{16\pi G}{c^4}_R^1\tau ^{\mu \nu },$$
(12)
where the retarded integral takes the standard form
$$(_R^1\tau )(𝐱,t)\frac{1}{4\pi }\frac{d^3𝐱^{}}{|𝐱𝐱^{}|}\tau (𝐱^{},t|𝐱𝐱^{}|/c).$$
(13)
Notice that since $`\tau ^{\mu \nu }`$ depends on $`h`$ and its derivatives, the equation (10) is to be viewed rather as an integro-differential equation equivalent to the Einstein equation (3) with no-incoming radiation.
### C Method and general physical picture
We want to describe an isolated system, for instance a “two-body system”, in Einstein’s theory. We expect (though this is not proved) that initial data sets $`g_{\mu \nu }`$, $`_tg_{\mu \nu }`$, $`\rho `$, v satisfying the constraint equations on the space-like hypersurface $`t=t_0`$ exist, and that this determines a unique solution of the field equations for any time $`t`$, which approaches in the case of two bodies a “scattering state” when $`t\mathrm{}`$, in which the bodies move on unbound (hyperbolic-like) orbits. We assume that the space-times generated by such data admit a past null infinity $`𝒥^{}`$ (or, if one uses approximate methods, $`𝒥_\mathrm{M}^{}`$) with no incoming radiation. (Note that in a situation with initial scattering the field might not satisfy the rigorous definitions of asymptotic flatness at $`𝒥^{}`$; see .) The point to make is that in this class of space-times there is no degree of freedom for the gravitational field (we could consider other situations where the motion is influenced by incoming radiation).
Both our technical assumptions of compact support for the matter source (with constant radius $`d`$) and stationarity before the time $`𝒯`$ contradict our expectation that a two-body system follows an unbound orbit in the remote past. We do not solve this conflict but argue as follows: (i) these technical assumptions permit to derive rigorously some results, for instance the expression \[given by (52) with (56) below\] of the far-field of an isolated past-stationary system; (ii) it is clear that these results do not depend on the constant radius $`d`$, and furthermore we check that they admit in the “scattering” situation a well-defined limit when $`𝒯\mathrm{}`$ ; (iii) this makes us confident that the results are actually valid for a more realistic class of physical systems which become unbound in the past and are never stationary (and, even, one can give a posteriori conditions under which the limit $`𝒯\mathrm{}`$ exists for a general system at some order of approximation).
Suppose that the system is “slowly-moving” \[in the sense of (12) below\], so that we can compute the field inside its compact support by means of a post-Newtonian method, say $`h_{\mathrm{in}}^{\mu \nu }\overline{h}^{\mu \nu }`$ where the overbar refers to the formal post-Newtonian series. The post-Newtonian iteration (say, for hydrodynamics) is not yet defined to all orders in $`1/c`$, but many terms are known: see the works of Lorentz and Droste , Einstein, Infeld and Hoffmann , Fock , Chandrasekhar and collaborators , Ehlers and followers , and many other authors .
On the other hand, outside the isolated system, the field is weak everywhere and it satisfies the vacuum equations. Therefore, the equations can be solved conjointly by means of a weak-field or post-Minkowskian expansion ($`G0`$), and, for each coefficient of $`G^n`$ in the latter expansion, by means of a multipole expansion (valid because we are outside). The general Multipolar-post-Minkowskian (MPM) metric was constructed in as a functional of two sets of “multipole moments” $`M_L(t)`$ and $`S_L(t)`$ which were left arbitrary at this stage (i.e. not connected to the source). The idea of combining the post-Minkowskian and multipole expansions comes from the works of Bonnor and Thorne . We denote by $`h_{\mathrm{ext}}^{\mu \nu }(h^{\mu \nu })`$ the exterior solution, where $``$ stands for the multipole expansion (as it will turn out, the post-Minkowskian expansion appears in this formalism to be somewhat less fundamental than the multipole expansion).
The key assumption is that the two expansions $`h_{\mathrm{in}}^{\mu \nu }=\overline{h}^{\mu \nu }`$ and $`h_{\mathrm{ext}}^{\mu \nu }=(h^{\mu \nu })`$ should match in a region of common validity for both the post-Newtonian and multipole expansions. Here is where our physical restriction to slow motion plays a crucial role, because such an overlap region exists (this is the so-called exterior near-zone) if and only if the system is slowly-moving. The matching is a variant of the well-known method of matching of asymptotic expansions, very useful in gravitational radiation theory . It consists of decomposing the inner solution into multipole moments (valid in the outside), re-expanding the exterior solution in the near zone ($`r/c0`$), and equating term by term the two resulting expansion series. From the requirement of matching we obtain in , and review in Sections 2 and 3 below, the general formula for the multipole expansion $`(h^{\mu \nu })`$ in terms of the “source” multipole moments (notably a mass-type moment $`I_L`$ and a current-type $`J_L`$), given as functionals of the post-Newtonian expansion of the pseudo-tensor, i.e. $`\overline{\tau }^{\mu \nu }`$. \[The previous moments $`M_L`$ and $`S_L`$ (referred below to as “canonical”) are deduced from the source moments after a suitable coordinate transformation.\] In addition the matching equation determines the radiation reaction contributions in the inner post-Newtonian metric .
To obtain the source multipole moments in terms of basic source parameters (mass density, pressure), it remains to replace $`\overline{\tau }^{\mu \nu }`$ by the result of an explicit post-Newtonian iteration of the inner field. This was done to 1PN order in , then to 2PN order in , and the general formulas obtained in permit recovering these results. See Section 6. On the other hand, if one needs the equations of motion of the source, simply one inserts the post-Newtonian metric into the conservation law $`_\nu \overline{\tau }^{\mu \nu }=0`$. (Note that we are speaking of the equations of motion, which take for instance the form of Euler-type equations with many relativistic corrections, but not of the solutions of these equations, which are typically impossible to obtain analytically.)
From the harmonic coordinates, one can perform to all post-Minkowskian orders a coordinate transformation to some radiative coordinates such that the metric admits a far-field expansion in powers of the inverse of the distance $`R`$ (without the powers of $`\mathrm{ln}R`$ which plague the harmonic coordinates). Considering the leading order $`1/R`$ one compares the exterior metric, which is parametrized by the source moments (connected to the source via the matching equation), to the metric defined with “radiative” multipole moments, say $`U_L`$ and $`V_L`$. This gives $`U_L`$ and $`V_L`$ in terms of the source moments, notably $`I_L`$ and $`J_L`$, and a fortiori of the source parameters. This solves, within approximate methods, the problem of the relation between the far field and the source. The radiative moments have been obtained with increasing precision reaching now 3PN , as reviewed in Section 5.
The previous scheme is developed for a general description of matter, however restricted to be smooth (we have in mind a general “hydrodynamical” $`T^{\mu \nu }`$). Thus the scheme a priori excludes the presence of singularities (no “point-particles” or black holes), but this is a serious limitation regarding the application to compact objects like neutron stars, which can adequately be approximated by point-masses when studying their dynamics. Fortunately, the formalism is applicable to a singular $`T^{\mu \nu }`$ involving Dirac measures, at the price of a further ansatz, that the infinite self-field of point-masses can be regularized in a certain way. By implementing consistently the regularization we obtain the multipole moments and the radiation field of a system of two point-masses at 2.5PN order , as well as their equations of motion at the same order in the form of ordinary differential equations (the result agrees with previous works ); see Section 7.
## II Multipole decomposition
In this section we construct the multipole expansion $`(h^{\mu \nu })h_{\mathrm{ext}}^{\mu \nu }`$ of the gravitational field outside an isolated system, supposed to be at once self-gravitating and slowly-moving. By slowly-moving we mean that the typical current and stress densities are small with respect to the energy density, in the sense that
$$\mathrm{max}\{\left|\frac{T^{0i}}{T^{00}}\right|,\left|\frac{T^{ij}}{T^{00}}\right|^{1/2}\}=O\left(\frac{1}{c}\right),$$
(14)
where $`1/c`$ denotes (slightly abusively) the small post-Newtonian parameter. The point about (12) is that the ratio between the size of the source $`d`$ and a typical wavelength of the gravitational radiation is of order $`d/\lambda =O(1/c)`$. Thus the domain of validity of the post-Newtonian expansion covers the source: it is given by $`r<b`$ where the radius $`b`$ can be chosen so that $`d<b=O(\lambda /c)`$.
### A The matching equation
The construction of the multipole expansion is based on several technical assumptions, the crucial one being that of the consistency of the asymptotic matching between the exterior and interior fields of the isolated system. In some cases the assumptions can be proved from the properties of the exterior field $`h_{\mathrm{ext}}^{\mu \nu }`$ as obtained in by means of a post-Minkowskian algorithm. However, since our assumptions are free of any reference to the post-Minkowskian expansion, we prefer to state them more generally, without invoking the existence of such an approximation (refer to for the full detailed assumptions). In many cases the assumptions have been explicitly verified at some low post-Newtonian orders .
The field $`h`$ (skipping space-time indices), solution in $`IR^4`$ of the relaxed field equations and the no-incoming radiation condition, is given as the retarded integral (10). We now assume that outside the isolated system, say, in the region $`r>`$ where $``$ is a constant radius strictly larger than $`d`$, we have $`h=(h)`$ where $`(h)`$ denotes the multipole expansion of $`h`$, a solution of the vacuum field equations in $`IR^4`$ deprived from the spatial origin $`r=0`$, and admitting a spherical-harmonics expansion of a certain structure (see below). Thus, in $`IR\times IR_{}^3`$ where $`IR_{}^3IR^3\{\mathrm{𝟎}\}`$,
$`_\nu (h^{\mu \nu })`$ $`=`$ $`0,`$ (16)
$`(h^{\mu \nu })`$ $`=`$ $`(\mathrm{\Lambda }^{\mu \nu }).`$ (17)
The source term $`(\mathrm{\Lambda })`$ is obtained from inserting $`(h)`$ in place of $`h`$ into (5), i.e. $`(\mathrm{\Lambda })\mathrm{\Lambda }((h))`$. \[Since the matter tensor has a compact support, $`(T)=0`$ so that $`(\tau )=\frac{c^4}{16\pi G}(\mathrm{\Lambda })`$.\] Of course, inside the source (when $`rd`$), the true solution $`h`$ differs from the vacuum solution $`(h)`$, the latter becoming in fact singular at the origin ($`r=0`$). We assume that the spherical-harmonics expansion of $`(h)`$ in $`IR\times IR_{}^3`$ reads
$$(h)(𝐱,t)=\underset{aN}{}\widehat{n}_Lr^a(\mathrm{ln}r)^p{}_{L}{}^{}F_{a,p}^{}(t)+R_N(𝐱,t).$$
(18)
This expression is valid for any $`NIN`$. The powers of $`r`$ are positive or negative, $`aZZ`$, and we have $`aN`$ (the negative powers of $`r`$ show that the multipole expansion is singular at $`r=0`$). For ease of notation we indicate only the summation over $`a`$, but there are two other summations involved: one over the powers $`pIN`$ of the logarithms, and one over the order of multipolarity $`lIN`$. The summations are considered only in the sense of formal series, as we do not control the mathematical nature of the series. The factor $`\widehat{n}_L`$ is a product of $`l`$ unit vectors, $`n_Ln^Ln^{i_1}\mathrm{}n^{i_l}`$, where $`Li_1\mathrm{}i_l`$ is a multi-index with $`l`$ indices, on which the symmetric and trace-free (STF) projection is applied: $`\widehat{n}_L\mathrm{STF}[n_L]`$. The decomposition in terms of STF tensors $`\widehat{n}_L(\theta ,\phi )`$ is equivalent to the decomposition in usual spherical harmonics. The functions $`{}_{L}{}^{}F_{a,p}^{}(t)`$ are smooth ($`C^{\mathrm{}}`$) functions of time, which become constant when $`t𝒯`$ because of our assumption (9). \[Of course, the $`{}_{L}{}^{}F_{a,p}^{}`$’s depend also on $`c`$: $`{}_{L}{}^{}F_{a,p}^{}(t,c)`$.\] Finally the function $`R_N(𝐱,t)`$ is defined by continuity throughout $`IR^4`$. Its two essential properties are $`R_NC^N(IR^4)`$ and $`R_N=O(r^N)`$ when $`r0`$ with fixed $`t`$. In addition $`R_N`$ is zero before the time $`𝒯`$. Though the function $`R_N(𝐱,t)`$ is given “globally” (as is the multipole expansion), it represents a small remainder $`O(r^N)`$ in the expansion of $`(h)`$ when $`r0`$, which is to be identified with the “near-zone” expansion of the field outside the source. It is convenient to introduce a special notation for the formal near-zone expansion (valid to any order $`N`$):
$$\overline{(h)}(𝐱,t)=\widehat{n}_Lr^a(\mathrm{ln}r)^p{}_{L}{}^{}F_{a,p}^{}(t),$$
(19)
where the summation is to be understood in the sense of formal series. \[Note that (14) and (15) are written for the field variable $`(h)`$, but it is easy to check that the same type of structure holds also for the source term $`(\mathrm{\Lambda })`$.\]
Our justification of the assumed structure (14) is that it has been proved to hold for metrics in the class of Multipolar-post-Minkowskian (MPM) metrics considered in , i.e. formal series $`h_{\mathrm{ext}}=G^nh_n`$ which satisfy the vacuum equations, are stationary in the past, and depend on a finite set of independent multipole moments. More precisely, from the theorem 4.1 in , the general MPM metric $`h_{\mathrm{ext}}`$, that we identify in this paper with $`(h^{\mu \nu })`$, is such that the property (14) holds for the $`h_n`$’s to any order $`n`$, with the only difference that to any finite order $`n`$ the integers $`a,p,l`$ vary into some finite ranges, namely $`a_{\mathrm{min}}(n)aN`$, $`0pn1`$ and $`0ll_{\mathrm{max}}(n)`$, with $`a_{\mathrm{min}}(n)\mathrm{}`$ and $`l_{\mathrm{max}}(n)+\mathrm{}`$ when $`n+\mathrm{}`$. The functions $`{}_{L}{}^{}F_{a,p}^{}`$ and the remainder $`R_N`$ in (14) should therefore be viewed as post-Minkowskian series $`G^n{}_{L}{}^{}F_{a,p,n}^{}`$ and $`G^nR_{N,n}`$. What we have done in writing (14) and (15) is to assume that one can legitimately consider such formal post-Minkowskian series. Note that because the general MPM metric represents the most general solution of the field equations outside the source (Theorem 4.2 in ), it is quite appropriate to identify the general multipole expansion $`(h)`$ with the MPM metric $`h_{\mathrm{ext}}`$. Actually we shall justify this assumption in Section 5 by recovering from $`(h)`$, step by step in the post-Minkowskian expansion, the MPM metric $`h_{\mathrm{ext}}`$. Because the properties are proved in for any $`n`$, and because we consider the formal post-Minkowskian sum, we see that (14)-(15), viewed as if it were “exact”, constitutes a quite natural assumption. In particular we have assumed in (14)-(15) that the multipolar series involves an infinite number of independent multipoles. In summary, we give to the properties (14)-(15) a scope larger than the one of MPM expansions (maybe they could be proved for exact solutions), at the price of counting them among our basic assumptions.
The multipole expansion $`(h)`$ is a mathematical solution of the vacuum equations in $`IR\times IR_{}^3`$, but whose “multipole moments” (the functions $`{}_{L}{}^{}F_{a,p}^{}`$) are not determined in terms of the source parameters. When the isolated system is slowly moving in the sense of (12), there exists an overlapping region between the domains of validity of the post-Newtonian expansion: the “near-zone” $`r<b`$, where $`d<b=O(\lambda /c)`$, and of the multipole expansion: the exterior zone $`r>`$. For this to be true it suffices to choose $``$, which is restricted only to be strictly larger than $`d`$, such that $`d<<b`$. We assume that the field $`h`$ given by (10) admits in the near-zone a formal post-Newtonian expansion, $`h=\overline{h}`$ when $`r<b`$. On the other hand, recall that $`h=(h)`$ when $`r>`$. Matching the two asymptotic expansions $`\overline{h}`$ and $`(h)`$ in the “matching” region $`<r<b`$ means that the (formal) double series obtained by considering the multipole expansion of (all the coefficients of) the post-Newtonian expansion $`\overline{h}`$ is identical to the double series obtained by taking the near-zone expansion of the multipole expansion. \[We use the same overbar notation for the post-Newtonian and near-zone expansions because the near-zone expansion ($`r/c0`$) of the exterior multipolar field is mathematically equivalent to the expansion when $`c\mathrm{}`$ with fixed multipole moments.\] The resulting matching equation reads
$$\overline{(h)}=(\overline{h}).$$
(20)
This equation should be true term by term, after both sides of the equation are re-arranged as series corresponding to the same expansion parameter. Though looking quite reasonable (if the theory makes sense), the matching equation cannot be justified presently with full generality; however up to 2PN order it was shown to determine a unique solution valid everywhere inside and outside the source . The matching assumption complements the framework of MPM approximations , by giving physical “pith” to the arbitrary multipole moments used in the construction of MPM metrics (see Section 4).
### B The field in terms of multipole moments
Let us consider the relaxed vacuum Einstein equation (13b), whose source term $`(\mathrm{\Lambda })`$, according to our assumptions, owns the structure (14) \[recall that (14) applies to $`(h)`$ as well as $`(\mathrm{\Lambda })`$\]. We obtain a particular solution of this equation (in $`IR\times IR_{}^3`$) as follows. First we multiply each term composing $`(\mathrm{\Lambda })`$ in (14) by a factor $`(r/r_0)^B`$, where $`B`$ is a complex number and $`r_0`$ a constant with the dimension of a length. For each term we can choose the real part of $`B`$ large enough so that the term becomes regular when $`r0`$, and then we can apply the retarded integral (11). The resulting $`B`$-dependent retarded integral is known to be analytically continuable for any $`BIC`$ except at integer values including in general the value of interest $`B=0`$. Furthermore one can show that the finite part (in short $`\mathrm{FP}_{B=0}`$) of this integral, defined to be the coefficient of the zeroth power of $`B`$ in the expansion when $`B0`$, is a retarded solution of the corresponding wave equation. In the case of a regular term in (14) such as the remainder $`R_N`$, this solution simply reduces to the retarded integral. Summing all these solutions, corresponding to all the separate terms in (14), we thereby obtain as a particular solution of (13b) the object $`\mathrm{FP}_{B=0}_R^1[(r/r_0)^B(\mathrm{\Lambda })]`$. This is basically the method employed in to solve the vacuum field equations in the post-Minkowskian approximation.
Now all the problem is to find the homogeneous solution to be added to the latter particular solution in order that the multipole expansion $`(h)`$ matches with the post-Newtonian expansion $`\overline{h}`$, solution within the source of the field equation (3) \[or, rather, (10)\]. Finding this homogeneous solution means finding the general consequence of the matching equation (16). The result is that the multipole expansion $`h^{\mu \nu }`$ satisfying the Einstein equation (10) together with the matching equation (16) reads
$$(h^{\mu \nu })=\mathrm{FP}_{B=0}_R^1[(r/r_0)^B(\mathrm{\Lambda }^{\mu \nu })]\frac{4G}{c^4}\underset{l=0}{\overset{+\mathrm{}}{}}\frac{()^l}{l!}_L\left\{\frac{1}{r}_L^{\mu \nu }(tr/c)\right\}$$
(21)
where the first term is the previous particular solution, and where the second term is a retarded solution of the source-free (homogeneous) wave equation, whose “multipole moments” are given explicitly by ($`utr/c`$)
$$_L^{\mu \nu }(u)=\mathrm{FP}_{B=0}d^3𝐱|𝐱/r_0|^Bx_L\overline{\tau }^{\mu \nu }(𝐱,u).$$
(22)
Here $`\overline{\tau }^{\mu \nu }`$ denotes the post-Newtonian expansion of the stress-energy pseudo-tensor $`\tau ^{\mu \nu }`$ appearing in the right side of (10). In (17) and (18) we denote $`L=i_1\mathrm{}i_l`$ and $`_L_{i_1}\mathrm{}_{i_l}`$, $`x_Lx_{i_1}\mathrm{}x_{i_l}`$.
It is important that the multipole moments (18) are found to depend on the post-Newtonian expansion $`\overline{\tau }^{\mu \nu }`$ of the pseudo-tensor, and not of $`\tau ^{\mu \nu }`$ itself, as this is precisely where our assumption of matching to the inner post-Newtonian field comes in. The formula is a priori valid only in the case of a slowly-moving source; it is a priori true only after insertion of a definite post-Newtonian expansion of the pseudo-tensor, where in particular all the retardations have been expanded when $`c\mathrm{}`$ \[the formulas (17)-(18) assume implicitly that one can effectively construct such a post-Newtonian expansion\].
Like in the first term of (17), the moments (18) are endowed with a finite part operation defined by complex analytic continuation in $`B`$. Notice however that the two finite part operations in the first term of (17) and in (18) act quite differently. In the first term of (17) the analytic continuation serves at regularizing the singularity of the multipole expansion at the spatial origin $`r=0`$. Since the pseudo-tensor is smooth inside the source, there is no need in the moments (18) to regularize the field near the origin; still the finite part is essential because it applies to the bound of the integral at infinity ($`|𝐱|\mathrm{}`$). Otherwise the integral would be (a priori) divergent at infinity, because of the presence of the factor $`x_L=O(r^l)`$ in the integrand, and the fact that the pseudo-tensor $`\overline{\tau }^{\mu \nu }`$ is non-compact supported. The two finite parts present in the two separate terms of (17) involve the same arbitrary constant $`r_0`$, but this constant can be readily checked to cancel out between the two terms \[i.e. the differentiation of $`(h^{\mu \nu })`$ with respect to $`r_0`$ yields zero\].
The formulas (17)-(18) were first obtained (in STF form) up to the 2PN order in by performing explicitly the matching. This showed in particular that the matching equation (16) is correct to 2PN order. Then the proof valid to any post-Newtonian order, but at the price of assuming (16) to all orders, was given in Section 3 of (see also Appendix A of for an alternative proof). The crucial step in the proof is to remark that the finite part of the integral of $`\overline{(\mathrm{\Lambda })}`$ over the whole space $`IR^3`$ is identically zero by analytic continuation:
$$\mathrm{FP}_{B=0}_{IR^3}d^3𝐱|𝐱/r_0|^Bx_L\overline{(\mathrm{\Lambda })}(𝐱,u)=0.$$
(23)
This follows from the fact that $`\overline{(\mathrm{\Lambda })}`$ can be written as a formal series of the type (15). Using (15) it is easy to reduce the computation of the integral (19) to that of the elementary radial integral $`_0^+\mathrm{}d|𝐱||𝐱|^{B+2+l+a}`$ (since the powers of the logarithm can be obtained by repeatedly differentiating with respect to $`B`$). The latter radial integral can be split into a “near-zone” integral, extending from zero to radius $``$, and a “far-zone” integral, extending from $``$ to infinity (actually any finite non-zero radius fits instead of $``$). When the real part of $`B`$ is a large enough positive number, the value of the near-zone integral is $`^{B+3+l+a}/(B+3+l+a)`$, while when the real part of $`B`$ is a large negative number, the far-zone integral reads the opposite, $`^{B+3+l+a}/(B+3+l+a)`$. Both obtained values represent the unique analytic continuations of the near-zone and far-zone integrals for any $`BIC`$ except $`3la`$. The complete integral $`_0^+\mathrm{}d|𝐱||𝐱|^{B+2+l+a}`$ is defined as the sum of the analytic continuations of the near-zone and far-zone integrals, and is therefore identically zero ($`BIC`$); this proves (19).
One may ask why the whole integration over $`IR^3`$ contributes to the multipole moment (18) – a somewhat paradoxical fact because the integrand is in the form of a post-Newtonian expansion, and is thus expected to be physically valid (i.e. to give accurate results) only in the near zone. This fact is possible thanks to the technical identity (19) which enables us to transform a near-zone integration into a complete $`IR^3`$-integration (refer to for details).
### C Equivalence with the Will-Wiseman multipole expansion
Recently a different expression of the multipole decomposition, with correlatively a different expression of the multipole moments, was obtained by Will and Wiseman , extending previous work of Epstein and Wagoner and Thorne . Basically, the multipole moments in are defined by an integral extending over a ball of finite radius $``$ (essentially the same $``$ as here), and thus do not require any regularization of the bound at infinity. By contrast, our multipole moments (18) involve an integration over the whole $`IR^3`$, which is allowed thanks to the analytic continuation \[leading to the identity (19)\]. Let us outline the proof of the equivalence between the Will-Wiseman formalism and the present one .
Will and Wiseman find, instead of (17)-(18),
$$(h^{\mu \nu })=_R^1[(\mathrm{\Lambda }^{\mu \nu })]_|\frac{4G}{c^4}\underset{l=0}{\overset{+\mathrm{}}{}}\frac{()^l}{l!}_L\left\{\frac{1}{r}𝒲_L^{\mu \nu }(tr/c)\right\}.$$
(24)
The first term is given by the retarded integral (11) acting on $`(\mathrm{\Lambda })`$, but truntated, as indicated by the subscript $``$, to extend only in the “far zone”: $`|𝐱^{}|>`$ in the notation (11). Thus, the near-zone part of the retarded integral, which contains the source, is removed, and there is no problem with the singularity of the multipole expansion at the origin. Then, the multipole moments $`𝒲_L`$ are given by an integral extending over the “near zone” only:
$$𝒲_L^{\mu \nu }(u)=_{|𝐱|<}d^3𝐱x_L\overline{\tau }^{\mu \nu }(𝐱,u).$$
(25)
The integral being compact-supported is well-defined. The multipole moments $`𝒲_L`$ look technically more simple than ours given by (18). On the other hand, practically speaking, the analytic continuation in (18) permits deriving many closed-form formulas to be used in applications . Of course, one is free to choose any definition of the multipole moments as far as it is used in a consistent manner.
We compute the difference between the moments $`_L`$ and $`𝒲_L`$. For the comparison we split $`_L`$ into far-zone and near-zone integrals corresponding to the radius $``$. Since the analytic continuation factor in $`_L`$ deals only with the bound at infinity, it can be removed from the near-zone integral, which is then clearly seen to agree with $`𝒲_L`$. So the difference $`_L𝒲_L`$ is given by the far-zone integral:
$$_L(u)𝒲_L(u)=\mathrm{FP}_{B=0}_{|𝐱|>}d^3𝐱|𝐱/r_0|^Bx_L\overline{\tau }(𝐱,u).$$
(26)
Next we transform the integrand. Successively we write $`\overline{\tau }=(\overline{\tau })`$ because we are in the far zone; $`(\overline{\tau })=\overline{(\tau )}`$ from the matching equation (16); and $`\overline{(\tau )}=\frac{c^4}{16\pi G}\overline{(\mathrm{\Lambda })}`$ because $`T`$ has a compact support. At this stage, the technical identity (19) allows one to transform the far-zone integration into a near zone integration (changing simply the overall sign in front of the integral). So,
$$_L(u)𝒲_L(u)=\frac{c^4}{16\pi G}\mathrm{FP}_{B=0}_{|𝐱|<}d^3𝐱|𝐱/r_0|^Bx_L\overline{(\mathrm{\Lambda })}(𝐱,u).$$
(27)
It is straightforward to check that the right side of this equation, when summed up over all multipolarities $`l`$, accounts exactly for the near-zone part that was removed from the retarded integral of $`(\mathrm{\Lambda })`$ \[first term in (20)\], so that the “complete” retarded integral as given by the first term in (17) is exactly reconstituted. In conclusion the two formalisms and are equivalent.
## III Source multipole moments
Quite naturally our source multipole moments will be closely related to the $`_L`$’s obtained in (18). However, before giving a precise definition, we need to find the equivalent of the multipole decomposition (17)-(18) in terms of symmetric and trace-free (STF) tensors, and we must reduce the number of independent tensors by imposing the harmonic gauge condition (13a). This leads to the definition of a “linearized” metric associated with the multipole expansion $`(h)`$, and parametrized by six sets of STF source multipole moments.
### A Multipole expansion in symmetric-trace-free form
The moments $`_L`$ given by (18) are non-trace-free because $`x_L`$ owns all its traces (i.e. $`\delta _{i_li_{l1}}x_L=𝐱^2x_{L2}`$, where $`L2=i_1\mathrm{}i_{l2}`$). Instead of $`_L`$, there are certain advantages in using STF multipole moments: indeed the STF moments are uniquely defined, and they often yield simpler computations in practice. It is not difficult, using STF techniques, to obtain the multipole decomposition equivalent to (17)-(18) but expressed in terms of STF tensors. We find
$$(h^{\mu \nu })=\mathrm{FP}_{B=0}_R^1[(r/r_0)^B(\mathrm{\Lambda }^{\mu \nu })]\frac{4G}{c^4}\underset{l=0}{\overset{+\mathrm{}}{}}\frac{()^l}{l!}_L\left\{\frac{1}{r}_L^{\mu \nu }(tr/c)\right\}$$
(28)
where the STF multipole moments are given by
$$_L^{\mu \nu }(u)=\mathrm{FP}_{B=0}d^3𝐱|𝐱/r_0|^B\widehat{x}_L_1^1𝑑z\delta _l(z)\overline{\tau }^{\mu \nu }(𝐱,u+z|𝐱|/c).$$
(29)
The notation for a STF product of vectors is $`\widehat{x}_L\mathrm{STF}(x_L)`$ (such that $`\widehat{x}_L`$ is symmetric in $`L`$ and $`\delta _{i_li_{l1}}\widehat{x}_L=0`$; for instance $`\widehat{x}_{ij}=x_ix_j\frac{1}{3}\delta _{ij}𝐱^2`$). As we see, the STF moments (25) involve an extra integration, over the variable $`z`$, with respect to the non-STF ones (18). The weighting function associated with the $`z`$-integration reads, for any $`l`$,
$$\delta _l(z)=\frac{(2l+1)!!}{2^{l+1}l!}(1z^2)^l;_1^1𝑑z\delta _l(z)=1.$$
(30)
In the limit of large $`l`$ the weighting function tends toward the Dirac delta measure (hence its name): $`lim_l\mathrm{}\delta _l=\delta `$. Remark that since (25) is valid only in the post-Newtonian approximation, the $`z`$-integration is to be expressed as a post-Newtonian series. Here is the relevant formula :
$$_1^1𝑑z\delta _l(z)\overline{\tau }(𝐱,u+z|𝐱|/c)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(2l+1)!!}{2^kk!(2l+2k+1)!!}\left(\frac{|𝐱|}{c}\frac{}{u}\right)^{2k}\overline{\tau }(𝐱,u).$$
(31)
In the limiting case of linearized gravity, one can neglect the first term in (24), and the pseudo-tensor $`\overline{\tau }^{\mu \nu }`$ in (25) can be replaced by the matter stress-energy tensor $`T^{\mu \nu }`$ (we have $`\overline{T}^{\mu \nu }=T^{\mu \nu }`$ inside the slowly-moving source). Since $`T^{\mu \nu }`$ has a compact support the finite part prescription can be removed, and we recover the known multipole decomposition corresponding to a compact-support source (see the appendix B in ).
### B Linearized approximation to the exterior field
Up to now we have solved the relaxed field equation (10) in the exterior zone, with result the multipole decomposition (24)-(25). In this section we further impose the harmonic gauge condition (13a), and from this we find a solution of the linearized vacuum equation, appearing as the first approximation in a post-Minkowskian expansion of the multipole expansion $`(h)`$.
Let us give a notation to the first term in (24):
$$u^{\mu \nu }\mathrm{FP}_{B=0}_R^1[(r/r_0)^B(\mathrm{\Lambda }^{\mu \nu })].$$
(32)
Applying on (24) the condition $`_\nu (h^{\mu \nu })=0`$, we find that the divergence $`w^\mu _\nu u^{\mu \nu }`$ is equal to a retarded solution of the source-free wave equation, given by
$$w^\mu =\frac{4G}{c^4}_\nu \left(\underset{l=0}{\overset{+\mathrm{}}{}}\frac{()^l}{l!}_L\left\{\frac{1}{r}_L^{\mu \nu }(tr/c)\right\}\right).$$
(33)
Now, associated to any $`w^\mu `$ of this type, there exists some $`v^{\mu \nu }`$ which is like $`w^\mu `$ a retarded solution of the source-free wave equation, $`(v^{\mu \nu })=0`$, and furthermore whose divergence is the opposite of $`w^\mu `$, $`_\nu v^{\mu \nu }=w^\mu `$. We refer to for the explicit formulas allowing the “algorithmic” construction of $`v^{\mu \nu }`$ once we know $`w^\mu `$. For definiteness, we adopt the formulas (2.12) in , which represent themselves a slight modification of the earlier formulas (4.13) in (see also the appendix B in ).
With $`v^{\mu \nu }`$ at our disposal we define what constitutes the linearized approximation to the exterior metric, say $`Gh_1^{\mu \nu }`$ where we factorize out $`G`$ in front of the metric in order to emphasize its linear character:
$$Gh_1^{\mu \nu }\frac{4G}{c^4}\underset{l=0}{\overset{+\mathrm{}}{}}\frac{()^l}{l!}_L\left\{\frac{1}{r}_L^{\mu \nu }(tr/c)\right\}v^{\mu \nu }.$$
(34)
The linearized metric satisfies the linearized vacuum equations in harmonic gauge: $`h_1^{\mu \nu }=0`$ since both terms in (30) satisfy the source-free wave equation, and $`_\nu h_1^{\mu \nu }=0`$ thanks to (29) and $`_\nu v^{\mu \nu }=w^\mu `$. Using the definition (30) one can re-write the multipole expansion of the exterior field as
$$(h^{\mu \nu })=Gh_1^{\mu \nu }+u^{\mu \nu }+v^{\mu \nu }.$$
(35)
Quite naturally the $`u^{\mu \nu }`$ and $`v^{\mu \nu }`$ will represent the non-linear corrections to be added to the “linearized” metric $`Gh_1^{\mu \nu }`$ in order to reconstruct the complete exterior metric (see Section 4).
Since $`h_1^{\mu \nu }`$ satisfies $`h_1^{\mu \nu }=0=_\nu h_1^{\mu \nu }`$, there is a unique way to decompose it into the sum of a “canonical” metric introduced by Thorne (see also ) plus a linearized gauge transformation,
$$h_1^{\mu \nu }=h_{\mathrm{can1}}^{\mu \nu }+^\mu \phi _1^\nu +^\nu \phi _1^\mu \eta ^{\mu \nu }_\lambda \phi _1^\lambda .$$
(36)
The canonical linearized metric is defined by
$`h_{\mathrm{can1}}^{00}`$ $`=`$ $`{\displaystyle \frac{4}{c^2}}{\displaystyle \underset{l0}{}}{\displaystyle \frac{()^l}{l!}}_L\left({\displaystyle \frac{1}{r}}I_L(u)\right),`$ (38)
$`h_{\mathrm{can1}}^{0i}`$ $`=`$ $`{\displaystyle \frac{4}{c^3}}{\displaystyle \underset{l1}{}}{\displaystyle \frac{()^l}{l!}}\{_{L1}\left({\displaystyle \frac{1}{r}}I_{iL1}^{(1)}(u)\right)`$ (40)
$`+{\displaystyle \frac{l}{l+1}}\epsilon _{iab}_{aL1}\left({\displaystyle \frac{1}{r}}J_{bL1}(u)\right)\},`$
$`h_{\mathrm{can1}}^{ij}`$ $`=`$ $`{\displaystyle \frac{4}{c^4}}{\displaystyle \underset{l2}{}}{\displaystyle \frac{()^l}{l!}}\{_{L2}\left({\displaystyle \frac{1}{r}}I_{ijL2}^{(2)}(u)\right)`$ (42)
$`+{\displaystyle \frac{2l}{l+1}}_{aL2}\left({\displaystyle \frac{1}{r}}\epsilon _{ab(i}J_{j)bL2}^{(1)}(u)\right)\},`$
where the $`I_L`$’s and $`J_L`$’s are two sets of functions of the retarded time $`u=tr/c`$ \[the subscript $`(n)`$ indicates $`n`$ time derivatives\], and which are STF with respect to all their indices $`L=i_1\mathrm{}i_l`$ (the symmetrization is denoted with parenthesis). As for the gauge vector $`\phi _1^\mu `$, it satisfies $`\phi _1^\mu =0`$ and depends in a way similar to (33) on four other sets of STF functions of $`u`$, denoted $`W_L`$, $`X_L`$, $`Y_L`$ and $`Z_L`$ (one type of function for each component of the vector). See for the expression of $`\phi _1^\mu =\phi _1^\mu [W_L,X_L,Y_L,Z_L]`$.
### C Derivation of the source multipole moments
The two sets of multipole moments $`I_L`$ and $`J_L`$ parametrizing the metric (33) constitute our definitions for respectively the mass-type and current-type multipole moments of the source. Actually, there are also the moments $`W_L`$, $`X_L`$, $`Y_L`$, $`Z_L`$, and we refer collectively to $`\{I_L,J_L,W_L,X_L,Y_L,Z_L\}`$ as the set of six source multipole moments.
With (32) it is easily seen (because $`\phi _1^\mu =0`$) that the gauge condition $`_\nu h_1^{\mu \nu }=0`$ imposes no condition on the source moments except the conservation laws appropriate to the gravitational monopole $`I`$ (having $`l=0`$) and dipoles $`I_i`$, $`J_i`$ ($`l=1`$): namely,
$$I^{(1)}=0;I_i^{(2)}=0;J_i^{(1)}=0.$$
(43)
The mass monopole $`I`$ and current dipole $`J_i`$ are thus constant, and agree respectively with the ADM mass and total angular momentum of the isolated system (later we shall denote the ADM mass by $`MI`$). According to (34) the mass dipole $`I_i`$ is a linear function of time, but since we assumed that the metric is stationary in the past, $`I_i`$ is in fact also constant, and equal to the (ADM) center of mass position.
The expressions of $`I_L`$ and $`J_L`$ (as well as of the other moments $`W_L,X_L,Y_L,`$ $`Z_L`$) come directly from (30) with (32)-(33) and the result of the matching, which is personified by the formula (25). To simplify the notation we define
$`\mathrm{\Sigma }`$ $``$ $`{\displaystyle \frac{\overline{\tau }^{00}+\overline{\tau }^{ii}}{c^2}},`$ (45)
$`\mathrm{\Sigma }_i`$ $``$ $`{\displaystyle \frac{\overline{\tau }^{0i}}{c}},`$ (46)
$`\mathrm{\Sigma }_{ij}`$ $``$ $`\overline{\tau }^{ij},`$ (47)
(where $`\overline{\tau }^{ii}\delta _{ij}\overline{\tau }^{ij}`$). The result is
$`I_L(u)`$ $`=`$ $`\text{FP}_{B=0}{\displaystyle }d^3𝐱|𝐱/r_0|^B{\displaystyle _1^1}dz\{\delta _l\widehat{x}_L\mathrm{\Sigma }{\displaystyle \frac{4(2l+1)}{c^2(l+1)(2l+3)}}\delta _{l+1}\widehat{x}_{iL}_t\mathrm{\Sigma }_i`$ (50)
$`+{\displaystyle \frac{2(2l+1)}{c^4(l+1)(l+2)(2l+5)}}\delta _{l+2}\widehat{x}_{ijL}_t^2\mathrm{\Sigma }_{ij}\}(𝐱,u+z|𝐱|/c),`$
$`J_L(u)`$ $`=`$ $`\epsilon _{ab<i_l}\text{FP}_{B=0}{\displaystyle }d^3𝐱|𝐱/r_0|^B{\displaystyle _1^1}dz\{\delta _l\widehat{x}_{L1>a}\mathrm{\Sigma }_b`$ (53)
$`{\displaystyle \frac{2l+1}{c^2(l+2)(2l+3)}}\delta _{l+1}\widehat{x}_{L1>ac}_t\mathrm{\Sigma }_{bc}\}(𝐱,u+z|𝐱|/c),`$
($`<>`$ refers to the STF projection). In a sense these expressions are exact, since they are formally valid up to any post-Newtonian order. \[See (68)-(69) below for explicit formulas at 2PN.\]
By replacing $`\overline{\tau }^{\mu \nu }`$ in (36) by the compact-support matter tensor $`T^{\mu \nu }`$ we recover the expressions of the multipole moments worked out in linearized gravity by Damour and Iyer (see also ). On the other hand the formulas (36) contain the results obtained by explicit implementation (“order by order”) of the matching up to the 2PN order .
## IV Post-Minkowskian approximation
In linearized gravity, the source multipole moments represent also the moments which are “measured” at infinity, using an array of detectors surrounding the source. However, in the non-linear theory, the gravitational source $`\mathrm{\Lambda }^{\mu \nu }`$ cannot be neglected and the first term in (24) plays a crucial role, notably it implies that the measured multipole moments at infinity differ from the source moments. Thus, we must now supplement the formulas of the source multipole moments (36) by the study of the “non-linear” term $`u^{\mu \nu }\mathrm{FP}_{B=0}_R^1[(r/r_0)^B(\mathrm{\Lambda }^{\mu \nu })]`$ in (24). For this purpose we develop following a post-Minkowskian approximation for the exterior vacuum metric.
### A Multipolar-post-Minkowskian iteration of the exterior field
The work started already with the formulas (31)-(33), where we expressed the exterior multipolar metric $`h_{\mathrm{ext}}^{\mu \nu }(h^{\mu \nu })`$ as the sum of the “linearized” metric $`Gh_1^{\mu \nu }`$ and the “non-linear” corrections $`u^{\mu \nu }`$, given by (28), and $`v^{\mu \nu }`$, algorithmically constructed from $`w^\mu =_\nu u^{\mu \nu }`$ \[see (29)\]. The linearized metric is a functional of the source multipole moments: $`h_1=h_1[I,J,W,X,Y,Z]`$. We regard $`G`$ as the book-keeping parameter for the post-Minkowskian series, and consider that $`Gh_1`$ is purely of first order in $`G`$, and thus that $`h_1`$ itself is purely of zeroth order. Of course we know from the previous section that this is untrue, because the source multipole moments depend on $`G`$; supposing $`h_1=O(G^0)`$ is simply a convention allowing the systematic implementation of the post-Minkowskian iteration.
Here we check that the non-linear corrections $`u^{\mu \nu }`$ and $`v^{\mu \nu }`$ in (31) generate the whole post-Minkowskian algorithm of . The detail demanding attention is how the post-Minkowskian expansions of $`u^{\mu \nu }`$ and $`v^{\mu \nu }`$ are related to a spliting of the gravitational source $`\mathrm{\Lambda }^{\mu \nu }`$ into successive non-linear terms. Let us pose, with obvious notation,
$$\mathrm{\Lambda }^{\mu \nu }=N^{\mu \nu }[h,h]+M^{\mu \nu }[h,h,h]+O(h^4),$$
(54)
where, from the exact formula (5), the quadratic-order piece reads (all indices being lowered with the Minkowski metric, and $`h`$ denoting $`\eta ^{\rho \sigma }h_{\rho \sigma }`$):
$`N^{\mu \nu }[h,h]=`$ $``$ $`h^{\rho \sigma }_{\rho \sigma }^2h^{\mu \nu }+{\displaystyle \frac{1}{2}}^\mu h_{\rho \sigma }^\nu h^{\rho \sigma }{\displaystyle \frac{1}{4}}^\mu h^\nu h`$ (55)
$``$ $`^\mu h_{\rho \sigma }^\rho h^{\nu \sigma }^\nu h_{\rho \sigma }^\rho h^{\mu \sigma }+_\sigma h^{\mu \rho }(^\sigma h_\rho ^\nu +_\rho h^{\nu \sigma })`$ (56)
$`+`$ $`\eta ^{\mu \nu }\left[{\displaystyle \frac{1}{4}}_\lambda h_{\rho \sigma }^\lambda h^{\rho \sigma }+{\displaystyle \frac{1}{8}}_\rho h^\rho h+{\displaystyle \frac{1}{2}}_\rho h_{\sigma \lambda }^\sigma h^{\rho \lambda }\right],`$ (57)
and where the cubic-order piece $`M[h,h,h]`$ and all higher-order terms can be obtained in a straightforward way.
First, reasoning ad absurdio, we prove (see for details) that both $`u`$ and $`v`$ indeed represent non-linear corrections to the linearized metric since they start at order $`G^2`$: $`u=G^2u_2+O(G^3)`$ and $`v=G^2v_2+O(G^3)`$. Next we obtain explicitly $`u_2`$ by substituting the linearized metric $`h_1`$ into (38) and applying the finite part of the retarded integral, i.e.
$$u_2^{\mu \nu }=\mathrm{FP}_{B=0}_R^1\left\{(r/r_0)^BN^{\mu \nu }[h_1,h_1]\right\}.$$
(58)
In this way we have a particular solution of the wave equation in $`IR\times IR_{}^3`$, $`u_2=N[h_1,h_1]`$. From $`u_2`$ one deduces $`v_2`$ by the same “algorithmic” equations as used when deducing $`v`$ from $`u`$ \[see after (29)\]. Then $`v_2=0`$ and the sum $`u_2+v_2`$ is divergenceless, so we can solve the quadratic-order vacuum equations in harmonic coordinates by posing
$$h_2^{\mu \nu }=u_2^{\mu \nu }+v_2^{\mu \nu }.$$
(59)
With this definition it is clear that the multipole expansion (31) reads to quadratic order:
$$(h^{\mu \nu })=Gh_1^{\mu \nu }+G^2h_2^{\mu \nu }+O(G^3).$$
(60)
Continuing in this fashion to the next order we find successively
$`u_3^{\mu \nu }=\mathrm{FP}_{B=0}_R^1\{(r/r_0)^B(M^{\mu \nu }[h_1,h_1,h_1]+N^{\mu \nu }[h_1,h_2]+N^{\mu \nu }[h_2,h_1])\};`$ (62)
(63)
$`h_3^{\mu \nu }=u_3^{\mu \nu }+v_3^{\mu \nu };`$ (64)
(65)
$`(h^{\mu \nu })=Gh_1^{\mu \nu }+G^2h_2^{\mu \nu }+G^3h_3^{\mu \nu }+O(G^4).`$ (66)
This process continues ad infinitum. The latter post-Minkowskian algorithm is exactly the one proposed in (see also Section 2 of ). That is, starting from $`h_1[I,J,W,X,Y,Z]`$ given by (32)-(33), one generates the infinite post-Minkowskian (MPM) series of , solving the vacuum (harmonic-coordinate) Einstein equations in $`IR\times IR_{}^3`$, and this formal series happens to be equal, term by term in $`G`$, to the general multipole decomposition of $`h^{\mu \nu }`$ given by (24). For any $`n`$, we have $`h_n^{\mu \nu }=u_n^{\mu \nu }+v_n^{\mu \nu }`$, and
$$(h^{\mu \nu })=\underset{n=1}{\overset{+\mathrm{}}{}}G^nh_n^{\mu \nu }.$$
(67)
This result is perfectly consistent with the fact that the MPM algorithm generates the most general solution of the field equations in $`IR\times IR_{}^3`$. Furthermore, the latter post-Minkowskian approximation is known to be reliable (existence of a one-parameter family of exact solutions whose Taylor expansion when $`G0`$ reproduces the approximation) – an interesting result which indicates that the multipole decomposition $`(h)`$ given by (24)-(25) might be proved within a context of exact solutions.
Recall that the source multipole moments $`I_L`$, $`J_L`$, $`W_L`$, $`X_L`$, $`Y_L`$, $`Z_L`$ entering the linearized metric $`h_1`$ at the basis of the post-Minkowskian algorithm are given by formulas like (36). Thus, in the present formalism, the source moments, including formally all post-Newtonian corrections \[and all possible powers of $`G`$\] as contained in (36), serve as “seeds” for the post-Minkowskian iteration of the exterior field, which as it stands leads to all possible non-linear interactions between the moments. As we can imagine, rapidly the formalism becomes extremely complicated when going to higher and higher post-Minkowskian and/or post-Newtonian approximations. Most likely the complexity is not due to the formalism but reflects the complexity of the field equations. It is probably impossible to find a different formalism in which things would be much simpler (except if one restricts to a particular type of source).
### B The “canonical” multipole moments
The previous post-Minkowskian algorithm started with $`h_1`$, a functional of six types of source multipole moments, $`I_L`$ and $`J_L`$ entering the “canonical” linearized metric $`h_{\mathrm{can1}}`$ given by (33), and $`W_L`$, $`X_L`$, $`Y_L`$, $`Z_L`$ parametrizing the gauge vector $`\phi _1`$ in (32). All these moments deserve their name of source moments, but clearly the moments $`W_L`$, $`X_L`$, $`Y_L`$ and $`Z_L`$ do not play a physical role at the level of the linearized approximation, as they simply parametrize a linear gauge transformation. But because the theory is covariant with respect to (non-linear) diffeomorphisms and not merely to linear gauge transformations, these moments do contribute to physical quantities at the non-linear level.
In practice, the presence of the moments $`W_L`$, $`X_L`$, $`Y_L`$, $`Z_L`$ complicates the post-Minkowskian iteration. Fortunately one can take advantage of the fact (proved in ) that it is always possible to parametrize the vacuum metric by means of two and only two types of multipole moments $`M_L`$ and $`S_L`$ (different from $`I_L`$ and $`J_L`$). The metric is then obtained by the same post-Minkowskian algorithm as in (39)-(43), but starting with the “canonical” linearized metric $`h_{\mathrm{can1}}[M,S]`$ instead of $`h_1[I,J,W,X,Y,Z]`$. The resulting non-linear metric $`h_{\mathrm{can}}`$ is isometric to our exterior metric $`h_{\mathrm{ext}}(h)`$, provided that the moments $`M_L`$ and $`S_L`$ are given in terms of the source moments $`I_L,J_L,\mathrm{},Z_L`$ by some specific relations
$`M_L`$ $`=`$ $`M_L[I,J,W,X,Y,Z],`$ (69)
$`S_L`$ $`=`$ $`S_L[I,J,W,X,Y,Z].`$ (70)
The two coordinate systems in which $`h_{\mathrm{can}}`$ and $`h_{\mathrm{ext}}`$ are defined satisfy the harmonic gauge condition in the exterior zone, but (probably) only the one associated with $`h_{\mathrm{ext}}`$ meshes with the harmonic coordinates in the interior zone. With the notation (32) the coordinate change reads $`\delta x^\mu =G\phi _1^\mu `$\+ non-linear corrections. We shall refer to the moments $`M_L`$ and $`S_L`$ as the mass-type and current-type canonical multipole moments. Of course, since at the linearized approximation the only “physical” moments are $`I_L`$ and $`J_L`$, we have
$`M_L`$ $`=`$ $`I_L+O(G),`$ (72)
$`S_L`$ $`=`$ $`J_L+O(G),`$ (73)
where $`O(G)`$ denotes the post-Minkowskian corrections. Furthermore, it can be shown that in terms of a post-Newtonian expansion the difference between both types of moments is very small: 2.5PN order, i.e.
$$M_L=I_L+O\left(\frac{1}{c^5}\right)$$
(74)
\[note that $`M=M_{\mathrm{ADM}}=I`$\]. Thus, from (46), the canonical moments are only “slightly” different from the source moments. Their usefulness is merely practical – in general they are used in place of the source moments to simplify a computation.
### C Retarded integral of a multipolar extended source
The previous post-Minkowskian algorithm has only theoretical interest unless we supply it with some explicit formulas for the computation of the coefficients $`h_n`$. Happily for us pragmatists, such formulas exist, and can be found in a rather elegant way thanks to the process of analytic continuation. Basically we need the retarded integral of an extended (non-compact-support) source with a definite multipolarity $`l`$. Here we present three exemplifying formulas; see the appendices A in and for more discussion.
Very often we meet a wave equation whose source term is of the type $`\widehat{n}_LF(tr/c)/r^k`$, where $`\widehat{n}_L`$ has multipolarity $`l`$ and $`F`$ denotes a certain product of multipole moments. \[Clearly, the near-zone expansion of such a term is of the form (15).\] When the power $`k`$ is such that $`3kl+2`$ (this excludes the scalar case $`l=0`$), we obtain the solution of the wave equation as
$`\text{FP}_{B=0}_R^1\left[(r/r_0)^B{\displaystyle \frac{\widehat{n}_L}{r^k}}F(tr/c)\right]={\displaystyle \frac{(k3)!(l+2k)!}{(l+k2)!}}\widehat{n}_L`$ (75)
$`\times {\displaystyle \underset{j=0}{\overset{k3}{}}}2^{k3j}{\displaystyle \frac{(l+j)!}{j!(lj)!}}{\displaystyle \frac{F^{(k3j)}(tr/c)}{c^{k3j}r^{j+1}}}.`$ (76)
As we see the (finite part of the) retarded integral depends in this case on the values of the extended source at the same retarded time $`tr/c`$ (for simplicity we use the same notation for the source and field points). But it is well known (see e.g. ) that this feature is exceptional; in most cases the retarded integral depends on the whole integrated past of the source. A chief example of such a “hereditary” character is the case with $`k=2`$ in the previous example, for which we find
$$_R^1\left[\frac{\widehat{n}_L}{r^2}F(tr/c)\right]=\frac{\widehat{n}_L}{r}_{\mathrm{}}^{ctr}𝑑sF(s/c)Q_l\left(\frac{cts}{r}\right)$$
(77)
where $`Q_l`$ denotes the Legendre function of the second kind, related to the usual Legendre polynomial $`P_l`$ by the formula
$$Q_l(x)=\frac{1}{2}P_l(x)\mathrm{ln}\left(\frac{x+1}{x1}\right)\underset{j=1}{\overset{l}{}}\frac{1}{j}P_{lj}(x)P_{j1}(x).$$
(78)
Since the retarded integral (48) is in fact convergent when $`r0`$, we have removed the factor $`(r/r_0)^B`$ and finite part prescription. When the source term itself is given by a “hereditary” expression such as the right side of (48), we get a more complicated but still manageable formula, for instance
$$_R^1\left[\frac{\widehat{n}_L}{r^2}_{\mathrm{}}^{ctr}𝑑sF(s/c)Q_p\left(\frac{cts}{r}\right)\right]=\frac{c\widehat{n}_L}{r}_{\mathrm{}}^{ctr}𝑑sF^{(1)}(s/c)R_{lp}\left(\frac{cts}{r}\right)$$
(79)
where $`F^{(1)}`$ denotes that anti-derivative of $`F`$ which is zero in the past \[from (9) we have restricted $`F`$ to be zero in the past\], and where
$$R_{lp}(x)=Q_l(x)_1^x𝑑yQ_p(y)\frac{dP_l}{dy}(y)+P_l(x)_x^+\mathrm{}𝑑yQ_p(y)\frac{dQ_l}{dy}(y).$$
(80)
Like in (48) we do not need a finite part operation. The function $`R_{lp}`$ is well-defined thanks to the behaviour of the Legendre function at infinity: $`Q_l(x)1/x^{l+1}`$ when $`x\mathrm{}`$.
The formulas (48)-(51) are needed to investigate the so-called tails of gravitational waves appearing at quadratic non-linear order, and even the tails generated by the tails themselves (“tails of tails”) which arise at cubic order . (These formulas do not show a dependence on the constant $`r_0`$, but other formulas do.)
## V Radiative multipole moments
In Section 2 we introduced the definition of a set of multipole moments $`\{I_L,J_L,W_L,X_L,Y_L,Z_L\}`$ for the isolated source, and in Section 3 we showed that the exterior field, and in particular the asymptotic field therein, is actually a complicated non-linear functional of the latter moments. Therefore, to define some source multipole moments is not sufficient by itself; this must be completed by a study of the relation between the adopted definition and some convenient far-field observables. The same is true of other definitions of source moments in different formalisms, such as in the Dixon local description of extended bodies , which should be completed by a connection to the far-zone gravitational field, for instance along the line proposed by in the case of the Dixon moments. In the present formalism, the connection rests on the relation between the so-called radiative multipole moments, denoted $`U_L`$ and $`V_L`$, and the source moments $`I_L`$, $`J_L`$,$`\mathrm{}`$, $`Z_L`$ \[in fact, for simplicity’s sake, we prefer using the two moments $`M_L`$ and $`S_L`$ instead of the more basic six source moments\].
### A Definition and general structure
The radiative moments $`U_L`$ (mass-type) and $`V_L`$ (current-type) are the coefficients of the multipolar decomposition of the leading $`1/R`$ part of the transverse-tracefree (TT) projection of the radiation field in radiative coordinates $`(T,𝐗)`$ (with $`R=|𝐗|`$ the radial distance to the source). Radiative coordinates are such that the metric coefficients admit an expansion when $`R\mathrm{}`$ in powers of $`1/R`$ (no logarithms of $`R`$). In radiative coordinates the retarded time $`TR/c`$ is light-like, or becomes asymptotically light-like when $`R\mathrm{}`$. By definition,
$`h_{ij}^{TT}(𝐗,T)`$ $`=`$ $`{\displaystyle \frac{4G}{c^2R}}𝒫_{ijab}(𝐍){\displaystyle \underset{l2}{}}{\displaystyle \frac{1}{c^ll!}}\{N_{L2}U_{abL2}`$ (82)
$`{\displaystyle \frac{2l}{c(l+1)}}N_{cL2}\epsilon _{cd(a}V_{b)dL2}\}+O({\displaystyle \frac{1}{R^2}}),`$
where $`N_i=X^i/R`$, $`N_{L2}=N_{i_1}\mathrm{}N_{i_{l2}}`$, $`N_{cL2}=N_cN_{L2}`$, and the TT algebraic projector reads $`𝒫_{ijab}=(\delta _{ia}N_iN_a)(\delta _{jb}N_jN_b)\frac{1}{2}(\delta _{ij}N_iN_j)(\delta _{ab}N_aN_b)`$. The radiative moments $`U_L`$ and $`V_L`$ depend on $`TR/c`$; from (52) they are defined $`l2`$. The radiative-coordinate retarded time differs from the corresponding harmonic-coordinate time by the well-known logarithmic deviation of light cones,
$$T\frac{R}{c}=t\frac{r}{c}\frac{2GM}{c^3}\mathrm{ln}\left(\frac{r}{r_0}\right)+O(G^2),$$
(83)
where we have introduced in the logarithm the same constant $`r_0`$ as in (39) (this corresponds simply to a choice of the origin of time in the far zone).
Now from the post-Minkowskian algorithm of Section 3, it is clear that the radiative moments $`U_L`$ and $`V_L`$ can be obtained to any post-Minkowskian order in principle, in the form of a non-linear series in the source or equivalently the canonical multipole moments $`M_L`$ and $`S_L`$. The practical detail (worked out in ) is to determine the transformation between harmonic and radiative coordinates, generalizing (53) to any post-Minkowskian order. The structure of e.g. the mass-type radiative moment is
$$U_L=M_L^{(l)}+\underset{n=2}{\overset{+\mathrm{}}{}}\frac{G^{n1}}{c^{3(n1)+2k}}X_{nL}.$$
(84)
The first term comes from the fact that the radiative moment reduces at the linearized approximation to the ($`l`$th time derivative of the) source or canonical moment. The second term represents the series of non-linear corrections, each of them is given by a certain $`X_{nL}`$ which is a $`n`$-linear functional of derivatives of multipole moments $`M_L`$ or $`S_L`$. Furthermore we know from e.g. (48) and (50) that each new non-linear iteration (which always involves a retarded integral) brings a priori a new “hereditary” integration with respect to the previous approximation. So we expect that $`X_{nL}`$ is of the form ($`UTR/c`$)
$$X_{nL}(U)=_{\mathrm{}}^U𝑑u_1\mathrm{}_{\mathrm{}}^U𝑑u_n𝒵_n(U,u_1,\mathrm{},u_n)M_{L_1}^{(a_1)}(u_1)\mathrm{}S_{L_n}^{(a_n)}(u_n)$$
(85)
where $`𝒵_n`$ denotes a certain kernel depending on time variables $`U,u_1,\mathrm{},u_n`$, and where the sum refers to all possibilities of coupling together the $`n`$ moments. \[See (56) below for examples of kernels $`𝒵_2`$ and $`𝒵_3`$.\] A useful information is obtained from imposing that $`𝒵_n`$ be dimensionless; this yields the powers of $`G`$ and $`1/c`$ in front of each non-linear term in (54), where $`k`$ is the number of contractions among the indices present on the $`n`$ moments (the current moments carrying their associated Levi-Civita symbol).
As an example of application of (54) let us suppose that one is interested in the 3PN or $`1/c^6`$ approximation. From (54) we have $`3(n1)+2k=6`$, and we deduce that the only possibility is $`n=3`$ (cubic non-linearity) and $`k=0`$ (no contractions between the moments). From this we infer immediately that the only possible multipole interaction at that order is between two mass monopoles and a multipole, i.e. $`M\times M\times M_L`$. This corresponds to the “tails of tails” computed explicitly in (56) below.
### B The radiative quadrupole moment to 3PN order
To implement the formula (54) a tedious computation is to be done, following in details the post-Minkowskian algorithm of Section 4 augmented by explicit formulas such as (47)-(51), and changing the coordinates from harmonic to radiative according to the prescription in . Here we present the result of the computation of the mass-type radiative quadrupole ($`l=2`$) up to the 3PN order:
$`U_{ij}(U)=M_{ij}^{(2)}(U)+2{\displaystyle \frac{GM}{c^3}}{\displaystyle _0^+\mathrm{}}𝑑\tau M_{ij}^{(4)}(U\tau )\left[\mathrm{ln}\left({\displaystyle \frac{c\tau }{2r_0}}\right)+{\displaystyle \frac{11}{12}}\right]`$ (86)
$`+{\displaystyle \frac{G}{c^5}}\{{\displaystyle \frac{2}{7}}{\displaystyle _0^+\mathrm{}}d\tau \left[M_{a<i}^{(3)}M_{j>a}^{(3)}\right](U\tau ){\displaystyle \frac{2}{7}}M_{a<i}^{(3)}M_{j>a}^{(2)}`$ (87)
$`{\displaystyle \frac{5}{7}}M_{a<i}^{(4)}M_{j>a}^{(1)}+{\displaystyle \frac{1}{7}}M_{a<i}^{(5)}M_{j>a}+{\displaystyle \frac{1}{3}}\epsilon _{ab<i}M_{j>a}^{(4)}S_b\}`$ (88)
$`+2\left({\displaystyle \frac{GM}{c^3}}\right)^2{\displaystyle _0^+\mathrm{}}𝑑\tau M_{ij}^{(5)}(U\tau )\left[\mathrm{ln}^2\left({\displaystyle \frac{c\tau }{2r_0}}\right)+{\displaystyle \frac{57}{70}}\mathrm{ln}\left({\displaystyle \frac{c\tau }{2r_0}}\right)+{\displaystyle \frac{124627}{44100}}\right]`$ (89)
$`+O\left({\displaystyle \frac{1}{c^7}}\right).`$ (90)
Recall that in this formula the moment $`M_{ij}`$ is the canonical moment which agrees with the source moment $`I_{ij}`$ up to a 2.5PN term \[see (46)\], and that the source moment $`I_{ij}`$ itself is given in terms of the pseudo-tensor of the source by (36a). See also the formulas (68)-(69) below for a more explicit expression of the source moment at the 2PN order \[of course, to be consistent, one should use (56) conjointly with 3PN expressions of the source moments\].
The “Newtonian” term in (56) corresponds to the quadrupole formalism. Next, there is a quadratic non-linear correction with multipole interaction $`M\times M_{ij}`$ representing the dominant effect of tails (scattering of linear waves off the space-time curvature generated by the mass $`M`$). This correction, computed in , is of order $`1/c^3`$ or 1.5PN and has the form of a hereditary integral with logarithmic kernel. The constant $`11/12`$ depends on the coordinate system chosen to cover the source, here the harmonic coordinates; for instance the constant would be $`17/12`$ in Schwarzschild-like coordinates . The next correction, of order $`1/c^5`$ or 2.5PN, is constituted by quadratic interactions between two mass quadrupoles, and between a mass quadrupole and a constant current dipole . This term contains a hereditary integral, of a type different from the tail integral, which is due to the gravitational radiation generated by the stress-energy distribution of linear waves . Sometimes this integral is referred to as the non-linear memory integral because it corresponds to the contribution of gravitons in the so-called linear memory effect . The non-linear memory integral can easily be found by using the effective stress-energy tensor of gravitational waves in place of the right side of (3); it follows also from rigorous studies of the field at future null infinity . Finally, at 3PN order in (56) appears the dominant cubic non-linear correction, corresponding to the interaction $`M\times M\times M_{ij}`$ and associated with the tails of tails of gravitational waves .
### C Tail contributions in the total energy flux
Observable quantities at infinity are expressible in terms of the radiative mass and current multipole moments. For instance the total gravitational-wave power emitted in all spatial directions (total gravitational flux or “luminosity” $``$) is given by the positive-definite multipolar series
$``$ $`=`$ $`{\displaystyle \underset{l=2}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{G}{c^{2l+1}}}\{{\displaystyle \frac{(l+1)(l+2)}{l(l1)l!(2l+1)!!}}U_L^{(1)}U_L^{(1)}`$ (92)
$`+{\displaystyle \frac{4l(l+2)}{c^2(l1)(l+1)!(2l+1)!!}}V_L^{(1)}V_L^{(1)}\}.`$
In the case of inspiralling compact binaries (a most prominent source of gravitational waves) the rate of inspiral is fixed by the flux $``$, which is therefore a crucial quantity to predict. Excitingly enough, we know that $``$ should be predicted to 3PN order for detection and analysis of inspiralling binaries in future experiments .
To 3PN order we can use the relation (56) giving the 3PN radiative quadrupole moment. Here we concentrate our attention on tails and tails of tails. The dominant tail contribution at 1.5PN order yields correspondingly a contribution in the total flux (with $`U=TR/c`$):
$$_{\mathrm{tail}}=\frac{4G^2M}{5c^8}I_{ij}^{(3)}(U)_0^+\mathrm{}𝑑\tau I_{ij}^{(5)}(U\tau )\left[\mathrm{ln}\left(\frac{c\tau }{2r_0}\right)+\frac{11}{12}\right].$$
(93)
Since we are interested in the dominant tail we have replaced using (46) the canonical mass quadrupole by the source quadrupole. Similarly there are some tail contributions due to the mass octupole, current quadrupole and all higher-order multipoles, but these are correlatively of higher post-Newtonian order \[see the factors $`1/c`$ in (57)\]. It has been shown that the work done by the dominant “hereditary” contribution in the radiation reaction force within the source – which arises at 4PN order in the equations of motion – agrees exactly with (58).
Next, because $``$ is made of squares of (derivatives of) radiative moments, it contains a term with the square of the tail integral at 1.5PN. This term arises at the 3PN relative order and reads
$$_{(\mathrm{tail})^2}=\frac{4G^3M^2}{5c^{11}}\left(_0^+\mathrm{}𝑑\tau I_{ij}^{(5)}(U\tau )\left[\mathrm{ln}\left(\frac{c\tau }{2r_0}\right)+\frac{11}{12}\right]\right)^2.$$
(94)
Finally, there is also the direct 3PN contribution of tails of tails in (56):
$`_{\mathrm{tail}(\mathrm{tail})}`$ $`=`$ $`{\displaystyle \frac{4G^3M^2}{5c^{11}}}I_{ij}^{(3)}(U){\displaystyle _0^+\mathrm{}}𝑑\tau I_{ij}^{(6)}(U\tau )`$ (96)
$`\times \left[\mathrm{ln}^2\left({\displaystyle \frac{c\tau }{2r_0}}\right)+{\displaystyle \frac{57}{70}}\mathrm{ln}\left({\displaystyle \frac{c\tau }{2r_0}}\right)+{\displaystyle \frac{124627}{44100}}\right].`$
By a control of all the hereditary integrals in $``$ up to 3PN we have checked that the terms (59)-(60) do exist. The two contributions (59) and (60) appear somewhat on the same footing – of course both should be taken into account in practical computations. Note that in a physical situation where the emission of radiation stops after a certain date, in the sense that the source multipole moments become constant after this date (assuming a consistent matter model which would do this at a given post-Newtonian order), the only contribution to $``$ which survives after the end of emission is the 3PN tail-square contribution (59).
## VI Post-Newtonian approximation
In Sections 2 and 3 we have reasoned upon the formal post-Newtonian expansion $`\overline{h}^{\mu \nu }`$ of the near-zone field to obtain the source multipole moments as functionals of the post-Newtonian pseudo-tensor $`\overline{\tau }^{\mu \nu }`$. We have also considered in Sections 4 and 5 the formal expansion $`c\mathrm{}`$ of the radiation field when holding the multipole moments fixed. Clearly missing in this scheme is an explicit algorithm for the computation of $`\overline{h}^{\mu \nu }`$ in the near zone. No such algorithm (say, in the spirit of the post-Minkowskian algorithm in Section 4) is known presently, but a lot is known on the first few post-Newtonian iterations .
The main difficulty in setting up a post-Newtonian algorithm is the appearance at some post-Newtonian order of divergent Poisson-like integrals. This comes from the fact that the post-Newtonian expansion is actually a near-zone expansion , which is valid only in the region where $`r=O(\lambda /c)`$, and that such an expansion blows up when taking formally the limit $`r+\mathrm{}`$. For instance, Rendall has shown that the post-Newtonian solution cannot be asymptotically flat starting at the 2PN or 3PN level, depending on the gauge. This is clear from the structure of the exterior near-zone expansion (15), which involves many positive powers of the radial distance $`r`$. Thus, one is not allowed in general to consider the limit $`r+\mathrm{}`$. In consequence, using the Poisson integral for solving a Poisson equation with non-compact-support source at a given post-Newtonian order is a priori meaningless. Indeed the Poisson integral not only extends over the near-zone but also over the regions at infinity. This means that the Poisson integral does not constitute the correct solution of the Poisson equation in this context. However, to the lowest post-Newtonian orders it works; for instance it was shown by Kerlick and Caporali that the post-Newtonian iteration (including the suggestion by Ehlers of an improvement with respect to previous work ) is well-defined up to the 2.5PN order where radiation reaction terms appear, but that some divergent integrals show up at the 3PN order.
Another difficulty is that the post-Newtonian approximation is in a sense not self-supporting, because it necessitates information coming from outside its own domain of validity. Of course we have in mind the boundary conditions at infinity which determine the radiation reaction in the source’s local equations of motion. Again, to the lowest post-Newtonian orders one can circumvent this difficulty by considering retarded integrals that are formally expanded when $`c\mathrm{}`$ as series of “instantaneous” Poisson-like integrals . However, this procedure becomes incorrect at the 4PN order, not to mention the problem of divergencies, because the near-zone field (as well as the source’s dynamics) ceases to be given by an instantaneous functional of the source parameters, due to the appearance of “tail-transported” hereditary integrals modifying the lowest-order radiation reaction damping .
Let us advocate here that the cure of the latter difficulty (and perhaps of all difficulties) is the matching equation (16). Indeed suppose that one knows a particular solution of the Poisson equation at some post-Newtonian order. This solution might be in the form of some “finite part” of a Poisson integral. The correct post-Newtonian solution will be the sum of this particular solution and of a homogeneous solution satisfying the Laplace equation, namely a harmonic solution, regular at the origin, which can always be written in the form $`A_L\widehat{x}_L`$, for some unknown constant tensors $`A_L`$. The homogeneous solution is associated with radiation reaction effects. Now the matching equation states that the multipole expansion of the post-Newtonian solution agrees with the near-zone expansion of the exterior field (which has been computed beforehand in Section 4). The multipole expansion of the known particular solution can be obtained by a standard method, and the multipole expansion of the homogeneous solution is simply itself, i.e. $`(A_L\widehat{x}_L)=A_L\widehat{x}_L`$. Therefore, we see that the matching equation determines in principle the homogeneous solution (i.e. all the unknown tensors $`A_L`$), and since the exterior field satisfies relevant boundary conditions at infinity, the $`A_L`$’s should correspond to the radiation reaction on a truly isolated system. See for implementation of this method to determine the radiation reaction force to 4PN order (1.5PN relative order).
### A The inner metric to 2.5PN order
Going to high post-Newtonian orders can become prohibitive because of the rapid proliferation of terms. Typically any allowed term (compatible dimension, correct index structure) does appear with a definite non-zero coefficient in front. However, high post-Newtonian orders can be manageable if one chooses some appropriate matter variables, and if one avoids expanding systematically the retardations due to the speed of propagation of gravity. Often it is sufficient, and clearer, to present a result in terms of matter variables still containing some $`c`$’s, and perhaps also in terms of some convenient retarded potentials (being clear that any retardation going to an order higher than the prescribed post-Newtonian order of the calculation is irrelevant). See for instance (65) and (68)-(69) below. Anyway, only in a final stage, when a result to the prescribed order is in hands, should we introduce the more basic matter variables (e.g. the coordinate mass density) and perform all necessary retardations. Then of course one does not escape to a profusion of terms, but at least we have been able to carry the post-Newtonian iteration using some reasonably simple expressions.
The matter variables are chosen in a way consistent with our earlier definitions (35), i.e.
$`\sigma `$ $``$ $`{\displaystyle \frac{T^{00}+T^{ii}}{c^2}};`$ (98)
$`\sigma _i`$ $``$ $`{\displaystyle \frac{T^{0i}}{c}};`$ (99)
$`\sigma _{ij}`$ $``$ $`T^{ij}.`$ (100)
To 2.5PN order one defines some retarded potentials $`V`$, $`V_i`$, $`\widehat{W}_{ij}`$, $`\widehat{X}`$ and $`\widehat{R}_i`$, with $`V`$ and $`V_i`$ looking like some retarded versions of the Newtonian and gravitomagnetic potentials, and $`\widehat{W}_{ij}`$ being associated with the matter and gravitational-field stresses:
$`V`$ $``$ $`_R^1\left\{4\pi G\sigma \right\},`$ (102)
$`V_i`$ $``$ $`_R^1\left\{4\pi G\sigma _i\right\},`$ (103)
$`\widehat{W}_{ij}`$ $``$ $`_R^1\left\{4\pi G(\sigma _{ij}\delta _{ij}\sigma _{kk})_iV_jV\right\},`$ (104)
$`\widehat{R}_i`$ $``$ $`_R^1\left\{4\pi G(V\sigma _iV_i\sigma )2_kV_iV_k{\displaystyle \frac{3}{2}}_tV_iV\right\},`$ (105)
$`\widehat{X}`$ $``$ $`_R^1\{4\pi GV\sigma _{ii}+2V_i_t_iV+V_t^2V`$ (107)
$`+{\displaystyle \frac{3}{2}}(_tV)^22_iV_j_jV_i+\widehat{W}_{ij}_{ij}^2V\},`$
where $`_R^1`$ denotes the retarded integral (11). All these potentials but $`V`$ and $`V_i`$ have a spatially non-compact support. The highest non-linearity entering them is cubic; it appears in the last term of $`\widehat{X}`$.
Based on the latter potentials one can show that the inner metric to order 2.5PN (in harmonic coordinates, $`_\nu (\sqrt{g}g^{\mu \nu })=0`$) takes the form
$`g_{00}`$ $`=`$ $`1+{\displaystyle \frac{2}{c^2}}V{\displaystyle \frac{2}{c^4}}V^2+{\displaystyle \frac{8}{c^6}}\left[\widehat{X}+V_iV_i+{\displaystyle \frac{V^3}{6}}\right]+O\left({\displaystyle \frac{1}{c^8}}\right),`$ (109)
$`g_{0i}`$ $`=`$ $`{\displaystyle \frac{4}{c^3}}V_i{\displaystyle \frac{8}{c^5}}\widehat{R}_i+O\left({\displaystyle \frac{1}{c^7}}\right),`$ (110)
$`g_{ij}`$ $`=`$ $`\delta _{ij}\left(1+{\displaystyle \frac{2}{c^2}}V+{\displaystyle \frac{2}{c^4}}V^2\right)+{\displaystyle \frac{4}{c^4}}\widehat{W}_{ij}+O\left({\displaystyle \frac{1}{c^6}}\right),`$ (111)
(writing $`\overline{g}_{\mu \nu }`$ would be more consistent with the notation of Section 2). With this form, we believe, the computational problems encountered in applications are conveniently divided into the specific problems associated with the computation of the various potentials (62), which constitute in this approach some appropriate computational “blocks” (having of course no physical signification separately). By expanding all powers of $`1/c`$ present into the matter densities (61) and into the retardations of the potentials (62), we find that the metric (63) becomes extremely complicated, as it really is (see e.g. ).
Because of our use of retarded potentials, the metric (63) involves explicitly only even post-Newtonian terms (using the post-Newtonian terminology that even terms correspond to even powers of $`1/c`$ in the equations of motion). We have checked that the odd post-Newtonian terms (responsible for radiation reaction), contained in (63) via the expansion of retardations, match, in the sense of the equation (16), to the exterior metric satisfying the no-incoming radiation condition (9).
The harmonic gauge condition implies some differential equations to be satisfied by the previous potentials. To 2.5PN order we find
$`_t\left\{V+{\displaystyle \frac{1}{c^2}}\left[{\displaystyle \frac{1}{2}}\widehat{W}_{ii}+2V^2\right]\right\}+_i\left\{V_i+{\displaystyle \frac{2}{c^2}}\left[\widehat{R}_i+VV_i\right]\right\}=O\left({\displaystyle \frac{1}{c^4}}\right),`$ (113)
$`_tV_i+_j\left\{\widehat{W}_{ij}{\displaystyle \frac{1}{2}}\delta _{ij}\widehat{W}_{kk}\right\}=O\left({\displaystyle \frac{1}{c^2}}\right),`$ (114)
where $`\widehat{W}_{ii}\delta _{ij}\widehat{W}_{ij}`$. These equations are in turn equivalent to the equation of continuity and the equation of motion for the matter system,
$`_t\sigma +_i\sigma _i`$ $`=`$ $`{\displaystyle \frac{1}{c^2}}\left(_t\sigma _{ii}\sigma _tV\right)+O\left({\displaystyle \frac{1}{c^4}}\right),`$ (116)
$`_t\sigma _i+_j\sigma _{ij}`$ $`=`$ $`\sigma _iV+O\left({\displaystyle \frac{1}{c^2}}\right).`$ (117)
Note that the precision is 1PN for the equation of continuity but only Newtonian for the equation of motion.
### B The mass-type source moment to 2.5PN order
From the 2.5PN metric (63) we obtain the pseudo-tensor $`\overline{\tau }`$ and the auxiliary quantities (35), that we replace into the formulas (36) to obtain the 2.5PN source multipole moments. Recall that the $`z`$-integration in the moments is to be carried out using the formula (27). Let us first see how this works at the 1PN order.
We need $`\mathrm{\Sigma }`$ to 1PN order and $`\mathrm{\Sigma }_i`$ to Newtonian order. The latter quantity reduces to the matter part, $`\mathrm{\Sigma }_i=\sigma _i+O(1/c^2)`$, and the former one reads after a simple transformation
$$\mathrm{\Sigma }=\sigma \frac{1}{2\pi Gc^2}\mathrm{\Delta }(V^2)+O\left(\frac{1}{c^4}\right).$$
(118)
The substitution into the moments $`I_L`$ given by (36a) leads to
$`I_L`$ $`=`$ $`\text{FP}_{B=0}{\displaystyle }d^3𝐱|𝐱/r_0|^B\{\widehat{x}_L\sigma {\displaystyle \frac{\widehat{x}_L}{2\pi Gc^2}}\mathrm{\Delta }(V^2)`$ (120)
$`+{\displaystyle \frac{|𝐱|^2\widehat{x}_L}{2c^2(2l+3)}}_t^2\sigma {\displaystyle \frac{4(2l+1)\widehat{x}_{iL}}{c^2(l+1)(2l+3)}}_t\sigma _i\}+O({\displaystyle \frac{1}{c^4}}).`$
The integrand is non-compact-supported because of the contribution of the second term, and accordingly we keep the regularization factor $`|𝐱/r_0|^B`$ and finite part operation. But let us operate by parts the second term, using the fact that $`|𝐱|^B\widehat{x}_L\mathrm{\Delta }(V^2)\mathrm{\Delta }(|𝐱|^B\widehat{x}_L)V^2=_i\{|𝐱|^B\widehat{x}_L_i(V^2)_i(|𝐱|^B\widehat{x}_L)V^2\}`$ is a pure divergence. When the real part of $`B`$ is a large negative number, we see thanks to the Gauss theorem that the latter divergence will not contribute to the moment, therefore by the unicity of the analytic continuation it will always yield zero contribution. Thus, using $`\mathrm{\Delta }\widehat{x}_L=0`$, we can replace $`|𝐱|^B\widehat{x}_L\mathrm{\Delta }(V^2)`$ in the second term of (67) by $`\mathrm{\Delta }(|𝐱|^B\widehat{x}_L)V^2=B(B+l+1)|𝐱|^{B2}\widehat{x}_LV^2`$, and because of the explicit factor $`B`$ we see that the second term can be non-zero only in the case where the factor $`B`$ multiplies an integral owning a simple pole $`1/B`$ due to the integration bound $`|𝐱|\mathrm{}`$. Expressing $`V^2`$ (to Newtonian order) in terms of source points $`𝐳_1`$ and $`𝐳_2`$, we obtain the integral $`d^3𝐱|𝐱|^{B2}\widehat{x}_L|𝐱𝐳_1|^1|𝐱𝐳_2|^1`$. When $`|𝐱|\mathrm{}`$ each $`|𝐱𝐳_{1,2}|^1`$ can be expanded as a series of $`\widehat{n}_{L_{1,2}}|𝐱|^{l_{1,2}1}`$; then performing the angular integration shows that the sum of “multipolarities” $`l+l_1+l_2`$ is necessarily an even integer. When this is realized the remaining radial integral reads $`d|𝐱||𝐱|^{B+ll_1l_22}`$ which develops a pole only when $`ll_1l_22=1`$. But that is incompatible with the previous finding. Thus the second term in (67) is identically zero, and we end up simply with a compact-support expression on which we no longer need to implement the finite part,
$$I_L=d^3𝐱\left\{\widehat{x}_L\sigma +\frac{|𝐱|^2\widehat{x}_L}{2c^2(2l+3)}_t^2\sigma \frac{4(2l+1)\widehat{x}_{iL}}{c^2(l+1)(2l+3)}_t\sigma _i\right\}+O\left(\frac{1}{c^4}\right).$$
(121)
This expression was first obtained in using a different method valid at 1PN order. Here we have recovered the same expression from the formula (36a) valid to any post-Newtonian order .
Only starting at the 2PN order does the mass multipole moment have a non-compact support (so the finite part becomes crucial at this order). By a detailed computation in we arrive at the following 2PN (or rather 2.5PN) expression:
$`I_L(t)`$ $`=`$ $`\mathrm{FP}_{B=0}{\displaystyle }d^3𝐱|𝐱/r_0|^B\{\widehat{x}_L[\sigma +{\displaystyle \frac{4}{c^4}}\sigma _{ii}V]+{\displaystyle \frac{|𝐱|^2\widehat{x}_L}{2c^2(2\mathrm{}+3)}}_t^2\sigma `$ (127)
$`+{\displaystyle \frac{|𝐱|^4\widehat{x}_L}{8c^4(2\mathrm{}+3)(2\mathrm{}+5)}}_t^4\sigma {\displaystyle \frac{2(2\mathrm{}+1)|𝐱|^2\widehat{x}_{iL}}{c^4(\mathrm{}+1)(2\mathrm{}+3)(2\mathrm{}+5)}}_t^3\sigma _i`$
$`+{\displaystyle \frac{2(2\mathrm{}+1)\widehat{x}_{ijL}}{c^4(\mathrm{}+1)(\mathrm{}+2)(2\mathrm{}+5)}}_t^2\left[\sigma _{ij}+{\displaystyle \frac{1}{4\pi G}}_iV_jV\right]`$
$`+{\displaystyle \frac{\widehat{x}_L}{\pi Gc^4}}\left[\widehat{W}_{ij}_{ij}^2V2V_i_t_iV+2_iV_j_jV_i{\displaystyle \frac{3}{2}}(_tV)^2V_t^2V\right]`$
$`{\displaystyle \frac{4(2\mathrm{}+1)\widehat{x}_{iL}}{c^2(\mathrm{}+1)(2\mathrm{}+3)}}_t[(1+{\displaystyle \frac{4V}{c^2}})\sigma _i`$
$`+{\displaystyle \frac{1}{\pi Gc^2}}(_kV[_iV_k_kV_i]+{\displaystyle \frac{3}{4}}_tV_iV)]\}+O({\displaystyle \frac{1}{c^6}}).`$
Recall that the canonical moment $`M_L`$ differs from the source moment $`I_L`$ at precisely the 2.5PN order \[see (46)\].
## VII Point-particles
So far the post-Newtonian formalism has been developed for smooth (i.e. $`C^{\mathrm{}}`$) matter distributions. As such, the source multipole moments (36) become ill-defined in the presence of singularities. We now argue that the formalism is in fact also applicable to singular sources (notably point-particles described by Dirac measures) provided that we add to our other basic assumptions a certain method for removing the infinite self-field of point-masses. Our main motivation is the inspiralling compact binary – a system of two compact objects (neutron stars or black holes) which can be described with great precision by two point-particles moving on a circular orbit, and whose orbital phase evolution should be computed prior to gravitational-wave detection with relative 3PN precision .
For this application we restrict ourselves to two point-masses $`m_1`$ and $`m_2`$ (constant Schwarzschild masses). The trajectories are $`𝐲_1(t)`$ and $`𝐲_2(t)`$ and the coordinate velocities $`𝐯_{1,2}=d𝐲_{1,2}/dt`$; we pose $`v_{1,2}^\mu =(c,𝐯_{1,2})`$. The symbol $`12`$ means the same term but with the labels of the two particles exchanged. A model for the stress-energy tensor of point-masses (say, at 2PN order) is
$`T_{\mathrm{point}\mathrm{mass}}^{\mu \nu }(𝐱,t)=\mu _1(t)v_1^\mu (t)v_1^\nu (t)\delta [𝐱𝐲_1(t)]+12;`$ (129)
$`\mu _1(t){\displaystyle \frac{m_1}{\sqrt{(gg_{\rho \sigma })_1{\displaystyle \frac{v_1^\rho v_1^\sigma }{c^2}}}}},`$ (130)
where $`\delta `$ denotes the three-dimensional Dirac measure, and $`g_{\mu \nu }`$ the metric coefficients in harmonic coordinates ($`g\mathrm{det}g_{\mu \nu }`$). The notation $`(gg_{\mu \nu })_1`$ means the value at the location of particle 1. However, due to the presence of the Dirac measure at particles 1 and 2, the metric coefficients will be singular at 1 and 2. Therefore, we must supplement the model (70) by a method of “regularization” able to give a sense to the ill-defined limit at 1 or 2. A priori the choice of one or another regularization constitutes a fully-qualified element of the model of point-particles. In the following we systematically employ the Hadamard regularization, based on the Hadamard “partie finie” of a divergent integral .
Let us discuss an example. The “Newtonian” potential $`U`$, defined by $`U=\mathrm{\Delta }^1(4\pi G\sigma )`$, where $`\sigma `$ is given by (61a) \[we have $`V=U+O(1/c^2)`$\], follows from (70a) as
$$U=\frac{G\mu _1}{r_1}\left[1+\frac{v_1^2}{c^2}\right]+12,$$
(131)
where $`r_1=|𝐱𝐲_1|`$. To Newtonian order $`U=Gm_1/r_1+O(1/c^2)+12`$. We compute $`U`$ at the 1PN order: from (70b) we deduce at this order $`\mu _1/m_1=1(U)_1/c^2+v_1^2/2c^2+O(1/c^4)`$, which involves $`U`$ itself taken at point 1, but of course this does not make sense because $`U`$ is singular at 1 and 2. Now, after applying the Hadamard regularization (described below), we obtain unambiguously the standard Newtonian result $`(U)_1=Gm_2/r_{12}+O(1/c^2)`$, where $`r_{12}=|𝐲_1𝐲_2|`$, that we insert back into $`\mu _1`$. So, $`U`$ at 1PN, and its regularized value at 1, read
$`U`$ $`=`$ $`{\displaystyle \frac{Gm_1}{r_1}}\left(1+{\displaystyle \frac{1}{c^2}}\left[{\displaystyle \frac{Gm_2}{r_{12}}}+{\displaystyle \frac{3}{2}}v_1^2\right]\right)+O\left({\displaystyle \frac{1}{c^4}}\right)+12,`$ (133)
$`(U)_1`$ $`=`$ $`{\displaystyle \frac{Gm_2}{r_{12}}}\left(1+{\displaystyle \frac{1}{c^2}}\left[{\displaystyle \frac{Gm_1}{r_{12}}}+{\displaystyle \frac{3}{2}}v_2^2\right]\right)+O\left({\displaystyle \frac{1}{c^4}}\right).`$ (134)
### A Hadamard partie finie regularization
We consider the class of functions of the field point $`𝐱`$ which are smooth on $`IR^3`$ except at the location of the two source points $`𝐲_{1,2}`$, around which the functions admit some power-like expansions in the radial distance $`r_1=|𝐱𝐲_1|`$, with fixed spatial direction $`𝐧_1=(𝐱𝐲_1)/r_1`$ (and idem for 2). Thus, for any $`F(𝐱)`$ in this class, we have
$`F`$ $`=`$ $`{\displaystyle \underset{a}{}}r_1^af_{1(a)}(𝐧_1)(\text{when }r_10);`$ (136)
$`F`$ $`=`$ $`{\displaystyle \underset{a}{}}r_2^af_{2(a)}(𝐧_2)(\text{when }r_20),`$ (137)
where the summation index $`a`$ ranges over values in $`ZZ`$ bounded from below, $`aa_0`$ (we do not need to be more specific), and where the coefficients of the various powers of $`r_{1,2}`$ depend on the spatial directions $`𝐧_{1,2}`$. In (73) we do not write the remainders for the expansions because we don’t need them; simply, we regard the expansions (73) as listings of the various coefficients $`f_{1(a)}`$ and $`f_{2(a)}`$. We assume also that the functions $`F`$ in this class decrease sufficiently rapidly when $`|𝐱|\mathrm{}`$, so that all integrals we consider are convergent at infinity.
The integral $`d^3𝐱F`$ is in general divergent because of the singular behaviour of $`F`$ near $`𝐲_{1,2}`$, but we can compute its partie finie ($`\mathrm{Pf}`$) in the sense of Hadamard . Let us consider two volumes surrounding the two singularities, of the form $`r_1s\rho _1(𝐧_1)`$ (and similarly for 2), where $`s`$ measures the size of the volume and $`\rho _1`$ gives its shape as a function of the direction $`𝐧_1`$ ($`\rho _1=1`$ in the case of a spherical ball). Using (73) it is easy to determine the expansion when $`s0`$ of the integral extending on $`IR^3`$ deprived from the two previous volumes, and then to subtract from the integral all the divergent terms when $`s0`$ in the latter expansion. The Hadamard partie finie is defined to be the limit when $`s0`$ of what remains. As it turns out, the result can be advantageously re-expressed in terms of an integral on $`IR^3`$ deprived from two spherical balls ($`\rho _{1,2}=1`$), at the price of introducing two constants $`s_{1,2}`$ which depend on the shape of the two regularizing volumes originally considered. With full generality the Hadamard partie finie of the divergent integral reads
$`\mathrm{Pf}{\displaystyle }d^3𝐱F\underset{s0}{lim}\{{\displaystyle _{\genfrac{}{}{0pt}{}{r_1>s}{r_2>s}}}d^3𝐱F`$ (138)
$`+{\displaystyle \underset{a+31}{}}{\displaystyle \frac{s^{a+3}}{a+3}}{\displaystyle }d\mathrm{\Omega }_1f_{1(a)}+\mathrm{ln}\left({\displaystyle \frac{s}{s_1}}\right){\displaystyle }d\mathrm{\Omega }_1f_{1(3)}+12\}`$ (139)
where $`s_1`$ is given by
$$\mathrm{ln}s_1=\frac{𝑑\mathrm{\Omega }_1f_{1(3)}\mathrm{ln}\rho _1}{𝑑\mathrm{\Omega }_1f_{1(3)}}.$$
(140)
Because of the two arbitrary constants $`s_{1,2}`$ the Hadamard partie finie is ambiguous, and one could think a priori that there is no point about defining a divergent integral by means of an ambiguous expression. Actually the point is that we control the origin of these constants: they come from the coefficients of $`1/r_{1,2}^3`$ in the expansions of $`F`$, which generate logarithmic terms in the integral. As we shall see the constants $`s_{1,2}`$ do not appear in the post-Newtonian metric up to the 2.5PN order (they are expected to appear only at 3PN order).
We can also give a meaning to the value of the function $`F`$ at the location of particle 1 for instance, by taking the average over all directions $`𝐧_1`$ of the coefficient of the zeroth power of $`r_1`$ in (73a), namely
$$(F)_1\frac{d\mathrm{\Omega }_1}{4\pi }f_{1(0)}.$$
(141)
We refer also to the definition (76) as the Hadamard partie finie (of the function $`F`$ at 1) because this definition is closely related to the definition (74) of the Hadamard partie finie of a divergent integral. To see this, apply (74) to the case where the function $`F`$ is actually a gradient, $`F=_iG`$, where $`G`$ satisfies (73) \[it is then clear that $`F`$ itself satisfies (73)\]. We find
$$\mathrm{Pf}d^3𝐱_iG=4\pi (n_1^ir_1^2G)_14\pi (n_2^ir_2^2G)_2$$
(142)
where in the right side the values at 1 and 2 are taken in the sense of the Hadamard partie finie (76). This nice connection between the Hadamard partie finie of a divergent integral and that of a singular function is clearly understood from applying the Gauss theorem on two surfaces $`r_{1,2}=s`$ surrounding the singularities (there is no dependence on the constants $`s_{1,2}`$).
### B Multipole moments of point-mass binaries
To compute the source moments (36) of two point-particles we insert (70) in place of the stress-energy tensor $`T^{\mu \nu }`$ of a continuous source, and we pick up the Hadamard partie finie \[in the sense of (74)\] of all integrals. This ansatz reads
$`(I_L)_{\mathrm{point}\mathrm{mass}}`$ $`=`$ $`\mathrm{Pf}\left\{I_L[T_{\mathrm{point}\mathrm{mass}}^{\mu \nu }]\right\};`$ (144)
$`(J_L)_{\mathrm{point}\mathrm{mass}}`$ $`=`$ $`\mathrm{Pf}\left\{J_L[T_{\mathrm{point}\mathrm{mass}}^{\mu \nu }]\right\}.`$ (145)
As we have seen in (69), the source multipole moments involve at high PN order many (non-compact-support) non-linear contributions which can be expressed in terms of retarded potentials such as $`V`$. The paradigm of such non-linear contributions is a term involving the quadratic product of two (derivatives of) potentials $`V`$, say $`VV`$, or, neglecting $`O(1/c^2)`$ corrections, $`UU`$. To Newtonian order $`U`$ is given by $`Gm_1/r_1+Gm_2/r_2`$ and it is easily checked that this paradigmatic term can be written as a certain derivative operator, say $``$, acting on the elementary integral (assuming for simplicity $`l=2`$)
$$Y_{ij}(𝐲_1,𝐲_2)\frac{1}{2\pi }\mathrm{FP}_{B=0}d^3𝐱|𝐱/r_0|^B\frac{\widehat{x}_{ij}}{r_1r_2}.$$
(146)
We see that the integral would be divergent at infinity without the finite part operation. However, it is perfectly well-behaved near 1 and 2 where there is no need of a regularization. The integral (79) can be evaluated in various ways; the net result is
$$Y_{ij}=\frac{r_{12}}{3}\left[y_1^{<ij>}+y_1^{<i}y_2^{j>}+y_2^{<ij>}\right]$$
(147)
where $`<ij>\mathrm{STF}(ij)`$. Starting at 3PN order we meet some elementary integrals which need the regularization at 1 or 2 in addition to involving the finite part at infinity. An example is
$$Z_{ij}(𝐲_1)\frac{1}{2\pi }\mathrm{Pf}\left\{\mathrm{FP}_{B=0}d^3𝐱|𝐱/r_0|^B\frac{\widehat{x}_{ij}}{r_1^3}\right\}.$$
(148)
To obtain this integral one splits it into a near-zone integral extending over the domain $`r_1<_1`$ (say), and a far-zone integral extending over $`_1<r_1`$. The Hadamard regularization at 1 applies only to the near-zone integral, while the finite part at $`B=0`$ is needed only for the far-zone integral. The result, found to be independent of the radius $`_1`$, reads
$$Z_{ij}=\left[2\mathrm{ln}\left(\frac{s_1}{r_0}\right)+\frac{16}{15}\right]y_1^{<ij>}.$$
(149)
In this case we find an explicit dependence on both the constants $`r_0`$ due to the finite part at infinity, and $`s_1`$ due to the Hadamard partie finie near 1 \[see (74)\]. However these constants do not enter the multipole moments before the 3PN order (collaboration with Iyer and Joguet ).
A long computation, done in , yields the mass-type quadrupole moment at the 2PN order fully reduced in the case of two point-masses moving on a circular orbit. The method is to start from (69) (issued from ) and to employ notably the elementary integral (79)-(80) (see also for the treatment of a cubically non-linear term). An equivalent result has been obtained by Will and Wiseman using their formalism . In a mass-centered frame the moment is of the form
$$I_{ij}=\mu \left(A\widehat{y}_{ij}+B\frac{\widehat{v}_{ij}}{\omega ^2}\right)+O\left(\frac{1}{c^5}\right),$$
(150)
where $`y_i=y_1^iy_2^i`$ and $`v_i=v_1^iv_2^i`$, where $`\omega `$ denotes the binary’s Newtonian orbital frequency \[$`\omega ^2=Gm/r_{12}^3`$ with $`m=m_1+m_2`$\], and where $`\mu =m_1m_2/m`$ is the reduced mass. The point is to obtain the coefficients $`A`$ and $`B`$ developed to 2PN order in terms of the post-Newtonian parameter $`\gamma =Gm/r_{12}c^2`$, where we recall that $`r_{12}`$ is the distance between the two particles in harmonic coordinates. Untill 2PN we find some definite polynomials in the mass ratio $`\nu =\mu /m`$ (such that $`0<\nu 1/4`$):
$`A`$ $`=`$ $`1+\gamma \left[{\displaystyle \frac{1}{42}}{\displaystyle \frac{13}{14}}\nu \right]+\gamma ^2\left[{\displaystyle \frac{461}{1512}}{\displaystyle \frac{18395}{1512}}\nu {\displaystyle \frac{241}{1512}}\nu ^2\right],`$ (152)
$`B`$ $`=`$ $`\gamma \left[{\displaystyle \frac{11}{21}}{\displaystyle \frac{11}{7}}\nu \right]+\gamma ^2\left[{\displaystyle \frac{1607}{378}}{\displaystyle \frac{1681}{378}}\nu +{\displaystyle \frac{229}{378}}\nu ^2\right].`$ (153)
The 2PN mass quadrupole moment (83)-(84) is part of a program aiming at computing the orbital phase evolution of inspiralling compact binaries to high post-Newtonian order (see Section 7.4). First-order black-hole perturbations, valid in the test-mass limit $`\nu 0`$ for one body, have already achieved the very high 5.5PN order . Recovering the result of black-hole perturbations in this limit constitutes an important check of the overall formalism. For the moment it passed the check to 2.5PN order ; this is quite satisfactory regarding the many differences between the present approach and the black-hole perturbation method.
### C Equations of motion of compact binaries
The equations of motion of two point-masses play a crucial role in accounting for the observed dynamics of the binary pulsar PSR1913+16 , and constitute an important part of the program concerning inspiralling compact binaries. The motivation for investigating rigorously the equations of motion came in part from the salubrious criticizing remarks of Jürgen Ehlers et al . Four different approaches have succeeded in obtaining the equations of motion of point-mass binaries complete up to the 2.5PN order (dominant order of radiation reaction): the “post-Minkowskian” approach of Damour, Deruelle and colleagues ; the “Hamiltonian” approach of Schäfer and predecessors ; the “extended-body” approach of Kopejkin et al ; and the “post-Newtonian” approach of Blanchet, Faye and Ponsot . The four approaches yield mutually agreeing results.
The post-Newtonian approach consists of (i) inserting the point-mass stress-energy tensor (70) into the 2.5PN metric in harmonic coordinates given by (63); (ii) curing systematically the self-field divergences of point-masses using the Hadamard regularization; and (iii) substituting the regularized metric into the standard geodesic equations. For convenience we write the geodesic equation of the particle 1 in the Newtonian-like form
$$\frac{d𝒫_1^i}{dt}=_1^i$$
(154)
where the (specific) linear momentum $`𝒫_1^i`$ and force $`_1^i`$ are given by
$$𝒫_1^i=c\left(\frac{v_1^\mu g_{i\mu }}{\sqrt{g_{\rho \sigma }v_1^\rho v_1^\sigma }}\right)_1;_1^i=\frac{c}{2}\left(\frac{v_1^\mu v_1^\nu _ig_{\mu \nu }}{\sqrt{g_{\rho \sigma }v_1^\rho v_1^\sigma }}\right)_1.$$
(155)
Crucial in this method, the quantities are evaluated at the location of particle 1 according to the rule (76). All the potentials (62) and their gradients are evaluated in a way similar to our computation of $`U`$ in (72), and then inserted into (85)-(86). We “order-reduce” the result, i.e. we replace each acceleration, consistently with the approximation, by its equivalent in terms of the positions and velocities as given by the (lower-order) equations of motion. After simplication we find, in agreement with other methods,
$`{\displaystyle \frac{dv_1^i}{dt}}=`$ $``$ $`{\displaystyle \frac{Gm_2}{r_{12}^2}}n_{12}^i+{\displaystyle \frac{Gm_2}{r_{12}^2c^2}}\{v_{12}^i[4(n_{12}v_1)3(n_{12}v_2)]`$ (156)
$`+`$ $`n_{12}^i[v_1^22v_2^2+4(v_1v_2)+{\displaystyle \frac{3}{2}}(n_{12}v_2)^2+5{\displaystyle \frac{Gm_1}{r_{12}}}+4{\displaystyle \frac{Gm_2}{r_{12}}}]\}`$ (157)
$`+`$ $`{\displaystyle \frac{Gm_2}{r_{12}^2c^4}}n_{12}^i\{[2v_2^4+4v_2^2(v_1v_2)2(v_1v_2)^2`$ (158)
$`+`$ $`{\displaystyle \frac{3}{2}}v_1^2(n_{12}v_2)^2+{\displaystyle \frac{9}{2}}v_2^2(n_{12}v_2)^26(v_1v_2)(n_{12}v_2)^2{\displaystyle \frac{15}{8}}(n_{12}v_2)^4]`$ (159)
$`+`$ $`{\displaystyle \frac{Gm_1}{r_{12}}}[{\displaystyle \frac{15}{4}}v_1^2+{\displaystyle \frac{5}{4}}v_2^2{\displaystyle \frac{5}{2}}(v_1v_2)`$ (160)
$`+`$ $`{\displaystyle \frac{39}{2}}(n_{12}v_1)^239(n_{12}v_1)(n_{12}v_2)+{\displaystyle \frac{17}{2}}(n_{12}v_2)^2]`$ (161)
$`+`$ $`{\displaystyle \frac{Gm_2}{r_{12}}}\left[4v_2^28(v_1v_2)+2(n_{12}v_1)^24(n_{12}v_1)(n_{12}v_2)6(n_{12}v_2)^2\right]`$ (162)
$`+`$ $`{\displaystyle \frac{G^2}{r_{12}^2}}[{\displaystyle \frac{57}{4}}m_1^29m_2^2{\displaystyle \frac{69}{2}}m_1m_2]\}`$ (163)
$`+`$ $`{\displaystyle \frac{Gm_2}{r_{12}^2c^4}}v_{12}^i\{v_1^2(n_{12}v_2)+4v_2^2(n_{12}v_1)5v_2^2(n_{12}v_2)4(v_1v_2)(n_{12}v_1)`$ (164)
$`+`$ $`4(v_1v_2)(n_{12}v_2)6(n_{12}v_1)(n_{12}v_2)^2+{\displaystyle \frac{9}{2}}(n_{12}v_2)^3`$ (165)
$`+`$ $`{\displaystyle \frac{Gm_1}{r_{12}}}[{\displaystyle \frac{63}{4}}(n_{12}v_1)+{\displaystyle \frac{55}{4}}(n_{12}v_2)]+{\displaystyle \frac{Gm_2}{r_{12}}}[2(n_{12}v_1)2(n_{12}v_2)]\}`$ (166)
$`+`$ $`{\displaystyle \frac{4G^2m_1m_2}{5c^5r_{12}^3}}\{n_{12}^i(n_{12}v_{12})[6{\displaystyle \frac{Gm_1}{r_{12}}}+{\displaystyle \frac{52}{3}}{\displaystyle \frac{Gm_2}{r_{12}}}+3v_{12}^2]`$ (167)
$`+`$ $`v_{12}^i[2{\displaystyle \frac{Gm_1}{r_{12}}}8{\displaystyle \frac{Gm_2}{r_{12}}}v_{12}^2]\}+O\left({\displaystyle \frac{1}{c^6}}\right),`$ (168)
\[where $`n_{12}^i=(y_1^iy_2^i)/r_{12}`$; $`v_{12}^i=v_1^iv_2^i`$; and e.g. $`(n_{12}v_1)`$ denotes the Euclidean scalar product\]. At the 1PN or $`1/c^2`$ level the equations were obtained before by Lorentz an Droste , and by Einstein, Infeld and Hoffmann . The 2.5PN or $`1/c^5`$ term represents the radiation damping in harmonic coordinates \[correct because the metric (63) we started with matches to the post-Minkowskian exterior field\]. In the case of circular orbits, the equations simplify drastically:
$$\frac{dv_{12}^i}{dt}=\omega _{2\mathrm{P}\mathrm{N}}^2y_{12}^i\frac{32G^3m^3\nu }{5c^5r_{12}^4}v_{12}^i+O\left(\frac{1}{c^6}\right),$$
(169)
where the orbital frequency $`\omega _{2\mathrm{P}\mathrm{N}}`$ of the 2PN circular motion reads
$$\omega _{2\mathrm{P}\mathrm{N}}^2=\frac{Gm}{r_{12}^3}\left[1+(3+\nu )\gamma +\left(6+\frac{41}{4}\nu +\nu ^2\right)\gamma ^2\right]$$
(170)
(the post-Newtonian parameter is $`\gamma =Gm/c^2r_{12}`$; and $`\nu =\mu /m`$).
### D Gravitational waveforms of inspiralling compact binaries
The gravitational radiation field and associated energy flux are given by (52) and (57) in terms of time-derivatives of the radiative multipole moments, themselves related to the source multipole moments by formulas such as (56). Furthermore, at a given post-Newtonian order, the source moments admit some explicit though complicated expressions such as (68)-(69), which, when specialized to (non-spinning) point-mass circular binaries, yield e.g. (83)-(84).
Now, for insertion into the radiation field and energy flux, one must compute the time-derivatives of the binary moments, with appropriate order-reduction using the binary’s equations of motion (87)-(89). This yields in particular the fully reduced (up to the prescribed post-Newtonian order) gravitational waveform of the binary, or more precisely the two independent “plus” and “cross” polarization states $`h_+`$ and $`h_\times `$. The result to 2PN order is written in the form
$$h_{+,\times }=\frac{2Gm\nu x}{c^2R}\left\{H_{+,\times }^{(0)}+x^{1/2}H_{+,\times }^{(1/2)}+xH_{+,\times }^{(1)}+x^{3/2}H_{+,\times }^{(3/2)}+x^2H_{+,\times }^{(2)}\right\},$$
(171)
where, for convenience, we have introduced a post-Newtonian parameter which is directly related to the orbital frequency: $`x=(Gm\omega _{2\mathrm{P}\mathrm{N}}/c^3)^{2/3}`$, where $`\omega _{2\mathrm{P}\mathrm{N}}`$ is given for circular orbits by (89). The various post-Newtonian coefficients in (90) depend on the cosine and sine of the “inclination” angle between the detector’s direction and the normal to the orbital plane ($`c_i=\mathrm{cos}i`$ and $`s_i=\mathrm{sin}i`$), and on the masses through the ratios $`\nu =\mu /m`$ and $`\delta m/m`$, where $`\delta m=m_1m_2`$. The result for the “plus” polarization (collaboration with Iyer, Will and Wiseman ) is
$`H_+^{(0)}`$ $`=`$ $`(1+c_i^2)\mathrm{cos}2\psi ,`$ (173)
$`H_+^{(1/2)}`$ $`=`$ $`{\displaystyle \frac{s_i}{8}}{\displaystyle \frac{\delta m}{m}}\left[(5+c_i^2)\mathrm{cos}\psi 9(1+c_i^2)\mathrm{cos}3\psi \right],`$ (174)
$`H_+^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left[19+9c_i^22c_i^4\nu (1911c_i^26c_i^4)\right]\mathrm{cos}2\psi `$ (176)
$`{\displaystyle \frac{4}{3}}s_i^2(1+c_i^2)(13\nu )\mathrm{cos}4\psi ,`$
$`H_+^{(3/2)}`$ $`=`$ $`{\displaystyle \frac{s_i}{192}}{\displaystyle \frac{\delta m}{m}}\{[57+60c_i^2c_i^42\nu (4912c_i^2c_i^4)]\mathrm{cos}\psi `$ (179)
$`{\displaystyle \frac{27}{2}}\left[73+40c_i^29c_i^42\nu (258c_i^29c_i^4)\right]\mathrm{cos}3\psi `$
$`+{\displaystyle \frac{625}{2}}(12\nu )s_i^2(1+c_i^2)\mathrm{cos}5\psi \}2\pi (1+c_i^2)\mathrm{cos}2\psi ,`$
$`H_+^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{120}}[22+396c_i^2+145c_i^45c_i^6+{\displaystyle \frac{5}{3}}\nu (706216c_i^2251c_i^4+15c_i^6)`$ (186)
$`5\nu ^2(98108c_i^2+7c_i^4+5c_i^6)]\mathrm{cos}2\psi `$
$`+{\displaystyle \frac{2}{15}}s_i^2[59+35c_i^28c_i^4{\displaystyle \frac{5}{3}}\nu (131+59c_i^224c_i^4)`$
$`+5\nu ^2(213c_i^28c_i^4)]\mathrm{cos}4\psi `$
$`{\displaystyle \frac{81}{40}}(15\nu +5\nu ^2)s_i^4(1+c_i^2)\mathrm{cos}6\psi `$
$`+{\displaystyle \frac{s_i}{40}}{\displaystyle \frac{\delta m}{m}}\{[11+7c_i^2+10(5+c_i^2)\mathrm{ln}2]\mathrm{sin}\psi 5\pi (5+c_i^2)\mathrm{cos}\psi `$
$`27[710\mathrm{ln}(3/2)](1+c_i^2)\mathrm{sin}3\psi +135\pi (1+c_i^2)\mathrm{cos}3\psi \}.`$
The “cross” polarization admits a similar expression (see ). Here, $`\psi `$ denotes a particular phase variable, related to the actual binary’s orbital phase $`\varphi `$ and frequency $`\omega \omega _{2PN}`$ by
$$\psi =\varphi \frac{2Gm\omega }{c^3}\mathrm{ln}\left(\frac{\omega }{\omega _0}\right);$$
(188)
$`\varphi `$ is the angle, oriented in the sense of the motion, between the vector separation of the two bodies and a fixed direction in the orbital plane (since the bodies are not spinning, the orbital motion takes place in a plane). In (92), $`\omega _0`$ denotes some constant frequency, for instance the orbital frequency when the signal enters the detector’s frequency bandwidth; see for discussion.
The previous formulas give the waveform of point-mass binaries whenever the frequency and phase of the orbital motion take the values $`\omega `$ and $`\varphi `$. To get the waveform as a function of time, we must replace $`\omega `$ and $`\varphi `$ by their explicit time evolutions $`\omega (t)`$ and $`\varphi (t)`$. Actually, the frequency is the time-derivative of the phase: $`\omega =d\varphi /dt`$. The evolution of the phase is entirely determined, for circular orbits, by the energy balance equation $`dE/dt=`$ relating the binding energy $`E`$ of the binary in the center of mass to the emitted energy flux $``$. $`E`$ is computed using the equations of motion (87), and $``$ follows from (57) and application of the previous formalism \[changing the radiative moments to the source moments, applying (83)-(84), etc…\]; the net result for the 2.5PN orbital phase is
$`\varphi `$ $`=`$ $`\varphi _0{\displaystyle \frac{1}{\nu }}\{\mathrm{\Theta }^{5/8}+({\displaystyle \frac{3715}{8064}}+{\displaystyle \frac{55}{96}}\nu )\mathrm{\Theta }^{3/8}{\displaystyle \frac{3}{4}}\pi \mathrm{\Theta }^{1/4}`$ (191)
$`+\left({\displaystyle \frac{9275495}{14450688}}+{\displaystyle \frac{284875}{258048}}\nu +{\displaystyle \frac{1855}{2048}}\nu ^2\right)\mathrm{\Theta }^{1/8}`$
$`+({\displaystyle \frac{38645}{172032}}{\displaystyle \frac{15}{2048}}\nu )\pi \mathrm{ln}\mathrm{\Theta }\},`$
where $`\varphi _0`$ is a constant phase (determined for instance when the frequency is $`\omega _0`$), and $`\mathrm{\Theta }`$ the convenient dimensionless time variable
$$\mathrm{\Theta }=\frac{c^3\nu }{5Gm}(t_ct),$$
(192)
$`t_c`$ being the instant of coalescence at which, formally, $`\omega (t)`$ tends to infinity (of course, the post-Newtonian method breaks down before the final coalescence). All the results are in agreement, in the limit $`\nu 0`$, with those of black-hole perturbation theory .
## VIII Conclusion
The formalism reviewed in this article permits investigating in principle all aspects of the problem of dynamics and gravitational-wave emission of a slowly-moving isolated system (with, say, $`v/c0.3`$ at most): the generation of waves, their propagation in vacuum, the back-reaction onto the system, the structure of the asymptotic field, and most importantly the relation between the far-field and the source parameters. Of course, the formalism is merely post-Newtonian and never “exact”, but in applications to astrophysical objects such as inspiralling compact binaries this should be sufficient provided that the post-Newtonian approximation is carried to high order.
Furthermore, there are several places in the formalism where some results are valid formally to any order of approximation. For instance, the source multipole moments are related to the infinite formal post-Newtonian expansion of the pseudo-tensor \[see (18) or (36)\], and the post-Minkowskian iteration of the exterior field is performed to any non-linear order \[see (43)\]. In such a situation, where an infinite approximate series can be defined, there is the interesting question of its relation to a corresponding element in the exact theory. For the moment the only solid work concerns the post-Minkowskian approximation of the exterior vacuum field, which has been proved to be asymptotic . Likewise it is plausible that the expressions of the source multipole moments could be valid in the case of exact solutions.
The most important part of the formalism where a general prescription for how to proceed at any approximate step is missing, is the post-Newtonian expansion for the field inside the isolated system. For instance, though the multipole moments are given in terms of the formal post-Newtonian expansion of the pseudo-tensor, no general algorithm for computing explicitly this post-Newtonian expansion is known. An interesting task would be to define such an algorithm, in a manner similar to the post-Minkowskian algorithm in Section 4. In the author’s opinion, the post-Newtonian algorithm should be defined conjointly with the post-Minkowskian algorithm, and should rely on the matching equation (16), so as to convey into the post-Newtonian field the information about the exterior metric.
Note that even if a general method for implementing a complete approximation series is defined, this method may be unworkable in practical calculations, because not explicit enough. For instance the post-Minkowskian series (43) is defined in terms of “iterated” retarded integrals, but needs to be suplemented by some formulas, to be used in applications, for the retarded integral of a multipolar extended source. In this respect it would be desirable to develop the formulas generalizing (50)-(51) to any non-linear order. This should permit in particular the study of the general structure of tails, tails of tails, and so on.
For the moment the only application of the formalism concerns the radiation and motion of point-particle binaries. Of course it is important to keep the formalism as general as possible, and not to restrict oneself to a particular type of source, but this application to point-particles offers some interesting questions. Indeed, it seems that the post-Newtonian approximation used conjointly with a regularization à la Hadamard works well, and that one is getting closer and closer to an exact (numerical) solution corresponding to the dynamics and radiation of two black-holes. So, in which sense does the post-Newtonian solution (corresponding to point-masses without horizons) approach a true solution for black-holes? Does the adopted method of regularizing the self-field play a crucial role? Is it possible to define a regularization consistently with the post-Newtonian approximation to all orders?
## Acknowledgments
The author is especially grateful to Bernd Schmidt for discussions and for remarks which led to improvement of this article. Stimulating discussions with Piotr Crusciel, Thibault Damour and Gerhard Schäfer are also acknowledged. The Max-Planck-Institut für Gravitationsphysik (Albert-Einstein-Institut) in Potsdam is thanked for an invitation during which the writing of the article was begun.
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# A supermassive scalar star at the Galactic Center?
## 1 Introduction
During the last years, the possible existence of a single large mass in the Galactic Center has been favored as the upper bound on its size tightens, and stability criteria rule out complex clusters. Although it is commonly believed that this central mass is a supermassive black hole, it is not yet established, as we discuss below, on a firm observational basis.
The aim of this paper is to present an alternative model for the supermassive dark object in the center of our Galaxy, formed by self-gravitating non-baryonic matter composed by bosons. This kind of objects, so-called boson stars, are well known to physicists, but up to now, observational astrophysical consequences have hardly been explored. The main characteristics of this model are:
1. It is highly relativistic, with a size comparable to (but slightly larger than) the Schwarzschild radius of a black hole of equal mass.
2. It has neither an event horizon nor a singularity, and after a physical radius is reached, the mass distribution exponentially decreases.
3. The particles that form the object interact between each other only gravitationally, in such a way that there is no solid surface to which falling particles can collide.
It is the purpose of this work to show that these features are able to produce a Galactic Center model which can be confronted with all known observational constraints, and to point the way in which such a center could be differentiated from a usual supermassive black hole.
### 1.1 Why scalar fields?
Interesting models for dark matter use weakly interacting bosons, and primordial nucleosynthesis show that most of the mass in the universe should be non-baryonic if $`\mathrm{\Omega }1`$. Most models of inflation make use of scalar fields. Scalar-tensor gravitation is the most interesting alternative to general relativity. Recent results from supernovae, which in principle were thought of to favor a cosmological constant (Perlmutter et al., 1999), can as well be supported by a variety of models, some of them with scalar fields too (Célérier, 1999). Particle physicist expect to detect the scalar Higgs particle in the next generation of accelerators. Scalar dilatons appear in low energy unified theories, where the tensor field $`g_{\mu \nu }`$ of gravity is accompanied by one or several scalar fields, and in string effective super-gravity. The axion is a scalar with a long history as a dark matter candidate, and Goldstone bosons have also already inferred masses. Symmetry arguments, which once led to the concept of neutron stars, may force to ask whether there could be stellar structures made up of bosons instead of fermions.
In recent works, Schunck and Liddle (1997), Schunck and Torres (2000), and Capozziello, Lambiase and Torres (2000) analyzed some of the observational properties of boson stars, and found them notoriously similar to black holes. In particular, the gravitational redshift of the radiation emitted within a boson star potential, and the rotational curves of accreted particles, were studied to assess possible boson star detection. The interest in observational properties of boson stars also led to investigate them as a possible lens in a gravitational lensing configuration (Da̧browski and Schunck, 2000). Recent studies are putting forward the gravitational lensing phenomenon in strong field regimes (Virbhadra et al., 1998; Virbhadra and Ellis, 2000; Torres et al., 1998a, b). This would be the case for boson stars, which are genuine relativistic objects.
This work is organized as follows. Sec. II is a brief summary of the main observational results concerning Sgr A, and the main hypotheses in order to explain it are outlined there. Sec. III analyzes what dynamical observational data show, and what kind of models can support them. Sec. IV gives the basic ingredients to theoretically construct scalar stars, shows mass and radius estimates, and study effective potentials and orbits of particles. In Sec. V we study the center of the Galaxy, show which are the stable scalar star models able to fit such a huge mass, and comment on the possible existence of boson particles with the required features. We also provide there an assessment of the disruption processes in boson star scenarios. Discussion and conclusions are given in Sections VI and VII.
## 2 The Galactic Center
### 2.1 Main observational facts
The Galactic Center is a very active region toward the Sagittarius constellation where, at least, six very energetic radio sources are present (Sgr A, A, B1, B2, C, and D). Furthermore, there are several supernova remnants, filaments, and very reach star clusters. Several observational campaigns (Genzel et al., 1994) have identified the exact center with the supermassive compact dark object in Sgr A, an extremely loud radio source. The mass and the size of the object has been established to be $`(2.61\pm 0.76)\times 10^6M_{}`$ concentrated within a radius of 0.016 pc (about 30 lds)(Ghez, 1998; Genzel et al., 1996).
More precisely, Ghez et al. (1998) have made extremely accurate velocity measurements in the central square arcsecs. From this bulk of data, it is possible to state that a supermassive compact dark object is present at the Galactic Center. It is revealed by the motion of stars moving within projected distances around 0.01 pc from the radio source Sgr A, at projected velocities in excess of 1000 km s<sup>-1</sup>. In other words, a high increase in the velocity dispersion of the stars toward the dynamical center is revealed. Furthermore, a large and coherent counter–rotation, especially of the early–type stars, is shown, supporting their origin in a well–defined epoch of star formation. Observations of stellar winds nearby Sgr A give a mass accretion rate of $`dM/dt=6\times 10^6M_{}`$ yr<sup>-1</sup> (Genzel et al., 1996). Hence, the dark mass must have a density $`10^9M_{}`$ pc<sup>-3</sup> or greater, and a mass–to–luminosity ratio of at least $`100M_{}/L_{}`$. The bottom line is that the central dark mass seems to be a single object, and that it is statistically very significant $`(68\sigma )`$.
### 2.2 Domain of the black hole?
Such a large density contrast excludes that the dark mass could be a cluster of almost $`2\times 10^6`$ neutron stars or white dwarfs. Detailed calculations of evaporation and collision mechanisms give maximal lifetimes of the order of $`10^8`$ years, much shorter than the estimated age of the Galaxy (Sanders, 1992; Maoz, 1995).
As a first conclusion, several authors state that in the Galactic Center there is either a single supermassive black hole or a very compact cluster of stellar size black holes (Genzel et al., 1996). We shall come back to this paradigm through the rest of the paper (particularly in Section III).
Dynamical evidence for central dark objects has been published for 17 galaxies, like M87 (Ford et al., 1994; Macchetto et al., 1997), or NGC4258 (Greenhill et al., 1995), but proofs that they really are black holes requires measurement of relativistic velocities near the Schwarzschild radius, $`r_s2M_{}/(10^8M_{})`$ AU (Kormendy and Richstone, 1995), or seeing the properties of the accretion disk. Due to the above mentioned mass accretion rate, if Sgr A is a supermassive black hole, its luminosity should be more than $`10^{40}`$erg s<sup>-1</sup>, provided the radiative efficiency is about 10%. On the contrary, observations give a bolometric luminosity less than $`10^{37}`$ erg s<sup>-1</sup>, already taking into account the luminosity extinction due to interstellar gas and dust. This discrepancy is the so–called “blackness problem” which has led to the notion of a “black hole on starvation”. Standard dynamics and thermodynamics of the spherical accretion onto a black hole must be modified in order to obtain successful models (Falcke and Heinrich, 1994). Recent observations, we recall, probe the gravitational potential at a radius larger than $`4\times 10^4`$ Schwarzschild radii of a black hole of mass $`2.6\times 10^6M_{}`$ (Ghez, 1998).
### 2.3 Domain of the massive neutrino?
An alternative model for the supermassive compact object in the center of our Galaxy has been recently proposed by Tsiklauri and Viollier (1998). The main ingredient of the proposal is that the dark matter at the center of the galaxy is non-baryonic (but in any case fermions, e.g. massive neutrinos or gravitinos), interacting gravitationally to form supermassive balls in which the degeneracy pressure balances their self–gravity. Such neutrino balls could have formed in early epochs, during a first–order gravitational phase transition and their dynamics could be reconciled, with some adjustments, to the Standard Model of Cosmology.
Several experiments are today running to search for neutrino oscillations. LSND (White, 1998) found evidence for oscillations in the $`\nu _e\nu _\mu `$ channel for pion decay at rest and in flight. On the contrary KARMEN (Zeitnitz, 1998) seems to be in contradiction with LSND evidence. CHORUS and NOMAD at CERN are just finishing the phase of 1994–1995 data analysis. It is very likely that exact predictions for and $`\nu _\mu \nu _\tau `$ oscillations will be available in next years. Thanks to the fact that it is possible to give correct values for the masses to the quarks up, charm, and top, it is possible to infer reasonable values of mass for $`\nu _e`$, $`\nu _\mu `$, and $`\nu _\tau `$. In order to explain the characteristics of Sgr A, we need fermions whose masses range between 10 and 25 keV which, cosmologically, fall into the category of warm dark matter. It is interesting to note that a good estimated value for the mass of $`\tau `$-neutrino is
$$m_{\nu _\tau }=m_{\nu _\mu }\left(\frac{m_t}{m_c}\right)^214.4\mathrm{keV}.$$
(1)
Choosing fermions like neutrinos or gravitinos in this mass range allows for the formation of supermassive degenerate objects from $`10^6M_{}`$ to $`10^9M_{}`$.
The theory of heavy neutrino condensates, bound by gravity, can be easily sketched considering a Thomas–Fermi model for fermions (Viollier et al., 1992). We can set the Fermi energy equal to the gravitational potential which binds the system, and see that the number density is a function of the gravitational potential. Such a gravitational potential will obey a Poisson equation where neutrinos (and anti-neutrinos) are the source term. This equation is valid everywhere except at the origin. By a little algebra, the Poisson equation reduces itself to the radial Lané–Emden differential equation with polytropic index $`n=3/2`$, which is equivalent to the Thomas–Fermi differential equation of atomic physics, except for a minus sign that is due to the gravitational attraction of the neutrinos. The Newtonian potential, which is the solution of such an equation, is $`\mathrm{\Phi }(r)r^4`$. Considering a standard accretion disk, if Sgr A is a neutrino star with radius $`R=30.3`$ lds ($`10^5`$ Schwarzschild radii), mass $`M_{GC}=2.6\times 10^6M_{}`$, and luminosity $`L_{GC}10^{37}`$erg sec<sup>-1</sup>, it should consist of neutrinos with masses $`m12.0`$ keV for $`g_\nu =4`$, or $`m14.3`$ keV for $`g_\nu =2`$ (Bilić, 1998; Viollier et al., 1992). It is specially appealing that with the same mass for neutrinos, several galactic centers can be modeled. For instance, similar results hold also for the dark object ($`M3\times 10^9M_{}`$) inside the center of M87. Due to the Thomas–Fermi theory, the model fails at the origin: we have to consider the effect of the surrounding baryonic matter which, in some sense, has to stabilize the neutrino condensate. In fact, the solution $`\mathrm{\Phi }(r)r^4`$ is clearly unbounded from below.
The model proceeds assuming a thermodynamical phase where a constant neutrino number density can be taken into consideration. This is quite natural for a Fermi gas at temperature $`T=0`$. The Poisson equation can be recast in a Lané–Emden form with polytropic index $`n=0`$. The solution of such an equation is $`\mathrm{\Phi }(r)r^2,`$ which is clearly bounded (Capozziello and Iovane, 1999). For $`T0`$, we get a solution of the form $`\mathrm{\Phi }(r)r^4`$. Matching these two results it is possible to confine the neutrino ball. On the other hand, a similar result is recovered using the Newton Theorem for a spherically symmetric distribution of matter of radius $`R`$ (Binney and Tremaine, 1987). In that case, the potential goes quadratic inside the sphere while it goes as $`\mathrm{\Phi }(r)r^1`$ matching on the boundary. In our case, the situation is similar assuming the matching with a steeper potential.
If such a neutrino condensate exists in the center of Galaxy, it could act as a spherical thick lens (a magnifying glass) for the stars behind it, so that their apparent velocities will be larger than in reality. In other words, depending on the line of sight, it should be possible to correct the projected velocities by a gravitational lensing contribution, so trying to explain the bimodal distribution (early and late type stars) actually observed (Ghez, 1998; Genzel et al., 1996). Since the astrophysical features of the object in Sgr A are quite well known, accurate observations by lensing could contribute to the exact determination of particle constituents. A detailed model and comparison with the data was presented by some of us (Capozziello and Iovane, 1999).
## 3 Observational status
### 3.1 Brief review on dynamical data and what are they implying
An important information on the central objects of galaxies (particularly in active galactic nuclei) is the short timescale of variability. This has the significance of putting an upper bound –known as causality constraint– on the size of the emitting region: If a system of size $`L`$ suddenly increases its emissivity at all points, the temporal width with which we receive it is $`L/c`$ and thus, a source can not fluctuate in a way that involves its entire volume in timescales shorter than this (unless $`c`$ is not a limiting velocity). This finally yields a corresponding maximum lenghscale, typically between $`10^4`$ to 10 pc (Krolik, 1999). Autocorrelation in the emission argue against the existence of a cluster of objects (unless only one member dominates the emissivity), and the small region in which the cluster should be allows two body encounters to be very common and produces the lost of the stability of complex clusters. These facts can be used to conclude that a single massive object must be in the center of most galaxies. What observations show in this case is that the size of the emitting regions are very small. This does not directly implicate black holes as such.
The way in which we expect to detect the influence of a very massive object is through its gravity. The “sphere of influence” of a large mass is defined by the distance at which its potential significantly affects the orbital motions of stars and gas, and is given by
$$R_{}=GM/\sigma _{}^24M_7\sigma _{,100}^2\mathrm{pc},$$
(2)
where the central mass is normalized to $`10^7M_{}`$ and $`\sigma _{,100}`$ is a rms orbital speed in units of 100 km s<sup>-1</sup>. Thus, even very large masses have a small sphere of influence. Within this sphere, the expected response to a large nuclear mass can be divided in two groups. Firstly, the response of interstellar gas can hardly be anything different from an isotropic motion in the center of mass frame. Random speeds are often thought to be much smaller than this overall speed, which far from the sphere of influence is just given by $`(GM/r)^{1/2}`$. If the energy produced by random motions is radiated away, the gas behaves as a whole and tend to flatten itself, conserving the angular momentum. Observations of the rotational velocity of gas as a function of the radius can thus provide a measure of the total mass at the center. Then, what observation searches is the existence of a Keplerian potential signature in the flattened gas velocity distribution. Examples of this are the radio galaxies M87 (Ford et al., 1994), NGC 4258 (Miyoshi et al., 1995), and many others. Then, any massive object producing a velocity distribution with a Keplerian decrease would be allowed by observation.
Secondly, we consider the response of the stars. In this case, peculiar velocities are larger or comparable to the bulk streaming and it is harder to actually differentiate the stellar response to a nuclear mass from the observable properties of a pure stellar potential. However, stellar motions, contrary to that of gas, is not affected by other forces (as those produced by magnetic fields), and are better tracers of mass distributions. What observations show in this case are the stellar density and the velocity distribution. Use of the collisionless Boltzmann equation allows one to get Jean’s relation (Binney and Tremaine, 1987; Krolik, 1999):
$$\frac{GM(r)}{r}=V_{rot}^2\sigma _r^2\left[\frac{d\mathrm{ln}n(r)}{d\mathrm{ln}r}\frac{d\mathrm{ln}\sigma _r^2}{d\mathrm{ln}r}+2\frac{\sigma _t^2}{\sigma _r^2}\right]$$
(3)
here, $`n(r`$) is the spherically symmetric distribution of stars, $`V_{rot}`$ is the mean rotational speed, $`\sigma _{r,\theta ,\varphi }`$ are velocity dispersions and $`\sigma _t^2=\sigma _\theta ^2+\sigma _\varphi ^2`$. Measurements of $`V`$ and random velocities determine that $`M(r)`$ decline inwards until a critical radius, where $`M`$ becomes constant. What really happens, however, if that even if this region is not scrutinized, the mass to light ratio becomes sufficiently large to suggest the presence of a large dark mass. A combination of this with other measurements strongly suggest that this dark mass is a single object.
To summarize, usual measurements point to inform us that there exist a single supermassive object in the center of some galaxies but do not definitively state its nature. As Kormendy and Richstone (1995) have stated, the black hole scenario has become our paradigm. But suggestions that the dark objects are black holes are based only on indirect astrophysical arguments, and surprises are possible on the way to the center.
### 3.2 The Galaxy
We follow Genzel et al. (1996) and parameterize the stellar density distribution as,
$$n(r)=\frac{\mathrm{\Sigma }_0}{R_0}\frac{1}{1+(R/R_0)^\alpha }.$$
(4)
Note that $`R_0`$ is related to the core radius through $`R_{\mathrm{core}}=b(\alpha )R_0`$. Genzel et al. found that the best fit parameters for the observed stellar cluster are a central density of $`4\times 10^6M_{}`$pc<sup>-3</sup>, a core radius of 0.38 pc, and a value $`\alpha =1.8`$ ($`b(\alpha )=2.19`$). With this distribution, they found that a dark mass of about 2.5 $`\times 10^6M_{}`$ was needed to fit the observational data.
The cluster distribution, and the cluster plus a black hole constant mass is shown in Fig. A supermassive scalar star at the Galactic Center?. Black boxes represent Genzel et al. 1996 (Table 10) and Eckart and Genzel 1997 (Fig. 5) data. A dashed line stands for the stellar cluster contribution, while a dot-dashed line represent an enclosed point-like black hole mass. The solid line in the right half of the figure stands for the mass distribution both for a black hole and a boson star plus the stellar cluster. The mass dependence we are plotting for the boson star is obtained in Section IV, and represents the mass distribution of a mini-boson star ($`\mathrm{\Lambda }=0`$, $`m[\mathrm{GeV}]=2.81\times 10^{26}`$) with dimensionless central density $`\sigma (0)`$ equal to 0.1. Other boson star configurations, with appropriate choice of the boson field mass $`m`$, yield to the same results. Note the break of more than three orders of magnitude in the x-axis, this is caused because the boson star distribution, further out of the equivalent Schwarzschild radius, behaves as a black hole. It begins to differ from the black hole case at radius more than three orders of magnitude less than the innermost data point, that is why the x-axis has to have a break. From the mass distribution, a boson star in the center of the galaxy is virtually indistinguishable from a black hole.
Tsiklauri and Viollier (1998) have shown that the same observational data can also be fitted using an extended neutrino ball. In that case, differences begin to be noticed just around the innermost data point. It is then hard to determine whether the central object is a black hole, a neutrino ball, or a boson star based only on dynamical data now at hand, and other problems concerning the accretion disk have to considered (see below). Even harder is the situation for deciding –using only this kind of data– if the supermassive object is a boson star instead of a black hole: as a boson star is a relativistic object, the decay of the enclosed mass curve happens close the center. This, however, will provide an equivalent picture than a black hole for disruption and accretion processes; we shall comment on it in the next sections.
To apply Eq. (3) to the observational data we have to convert intrinsic velocity dispersions ($`\sigma _r(r)`$) and volume densities ($`n(r)`$) to projected ones (Binney and Tremaine, 1987), these are the ones we observe. We shall also consider that $`\delta =1\sigma _\theta ^2/\sigma _r^2`$, the anisotropy parameter, is equal to 0, and we are assuming implicitly that $`\sigma _\theta =\sigma _\varphi `$. We take into account the following Abel integrals,
$$\mathrm{\Sigma }(p)=2_p^{\mathrm{}}\frac{n(r)r}{\sqrt{r^2p^2}}𝑑r,$$
(5)
$$\mathrm{\Sigma }(p)\sigma _r(p)=2_p^{\mathrm{}}\frac{n(r)\sigma _r(r)^2r}{\sqrt{r^2p^2}}𝑑r.$$
(6)
$`\mathrm{\Sigma }(p)`$ denotes surface density, $`\sigma _r(p)`$ is the projected velocity dispersion, and $`p`$ is the projected distance. We adopt Genzel et al.’s (1996) and Tsiklauri and Viollier’s (1998) parameterization for $`\sigma _r(r)`$,
$$\sigma _r(r)=\sigma (\mathrm{})^2+\sigma (2^{\prime \prime })^2\left(\frac{R}{2^{\prime \prime }}\right)^{2\beta }.$$
(7)
What one usually does is to numerically integrate Eqs. (5,6) and fit the observational data ($`\sigma _r(p)vs.p`$). To do so one also has to assume a density dependence for the cluster (as in Eq. (4)), and the dark mass. With a point-like dark mass of 2.5 $`\times 10^6M_{}`$, the parameters of the fit result to be $`\sigma (\mathrm{})=55`$ km s<sup>-1</sup>, $`\sigma (2^{\prime \prime })=350`$ km s<sup>-1</sup>, and $`\beta =0.95`$.
If one now changes the central black hole for a neutrino ball, or a boson star, one has to consider the particular density dependence for these objects. In this way, $`n(r)=n_{\mathrm{stellar}\mathrm{cluster}}(r)+n_{\mathrm{dark}\mathrm{mass}}(r)`$. We obtained $`n_{\mathrm{boson}\mathrm{star}}(r)`$ as $`M(r)/(4/3\pi r^3)`$, where $`M(r)`$ is the fitted mass dependence of the boson star as explained in the next section. It is now useful to consider that the fitting of $`\sigma _r(p)`$ is made taking into account observational data points in regions where the boson star generated space-time is practically indistinguishable from a black hole. Then, we may expect that the actual parameters, $`\sigma (\mathrm{})`$, $`\sigma (2^{\prime \prime })`$, and $`\beta `$, will be very close to those obtained for the black hole. For our purposes, it is enough to take the same $`\sigma _r(r)`$ as in the black hole case, and compute $`\sigma _{r\mathrm{boson}\mathrm{star}}(p)`$ using Eqs. (5,6), with the adequate total $`n(r)`$.
In Fig. A supermassive scalar star at the Galactic Center? we show the observational data of Eckart and Genzel (1997) (filled black boxes) and Genzel et al. (1996) (hollow circles) super-imposed with the curve for $`\sigma _r(p)`$ that we obtain with a mini-boson star ($`\mathrm{\Lambda }=0`$, $`\sigma (0)=0.1`$, $`m[\mathrm{GeV}]=2.81\times 10^{26}`$) in the galactic center. Other boson star configurations, with appropriate choice of the boson field mass $`m`$, yield to the same results. For this configuration, we used, as will be explained in Section IV, a mass distribution given by a Boltzmann-like equation with $`A_2=2.5\times 10^6M_{}`$, $`A_1=0.237\times 10^6M_{}`$, $`R_0=1.19210^6`$ pc, and $`\mathrm{\Delta }R=4.16310^6`$ pc. In the range plotted, and in which data is available, the differences between boson and black hole theoretical curves is undetectable. They only begin to deviate from each other at $`p10^4`$ pc, well beyond the last observational data point. Even the deviation in such a region is as slight as 1 km s<sup>-1</sup>, and it only becomes more pronounced when $`p`$ values are closer to the center. However, we should recall that it has no sense to go to such extreme values of $`p`$: the stars will be disrupted by tidal forces (see the discussion in Section VI) in those regions.
## 4 Boson stars
### 4.1 Basic concepts and configurations
Let us study the Lagrangian density of a massive complex self-gravitating scalar field (taking $`\mathrm{}=c=1`$),
$$=\frac{1}{2}\sqrt{g}\left[\frac{m_{\mathrm{Pl}}^2}{8\pi }R+_\mu \psi ^{}^\mu \psi U(|\psi |^2)\right],$$
(8)
where $`R`$ is the scalar of curvature, $`|g|`$ the modulus of the determinant of the metric $`g_{\mu \nu }`$, and $`\psi `$ is a complex scalar field with potential $`U`$. Using this Lagrangian as the matter sector of the theory, we get the standard field equations,
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R`$ $`=`$ $`{\displaystyle \frac{8\pi }{m_{\mathrm{Pl}}^2}}T_{\mu \nu }(\psi ),`$ (9)
$`\mathrm{}\psi +{\displaystyle \frac{dU}{d|\psi |^2}}\psi `$ $`=`$ $`0,`$ (10)
where the stress energy tensor is given by,
$`T_{\mu \nu }`$ $`=`$ $`(_\mu \psi ^{})(_\nu \psi )+`$ (11)
$`{\displaystyle \frac{1}{2}}g_{\mu \nu }\left[g^{\alpha \beta }(_\alpha \psi ^{})(_\beta \psi )U(|\psi |^2)\right]`$
and
$$\mathrm{}=_\mu \left[\sqrt{g}g^{\mu \nu }_\nu \right]/\sqrt{g}$$
(12)
is the covariant d’Alembertian. Because of the fact that the potential is a function of the square of the modulus of the field, we obtain a global $`U(1)`$ symmetry. This symmetry, as we shall later discuss, is related with the conserved number of particles. The particular form of the potential is what makes the difference between mini-boson, boson, and soliton stars. Conventionally, when the potential is given by
$$U=m^2|\psi |^2+\frac{\lambda }{2}|\psi |^4,$$
(13)
where $`m`$ is the scalar mass and $`\lambda `$ a dimensionless constant measuring the self-interaction strength, mini-boson stars are those spherically symmetric equilibrium configurations with $`\lambda =0`$. Boson stars, on the contrary, have a non-null value of $`\lambda `$. The previous potential with $`\lambda 0`$ was introduced by Colpi et al. (1986), who numerically found that the masses and radius of the configurations were deeply enlarged in comparison to the mini-boson case.
Soliton (also called non-topological soliton) stars are different in the sense that, apart from the requirement that the Lagrangian must be invariant under a global $`U(1)`$ transformation, it is required that –in the absence of gravity– the theory must have non-topological solutions; i.e. solutions with a finite mass, confined to a finite region of space, and non-dispersive. An example of this kind of potentials is the one introduced by Lee and his coworkers (Friedberg et al., 1987),
$$U=m^2|\psi |^2\left(1\frac{|\psi |^2}{\mathrm{\Phi }_0^2}\right)^2,$$
(14)
where $`\mathrm{\Phi }_0`$ is a constant. In general, boson stars accomplish the requirement of invariance under a $`U(1)`$ global transformation but not the solitonic second requirement. To fulfill it, it is necessary that the potential contains attractive terms. This is why the coefficient of $`(\psi ^{}\psi )^2`$ of the Lee form has a negative sign. Finally, when $`|\psi |\mathrm{}`$, $`U`$ must be positive, which leads, minimally, to a sixth order function of $`\psi `$ for the self-interaction. It is usually assumed, because of the range of masses and radius for soliton stars in equilibrium, that they are huge and heavy objects, although this finally depends on the choice of the different parameters.
We shall now briefly explain how these configurations can be obtained (Kaup, 1968; Ruffini and Bonazzola, 1969; Lee and Pang, 1992; Liddle and Madsen, 1992; Mielke and Schunck, 1998). We adopt a spherically symmetric line element
$$ds^2=e^{\nu (r)}dt^2e^{\mu (r)}dr^2r^2(d\vartheta ^2+\mathrm{sin}^2\vartheta d\phi ^2),$$
(15)
with a scalar field time-dependence ansatz consistent with this metric:
$$\psi (r,t)=\sigma (r)e^{i\omega t}$$
(16)
where $`\omega `$ is the (eigen-)frequency. This form of the field ensures us to be working in the configurations of minimal energy (Friedberg et al., 1987).
The non-vanishing components of the energy-momentum tensor are
$`T_0{}_{}{}^{0}=\rho ={\displaystyle \frac{1}{2}}[\omega ^2\sigma ^2(r)e^\nu +\sigma ^2(r)e^\mu +U],`$ (17)
$`T_1{}_{}{}^{1}=p_r={\displaystyle \frac{1}{2}}[\omega ^2\sigma ^2(r)e^\nu +\sigma ^2(r)e^\mu U],`$ (18)
$`T_2{}_{}{}^{2}=T_3{}_{}{}^{3}=p_{}={\displaystyle \frac{1}{2}}[\omega ^2\sigma ^2(r)e^\nu \sigma ^2(r)e^\mu U]`$ (19)
where $`{}_{}{}^{}=d/dr`$. One interesting characteristic of this system is that the pressure is anisotropic; thus, there are two equations of state $`p_r=\rho U`$ and $`p_{}=\rho U\sigma ^2(r)e^\mu `$. The non-vanishing independent components of the Einstein equation are
$`\nu ^{}+\mu ^{}`$ $`=`$ $`{\displaystyle \frac{8\pi }{m_{\mathrm{Pl}}^2}}(\rho +p_r)re^\mu ,`$ (20)
$`\mu ^{}`$ $`=`$ $`{\displaystyle \frac{8\pi }{m_{\mathrm{Pl}}^2}}\rho re^\mu {\displaystyle \frac{1}{r}}(e^\mu 1).`$ (21)
Finally, the scalar field equation is
$$\sigma ^{\prime \prime }+\left(\frac{\nu ^{}\mu ^{}}{2}+\frac{2}{r}\right)\sigma ^{}+e^{\mu \nu }\omega ^2\sigma e^\mu \frac{dU}{d\sigma ^2}\sigma =0.$$
(22)
To do numerical computations and order of magnitude estimates, it is useful to have a new set of dimensionless variables. We adopt here
$$x=mr,$$
(23)
for the radial distance, we redefine the radial part of the boson field as
$$\sigma =\sqrt{4\pi }\sigma /m_{\mathrm{Pl}},$$
(24)
and introduce
$$\mathrm{\Lambda }=\lambda m_{\mathrm{Pl}}^2/4\pi m^2,\mathrm{\Omega }=\frac{\omega }{m}.$$
(25)
In order to obtain solutions which are regular at the origin, we must impose the following boundary conditions $`\sigma ^{}(0)=0`$ and $`\mu (0)=0`$. These solutions have two fundamental parameters: the self-interaction and the central density (represented by the value of the scalar field at the center of the star). The mass of the scalar field fixes the scale of the problem. Boundary conditions representing asymptotic flatness must be applied upon the metric potentials, these determine –what is actually accomplished via a numerical shooting method– the initial value of $`\nu =\nu (0)`$. Then, having defined the value of the self interaction, or alternatively, the form of the soliton potential, the equilibrium configurations are parameterized by the central value of the boson field. As this central value increases, so does the mass and radius of the the star. This happens until a maximum value is reached in which the star looses its stability and disperses away (the binding energy being positive). Up to this value of $`\sigma _0`$, catastrophe theory can be used to show that these equilibrium configurations are stable (Kusmartsev et al., 1991). As an example of boson star configurations, we show in Fig. A supermassive scalar star at the Galactic Center?, the mass and number of particles (see below) for a $`\mathrm{\Lambda }=10`$ Colpi et al.’s potential, and in Fig. A supermassive scalar star at the Galactic Center?, the stability analysis.
When $`\mathrm{\Lambda }1`$, as in some of the cases we explore in the next sections, we must follow an alternative adimensionalization (Colpi et al., 1986). For large $`\mathrm{\Lambda }`$, we shift to the following set of variables,
$$\sigma _{}=\sigma \mathrm{\Lambda }^{1/2},$$
(26)
$$x_{}=x\mathrm{\Lambda }^{1/2},$$
(27)
$$M_{}=M\mathrm{\Lambda }^{1/2}.$$
(28)
Here, $`M_{}`$ is defined by
$$e^\mu =\left(12\frac{M_{}}{x_{}}\right)^1,$$
(29)
which corresponds to the Schwarzschild mass (see below). Ignoring terms $`𝒪(\mathrm{\Lambda }^1)`$, the scalar wave equation is solved algebraically to yield,
$$\sigma _{}=(\mathrm{\Omega }^2e^\nu 1)^{1/2},$$
(30)
and up to the same accuracy, the field equations are
$$\frac{dM_{}}{dx_{}}=\frac{1}{4}x_{}^2(3\mathrm{\Omega }^2e^\nu +1)(\mathrm{\Omega }^2e^\nu 1),$$
(31)
$$\frac{d\nu }{dx_{}}\frac{e^\mu }{x_{}}\frac{1}{x}_{}^2(1e^\mu )=\frac{1}{2}(\mathrm{\Omega }^2e^\nu 1)^2.$$
(32)
The system now depends only on one free parameter $`\mathrm{\Omega }^2e^{\nu (0)}`$. Numerical solutions show that the maximum mass corresponding to a stable star is given by $`M_{\mathrm{max}}0.22\mathrm{\Lambda }^{1/2}m_{\mathrm{Pl}}^2/m`$.
### 4.2 Masses estimates
The invariance of the Lagrangian density under a global phase transformation $`\psi \psi e^{i\vartheta }`$ of the complex scalar field gives (via the Noether’s theorem) a locally conserved current $`_\mu j^\mu =0`$, and a conserved charge (number of particles). We need to study the number of particles because it is essential to determine whether the configurations are stable or not. A necessary requirement towards the stability of the configurations is a negative binding energy ($`BE=MmN`$), i.e. the star must be energetically more favorable than a group of unbound particles of equal mass. From the Noether theorem, the current $`j^\mu `$ is given by
$$j^\mu =\frac{i}{2}\sqrt{g}g^{\mu \nu }[\psi ^{}_\nu \psi \psi _\nu \psi ^{}].$$
(33)
and the number of particles is
$$N:=j^0d^3x.$$
(34)
For the total gravitational mass of localized solutions, we may use Tolman’s expression (Tolman, 1934), or equivalently, the Schwarzschild mass:
$$M=(2T_0^{\mathrm{\hspace{0.33em}0}}T_\mu ^\mu )\sqrt{g}d^3x.$$
(35)
In Fig. A supermassive scalar star at the Galactic Center?, we show the mass of a boson star with $`\mathrm{\Lambda }=0`$ as a function of $`x`$. We have fitted this curve with a Boltzmann-like function
$$M_{\mathrm{fit}}(x)=A_2+\frac{A_1A_2}{1+e^{(xx_0)/\mathrm{\Delta }x}},$$
(36)
which is reliable except in regions very near the center, where $`M_{\mathrm{fit}}(x)`$ becomes slightly negative. The $`\chi ^2`$-parameter of the fitting is around $`10^5`$, and values for $`A_1,A_2,`$ and $`\mathrm{\Delta }x`$ are given in the figure. It is this formula what we have used in Fig. A supermassive scalar star at the Galactic Center? to analytically get the $`\sigma _r(p)`$ dependence of the boson star (in the range we use the approximation, the actual mass and the fitting differ negligibly). Note also what the fitting is physically telling us: it represents a black hole of mass $`A_1`$ plus an inner exponentially decreasing correction. This is the first time that such a Boltzmann-like fitting is done, and we think it could be usefully applied in other analytical computations. Models with different $`\mathrm{\Lambda }`$ and $`\sigma `$ can be equally well fitted.
Since boson stars are prevented from gravitational collapse by the Heisenberg uncertainty principle, we may make some straightforward mass estimates (Mielke and Schunck, 1998): For a boson to be confined within the star of radius $`R_0`$, the Compton wavelength has to satisfy $`\lambda _\psi =(2\pi \mathrm{}/mc)2R_0`$. In addition, the star radius must be of the order of the last stable Kepler orbit $`3R_\mathrm{S}`$ around a black hole of Schwarzschild radius $`R_\mathrm{S}:=2GM`$. In the case of a mini-boson star of effective radius $`R_0(\pi /2)^2R_\mathrm{S}`$ close to its Schwarzschild radius one obtains the estimate
$$M_{\mathrm{crit}}(2/\pi )m_{\mathrm{Pl}}^2/m0.633M_{\mathrm{Pl}}^2/m.$$
(37)
The exact value in the second expression was found only numerically.
For a mass of $`m=30`$ GeV, one can estimate the total mass of this mini–boson star to be $`M10^{10}`$ kg and its radius $`R_010^{17}`$ m, amounting a density 10<sup>48</sup> times that of a neutron star. In the case of a boson star ($`\lambda 0`$), since $`|\psi |m_{\mathrm{Pl}}/\sqrt{8\pi }`$ inside the boson star (Colpi et al., 1986), the energy density is
$$\rho m^2m_{\mathrm{Pl}}^2\left(1+\mathrm{\Lambda }/8\right).$$
(38)
Equivalently, we may think that this corresponds to a star formed from non–interacting bosons with re-scaled mass $`mm/\sqrt{1+\mathrm{\Lambda }/8}`$ and consequently, the maximal mass scales with the coupling constant $`\mathrm{\Lambda }`$ as,
$$M_{\mathrm{crit}}\frac{2}{\pi }\sqrt{1+\mathrm{\Lambda }/8}\frac{m_{\mathrm{Pl}}^2}{m}.$$
(39)
There is a large range of values of mass and radius that can be covered by a boson star, for different values of $`\mathrm{\Lambda }`$ and $`\sigma _0`$. For instance, if $`m`$ is of the order of the proton mass and $`\lambda 1`$, this is in the range of the Chandrasekhar limiting mass $`M_{\mathrm{Ch}}:=M_{\mathrm{Pl}}^3/m^21.5M_{}`$.
Larger than these estimates is the range of masses that a non-topological soliton star produces, this is because the power law dependence on the Planck mass is even higher: $`10^2(m_{\mathrm{Pl}}^4/m\mathrm{\Phi }_0^2)`$. These configurations are static and stable with respect to radial perturbations. It is by no means true, however, that these are the only stable equilibrium configurations that one can form with scalar fields. Many extensions of this formalism can be found. It is even possible to see that boson stars are a useful setting where to study gravitation theory in itself (Torres, 1997; Torres et al., 1998c, d; Whinnett and Torres, 1999). Most importantly for our hypothesis is that rotating stable relativistic boson stars can also be found with masses and radii comparable in magnitude to their static counterparts (Schunck and Mielke, 1996; Yoshida and Eriguchi, 1997). In astrophysical settings it is usual to expect some induced rotation of stellar objects, and it is important that these rotation may not destabilize the structure. Other interesting generalization is that of electrically charged boson stars introduced by Jetzer and van der Bij (1989). Although it is usually assumed that selective accretion will quickly discharge any astrophysical object, some recent results by Punsly (1998) suggest this may not always be the case.
From these simple considerations, we can set a difference between boson and fermion condensations (e.g. neutrino condensation): in the second case we can have an extended object (in the case of Sgr A, its size can be of about $`30`$ld) and it can be quite diluted. In the case of a boson star, the object is “strongly” relativistic, extremely compact, and its size is comparable with its Schwarzschild radius. We shall discuss the consequences of this point widely in the next section.
### 4.3 Effective potentials
The motion of test particles can be obtained from the Euler-Lagrange equations taking into account the conserved canonical momenta (Shapiro and Teukolsky, 1983). In a general spherically symmetric potential, the invariant magnitude squared of the four velocity ($`u^2g_{\mu \nu }\dot{x}^\mu \dot{x}^\nu `$), which is 1 for particles with non-zero rest mass and 0 for massless particles, yields to
$$\dot{r}^2=\frac{1}{g_{rr}}\left[\frac{E_{\mathrm{}}^2}{g_{tt}}u^2\frac{l^2}{g_{\varphi \varphi }}\right],$$
(40)
where $`l`$ and $`E_{\mathrm{}}`$ are constants of motion given by $`g_{\varphi \varphi }\mathrm{sin}^2\theta \dot{\varphi }`$ and $`g_{tt}\dot{t}`$ (angular momentum and energy at infinity, both per unit mass) respectively, and a dot stands for derivation with respect to an affine parameter. This equation can be transformed to
$$\frac{1}{2}\dot{r}^2+V_{\mathrm{eff}}=\frac{1}{2}E_{\mathrm{}}^2e^{\mu \nu },$$
(41)
where $`V_{\mathrm{eff}}`$ is an effective potential. This name comes from the fact that in the Schwarzschild solution $`e^{\mu \nu }1`$ and thus, the previous equation can be understood as a classical trajectory of a particle of energy $`E_{\mathrm{}}^2/2`$ moving in a central potential. This is not so in a more general spherically symmetric case, like these non-baryonic stars. In the boson star case, for instance, typical metric potentials are shown in Fig. A supermassive scalar star at the Galactic Center?; note that for them $`e^{\mu \nu }1`$ We can see, however, that $`e^{\mu \nu }<e^{\mu (0)\nu (0)}=C`$, where $`C`$ is a constant, and then, usual classical trajectories can be looked at, in the sense that we may construct an equation of the form $`\dot{r}^2/2+V_{\mathrm{eff}}<1/2E_{\mathrm{}}^2e^{\mu (0)\nu (0)}`$, and it will be always satisfied.
The effective potential that massive or massless particles would feel must be very different from the black hole case. In the case of a massless particle, the effective potential is given by
$$V_{\mathrm{eff}}=e^\mu \frac{l^2}{2r^2},$$
(42)
while for a massive test particle it is
$$V_{\mathrm{eff}}=\frac{1}{2}e^\mu \left(1+\frac{l^2}{r^2}\right).$$
(43)
We show both, boson (for the case of $`\mathrm{\Lambda }=0`$) and black hole potentials in Figs. A supermassive scalar star at the Galactic Center? and A supermassive scalar star at the Galactic Center?. The mass of the central object is fixed to be the same in both cases (see captions) and the curves represent a fourth order Runge-Kutta numerical integration of the equations we previously derived with $`\mathrm{\Lambda }=0`$ and $`\sigma (0)=0.1`$. In the case of massive test particles we use $`l^2/M^2=0,12`$ and 15; thus explicitly showing the change in the behavior of $`V_{\mathrm{eff}}`$ for the black hole case. As $`l`$ increases, this shape changes from a monotonous rising curve to one that has a maximum and a minimum before reaching its asymptotic limit. These extrema disappear for $`l^2/M^2<12`$. In the case of boson stars, however, if $`l0`$ we have a divergence in the center at $`r=0`$ and only one extremum -a minimum-, which occurs at rising values of $`r/M`$ as $`l`$ grows. The curve $`V_{\mathrm{eff}}`$ for $`l=0`$ -radially moving objects- is not divergent, and we can see that the particles may reach the center of the star with a non-null velocity, given by $`1/2\dot{r}^2=1/2E_{\mathrm{}}^2e^{\mu (0)\nu (0)}V_{\mathrm{eff}}(0)>0`$, and will then traverse the star unaffected.
For massless particles, differences are also notorious: a black hole produces a negative divergence and a boson star a positive one. Radial motion of massless particles is insensitive to $`V_{\mathrm{eff}}`$, being this equal to zero, as in the Schwarzschild case. In both cases, for massive and massless particles, outside the boson star the potential mimics the Schwarzschild one.
### 4.4 Particle orbits
In the case of a black hole, we can use the Newtonian analogy. Orbits can be of three types: if the energy is bigger than the effective potential at all points, particles are captured. If the energy is such that the energy equals the effective potential just once, particles describe an unbound orbit, and the point of equality is known as turning point. If there are two such points, orbits are bound around the black hole. Orbits in which the energy equals the potential in a minimum of the latter are circular and stable ($`\dot{r}=0`$, $`V_{\mathrm{eff}}^{\prime \prime }<0`$).
We can not, because of the different relationship between the metric potentials we commented above, do the same analysis using $`V_{\mathrm{eff}}`$ for the boson star. We can, however, note that in most cases the total equation for the derivative of $`r`$ will be modified in a trivial way: If we look at the cases where the effective potential has a divergence at the center (and because the metric coefficients do not diverge), there is no other possibility for the particles more than found a turning point. Orbits can then be bound or unbound depending on the energy, but we can always find a place where $`\dot{r}=0`$ and then it has to reverse its sign. In the particular cases in which the equality happens at the minimum of the potential we have again stable circular orbits. Only in the case where $`l=0`$ particles can traverse the scalar star unaffected. For $`m0`$ there is still the possibility of finding a turning point, if the energy is low enough. However, if the particle is freely falling from infinity, with $`E_{\mathrm{}}=1`$, all energy is purely rest mass, it will radially traverse the star, as would do a photon.
We conclude that all orbits are not of the capture type. They can be circular, or unbound, and they all have at least one turning point. This helps to explain why a non-baryonic object will not develop a singularity while still being a relativistic object (comparable effective potentials, equivalently, comparable relativisticity coefficient $`GM(r)/r`$.).
## 5 Galactic parameters and mass
One interesting fact, which seems to be not referred before, is the point that for all these scalar stars, their radius is always related with the mass in the same way: $`RMm_{\mathrm{Pl}}^2`$. This is indeed the statement -contrary to what it is usually assumed- that not all interesting astrophysical ranges of mass and radius can be modeled with scalar fields. In the scalar star models, from the given central mass, the radius we obtain for the star is comparable to that of the horizon, $`Rm_{\mathrm{Pl}}^2\times 2.61\times 10^6M_{}3.9\times 10^{11}`$cm.
The question now is for which values of the parameters we can obtain a scalar object of such a huge mass. For the case of mini-boson star we need an extremely light boson:
$$m[\mathrm{GeV}]=1.33\times 10^{25}\frac{M(\mathrm{})}{M_{\mathrm{BH}}}.$$
(44)
Given a central density $`\sigma (0)`$, $`M(\mathrm{})`$ stands for the dimensionless value of the boson star mass as seen by an observer at infinity. $`M_{\mathrm{BH}}`$ is the value of the black hole mass (in millions of solar masses), obtained by fitting observational data. Then, we are requiring that the total mass of the boson star equals that of the black hole. For instance, in Fig. A supermassive scalar star at the Galactic Center?, we have taken Eckart and Genzel’s (1997) and Genzel et al.’s (1996) data, and fit them with a mini-boson star with $`\sigma (0)=0.1`$, which yields to a boson mass given by $`m[\mathrm{GeV}]=2.81\times 10^{26}`$; the total mass of the star (without the cluster contribution) is 2.5 $`\times 10^6`$ M.
In the case of boson stars, and using the critical mass dependence, $`\sqrt{\lambda }m_{\mathrm{Pl}}^3/m^2`$, the requirement of a 2.5 million mass star yields to the following constraint,
$$m[\mathrm{GeV}]=7.9\times 10^4\left(\frac{\lambda }{4\pi }\right)^{1/4}.$$
(45)
It is possible to fulfill the previous relationship, for instance, with a more heavy boson of about 1 MeV and $`\lambda 1`$. A plot of this relation -for some values of $`\lambda `$\- is shown in Fig. A supermassive scalar star at the Galactic Center?. Note that in this case, the value of the dimensionless parameter $`\mathrm{\Lambda }`$ is huge, and special numerical procedures, as explained above, must be used to obtain solutions. The characteristics of these solutions have proven to be totally similar to those with $`\mathrm{\Lambda }=0`$ which were used in Fig. A supermassive scalar star at the Galactic Center?, just the adimensionalization differs.
Finally, in the case of a non-topological soliton we obtain the following constraint,
$$m[\mathrm{GeV}]=\frac{7.6\times 10^{12}}{\mathrm{\Phi }_0^2[\mathrm{GeV}]^2}.$$
(46)
For the usually assumed case, in which the order parameter $`\mathrm{\Phi }_0`$ is of equal value than the boson mass, we need very heavy bosons of unit mass $`m=1.2\times 10^4`$GeV. Other possible pairs are shown in Fig. A supermassive scalar star at the Galactic Center?.
### 5.1 Boson candidates?
Based only on the constraints imposed by the mass–radius relationship valid for the scalar stars analyzed, we may conclude that:
1. if the boson mass is comparable to the expected Higgs mass (hundreds of GeV), then the Center of Galaxy could be a non-topological soliton star;
2. an intermediate mass boson could produce a super-heavy object in the form of a boson star;
3. for a mini-boson star to be used as central objects for galaxies it is needed the existence of an ultra-light boson.
These conclusions should be considered as order of magnitude estimations. Several reasons force us to make this warning. Firstly, we are just considering static and uncharged stars, this is just a model (the simplest), but more complicated ones can modify the actual constraints. Secondly, we do not know the exact form of the self-interaction, or in the case of non-topological stars, the value of $`\mathrm{\Phi }_0`$. For instance, consider the Higgs mass. In the electro-weak theory a Higgs boson doublet $`(\mathrm{\Phi }^+,\mathrm{\Phi }^0)`$ and its anti-doublet $`(\mathrm{\Phi }^{},\overline{\mathrm{\Phi }}^0)`$ are necessary ingredients in order to generate masses for the $`W^\pm `$ and $`Z^0`$ gauge vector bosons. Calculations of two–loop electro-weak effects have lead to an indirect determination of the Higgs mass (Gambino, 1998). For a top quark mass of $`M_\mathrm{t}=173.8\pm 5`$ GeV, the Higgs mass is $`m_\mathrm{h}=104_{49}^{+93}`$ GeV. However, experimental constraints are weak. Fermilab’s tevatron (Han and Zhang, 1999) has a mass range of $`135<m_\mathrm{h}<186`$ GeV, and together with LHC at Cern, they could decide if these Higgs particles (in the given range) exist in nature. Interesting is to note that as a free particle, the Higgs boson is unstable with respect to the decays $`hW^++W^{}`$ and $`hZ^0+Z^0`$. However, as remarked by Mielke and Schunck (2000), in a compact object, and in full analogy with neutron stars –where there is equilibrium of $`\beta `$ and inverse $`\beta `$ decay– these decay channels are expected to be in equilibrium with the inverse process $`Z^0+Z^0h`$.
We should also mention the possible dilatons appearing in low energy unified theories, where the tensor field $`g_{\mu \nu }`$ of gravity is accompanied by one or several scalar fields. In string effective super-gravity (Ferrar et al., 1994), for instance, the mass of the dilaton can be related to the super-symmetry breaking scale $`m_{\mathrm{susy}}`$ by $`m_\phi 10^3(m_{\mathrm{susy}}/`$ TeV)<sup>2</sup> eV. Finally, a scalar with a long history as a dark matter candidate is the axion, which has an expected light mass $`m_\sigma =7.4\times (10^7\mathrm{GeV}/f_\sigma )`$ eV $`>10^{11}`$ eV with decay constant $`f_\sigma `$ close to the inverse Planck time. Goldstone bosons have also inferred mass in the range of eV and less, $`m_g<0.060.3`$ eV (Umeda et al., 1998).
If boson stars really exist, they could be the remnants of first-order gravitational phase transitions and their mass should be ruled by the epoch when bosons decoupled from the cosmological background. The Higgs particle, besides its leading role in inflationary theories, should be the best and natural candidate as constituent of a boson condensation if the phase transition occurred in early epochs. A boson condensation should be considered as a sort of topological defect relic. In this case, as we have seen, Sgr A could be a soliton star. If soft phase-transitions took place during cosmological evolution (e.g. soft inflationary events), the leading particles could have been intermediate mass bosons and so our supermassive objects should be genuine boson stars. If the phase transitions are very recent, the ultra-light bosons could belong to the Goldstone sector giving rise to mini-boson stars.
We shall not discuss here any further the formation processes of boson stars. The reader is referred to the recent paper by Mielke and Schunck (2000) and references therein. It is apparent that for every possible boson mass in the particle spectrum there is a boson star model able to fit the galactic center constraints, at least in order of magnitude.
## 6 Accretion and luminosity
### 6.1 Relativistic rotational velocities
For the static spherically symmetric metric considered here circular orbit geodesics obey,
$`v_\phi ^2`$ $`=`$ $`{\displaystyle \frac{r\nu ^{}e^\nu }{2}}=e^\nu {\displaystyle \frac{e^\mu 1}{2}}+{\displaystyle \frac{8\pi }{m_{\mathrm{Pl}}^2}}p_rr^2{\displaystyle \frac{e^{\mu +\nu }}{2}}`$ (47)
$``$ $`{\displaystyle \frac{M(r)}{r}}+{\displaystyle \frac{8\pi }{m_{\mathrm{Pl}}^2}}p_rr^2{\displaystyle \frac{e^{\mu +\nu }}{2}}.`$
These curves increase up to a maximum followed by a Keplerian decrease. Liddle and Schunck (1997) found that the possible rotation velocities circulating within the gravitational boson star potential are quite remarkable: their maximum reaches more than one-third of the velocity of light. Schunck and Torres (2000) proved that these high velocities are quite independent of the particular form of the self-interaction and are usually found in general models of boson stars. For instance, for $`\mathrm{\Lambda }=0,300`$ of the Colpi et al. ’s (1986) choice, $`U_{\mathrm{cosh}}=\alpha m^2\left[\mathrm{cosh}(\beta \sqrt{|\psi |^2})1\right]`$, and $`U_{\mathrm{exp}}=\alpha m^2\left[\mathrm{exp}(\beta ^2|\psi |^2)1\right]`$, the maximal velocities are: 122990 km/s at $`x=20.1`$ for $`\mathrm{\Lambda }=300`$, 102073 km/s at $`x=4.1`$ for $`\mathrm{\Lambda }=0`$, 104685 km/s at $`x=4.2`$ for $`U_{\mathrm{exp}}`$, and 102459 km/s at $`x=5.9`$ for $`U_{\mathrm{cosh}}`$ (Schunck and Torres, 2000).
With such high velocities, the matter possesses an impressive kinetic energy, of about 6% of the rest mass; i.e. to obtain the required luminosity we would need that about $`10^8`$ to $`10^7`$ solar masses per year be transformed into radiation. Note that the required matter-radiation transfer is at least two orders of magnitude smaller than the accretion rate towards Sgr A.
The maximum rotational speed is attained well outside the physical radius of the star, as can be seen by computing the dimensionless $`x`$ value for the star radius (e.g. it happens between $`x5`$ and 15 for $`\mathrm{\Lambda }`$ going from 0 to 300). It is interesting to note also that the dependence of the maximum velocity on $`\mathrm{\Lambda }`$ is not very critical, and the same process can be operative with mini-boson stars. The rotational velocity is dependent on the central density, increasing with a higher value of $`\sigma _0`$. To obtain large rotational velocities, it is needed that the central density of the star be highly relativistic, for Newtonian solutions velocities are low and quite constant over a larger interval. This is consistent with the density constraint of the dark object in Sgr A.
### 6.2 The black hole danger
How can one justify that the accretion onto the central object –a neutrino ball or a boson star– will not create a black hole in its center anyway. Interstellar gas and stars, while spiraling down towards the center of the object, will begin to collide with each other, and may glue together at the center, what could be the seed for a very massive black hole. Even a black hole of small mass can spiral inwards, and if it remains at the center, that black hole itself could be the seed. On the other hand, stellar formation of massive stars would yield, after evolution, to a black hole. Then, we need to consider whether there is a mechanism that prevents the formation of a very massive baryonic object -leading inevitably to a black hole- in the center of the galaxy. Below we provide an answer to this question.
The key aspect to consider here is disruption. A star interacting with a massive object can not be treated as a point mass when it is close enough to the object such that it becomes vulnerable to tidal forces. Such effects become important when the pericentre $`r_{\mathrm{min}}`$ is comparable to the tidal radius (Rees, 1988),
$$r_t=5\times 10^{12}M_6^{1/3}\left(\frac{R_{}}{R_{}}\right)\left(\frac{M_{}}{M_{}}\right)^{1/3}\mathrm{cm}.$$
(48)
Here, $`R_,`$ stands for the radius of the star and the sun respectively, while $`M_,`$ for the masses. $`M_6`$ is the mass of the central object in millions of solar masses. $`r_t`$ is the distance from the center object at which $`M/r^3`$ equals the mean internal energy of the passing star. Alternatively, we may compute the Roche criterion, a comparison between the smoothed out mass of the disruptor and the internal mean energy of the disruptee (Krolik, 1999). This catastrophe radius is
$$R_cR_g\mathrm{\hspace{0.17em}2}\left(\frac{M}{10^8M_{}}\right)^{2/3}(\rho _{}/\rho _{})^{1/3},$$
(49)
where $`R_c`$ is given in units of the event horizon, $`R_g`$, of the black hole, and the falling star has density $`\rho _{}`$. Only for black holes with masses smaller than $`10^8M_{}`$ and for certain low density stars -red giants- we can expect that the disruption happens outside the event horizon. This is the reason why supplying the material at lower densities, with the black hole gravity dominating the situation, can generate more power. Rees has also given an estimate of how frequently a star enters this zone. When star velocities are isotropic, the frequency with which a solar-like star pass within a distance $`r_{\mathrm{min}}`$ is
$$10^4M_6^{4/3}\left(\frac{N_{}}{10^5\mathrm{pc}^3}\right)\left(\frac{\sigma }{100\mathrm{k}\mathrm{m}\mathrm{s}^1}\right)\left(\frac{r_{\mathrm{min}}}{r_t}\right)\mathrm{yr}^1,$$
(50)
where $`N_{}`$ is the star density and $`\sigma `$ the velocity distribution. Disruptions are rare events that happens once in about 10000 years.
#### 6.2.1 Scalar stars
Because of the similar metric potentials, far from the center of the non-baryonic star, the accretion mechanism will be the same than that operative in the Schwarzschild case. If a boson star is in the center of the galaxy, the characteristics of the tidal radius and the timescale of disruption occurrence will be similar to those of a black hole of equal mass. Stars falling inwards will all be disrupted after they approach a minimum radius. Contrary to black holes of big masses, which can swallows stars as a whole and disrupt them behind their event horizons, boson stars will disrupt all stars –most stars outside and a few inside the Schwarzschild radius of a black hole of equal mass– at everyone sight.
In the case of black holes, the most recent simulations (Ayal et al., 2000) show that up to 75% of the mass that once formed the disrupted star become unbound. For boson stars, we argue that once the star is disrupted to test particles (with masses absolutely negligible to that of the central object), and because there is no capture orbits, all particles follow unbound trajectories. In this sense, all material is diverted from the center and the formation of a black hole is avoided.
It could be worth of interest to perform numerical simulations changing the central object from a black hole to a boson star, to see the aftermath and the fate of the debris in an actual boson-star-generated disruption. This, however, could well not be an easy task: In black hole simulations, the minimum radius is maintained still far from the center ($`10`$ Schwarzschild radius), where Newtonian or Post-Newtonian approximations are valid. To actually see the difference between a black hole and a boson star it could be necessary to attain inner values of radii, where the behavior is completely relativistic.
One case could merit further attention: the possible spiraling of a black hole of stellar size. This case may complicate the situation since a black hole can not be disrupted. However, differences in masses are so large that it will behave as a test particle for the boson central potential, and will be also diverted from the center. Moreover, being of stellar size, it will be appreciably influenced by other intruding stars, also making it to left an static position at the boson star center.
Finally, it is worth of interest to study other possible observational consequences of boson star disruption. The ejecta now is 100% of the star mass, and may possible produce some long term effects in the surrounding medium. In black hole scenarios, the bound debris will create a flare as it accretes. Does the same happen here? While we defer this kind of analysis to another work, we mention that if boson stars exists, disruption, what was once thought of as an inevitable concomitant of black holes, could now happen in non-baryonic environments.
#### 6.2.2 Neutrino balls
If the center of the galaxy is a neutrino ball, one also has to obtain a mechanism that prevents the formation of a very massive baryonic object. However, we have to expect very crucial differences with a boson star case, caused by the fact that a neutrino ball is an extended object and that the gravitational potential is shallower. The first thing to note is that $`r_t`$ is well within the neutrino ball, and then, stars will traverse the exterior parts of the ball without being disrupted. In doing so, however, the central mass that they see at the center will be less than the total mass of the ball, and at a distance $`r=r_t`$ the mass enclosed is negligible. Disruption can not proceed and other mechanism have to devised. We note then that the observation of disruption processes in the center of a galaxy is then indicative that a neutrino ball is not there.
When the mass enclosed by the neutrino ball is small enough (say, $`𝒪(10^3)M_{}`$), the accretion disk will be unstable. This happens about 0.1 – 1 light years from the center. There, stars which could actually form at a rate of 1 per hundred thousand years (the actual number will depend on the mass of the star) will be probably kicked off by intruding stars (Viollier, 2000). The absence of the disruption mechanism makes this problem worse, since given an enough amount of time, it is hard to think in compelling reasons by which gas and stars are expelled from the center.
In a published paper, Tsiklauri and Viollier (1999) suggested that matter arriving at the center would be diverted in the form of non-radiating jets generated by pressure of the inner accretion disk. However, it is not clear that gravity attractive forces of the spiraling objects will be smaller than gas pressure exerted by the disk. Moreover, even when the baryonic mass acquired by the central object during the entire age of the universe is at least two orders of magnitude smaller than the mass of the central neutrino ball, one should be worried if this yields to a black hole of that mass, since it may be the seed for a further collapse, or appreciably influence the dynamics. The mechanism considered in the previous paragraph seems more adequate than this for solving this question, if that is finally possible. To us, it is yet an unclear issue in the neutrino ball scenario.
## 7 Discussion: how to differentiate among these models?
One of the easiest things one may think of is to follow the trajectory of a particular star. This has been done by Munyaneza et al. (Munyaneza et al., 2000) in the case of a neutrino ball. The trajectory of S1, a fast moving star near Sgr A, offers the possibility of distinguishing between a black hole and a neutrino condensate, since Newtonian orbits deviate from each other by several degrees in a period of some years. However, as soon as the central object is not so extended, as in the boson star case, this technique is useless (in every case in which the pericentre is far than the tidal radius), and other forms of detecting their possible presence have to be devised.
It has been already noted that X-ray astronomy can probe regions very close to the Schwarzschild radius. It is only observations in the X-ray band that can study the inner accretion disk as close to the center as an event horizon would be. Recent results from the Japanese-US ASCA mission have revealed a broadened iron line feature that comes from so close to the event horizon that a gravitationally redshift is observed. This is 10000 times closer into the black hole than what can be pictured by HST. In particular, Iwasawa et al. (1996) claimed that ASCA observations to Seyfert 1 galaxy MCG-6-30-15 got data from 1.5 gravitational radii, and conclude that the peculiar line profile suggests that the line-emitting region is very close to a central spinning (Kerr) black hole where enormous gravitational effects operate. By the way, this is stating that a neutrino ball can not be the center of that galaxy. However, as was already noted (Schunck and Liddle, 1997), a boson star could well be a possible alternative, and X-ray could be used to map out in detail the form of the potential well. The NASA Constellation-X (Constellation-X web page, 2000) mission, to be launched in 2008, is optimized to study the iron K line feature discovered by ASCA and, if they are there, will determine the black hole mass and spin for a large number of systems. Still, Constellation-X will provide an indirect measure of the properties of the region within a few event horizon radii. A definite answer in this sense will probably be given by NASA-planned MAXIM mission (Maxim web page, 2000), a $`\mu `$-arcsec X-ray imaging mission, that would be able to take direct X-ray pictures of regions of the size of a black hole event horizon. Both of these space mission will have the ability to give us proofs of black hole existence, or to provide evidence for more strange objects, like boson stars.
Very recently, Falcke et al. (2000) have noted that gravitational lensing observations of very large baseline interferometry (VLBI) could give the signature to discriminate among these models. Falcke et al. assumed that the overall specific intensity observed at infinity is an integration of the emissivity (taken as independent of the frequency or falling as $`r^2`$) times the path length along geodesics. Defining the apparent boundary of a black hole as the curve on the sky plane which divides a region where geodesics intersects the horizon from a region whose geodesics miss the horizon, they noted that photons on geodesics located within the apparent boundary that can still escape to the observer will experience strong gravitational redshift and a shorter total path length, leading to a smaller integrated emissivity. On the contrary, photons just outside the apparent boundary could orbit the black hole near the circular photon radius several times, adding to the observed intensity. This is what produces a marked deficit of the observed intensity inside the apparent boundary, which they refer to as the “shadow” of the black hole. The apparent boundary of the black hole is a circle of radius $`27R_g`$ in the Schwarzschild case, which is much larger than the event horizon due to strong bending of light by the black hole. This size is enough to consider the imaging of it as a feasible experiment for the next generation of mm and sub-mm VLBI. While the observation of this shadow would confirm the presence of a single relativistic object, a non-detection would be a major problem for the current black hole paradigm.
Concerning the ideas put forward in this work, Falcke et al.’s shadow concept is appealing. In the case of a boson star, we might expect some diminishing of the intensity right in the center, this would be provided by the effect on relativistic orbits, however, this will not be as pronounced as if a black hole is present: for that case, many photons are really gone through the horizon and this deficit also shows up in the middle. If a boson star is there, some photons will traverse it radially, and the center region will not be as dark as in the black hole case. A careful analysis of Falcke et al.’s shadow behavior replacing the central black hole with a boson star model would be necessary to get any further detail, and eventually an observable prediction.
We also mention that the project ARISE (Advanced Radio Interferometry between Space and Earth) is going to use the technique of Space VLBI to increase our understanding of black holes and their environments. The mission, to be launched in 2008, will be based on a 25-meter inflatable space radio telescope working between 8 and 86 GHz (Ulvestad, 1999). It will study gravitational lenses at resolutions of tens of $`\mu `$arcsecs, yielding information on the possible existence (and signatures) of compact objects with masses between $`10^3M_{}`$ and $`10^6M_{}`$.
Another possible technique for detecting boson stars from other relativistic objects could be gravitational wave measurements (Ryan, 1997). If a particle with stellar mass is observed to spiral into a spinning object with a much larger mass and a radius comparable to its Schwarzschild length, from the emitted gravitational waves, one could in principle obtain the lowest multipole moments. The black hole no-hair theorem says that all moments are determined by its lowest two, the mass and angular momentum (assuming the charge equal to zero). Should this not be so, the central object would not be a black hole, and as far as we know, the only remaining viable candidate would be a boson star. In ten years time, perhaps, a combination of gravitational wave measurements, better determination of stellar motions, and mm and sub-mm VLBI techniques could give us a definite picture of the single object at the center of our Milky Way.
From a theoretical point of view, developments on gravitational lensing theory in very strong field regimes will be of extreme importance since, for objects like Sgr $`A^{}`$, the standard weak field theory does not hold and new effects have to be expected (Virbhadra and Ellis, 2000).
## 8 Conclusions
We have shown that boson stars provide the basic necessary ingredients to fit dynamical data and observed luminosity of the center of the Galaxy. They constitute viable alternative candidates for the central supermassive object, producing a theoretical curve for the projected stellar velocity dispersion consistent with Keplerian motion, relativistic rotational velocities, and having an extremely small size.
Other singularity-free models were as well considered, as non-baryonic fermion stars (e.g. neutrino, or gravitino condensations). In this case, the object is sustained by its Fermi energy, while in the boson star case, it is the Heisenberg’s Uncertainty Principle which prevents the system from collapsing to a singularity. Due to this fact, boson stars are genuine relativistic objects where a strong gravitational field regime holds.
This difference in the relativistic status of both objects is not trivial. While fermion neutrino balls are extended objects, boson stars mimic a black hole. Disruption processes can not happen in a fermion condensation, and it has the unpleasant consequence of not providing with a straightforward mechanism by which stars could be diverted from the center, and through which finally avoid the formation of a massive black hole inside the condensate.
The formation of boson stars, neutrino balls, and black holes can all be competitive processes. Then, it might well be that even if we discover that a black hole is in the center of the galaxy, others galaxies could harbor non-baryonic centers. In the case of boson stars, only after the discovery of the boson mass spectrum we shall be in position to determine a priori which galaxies could be modeled by such a center. Observations of galactic centers could then suggest the existence of boson scalars much before than their discovery in particle physicists labs.
We acknowledge insightful comments by H. Falcke, F. Schunck, D. Tsiklauri, R. Viollier, and A. Whinnett. Also, we acknowledge informal talks with D. Bennet, R. di Stefano, C. Kochanek, and A. Zaharov at the Moriond 2000 meeting. D.F.T. was supported by CONICET as well as by funds granted by Fundación Antorchas, and particularly thanks G. E. Romero for his encouragement and criticism. S. C. and G. L.’s research was supported by MURST fund (40%) and art. 65 D.P.R. 382/80 (60%). G.L. further thanks UE (P.O.M. 1994/1999).
## Appendix A Appendix
For reader’s convenience, we quote here the dimensional conversion for the radius and the mass of a boson star. Using the value of 1 GeV in cm<sup>-1</sup>, and taking into account the dimensionless parameter $`x=mr`$, we get
$$r[\mathrm{pc}]=\frac{x}{m[\mathrm{GeV}]}6.38\times 10^{33}.$$
(A1)
For the mass, recalling that $`M=M(x)m_{\mathrm{Pl}}^2/m`$, we get
$$M[10^6M_{}]=\frac{M(x)}{m[\mathrm{GeV}]}1.33\times 10^{25}.$$
(A2)
In the case where $`\mathrm{\Lambda }1`$, both right hand sides of the previous formulae get multiplied by $`\mathrm{\Lambda }^{1/2}`$.
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# Self-tuning in an Outgoing Brane Wave Model
## I Introduction
The idea that our universe is embedded in a higher dimensional world has received a great deal of renewed attention over the last two years. Of particular interest is the fact that this opens up new approaches to solving the cosmological constant problem (see, e.g., ) though related ideas were pursued some time ago. In particular, show that in a five dimensional theory of gravity coupled to a scalar field, the brane geometry can remain flat, independent of the value of the cosmological constant on the brane. This is known as “self-tuning”. There has been extensive discussion of this idea for static branes . However, a natural question to ask is whether the same self-tuning feature will remain in a dynamical context where the vacuum energy on the brane changes, e.g., due to a phase transition.
Dynamical discussions of our universe as a brane suffer from the awkward feature that evolution requires not just initial data on the brane, but also initial data for the bulk fields as well. Thus, there are a large number of degrees of freedom that we do not have direct access to. The simplest, and most natural assumption is that the these bulk degrees of freedom simply respond to motion on the brane. In other words, there are no incoming bulk waves, but only outgoing waves. We construct and study this “outgoing brane wave model”. It has a number of desirable features: (1) The metric in the bulk is an exact plane wave. Since the curvature is null, all higher order corrections to Einstein’s equation constructed from higher powers of the Riemann tensor automatically vanish. Hence, these metrics are exact solutions to string theory to all order in $`\alpha ^{}`$ . (2) The geometry on the brane is not flat, but is typically given by an expanding Robertson-Walker cosmology<sup>*</sup><sup>*</sup>*For other discussions of brane-world cosmology see .. However, the rate of expansion is independent of the value of the cosmological constant, so the self-tuning feature is incorporated. This continues to hold even when the vacuum energy on the brane changes due to a phase transition. (3) There are no naked singularities. Instead, there are null singularities in the past which extend from the big bang on the brane. (4) As mentioned above, this is a deterministic model in which the bulk geometry is determined solely by the brane.
Even though one does not need to specify boundary conditions at the null singularity, it still cuts off space in the fifth dimension. So the proper distance to the singularity on appropriate spacelike surfaces is finite, and the local four dimensional Newton’s constant will be nonzero. However, it is also time dependent. This is a problem with this model which is currently under investigation. Perhaps by adding a bulk cosmological constant one could localize gravity near the brane and stabilize the effective Newton’s constant.
We also find exact, time dependent solutions to the bulk equations which are more general than plane waves. These solutions can be used to study a phase transition in which the vacuum energy changes on an initially static, Poincare invariant brane. We find a solution in which the brane becomes time dependent after the transition. Nevertheless, it remains an open question whether or not the brane becomes time dependent in all solutions with an initially static brane.
After discussing the basic equations of motion in the next section, we introduce our outgoing brane wave model in section III and study its properties. We then consider more general time dependent solutions in section IV and use them to study phase transitions on initially static branes in section V.
## II Equations of Motion
Following , we take the low energy effective action to be
$$S=d^5x\sqrt{G}\left[R\frac{4}{3}(\phi )^2\right]+d^4x\sqrt{g}e^{b\phi }_{4D}$$
(II.1)
Gravity and the dilaton field $`\phi `$ live in a 5-dimensional world (with coordinates $`t,`$ $`x^i,`$ and $`y)`$ and are coupled to a thin 4-dimensional brane whose position is taken to be at $`y=0.`$ The metric $`g_{\mu \nu }=\delta _\mu ^M\delta _\nu ^NG_{MN}`$ is the 4-dimensional metric induced on the brane, and $`_{4D}`$ denotes the Lagrangian of the 4-dimensional world. (Here $`M,N=0,1,2,3,5`$; $`\mu ,\nu =0,1,2,3`$; and we have set $`16\pi G_5=1`$.) We take $`b`$ as a free phenomenological parameter which ultimately may be determined by string theory. We will assume the geometry on the brane is homogeneous and isotropic, so the stress energy tensor must take the form of a perfect fluid:
$`T_\nu ^\mu `$ $`=`$ $`e^{b\phi }\text{diag}(\rho _T,P_T,P_T,P_T)`$ (II.2)
$`=`$ $`e^{b\phi }V\delta _\nu ^\mu +e^{b\phi }\text{diag}(\rho ,P,P,P)`$ (II.3)
where we have set $`\rho _T=\rho +V`$ and $`P_T=PV.`$ With this form we describe a 4-dimensional world containing a vacuum energy $`V`$ and a perfect fluid with energy density $`\rho `$ and pressure $`P`$. We will treat $`V,`$ $`\rho ,`$ and $`P`$ as phenomenological functions of $`t.`$
Einstein’s equations read
$$R_{MN}\frac{1}{2}G_{MN}R=\frac{4}{3}\left[_M\phi _N\phi \frac{1}{2}G_{MN}(\phi )^2\right]+\frac{1}{2}\sqrt{\frac{g}{G}}T_{\mu \nu }\delta _M^\mu \delta _N^\nu \delta (y)$$
(II.4)
with $`g`$ and $`G`$ the determinant of $`g_{\mu \nu }`$ and $`G_{MN}`$ respectively. The equation of motion for the scalar field has a source involving the the four dimensional Lagrangian $`_{4D}`$. The Lagrangian for a perfect fluid (with pressure determined in terms of the energy density) is simply $`=\rho `$ . So if we model the matter on the brane by a vacuum energy and perfect fluid, we have $`_{4D}=(V+\rho )`$. The dilaton equation is thus
$$\frac{8}{3}^2\phi =\sqrt{\frac{g}{G}}be^{b\phi }(V+\rho )\delta (y)$$
(II.5)
We wish to consider solutions depending on both $`t`$ and $`y`$, which are homogeneous and isotropic in the three transverse directions. It will be convenient to adopt “conformal gauge” in the $`t,y`$ directions, and write the metric in the form
$`ds^2`$ $`=`$ $`e^{2A(t,y)}(dt^2+dy^2)+e^{2B(t,y)}(dx_1^2+dx_2^2+dx_3^2)`$ (II.6)
$`=`$ $`e^{2A(u,v)}dudv+e^{2B(u,v)}(dx_1^2+dx_2^2+dx_3^2)`$ (II.7)
where as usual $`uty`$ and $`vt+y.`$
The easiest way to solve Einstein’s equation (II.4) and the dilaton equation (II.5) is to first solve them in the bulk, that is, away from the brane at $`y=0.`$ We will solve these equations for $`A(t,y),`$ $`B(t,y),`$ and $`\phi (t,y)`$ for $`y>0`$ and $`y<0`$ separately and then match the solutions across the brane. We find it convenient to use the $`(u,v)`$ light cone coordinates to solve the equations in the bulk, and the $`(t,y)`$ coordinates to do the matching across the brane.
Using the metric (II.6) we obtain the following bulk equations of motion
$`B_{,uv}+3B_{,u}B_{,v}`$ $`=`$ $`0`$ (II.8)
$`2\phi _{,uv}+3(B_{,u}\phi _{,v}+B_{,v}\phi _{,u})`$ $`=`$ $`0`$ (II.9)
$`A_{,u}B_{,u}{\displaystyle \frac{1}{2}}(B_{,u}{}_{}{}^{2}+B_{,uu}){\displaystyle \frac{2}{9}}\phi _{,u}^2`$ $`=`$ $`0`$ (II.10)
$`A_{,v}B_{,v}{\displaystyle \frac{1}{2}}(B_{,v}{}_{}{}^{2}+B_{,vv}){\displaystyle \frac{2}{9}}\phi _{,v}^2`$ $`=`$ $`0`$ (II.11)
$`2\phi _{,u}\phi _{,v}+3(A_{,uv}+2B_{,uv}+3B_{,u}B_{,v})`$ $`=`$ $`0`$ (II.12)
Here (II.8) is the $`uv`$ component of Einstein’s equation. It is independent of $`\phi `$ since the dilaton is essentially a two dimensional scalar field and the $`uv`$ component of the stress energy tensor is like the trace of the two dimensional stress tensor which vanishes. (II.9) is the dilaton equation of motion. Both of these first two equations are independent of $`A`$ essentially because the two dimensional wave equation is conformally invariant. (II.10) is the $`uu`$ component of the Einstein equation, (II.11) is the $`vv`$ component, and (II.12) the $`ii`$ component. This last equation is a consequence of the first four by virtue of the Bianchi identity.
The jump conditions for the metric at the brane can be obtained from the usual thin shell conditions of general relativityThey can also be obtained by writing out (II.4) and integrating across $`y=0`$ at fixed $`t`$. . The unit outward normal to the surface $`y=0`$ is $`n=e^A/y`$. The extrinsic curvature $`K_{\mu \nu }=_\mu n_\nu `$ is
$$K_{tt}=e^AA_{,y}K_{ij}=e^{2BA}B_{,y}\delta _{ij}$$
(II.13)
and its trace is $`K=K_{\mu }^{}{}_{}{}^{\mu }=e^A(A_{,y}+3B_{,y})`$. In general, if $`K_{\mu \nu }]`$ is the jump in the extrinsic curvature across the surface $`y=0`$, the stress tensor on the brane is given by
$$T_{\mu \nu }=\frac{1}{8\pi G_5}(K_{\mu \nu }]g_{\mu \nu }K])$$
(II.14)
(For any function $`F(t,y)`$, we use the notation $`F]=F(t,0^+)F(t,0^{})`$ and $`F|=F(t,0)`$.) Using the fact that we have set $`16\pi G_5=1`$, and the form of the stress tensor (II.2), we obtain the matching conditions
$`6{\displaystyle \frac{B}{y}}]`$ $`=`$ $`e^{A+b\phi }|\rho _T=e^{A+b\phi }|(V+\rho ),`$ (II.15)
$`6{\displaystyle \frac{A}{y}}]`$ $`=`$ $`e^{A+b\phi }|(2\rho _T+3P_T)`$ (II.16)
$`=`$ $`e^{A+b\phi }|(2\rho +3PV)`$ (II.17)
$`{\displaystyle \frac{8}{3}}{\displaystyle \frac{\phi }{y}}]`$ $`=`$ $`be^{A+b\phi }|(V+\rho ).`$ (II.18)
where the matching condition for the dilaton is obtained by integrating (II.5) at fixed $`t`$ across $`y=0`$.
It is worth emphasizing that the matching conditions have to be satisfied at any instant in $`t`$. In other words, both sides in the matching conditions (II.15), (II.16) and (II.18) are functions of $`t`$, as opposed to be just constants in the static case. Comparing (II.18) and (II.15) we find
$$\phi _{,y}]=\frac{9}{4}bB_{,y}]$$
(II.19)
which is a rather strong constraint since $`b`$ is a constant. In many of our solutions this condition forces $`\phi `$ to be proportional to $`B.`$
Note that in the 4-dimensional universe on the brane, the metric is given by
$$ds^2=e^{2A(t,0)}dt^2+e^{2B(t,0)}(dx_1^2+dx_2^2+dx_3^2)$$
(II.20)
which can be written in the standard form of a Robertson-Walker universe
$$ds^2=d\tau ^2+a(\tau )^2(dx_1^2+dx_2^2+dx_3^2)$$
(II.21)
with
$`\tau `$ $`=`$ $`{\displaystyle 𝑑te^{A(t,0)}},`$ (II.22)
$`a(\tau )`$ $`=`$ $`e^{B(t(\tau ),0)}.`$ (II.23)
## III Outgoing brane wave model
There is a simple class of solutions to the bulk equations (II.8) – (II.12) which consists of taking $`A,`$ $`B,`$ and $`\phi `$ to depend only on $`u`$. The five field equations then collapse to just one
$$A_{,u}B_{,u}\frac{1}{2}(B_{,u}{}_{}{}^{2}+B_{,uu})\frac{2}{9}\phi _{,u}{}_{}{}^{2}=0$$
(III.1)
This corresponds to a plane wave propagating to the right. It may appear that this class of solutions has two arbitrary functions of $`u`$, but one of them can be absorbed by reparametrizing $`u`$. The remaining free function can be taken to be the amplitude of the dilaton field $`\phi `$. There are no independent gravitational degrees of freedom since we have assumed isotropy in the transverse $`x^i`$ space. (Gravitational plane waves would expand some directions while contracting others.) There are clearly analogous solutions depending only on $`v`$ representing plane waves moving to the left.
We now construct a solution of the full equations (II.4) and (II.5) by matching a solution depending only on $`u`$ (for $`y>0`$) to a solution depending only on $`v`$ (for $`y<0`$). The result is a solution in which the bulk spacetime consists of plane waves moving away from the brane on both sides. We will call this the outgoing brane wave model.
The matching conditions at $`y=0`$ can be solved as follows. It will be convenient to write
$$B(u,v)=\mathrm{log}h(u)$$
(III.2)
on the $`y>0`$ side. The continuity of $`B`$ implies that $`B(u,v)=\mathrm{log}h(v)`$ on the $`y<0`$ side. Since, according to (II.19), the jump in $`\phi _{,y}`$ must be proportional to the jump in $`B_{,y}`$ for all $`t`$, $`\phi `$ itself must be proportional to $`B`$.
$$\phi (u,v)=\frac{9}{4}b\mathrm{log}h(u)$$
(III.3)
for $`y>0`$ and $`\phi (u,v)=\frac{9}{4}b\mathrm{log}h(v)`$ for $`y<0`$. Note that an additive constant in $`\phi (u,v)`$ can be absorbed by scaling $`h`$ and the consequent additive constant in $`B`$ can be absorbed by scaling $`x^i`$. We can now immediately integrate (III.1) to obtainThroughout this paper it is understood that all functions inside the logarithm have absolute value signs.
$`A(u,v)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{9}{4}}b^2\right)B(u,v)+{\displaystyle \frac{1}{2}}\mathrm{log}B_{,u}(u,v)`$ (III.4)
$`=`$ $`{\displaystyle \frac{9}{8}}b^2\mathrm{log}h(u)+{\displaystyle \frac{1}{2}}\mathrm{log}h^{}(u)`$ (III.5)
for $`y>0`$ and $`A(u,v)=\frac{9}{8}b^2\mathrm{log}h(v)+\frac{1}{2}\mathrm{log}h^{}(v)`$ for $`y<0`$. An additive constant can be absorbed by scaling $`u`$ and $`v.`$
The function $`h`$ is determined once one picks an equation of state for the matter on the brane. In other words, the equation of states fixes both the amplitude of the bulk waves, and the dynamics of the brane geometry. Let
$$P_T=\gamma \rho _T$$
(III.6)
where $`\gamma `$ is a constant, and $`P_T`$ and $`\rho _T`$ are the total pressure and energy density. (A cosmological constant corresponds to $`\gamma =1`$.) Then the matching conditions (II.15) and (II.16) imply that the jump in $`A_{,y}`$ is proportional to the jump in $`B_{,y}`$. Since this must hold for all time, this implies
$$A=(2+3\gamma )B+k$$
(III.7)
for some constant $`k`$. This yields a first order equation for $`h`$ which can be integrated explicitly.
To illustrate some of the features of this model we start with a particularly simple special case. Inspection of (III.1) shows that we can choose $`\phi `$, $`A,`$ and $`B`$to be linear functions of $`u`$ on the $`y>0`$ side and linear functions of $`v`$ on the $`y<0`$ side respectively. This corresponds to setting
$$h(t)=e^{\lambda t}$$
(III.8)
where we will assume the constant $`\lambda `$ is positive. Thus, on the $`y>0`$ side, $`B=\lambda u,`$ $`\phi =\frac{9}{4}b\lambda u,`$ and
$$A=\frac{1}{2}(1+\frac{9}{4}b^2)\lambda u+\frac{1}{2}\mathrm{log}\lambda $$
(III.9)
and similarly on the $`y<0`$ side. We now set $`b=\pm \frac{2}{3}`$. (The case of general $`b`$ will be considered below.) Then $`A=B`$ \+ const, and it follows from (III.7) that $`\gamma =1`$ so the stress energy on the brane is a pure cosmological constant. From (II.15), the vacuum energy is
$$V=12\lambda ^{\frac{1}{2}}$$
(III.10)
The bulk metric is
$$ds^2=e^{2\lambda (ty)}[\lambda dt^2+\lambda dy^2+dx_idx^i]$$
(III.11)
for $`y>0`$ and
$$ds^2=e^{2\lambda (t+y)}[\lambda dt^2+\lambda dy^2+dx_idx^i]$$
(III.12)
for $`y<0.`$ Changing to cosmological time $`\tau =e^{\lambda t}/\sqrt{\lambda }`$, we see that the metric on the brane at $`y=0`$ has the Robertson-Walker form $`ds^2=d\tau ^2+\lambda \tau ^2dx_idx^i`$ which after scaling $`x^i`$ gives
$$ds^2=d\tau ^2+\tau ^2dx_idx^i$$
(III.13)
Thus, with a constant vacuum energy $`V`$, we have a universe with the scale factor growing linearly $`a(\tau )=\tau `$ rather the usual exponential growth. It is important to note that the expansion rate is independent of the value of the vacuum energy (as long as it is nonzero). Thus this solution has the self-tuning feature described in .
Unlike the static solutions discussed in , this solution has no naked timelike singularity. However there is a null singularity at $`u=\mathrm{}`$ on the right of the brane and at $`v=\mathrm{}`$ on the left. Null geodesics from these singularities reach the brane in finite affine parameter so they are not really at infinity. This is most easily seen by introducing a new coordinate $`U=e^{2\lambda u}/2`$ so the metric for $`y>0`$ becomes
$$ds^2=dUdv+2Udx_idx^i$$
(III.14)
The singularity is now at $`U=0`$. The brane is located at $`u=v`$ or $`U=e^{2\lambda v}/2`$, so the brane never actually hits the singularity, but instead becomes asymptotically null as $`v\mathrm{}`$ (see Fig. 1). The geometry on the left of the brane is similar with the roles of $`u`$ and $`v`$ interchanged.
Since the null singularity cuts off space in the fifth direction, one has a nonzero effective four dimensional Newton’s constant $`G_4`$. But since the distance of the singularity to the brane changes with time, $`G_4`$ will be time dependent. This is the main drawback of the outgoing brane wave model. It can perhaps be cured by modifying the bulk Lagrangian to try to localize gravity as in the Randall-Sundrum scenario.
As discussed in , the value of $`b`$ that we have chosen, $`b=\frac{2}{3}`$, arises naturally in string theory since then (II.1) is just the Einstein frame action arising from a string frame source proportional to $`e^{2\phi }`$. For our solution, the corresponding string metric $`\stackrel{~}{G}_{MN}=e^{4\phi /3}G_{MN}`$ is flat! The bulk solution is simply a linear dilaton vacuum where the dilaton is proportional to a null coordinate. This is well known to be an exact solution to string theory . The brane geometry is also flat in the string frame.
The general solution for $`h`$ given an arbitrary dilaton coupling to the brane $`b`$ and arbitrary equation of state $`P_T=\gamma \rho _T`$ is easy to construct. From (III.5) and (III.7) we get a first order equation for $`h`$ with solution
$$h(u)=(\lambda u)^{4/(20+24\gamma +9b^2)}$$
(III.15)
for $`y>0`$ where $`\lambda `$ is a positive constant related to $`k`$ in (III.7). The solution for $`y<0`$ is identical with $`u`$ replaced by $`v`$. This is valid whenever the exponent is finite. For the special case considered above ($`\gamma =1`$ and $`b=\pm 2/3`$) the exponent diverges and the solution for $`h`$ is not a power law, but rather an exponential as we saw. The metric on the brane turns out to be
$$ds^2=d\tau ^2+\tau ^{8/[12(1+\gamma )+9b^2]}dx_idx^i$$
(III.16)
after rescaling the $`x_i`$. In particular, for a pure cosmological constant ($`\gamma =1`$) the brane geometry is
$$ds^2=d\tau ^2+\tau ^{8/(9b^2)}dx_idx^i$$
(III.17)
Thus, for any nonzero coupling $`b`$, the expansion is a power law which is independent of the value of the vacuum energy. The value of this vacuum energy can be computed from (II.15). When there is no other matter on the brane ($`\gamma =1`$) one finds that the vacuum energy is indeed constant on the brane and given by
$$V=\pm \frac{24}{|9b^24|^{\frac{1}{2}}}\lambda ^{\frac{1}{2}}.$$
(III.18)
where the sign is the same as the sign of $`9b^24`$. The causal structure of the resulting spacetime is shown in Fig. 2. In this case, the singularity is at $`u=0`$ and the brane hits the singularity at the time of the big bang.
The properties of this solution are similar to the case $`b=\frac{2}{3}`$ discussed above. (In fact, Fig. 2 accurately represents the causal structure of the complete spacetime in this case also.) It is even true that the bulk spacetime is still an exact solution to string theory. In terms of a new null coordinate $`U`$, the bulk metric becomes
$$ds^2=dUdv+U^{8/(4+9b^2)}dx_idx^i$$
(III.19)
Surprisingly, this metric is independent of the equation of state parameter $`\gamma `$. However, the position of the brane in $`U,v`$ coordinates now depends on $`\gamma `$. The metric (III.19) has a covariantly constant null Killing vector $`\mathrm{}=/v`$ in addition to the translational symmetries in $`x^i`$. So it belongs to a general class of metrics known as exact plane waves. The Riemann tensor involves two powers of $`\mathrm{}`$ so all scalar curvature invariants vanish. Nevertheless, the singularity is real and results in infinite tidal forces. From the standpoint of string theory, since the curvature is null, all $`\alpha ^{}`$ corrections automatically vanish . Moreover, one can show that the exact two dimensional beta functions of a sigma-model with this target space vanish, so this is an exact solution to string theory even nonperturbatively .
When $`b=0`$ (and $`\gamma =1`$) (III.15) shows that $`h=(\lambda u)^1`$. The bulk metric is then $`ds^2=(\lambda u)^2(\lambda dudv+dx_idx^i)`$ which is actually flat space in disguise. The metric on the brane is $`d\tau ^2+e^{2\tau \sqrt{\lambda }}dx_idx^i`$ and we recover the usual de Sitter metric.
We now ask what happens if there is a transition from one value of the vacuum energy $`V_1`$ to another value $`V_2.`$ Since a pure vacuum energy is always constant, to achieve this we must assume that the matter on the brane consists of both a vacuum energy and another component which we take to be radiation. Physically, some of the initial vacuum energy is converted to radiation which is then redshifted away leaving a smaller vacuum energy in the future. Since the two independent components of the stress energy tensor (II.15) and (II.16) are determined in terms of $`A,B`$ and $`\phi `$ and these three functions are determined in terms of $`h`$ via (III.2), (III.3), and (III.5), we can compute $`V(t)`$ and $`\rho (t)`$ given any $`h(t)`$. Setting $`P=\rho /3`$ we find
$$V=\frac{3h^{\frac{9}{8}b^2}}{|h^{}|^{\frac{1}{2}}}\left[\left(3+\frac{9}{8}b^2\right)\frac{h^{}}{h}+\frac{h^{\prime \prime }}{2h^{}}\right]$$
(III.20)
and
$$\rho =\frac{3h^{\frac{9}{8}b^2}}{|h^{}|^{\frac{1}{2}}}\left[\left(1\frac{9}{8}b^2\right)\frac{h^{}}{h}\frac{h^{\prime \prime }}{2h^{}}\right]$$
(III.21)
Here everything is evaluated on the brane so that $`h`$ stands for $`h(t),`$ $`h^{}`$ for $`\frac{dh(t)}{dt},`$ and so forth.
To show that the brane geometry remains independent of the value of the vacuum energy even after a phase transition, one must find a function $`h(t)`$ for which $`\rho 0`$ everywhere, $`\rho 0`$ in the past and future, and $`V`$ approaches different constants in the past and future. Since (III.17) is the only solution for pure vacuum energy, the geometry will expand at this rate both in the past and future. The only difference will be a constant rescaling of the $`x^i`$ coordinates. It is not hard to find a suitable function $`h(t)`$. For example, in the case $`b=\pm \frac{2}{3}`$, we can choose
$$h(t)=\frac{e^{\lambda _1t}}{1+e^{(\lambda _1\lambda _2)t}}$$
(III.22)
with $`\lambda _1>\lambda _2`$. This describes a universe with pure vacuum energy $`V_1=12\lambda _1^{1/2}`$ in the past and $`V_2=12\lambda _2^{1/2}`$ in the future. During intermediate times there is radiation with energy density given by substituting (III.22) into (III.21). In Fig. 3 we show the time dependence of $`V(t)`$ and $`\rho (t)`$, assuming radiation only, for the case $`\lambda _1=2.5`$, $`\lambda _2=1`$ and $`b=2/3`$.
## IV General Time Dependent Solutions
We now turn to more general time dependent solutions of the bulk equations of motion (II.8) -(II.12). Remarkably, one can find explicit solutions which are much more general than the plane waves discussed so far. One can then match these solutions on the brane to obtain solutions to the full field equations.
The key is to observe that (II.8) involves only $`B`$ and can be solved to give
$$B(u,v)=\frac{1}{3}\mathrm{log}(f(u)+g(v))$$
(IV.1)
where $`f`$ and $`g`$ are two arbitrary functions. (A possible additive integration constant can always be absorbed by an overall scaling of $`f`$ and $`g`$.)
Given $`B(u,v),`$ we can solve the dilaton equation of motion (II.9). Using separation of variables, the general solution can be written
$$\phi (u,v)=𝑑k\frac{c(k)}{\sqrt{(f(u)k)(g(v)+k)}}$$
(IV.2)
A special case is obviously
$$\phi (u,v)=\frac{c}{\sqrt{(f(u)k)(g(v)+k)}}$$
(IV.3)
For this case, we can solve (II.10) and (II.11) explicitly to find $`A(u,v)`$ with the result
$$A(u,v)=\frac{1}{3}\mathrm{log}(f(u)+g(v))+\frac{1}{2}\mathrm{log}(f^{}(u)g^{}(v))\frac{c^2}{12}\frac{(f(u)+g(v))^2}{((f(u)k)(g(v)+k))^2}+a$$
(IV.4)
Another special solution for $`\phi `$ which will be more useful in the following is
$$\phi (u,v)=\frac{\alpha }{3}\mathrm{log}(f(u)+g(v))$$
(IV.5)
and the associated solution for $`A`$ is
$$A(u,v)=\left(\frac{2\alpha ^2}{9}1\right)\frac{1}{3}\mathrm{log}(f(u)+g(v))+\frac{1}{2}\mathrm{log}(f^{}(u)g^{}(v))+a$$
(IV.6)
with $`\alpha `$ and $`a`$ integration constants.
As an aside we mention one other special solution
$$\phi (u,v)=f(u)g(v)$$
(IV.7)
and
$$A(u,v)=\frac{1}{3}\mathrm{log}(f(u)+g(v))+\frac{1}{2}\mathrm{log}(f^{}(u)g^{}(v))+\frac{1}{3}(f(u)+g(v))^2+a$$
(IV.8)
This solution is not interesting for our purposes since the matching conditions (II.15), (II.16), and (II.18) can only be satisfied if the stress-energy tensor on the brane vanishes<sup>§</sup><sup>§</sup>§Continuity of $`\phi `$ and $`B`$ requires $`f_+=f_{}`$ and $`g_+=g_{}`$ (where the subscripts denote the solution for $`y>0`$ or $`y<0`$). From (II.15) - (II.18) we see that this leads to the conclusion that on the brane $`\rho _T=P_T=0.`$.
## V Matching to an initially static solution
We can now address the question of what happens if one starts with the static solutions discussed in and then changes the vacuum energy. We begin by writing down the static solution with our form of metric. For our solutions to be independent of $`t`$, the functions $`A,`$ $`B,`$ and $`\phi `$ can depend on $`(vu)`$ only. Since we want the metric on the brane to be Poincare invariant, we set $`A=B`$. With these restrictions, the bulk equations (II.8) -(II.12) reduce to two equations:
$`A^{\prime \prime }(y)+3A^{}(y)^2`$ $`=`$ $`0,`$ (V.1)
$`9A^{}(y)^2\phi ^{}(y)^2`$ $`=`$ $`0,`$ (V.2)
We wish to solve these equations for $`y>0`$ and $`y<0`$ separately and then match at $`y=0`$. Denoting the solutions in the two regions by subscripts $`+`$ or $``$, we have for $`y>0`$
$$A(y)=\frac{1}{3}\mathrm{log}(y_+y)+a_+,$$
(V.3)
and for $`y<0`$
$$A(y)=\frac{1}{3}\mathrm{log}(y+y_{})+a_{}$$
(V.4)
where $`y_\pm `$ and $`a_\pm `$ are constants of integration. In order to have a finite 4D Planck mass are forced to have naked singularities at $`y=y_+`$ and $`y=y_{}`$ to cut off spacetime. So the relevant region of spacetime is $`y_{}<y<y_+`$. Making the appropriate coordinate transformation we see that these agree with the solutions given in which we will refer to as KSS. From the dilaton equation (V.2), we have the freedom to choose on either side of the brane $`\phi =+3A`$ or $`\phi =3A`$ up to a constant. Let us fix
$$\phi (y)=\mathrm{log}(y_+y)+d_+$$
(V.5)
for $`y_+>y>0.`$ In the nomenclature of KSS, for $`0>y>`$ $`y_{}`$ the choice
$$\phi (y)=\mathrm{log}(y+y_{})+d_{}$$
(V.6)
is known as a type I solution, while the choice
$$\phi (y)=\mathrm{log}(y+y_{})+d_{}$$
(V.7)
is known as type II. The continuity of $`\phi `$ and $`A`$ at $`y=0`$ determines $`d_+`$ and $`a_+`$ in terms of the other constants.
Since we have set $`A=B`$, the matching conditions (II.15) and (II.16) imply that the stress-energy tensor on the brane takes the form of a pure vacuum energy density, with
$$6A_{,y}]=2(\frac{1}{y_+}+\frac{1}{y_{}})=e^{A+b\phi }|V$$
(V.8)
The matching condition for $`\phi `$ (II.18) implies
$$\frac{8}{3}\phi _{,y}]=\frac{8}{3}(\frac{1}{y_+}\frac{1}{y_{}})=be^{A+b\phi }|V$$
(V.9)
for type I and
$$\frac{8}{3}\phi _{,y}]=\frac{8}{3}(\frac{1}{y_+}+\frac{1}{y_{}})=be^{A+b\phi }|V$$
(V.10)
for type II. Taking the ratio of $`\phi ^{}]`$ and $`A^{}]`$ we see that type I is allowed only if $`b\pm \frac{4}{3}`$ (assuming that both $`y_+`$ and $`y_{}`$ are finite) while type II is allowed only if $`b=\frac{4}{3}`$. The solution for $`b=+\frac{4}{3}`$ is obtained from the type II solution by simply changing the sign of $`\phi `$ on both sides of the brane so $`\phi =3A`$ everywhere.
Suppose now that there is a phase transition in the microphysics so that $`V`$ changes from one constant to another. Will the brane remain flat? We found that it is most easy to use the special bulk solution (IV.5) and (IV.6) to match onto the type II solution in KSS. By causality, the static solution must persist for the region of the bulk defined by $`u=ty<0`$ and $`v=t+y<0`$. We thus take the solution (IV.1) for $`B`$ to be the following form
$$B(u,v)=\frac{1}{3}\mathrm{log}(y_{}(u)+uv)$$
(V.11)
for $`y>0`$ and, for $`y<0`$,
$$B(u,v)=\frac{1}{3}\mathrm{log}(y_{}(v)+vu),$$
(V.12)
In writing down (V.11) and (V.12) we have used the continuity condition of $`B(u,v)`$ at the location of the 3-brane $`u=v`$. $`y_{}`$ describes the time evolution of the location of the naked singularity on both sides. Before the phase transition, the universe is static and we have $`y_{}=`$constant. Given $`B(u,v)`$ we obtain $`\phi (u,v)`$ and $`A(u,v)`$ for $`y>0`$
$`\phi (u,v)`$ $`=`$ $`ϵ\mathrm{log}(y_{}(u)+uv)+d,`$ (V.13)
$`A(u,v)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{log}(y_{}(u)+uv)+{\displaystyle \frac{1}{2}}\mathrm{log}(1+y_{}^{}(u))+a`$ (V.14)
and for $`y<0`$
$`\phi (u,v)`$ $`=`$ $`ϵ\mathrm{log}(y_{}(v)+vu)+d,`$ (V.15)
$`A(u,v)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{log}(y_{}(v)+vu)+{\displaystyle \frac{1}{2}}\mathrm{log}(1+y_{}^{}(v))+a.`$ (V.16)
$`d`$ and $`a`$ are integration constants and $`ϵ=\pm 1`$. This solution was determined as follows. Demanding Poincare invariance on the brane for the initially static solution sets $`A=B`$ and thus $`\alpha ^2=9`$ in (IV.6). The continuity condition for $`\phi (u,v)`$ requires $`ϵ`$ to take the same sign on both sides of the brane. (II.19) determines $`b=ϵ\frac{4}{3}`$. It is clear that, for $`y_{}=`$constant$`>0`$, this solution reduces to the type II static solution discussed in KSS.
Using the matching conditions (II.15) and (II.16) we obtain
$`\rho _T(t)`$ $`=`$ $`4e^{abd}{\displaystyle \frac{y_{}^{}(t)+2}{\left|y_{}^{}(t)+1\right|^{\frac{1}{2}}}},`$ (V.17)
$`\rho _T(t)+P_T(t)`$ $`=`$ $`2e^{abd}{\displaystyle \frac{y_{}(t)y_{}^{\prime \prime }(t)}{\left|y_{}^{}(t)+1\right|^{\frac{1}{2}}(y_{}^{}(t)+1)}}.`$ (V.18)
From (V.17) we see immediately that $`\rho _T`$ depends on $`y_{}^{}`$ only, but not on the value of $`y_{}`$. This means that, in order for the vacuum energy density $`V`$ to settle to a constant value after the phase transition, $`y_{}^{}(t)`$ has to be a constant in the future.
There are two possible scenarios following from this observation. The first is that $`y_{}(t)`$ settles to a constant value and thus $`y_{}^{}=0`$. (V.17) and (V.18) implies the vacuum energy density $`V`$ would always go back to its original value before the phase transition and the universe becomes static again. In this scenario the vacuum energy density is determined by the constants of integration $`a`$ and $`d`$ and is independent of the value of $`y_{}`$.
The second scenario is that $`y_{}(t)`$ becomes linear in $`t`$ for $`t>0`$. In this case, the vacuum energy density $`V`$ could change to another constant value. However, the universe is not static even after the phase transition and would start evolving thereafter. As an example we take $`y_{}(t)=y_0\xi t`$ with $`0<\xi <1`$, then the vacuum energy density changes to
$`V_2`$ $`=`$ $`4e^{abd}{\displaystyle \frac{2\xi }{\sqrt{1\xi }}}`$ (V.19)
$`=`$ $`{\displaystyle \frac{V_1}{2}}{\displaystyle \frac{2\xi }{\sqrt{1\xi }}},`$ (V.20)
where $`V_1`$ is the initial value of the vacuum energy density. One can check that the transition can occur with positive energy density in the radiation. After the transition the universe is no longer static, but evolves toward a ‘big crunch’ singularity.
Here we have constructed an example which matches onto the initially static type II solution of KSS by using the special bulk solution (IV.5) and (IV.6). The continuity condition of $`\phi `$ requires $`ϵ`$ to take the same sign in (V.13) and (V.15). We are thus unable to match onto the type I solution of KSS using this special solution since the type I solution needs $`ϵ`$ to take opposite sign. It would be interesting to try to use the general solutions in Section IV to match onto the type I solution of KSS and study scenarios with phase transitions.
###### Acknowledgements.
We thank Lisa Randall, Eva Silverstein and Erik Verlinde for informative discussions. We also would like to thank Daniel Holz for help making the figures. This work was supported in part by the NSF under grant numbers PHY89-04035 and PHY95-07065, and also by Department of Energy under grant number DOE-ER-40682-143. I. L. is supported in part by an ITP Graduate Fellowship.
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# Dynamical Mass and Parity Condensate in Varying Topological Mass
## 1 Introduction
In a sense, a space-time with dimensions ”2+1” is mysterious because the mathematics tells us that there exists a specific term called a Chern-Simons term . This term is allowed just in (2+1) dimensions (generally in odd space-time dimensions). As is well known, a Lagrangian density for an electromagnetic field is given by a Maxwell term, which is 1) gauge invariant, 2) Lorentz invariant, and 3) bilinear for the gauge field. The (abelian) Chern-Simons term also satisfies all of 1) $``$ 3). Therefore the Maxwell theory in the (2+1) dimensions has a natural extension which is defined by adding the Chern-Simons term to the Maxwell Lagrangian. This extended version is called a Maxwell-Chern-Simons theory.
The modern technology of engineering makes possible to produce low-dimensional electron systems in realistic electronic devises. Especially, in (2+1)-dimensional systems, novel phenomena as the quantum Hall effect and the high-$`T_C`$ superconductivity were discovered. It is plausible that these phenomena may have their origins in the dimensions of the space-time.
In fact, there appeared many approaches to understand these macroscopic quantum effects by using (2+1)-dimensional quantum field theories with and without the Chern-Simons term. For example, the (2+1)-dimensional quantum electrodynamics ($`QED_3`$) was used to explain the quantum Hall effect and $`QED_3`$ with the Chern-Simons term (Maxwell-Chern-Simons $`QED_3`$) provided an anyon model which was expected to give us an essential mechanism for the high-$`T_C`$ superconductivity. These investigations produced important results and are still in progress.
What is the physical meaning of the Chern-Simons term? The Chern-Simons term gives the gauge field a mass without breaking the gauge symmetry. This mass is called a topological mass because the Chern-Simons term has a topological meaning as the secondary characteristic class. Whether the gauge field is massless or massive affects the nature of interactions, e.g., the range of interactions. In this case, the most important effect of the massive gauge field is to rescue the (2+1)-dimensional Maxwell theory from the infrared catastrophe which appears in a self-energy of fermion when the Maxwell field interacts with matters.
On the other hand, it is known that nonperturbative radiative corrections can produce a mass of fermion called a dynamical mass. The dynamical mass generation of four-component fermions in $`QED_3`$ without the Chern-Simons term has been studied in Ref.. One four-component fermion is equivalent to two two-component fermions. The mass term of the two-component fermion breaks the parity ($`P`$) symmetry while the one of the four-component fermion breaks $`P\times Z_2(\mathrm{flavour})`$ symmetry. In Ref., one of the present authors and others have investigated the dynamical mass generation of a single two-component fermion in $`QED_3`$ without the Chern-Simons term. A parity-breaking solution which generates the dynamical mass has been found .
Both analyses have been extended to the cases with the Chern-Simons term. The study of the four-component dynamical mass in the Maxwell-Chern-Simons $`QED_3`$ has been done in Ref.. The dynamical mass generation of a single two-component fermion in the Maxwell-Chern-Simons $`QED_3`$ has been studied in Ref..
The study of Ref. is motivated to clarify a role of the topological mass in the nonperturbative dynamics. The dependence of the dynamical mass on the topological mass has been investigated. However, as was pointed out there, the numerical estimation of the dynamical mass for a very small but non-zero value of the topological mass is very difficult technically. In this paper, we devise a way of estimation and extend the previous analysis further to the case in which the topological mass has much more smaller value. In addition, we evaluate a parity condensate which is an important quantity of indicating an amount of the parity breaking. We make clear a linking between theories with and without the Chern-Simons term in a nonperturbative level.
This paper is organized as follows. In Sec. 2, we explain the Maxwell-Chern-Simons $`QED_3`$ briefly. By examining a fermion self-energy in a perturbation, we demonstrate why we need a nonperturbative analysis for our purpose in Sec. 3. The Schwinger-Dyson equations are derived and studied by an approximated analytical method in Sec. 4. We solve the equations by a numerical method in Sec. 5, where the dynamical mass and a parity condensate are estimated. In Sec. 6, we summarize our results.
## 2 Maxwell-Chern-Simons $`QED_3`$
We consider the Maxwell-Chern-Simons $`QED_3`$ with the two-component Dirac fermion. The Lagrangian density of the theory is given by
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{\mu }{2}}\epsilon ^{\mu \nu \rho }A_\mu _\nu A_\rho {\displaystyle \frac{1}{2\alpha }}(_\mu A^\mu )^2+\overline{\psi }(i\overline{)}e\overline{)}A)\psi ,`$ (1)
where $`e`$ is the gauge coupling constant and $`\alpha `$ is the gauge-fixing parameter. The second term in the right-hand side of Eq.(1) is the so-called Chern-Simons term. It is well-known that the term gives the gauge field the mass $`\mu `$ without breaking the gauge symmetry. In fact, a free propagator of the gauge field $`iD_{\mu \nu }(pk)`$ derived from Eq.(1) is written as
$$iD_{\mu \nu }(p)=i\frac{1}{p^2\mu ^2}\left(g_{\mu \nu }\frac{p_\mu p_\nu }{p^2}\right)+\mu \frac{1}{p^2\mu ^2}\frac{1}{p^2}\epsilon _{\mu \nu \rho }p^\rho i\alpha \frac{p_\mu p_\nu }{p^4},$$
(2)
in which we find a massive pole at $`p^2=\mu ^2`$ so that $`\mu `$ is called the topological mass of the gauge field.
$`\psi `$ is the two-component fermion field which belongs to the irreducible spinor representation in (2+1)-dimensions. The Dirac matrices are defined by $`\gamma ^0=\sigma _3,\gamma ^1=i\sigma _1,\gamma ^2=i\sigma _2`$ with diag$`(g^{\mu \nu })=(1,1,1)`$ where $`\sigma _i`$’s (i=1, 2, 3) are the Pauli matrices. The $`\gamma ^\mu `$’s satisfy relations as $`\{\gamma ^\mu ,\gamma ^\nu \}=2g^{\mu \nu }`$, $`\gamma ^\mu \gamma ^\nu =iϵ^{\mu \nu \rho }\gamma _\rho +g^{\mu \nu }`$ and $`tr[\gamma ^\mu \gamma ^\nu ]=2g^{\mu \nu }`$. In this representation, there does not exist a matrix which anti-commutes with all of $`\gamma ^\mu `$’s so that we cannot define the chiral transformation. This is a specific aspect of the odd-dimensional space-time. In even-dimensions, the chiral symmetry requires that a fermion is massless. In odd-dimensions, the chiral symmetry itself does not exist. Instead, the mass term of the fermion is forbidden by parity symmetry. The parity transformation in (2+1)-dimensions is defined as $`x=(t,x,y)x^{}=(t,x,y),\psi (x)\gamma ^1\psi (x^{}),A^0(x)A^0(x^{}),A^1(x)A^1(x^{}),A^2(x)A^2(x^{}).`$ Under the parity transformation, the mass term of the fermion and the Chern-Simons term change their signs. Thus the mass terms of both the fermion and the gauge field are forbidden by the parity symmetry. We study how the breaking of parity by the topological mass affects the mass generation of the fermion.
## 3 Fermion self-energy in perturbation
Before proceeding to a nonperturbative analysis, it would be useful to see a fermion self-energy in the lowest order of perturbation in order to make prominent why we need a nonperturbative analysis to see the behaviour in the $`\mu 0`$ limit. The fermion self-energy in the one-loop approximation, $`\mathrm{\Sigma }^{(1)}(p)`$, is expressed as
$$\mathrm{\Sigma }^{(1)}(p)=\frac{d^3k}{(2\pi )^3}(ie\gamma ^\mu )iS_F(k)(ie\gamma ^\nu )iD_{\mu \nu }(pk),$$
(3)
where $`iS_F(p)`$ is a free fermion propagator as
$$iS_F(p)=\frac{i}{\mathit{}},$$
(4)
and $`iD_{\mu \nu }(pk)`$ is a free propagator of the gauge field given in Eq. (2).
The allowed form of the fermion propagator in the relativistic theory is written as
$$iS_F^{(1)}(p)=\frac{i}{A^{(1)}(p)\mathit{}B^{(1)}(p)}=\frac{i}{\mathit{}i\mathrm{\Sigma }^{(1)}(p)},$$
(5)
where $`A^{(1)}(p)`$ and $`B^{(1)}(p)`$ are functions of $`\sqrt{p_\mu p^\mu }`$, while $`\mathrm{\Sigma }^{(1)}(p)`$ depends on each $`p_\mu `$. $`A^{(1)}(p)^1`$ is the wave function renormalization and $`B^{(1)}(p)/A^{(1)}(p)`$ is a mass induced by dynamical effects at the momentum scale $`p`$. The so-called dynamical mass $`m_{phys}`$ is defined by $`m_{phys}=B^{(1)}(0)/A^{(1)}(0)`$ as usual. It is useful to notice the relations as
$$tr\left[\mathrm{\Sigma }^{(1)}(p)\right]=2iB^{(1)}(p),tr\left[\mathit{}\mathrm{\Sigma }^{(1)}(p)\right]=2i\{A^{(1)}(p)1\}p^2.$$
(6)
We substitute Eqs.(4) and (2) into Eq.(3) and use Eq.(6). We change the metric to the Euclidean one by the Wick rotation as $`(k^0,\stackrel{}{k})(ik^0,\stackrel{}{k})`$ and $`(p^0,\stackrel{}{p})(ip^0,\stackrel{}{p})`$. Then $`k^2`$ and $`p^2`$ are replaced by $`k^2=(k^0)^2(k^1)^2`$ and $`p^2=(p^0)^2(p^1)^2`$. After that, we transform the integral variables $`k^\mu `$’s to the polar coordinates $`(k,\theta ,\varphi )`$. The angular integration on $`\theta `$ and $`\varphi `$ can be done explicitly. Finally we obtain
$`B^{(1)}(p)`$ $`=`$ $`{\displaystyle \frac{e^2}{8\pi ^2p}}{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{1}{k}}[{\displaystyle \frac{1}{\mu }}(p^2k^2)\mathrm{ln}{\displaystyle \frac{(p+k)^2}{(pk)^2}}`$ (7)
$`+`$ $`{\displaystyle \frac{1}{\mu }}(p^2k^2+\mu ^2)\mathrm{ln}{\displaystyle \frac{(p+k)^2+\mu ^2}{(pk)^2+\mu ^2}}],`$
$`A^{(1)}(p)`$ $`=`$ $`1+{\displaystyle \frac{e^2}{8\pi ^2p^3}}{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{1}{k}}[2pk(\alpha +1)`$ (8)
$`+`$ $`\left\{{\displaystyle \frac{1}{2\mu ^2}}(p^2k^2)^2+{\displaystyle \frac{1}{2}}\alpha (p^2+k^2)\right\}\mathrm{ln}{\displaystyle \frac{(p+k)^2}{(pk)^2}}`$
$`+`$ $`\{{\displaystyle \frac{1}{2}}\mu ^2{\displaystyle \frac{1}{2\mu ^2}}(p^2k^2)^2\}\mathrm{ln}{\displaystyle \frac{(p+k)^2+\mu ^2}{(pk)^2+\mu ^2}}].`$
The dynamical mass of fermion is defined in the infrared limit so that we are interested in the behaviour of $`A^{(1)}(p)`$ and $`B^{(1)}(p)`$ in this limit. In the region of $`p1`$, Eqs.(7) and (8) are written as
$`B^{(1)}(p)`$ $`=`$ $`{\displaystyle \frac{e^2}{\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑k\left[{\displaystyle \frac{\mu }{k^2+\mu ^2}}+O(p^2)\right]`$ (9)
$`A^{(1)}(p)`$ $`=`$ $`1+{\displaystyle \frac{e^2}{\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑k\left[{\displaystyle \frac{1}{3}}\left\{{\displaystyle \frac{1}{k^2}}\alpha 2{\displaystyle \frac{\mu ^2}{(k^2+\mu ^2)^2}}\right\}+O(p^2)\right].`$ (10)
This infrared approximation makes the integration on k possible and we have
$`B^{(1)}(0)={\displaystyle \frac{e^2}{2\pi }}{\displaystyle \frac{|\mu |}{\mu }},A^{(1)}(0)=1{\displaystyle \frac{e^2}{6\pi }}{\displaystyle \frac{|\mu |}{\mu ^2}}+{\displaystyle \frac{e^2}{3\pi ^2}}{\displaystyle \frac{\alpha }{ϵ}},`$ (11)
where $`ϵ`$ is the infrared cutoff in the integration on $`k`$. We notice that $`A^{(1)}(0)`$ is free from the infrared divergence only in the Landau gauge.
At first sight, we can see the non-analytic $`\mu `$-dependence of $`A^{(1)}(0)`$ and $`B^{(1)}(0)`$ in Eq. (11): $`B^{(1)}(0)`$ depends on the sign of $`\mu `$ so that the value of $`B^{(1)}(0)`$ at $`\mu =0`$ is not well defined. $`A^{(1)}(0)`$ is singular at $`\mu =0`$. Therefore the theory with the Chern-Simons term may not be connected to the theory without the Chern-Simons term at least in the lowest order perturbation. However this consideration is not enough.
The theory has the two dimensional parameters $`e^2`$ and $`\mu `$ which have the dimension of mass. We take the $`p0`$ limit to define the dynamical mass and also the $`\mu 0`$ limit to see the linking of the theories with and without the Chern-Simons term. Then we have to fix the scale of $`e^2`$ in taking these limits. Thus the correct procedures of taking the limits are given by $`\frac{|\stackrel{}{p}|}{e^2}0`$ and $`\frac{\mu }{e^2}0`$. Here we notice that the $`\mu 0`$ limit in fixing $`e^2`$ is equivalent to the strong coupling limit. It means that Eq. (11) obtained in the perturbation is not available any more and the non-perturbative analysis is needed to investigate the limit. This situation motivates us to study the $`\mu 0`$ limit of the dynamical fermion mass by a non-perturbative method. This issue is extensively studied in the successive sections.
## 4 Schwinger-Dyson equation
In this section, we proceed to a nonperturbative analysis, where we use the Schwinger-Dyson technique to evaluate the dynamical mass of the fermion. The Schwinger-Dyson equation for the fermion self-energy $`\mathrm{\Sigma }(p)`$ is written as
$`\mathrm{\Sigma }(p)=(ie)^2{\displaystyle \frac{d^3k}{(2\pi )^3}\gamma ^\mu iS_F^{}(k)\mathrm{\Gamma }^\nu (k,pk)iD_{\mu \nu }^{}(pk)}.`$ (12)
$`\mathrm{\Gamma }^\nu (k,pk)`$ is a full vertex function and $`D_{\mu \nu }^{}(pk)`$ is a full propagator of the gauge field. $`S_F^{}(p)`$ is the full propagator of the fermion field which is written as
$`iS_F^{}(p)={\displaystyle \frac{i}{A(p)\mathit{}B(p)}}={\displaystyle \frac{i}{\mathit{}i\mathrm{\Sigma }(p)}},`$ (13)
which includes the full correction beyond the perturbative fermion propagator given in Eq. (5).
To analyze Eq.(12) further, we need to introduce any suitable approximation. In this paper, we limit ourselves to use the lowest ladder approximation where the full propagator of the gauge field and the full vertex are replaced by the free propagator and the bare vertex respectively as
$`iD_{\mu \nu }^{}(pk)iD_{\mu \nu }(pk),\mathrm{\Gamma }^\nu (k,pk)\gamma ^\nu ,`$ (14)
where $`iD_{\mu \nu }`$ has been given in Eq.(2). Thus the Schwinger-Dyson equation in the lowest ladder approximation becomes
$`\mathrm{\Sigma }(p)=(ie)^2{\displaystyle \frac{d^3k}{(2\pi )^3}\gamma ^\mu iS_F^{}(k)\gamma ^\nu iD_{\mu \nu }(pk)}.`$ (15)
We substitute Eqs.(2) and (13) into Eq.(15). Following the same steps used in getting Eqs.(7) and (8), we finally obtain the coupled integral equations as
$`B(p)`$ $`=`$ $`{\displaystyle \frac{e^2}{8\pi ^2p}}{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{k}{A(k)^2k^2+B(k)^2}}[\{\alpha B(k){\displaystyle \frac{1}{\mu }}(p^2k^2)A(k)\}\mathrm{ln}{\displaystyle \frac{(p+k)^2}{(pk)^2}}`$ (16)
$`+`$ $`\{{\displaystyle \frac{1}{\mu }}(p^2k^2)A(k)+\mu A(k)+2B(k)\}\mathrm{ln}{\displaystyle \frac{(p+k)^2+\mu ^2}{(pk)^2+\mu ^2}}],`$
$`A(p)`$ $`=`$ $`1+{\displaystyle \frac{e^2}{8\pi ^2p^3}}{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{k}{A(k)^2k^2+B(k)^2}}[2pk(\alpha +1)A(k)`$ (17)
$`+`$ $`\left\{{\displaystyle \frac{1}{2\mu ^2}}(p^2k^2)^2A(k)+{\displaystyle \frac{1}{\mu }}(p^2k^2)B(k)+{\displaystyle \frac{1}{2}}\alpha (p^2+k^2)A(k)\right\}\mathrm{ln}{\displaystyle \frac{(p+k)^2}{(pk)^2}}`$
$`+`$ $`\left\{{\displaystyle \frac{1}{2}}\mu ^2A(k){\displaystyle \frac{1}{2\mu ^2}}(p^2k^2)^2A(k)+\mu B(k){\displaystyle \frac{1}{\mu }}(p^2k^2)B(k)\right\}`$
$`\times `$ $`\mathrm{ln}{\displaystyle \frac{(p+k)^2+\mu ^2}{(pk)^2+\mu ^2}}],`$
which contain only the integration on the radial variable $`k`$. We solve these equations by an approximated analytical method in this section and also numerically by using an iteration method in the following section.
The limitation to the lowest ladder approximation is really unsatisfactory. This is mainly because the Schwinger-Dyson equation in the theory with the Chern-Simons term is complicated than the one in the theory without the Chern-Simons term, which will be explained soon later. In the situation, we think that it would be better to accumulate our experience in analysing the theory under a rather simpler approximation and to make a confidential first step as a springboard which is useful for a more extended next analysis. Even in the lowest ladder approximation, we can get many important informations. The results so found within the approximation might be denied in a more sophisticated approximation. But the analysis in the lowest ladder approximation will tell us what we should be careful for. Of course, there is no guarantee that the lowest ladder approximation is enough for our purpose. We will proceed our analysis beyond the lowest ladder approximation in future.
We can check easily that Eqs.(16) and (17) reduce to the Schwinger-Dyson equations in $`QED_3`$ without the Chern-Simons term if we put the topological mass $`\mu `$ equal to zero. In fact, taking the limit as $`\mu 0`$ in Eqs. (16) and (17), we obtain
$`B(p)`$ $`=`$ $`(\alpha +2){\displaystyle \frac{e^2}{8\pi ^2p}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle \frac{kB(k)}{A(k)^2k^2+B(k)^2}}\mathrm{ln}{\displaystyle \frac{(p+k)^2}{(pk)^2}},`$ (18)
$`A(p)`$ $`=`$ $`1\alpha {\displaystyle \frac{e^2}{4\pi ^2p^3}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle \frac{kA(k)}{A(k)^2k^2+B(k)^2}}\left[pk{\displaystyle \frac{p^2+k^2}{4}}\mathrm{ln}{\displaystyle \frac{(p+k)^2}{(pk)^2}}\right],`$ (19)
which is the Schwinger-Dyson equations in the lowest ladder approximation derived in $`QED_3`$ without Chern-Simons term. We can see that there exists the specific gauge where the wave function renormalization is absent. Thus in the Landau gauge$`(\alpha =0)`$, Eq.(19) gives us the simple solution as $`A(p)=1`$ and the problem reduces to solve Eq. (18) with $`A(p)=1`$.
In the case with the Chern-Simons term, as is seen in Eqs.(16) and (17), there does not exist such a specific gauge where the wave function is not renormalized. So far we cannot find a self-evident reason that the Landau is still specific in $`QED_3`$ with the Chern-Simons term, it must be fair to study Eqs. (16) and (17) for various values of the gauge-fixing parameter $`\alpha `$.
While the Schwinger-Dyson equations (16) and (17) in the Maxwell-Chern-Simons $`QED_3`$ reduce to the equations (18) and (19) in $`QED_3`$ without the Chern-Simons term under the ”naive” $`\mu 0`$ limit before integration, it is not so obvious whether the solutions of Eqs. (16) and (17) tend to the ones of Eq. (18) and (19). We cannot exclude a possibility that the interchange of the $`\mu 0`$ limit and the integration is not allowed because of a nontrivial nature of the integration kernels as distribution functions. In addition, we know that the Schwinger-Dyson equations (18) and (19) have two solutions. One is trivial ($`B(p)=0`$) and the other is non-trivial. It is interesting to find how the solution of Eqs. (16) and (17) behaves in the $`\mu 0`$ limit.
It is very useful if we can estimate $`A(0)`$ and $`B(0)`$ analytically even under a fairly crude approximation. The kernels of these integral equations are dumped rapidly as the integral variable $`k`$ increases so that the contribution from $`k0`$ is the most dominant one in the integrals. We approximate $`A(k)`$ and $`B(k)`$ by $`A(0)`$ and $`B(0)`$ in the integrals. We call this approximation ”the $`constant`$ approximation”. Of course this approximation might be too crude for our purpose and we only use the result as reference in the numerical analysis. Under this approximation, we can perform the remaining radial integration and obtain
$`B(0)={\displaystyle \frac{e^2}{2\pi }}+{\displaystyle \frac{e^2}{12\pi }}\alpha ,A(0)=1+{\displaystyle \frac{2\alpha }{\alpha +6}},`$ (20)
where we have considered the case of $`\mu >0`$.
From Eq.(20), we can see that the dependence of $`B(0)`$ and $`A(0)`$ on the gauge-fixing parameter, the coupling constant and the topological mass has the following peculiar features: $`B(0)`$ depends linearly on $`\alpha `$. In the Landau gauge ($`\alpha =0`$), $`A(0)=1`$ and $`B(0)=e^2/2\pi `$. $`A(0)=1`$ is favourable for us because $`A(p)=1`$ means that the Ward-Takahashi identity is satisfied. $`A(0)`$ does not depend on $`e^2`$. It means that the deviation of $`A(0)`$ from 1 is independent of the coupling constant. This is crucially different from the perturbative result given by Eq.(11) where the deviation is proportional to $`e^2`$. On the other hand, $`B(0)`$ is proportional to $`e^2`$. We recognize that there is no dependence on the topological mass $`\mu `$ in Eq.(20). In fact, if we apply the constant approximation to the case without the Chern-Simons term, we obtain the same results as Eq.(20). It means that the amount of the explicit parity breaking in the gauge sector by the topological mass does not affect the dynamical mass in the fermion sector in the constant approximation.
Now we proceed to a more precise numerical evaluation in the next section.
## 5 Numerical analysis
### 5.1 Nontrivial solutions and gauge dependence
We solve the two coupled integral equations (16) and (17) numerically by using a method of iteration. First we substitute trial functions into $`A(k)`$ and $`B(k)`$ on the right-hand sides of Eqs. (16) and (17) and then calculate the integrals numerically. The outputs so obtained, $`A(p)`$ and $`B(p)`$, are substituted back to the right-hand sides until the outputs coincide with the inputs. Finally we obtain convergent functions $`A(p)`$ and $`B(p)`$, which satisfy the integral equations, if there exist any solutions in Eqs. (16) and (17).
We have obtained the nontrivial solutions for the various values of the gauge parameter $`\alpha `$. We have found that $`A(p)`$ is almost a constant and its value is fairy close to 1 in the Landau gauge ($`\alpha =0`$) even for a very small $`\mu `$. Within the ladder approximation, no vertex correction is taken into account so that the wave function also should not be renormalized and $`A(p)`$ should be equal to 1. This means that the Ward-Takahashi identity is satisfied and the approximation is consistent with the gauge invariance. While we do not show the details here since the situation on this point is same as the one in ref. , the Landau gauge is still the best gauge. Therefore we present the results obtained in the Landau gauge hereafter.
### 5.2 Dependence on the topological mass
We are most interested in the dependence of the dynamical fermion mass on the topological mass of the gauge field. In the constant approximation, both $`A(0)`$ and $`B(0)`$ do not depend on the topological mass as seen in Eq. (20). We estimate them by the more precise numerical method.
We have studied the dependence of $`A(0)`$ on the dimensionless parameter $`\widehat{\mu }`$ which is defined by $`\widehat{\mu }=\mu /e^2`$. ($`A(0)`$ has no dimension. $`\mu `$ and $`e^2`$ have the dimension of mass.) We have found that the deviation of $`A(0)`$ from $`1`$ is less than 1 %. We may say that $`A(0)`$ is almost equal to 1 even in the extended region of $`\widehat{\mu }`$ which is wider than the one in Ref. .
The $`\widehat{\mu }`$-dependence of $`B(0)`$ is nontrivial. One of our motivations is to see whether or not the Maxwell-Chern-Simons $`QED_3`$ is smoothly connected to $`QED_3`$ without the Chern-Simons term in the $`\widehat{\mu }0`$ limit. Our numerical calculations show that very small meshes are needed to obtain reliable values of $`B(0)`$ in the region $`\widehat{\mu }1`$. Because of the limitation of our machine ability, we take another strategy different from our previous work . We estimate $`B(0)`$ on meshes of zero-width by an extrapolation from some $`B(0)`$’s on meshes of different finite-widths.
In Fig.1, we show the $`\widehat{\mu }`$-dependence of $`B(0)`$ in the region of $`10^5\widehat{\mu }10^2`$. ($`B(0)`$ has the dimension of mass so that we plot $`B(0)/e^2`$ with no mass dimension.) $`B(0)`$’s on meshes of finite-width show an abnormal behaviour in the region $`\widehat{\mu }1`$ that $`B(0)`$’s depart from $`B(0)`$ of $`QED_3`$ without the Chern-Simons term as $`\widehat{\mu }`$ decreases. This behaviour is improved by the extrapolation. We have estimated $`B(0)`$ on meshes of zero-width by the curve fitting of three $`B(0)`$’s on meshes of finite-widths. $`B(0)`$’s obtained by the extrapolation smoothly tend to the value of $`B(0)`$ in $`QED_3`$ without the Chern-Simons term.
The whole shape of $`B(0)`$ in the region of $`10^5\widehat{\mu }10^4`$, combining with the result in Ref. , is given in Fig.2. $`B(0)`$ is almost a constant in the region of $`\widehat{\mu }>10`$ and decreases rapidly in the region $`\widehat{\mu }=1.00.01`$. In the region of $`\widehat{\mu }<0.001`$, $`B(0)`$ becomes almost a constant again. The upper dotted line in Fig.2 is the result obtained in the constant approximation (Eq. (20)) and also in the lowest order perturbation (Eq. (11)). The lower dotted line shows the value obtained by a nonperturbative calculation in the case without Chern-Simons term. In other words, B(0) reproduces the result of perturbation in the region of $`e^2\mu `$ while it is close to the nonperturbative result of $`QED_3`$ without the Chern-Simons term in the region of $`e^2\mu `$.
### 5.3 Parity condensate
Another important quantity of indicating to what extent the parity symmetry is broken by the topological mass is a parity condensate, which is a gauge invariant order parameter of a vacuum. We evaluate the parity condensate as a function of the topological mass in order to know the nature of the vacuum. The parity condensate is defined by
$`<\overline{\psi }\psi >=\underset{x0}{lim}tr[iS_F^{}(x)],`$ (21)
where $`iS_F^{}(x)`$ is a propagator in real space-time coordinates, which is related to Eq.(13) by the Fourier transformation as
$`iS_F^{}(x)={\displaystyle \frac{d^3p}{(2\pi )^3}iS_F^{}(p)\mathrm{e}^{ipx}}.`$ (22)
By combining Eqs.(13), (21) and (22) and using the Wick rotation, we obtain
$`<\overline{\psi }\psi >={\displaystyle \frac{1}{\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑k{\displaystyle \frac{k^2B(k)}{A(k)^2k^2+B(k)^2}}.`$ (23)
We have known the numerical data of $`A(p)`$ and $`B(p)`$ for various values of $`\widehat{\mu }`$, which has been obtained by solving Eqs. (16) and (17), so that the $`\widehat{\mu }`$-dependence of $`<\overline{\psi }\psi >`$ is calculable by a numerical integral.
In Fig. 3, we show the $`\widehat{\mu }`$-dependence of $`<\overline{\psi }\psi >/e^4`$ in the region of $`10^5<\widehat{\mu }<10^2`$. ($`<\overline{\psi }\psi >`$ has the dimension of (mass)<sup>2</sup>. Therefore we plot a dimensionless quantity $`<\overline{\psi }\psi >/e^4`$.) To inspect the $`\widehat{\mu }0`$ limit, we extrapolate the $`B(0)`$’s in that region to $`\widehat{\mu }=0`$ by a least mean square method. The value obtained by the extrapolation is $`<\overline{\psi }\psi >/e^4=1.16187\times 10^3`$. On the other hand, we also evaluate $`B(p)`$ in $`QED_3`$ without the Chern-Simons term by solving Eqs. (18) in the Landau gauge ($`A(p)=1`$). Then we calculate the condensate of Eq. (23) numerically and obtain $`<\overline{\psi }\psi >/e^4=1.15222\times 10^3`$. Both values coincide within an error less than 1 %. Therefore we may say that the behaviour of $`<\overline{\psi }\psi >`$ also supports that $`QED_3`$ with the Chern-Simons term reduces to $`QED_3`$ without the Chern-Simons term smoothly in the $`\widehat{\mu }0`$ limit.
The $`\widehat{\mu }`$-dependence of $`<\overline{\psi }\psi >/e^4`$ in the region of $`10^5\widehat{\mu }10^4`$ is shown in Fig. 4. In the region of $`\widehat{\mu }1`$, the condensate is almost a constant. Around $`\widehat{\mu }0.001`$, it starts to increase. For $`\widehat{\mu }1`$, the increasing of the condensate is almost linear. The parity condensate increases more and more as $`\widehat{\mu }`$ does. Notice that there is no saturation for the increasing of $`<\overline{\psi }\psi >/e^4`$.
## 6 Conclusion
We have studied the dependence of the dynamical fermion mass and the parity condensate on the topological mass in the Maxwell-Chern-Simons $`QED_3`$ nonperturbatively by using the Schwinger-Dyson method.
When the topological mass is larger than the square of coupling constant, the value of the topological mass remains to be the one obtained by the perturbation. As the topological mass decreases, the value is changed to a nonperturbative value rapidly. The transition from the perturbative value to the nonperturbative one is sharp but not critical. Though it does not seem to be a phase transition, the inclusion of the Chern-Simons changes the nature of the theory drastically.
Motivated by the behaviour of the fermion self-energy in the perturbation, we have studied whether or not $`QED_3`$ with the Chern-Simons term reduces to $`QED_3`$ without the Chern-Simons term in the zero limit of the topological mass. We have checked the behaviour of the dynamical mass and the parity condensate for the extremely small topological mass in detail. The result shows that both quantities tend to the ones of $`QED_3`$ without the Chern-Simons term. Thus we conclude that $`QED_3`$ with the Chern-Simons term reduces to $`QED_3`$ without the Chern-Simons term smoothly in the nonperturbative level.
In general, an addition of a topological term to an action can give us a highly nontrivial deformation of a theory. Then it is not self-evident how the theories with or without the topological term are connected each other. The present work suggests that the linking of the theories should be considered in a nonperturbative level.
## Acknowledgment
One of the authors (T. M.) would like to thank Professor M. Kenmoku for his hospitality at the Nara Women’s University.
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# NEUTRINO OSCILLATIONS AND COSMOLOGY
## 1 Neutrino oscillations in vacuum. Basic concepts.
As is very well known, neutrinos possibly oscillate if (and because) their mass eigenstates are different from their interaction eigenstates <sup>-</sup> (for a review and more references see ). In other words, the mass matrix of different neutrino species is not diagonal in the basis of neutrino flavors: $`[\nu _e,\nu _\mu ,\nu _\tau ]`$. The latter is determined by the interaction with charged leptons, so that a beam of e.g. electrons would produce $`\nu _e`$ which is a mixture of several different mass states. An important condition is that the masses are different, otherwise oscillations would be unobservable. Indeed, if all the masses are equal, the mass matrix would be proportional to the unit matrix which is diagonal in any basis.
Of course not only neutrinos are capable to oscillate. All particles that are produced in the same reactions will do that, but usually the oscillation frequency, $`\omega _{osc}\delta m^2/2E`$ is so huge and correspondingly the oscillation length
$`l_{osc}=2p/\delta m^2`$ (1)
is so small, that the effect is very difficult to observe. Here $`E`$ and $`p`$ are respectively the energy and momentum of the particles in consideration and $`\delta m^2=m_1^2m_2^2`$. The expression is written in relativistic limit. Only for $`K`$-mesons and hopefully for neutrinos the mass difference is so small that $`l_{osc}`$ is, or may be, macroscopically large.
The neutrino Lagrangian can be written as follows:
$`_\nu =i\overline{\nu }/\nu +\overline{\nu }\nu +\overline{\nu }Z/\nu +\overline{\nu }W/l`$ (2)
where the vector-column $`\nu =[\nu _e,\nu _\mu ,\nu _\tau ]^T`$ is the operator of neutrino field in interaction basis, $`l=[e,\mu ,\tau ]^T`$ is the vector of charged lepton operators; the last two terms describe respectively neutral and charge current interactions (with $`Z`$ and $`W`$ bosons). The upper index “T” means transposition. The matrix $``$ is the mass matrix and by assumption it is non-diagonal in the interaction basis. This is not necessary but quite natural because normally masses know nothing about interactions.
Transformation between mass and interaction eigenstates is realized by an orthogonal matrix $`U`$ with the entries that are parameters of the theory to be determined by experiment. In the simplest case of only two mixed particles the matrix $`U`$ has the form
$`U=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)`$ (5)
If for example the only noticeable mixing is between electronic and muonic neutrinos, then the flavor eigenstates are related to the mass eigenstates $`\nu _{1,2}`$ as:
$`\nu _e`$ $`=`$ $`\nu _1\mathrm{cos}\theta +\nu _2\mathrm{sin}\theta ,`$
$`\nu _\mu `$ $`=`$ $`\nu _1\mathrm{sin}\theta +\nu _2\mathrm{cos}\theta `$ (6)
Thus if electronic neutrinos are produced on a target by a beam of electrons, their propagating wave function would have the form
$`\psi _{\nu _e}(\stackrel{}{r},t)=\mathrm{cos}\theta |\nu _1e^{ik_1x}+\mathrm{sin}\theta |\nu _2e^{ik_2x}`$ (7)
where $`kx=\omega t\stackrel{}{k}\stackrel{}{r}`$ and sub-$`\nu _e`$ means that the initial state was pure electronic neutrino. Below (in this section only) we denote neutrino energy by $`\omega `$ to distinguish it from the energies of heavy particles that are denoted as $`E`$. We assume, as is normally done, plane wave representation of the wave function.
If such a state hits a target, what is the probability of producing an electron or a muon? This probability is determined by the fraction of $`\nu _e`$ and $`\nu _\mu `$ components in the wave function $`\psi _{\nu _e}`$ at space-time point $`x`$. The latter can be found by re-decomposition of $`\nu _{1,2}`$ in terms of $`\nu _{e,\mu }`$:
$`\psi _{\nu _e}(\stackrel{}{r},t)=\mathrm{cos}\theta e^{ip_1x}\left(\mathrm{cos}\theta |\nu _e\mathrm{sin}\theta |\nu _\mu \right)+\mathrm{sin}\theta e^{ip_2x}\left(\mathrm{sin}\theta |\nu _e+\mathrm{cos}\theta |\nu _\mu \right)`$ (8)
One can easily find from that expression that the probability to register $`\nu _e`$ or $`\nu _\mu `$ are respectively:
$`P_{\nu _e}(\stackrel{}{r},t)\mathrm{cos}^4\theta +\mathrm{sin}^4\theta +2\mathrm{sin}^2\theta \mathrm{cos}^2\theta \mathrm{cos}\delta \mathrm{\Phi },`$ (9)
$`P_{\nu _\mu }(\stackrel{}{r},t)2\mathrm{sin}^2\theta \mathrm{cos}^2\theta \left(1\mathrm{cos}\delta \mathrm{\Phi }\right)`$ (10)
Here $`\delta \omega =\omega _1\omega _2`$, $`\delta \stackrel{}{k}=\stackrel{}{k}_1\stackrel{}{k}_2`$, and
$`\delta \mathrm{\Phi }(\stackrel{}{r},t)=\delta \omega t\stackrel{}{\delta }k\stackrel{}{r}`$ (11)
The energy difference between the mass eigenstates is
$`\delta \omega ={\displaystyle \frac{\omega }{m^2}}\delta m^2+{\displaystyle \frac{\omega }{\stackrel{}{k}}}\delta \stackrel{}{k}`$ (12)
Using this expression we find for the phase difference
$`\delta \mathrm{\Phi }(\stackrel{}{r},t)={\displaystyle \frac{\delta m^2}{2\omega }}t+\delta \stackrel{}{k}\left({\displaystyle \frac{\stackrel{}{k}}{\omega }}t\stackrel{}{r}\right)`$ (13)
The standard result of the neutrino oscillation theory are obtained if one assumes that: 1) $`\delta \stackrel{}{k}=0`$, 2) $`\stackrel{}{k}=\omega \stackrel{}{v}`$, and 3) $`t=r/v`$:
$`\delta \mathrm{\Phi }={\displaystyle \frac{\delta m^2r}{2k}}`$ (14)
Each of these assumptions is difficult to understand and, even more, some of them, in particular, $`\delta k=0`$ may be explicitly incorrect (see below). Both the second and the third conditions are fulfilled for a classical motion of a point-like body, however their validity should be questioned for a quantum mechanical particle (for a wave). Despite all that the final result (14) is true and if there are some corrections, they can be trivially understood.
Basic features of neutrino oscillations were discussed in many papers. An incomplete list of references includes <sup>-</sup>. One can find more citations and discussion in the above quoted papers and in the books <sup>-</sup>. Still some confusion and indications on possible controversies reappear from time to time in the literature, so it seems worthwhile to present a consistent derivation of eq. (14) from the first principles. A large part of this section is based on the discussions (and unpublished work) with A.Yu. Morozov, L.B. Okun, and M.G. Schepkin.
Let us consider a localized source that produces oscillating neutrinos; we keep in mind, for example, a pion decaying through the channel $`\pi \mu +\nu _\mu `$. The wave function of the source $`\psi _s(\stackrel{}{r},t)`$ can be Fourier decomposed in terms of plane waves:
$`\psi _s(\stackrel{}{r},t)={\displaystyle d^3pC(\stackrel{}{p}\stackrel{}{p}_0)e^{iEti\stackrel{}{p}\stackrel{}{r}}}`$ $``$
$`e^{iE_0ti\stackrel{}{p}_0\stackrel{}{r}}{\displaystyle d^3qC(\stackrel{}{q})\mathrm{exp}\left[i\stackrel{}{q}\left(\stackrel{}{r}\stackrel{}{V}_0t\right)\right]}`$ $`=`$ $`e^{iE_0ti\stackrel{}{p}_0\stackrel{}{r}}\stackrel{~}{C}\left(\stackrel{}{r}\stackrel{}{V}_0t\right)`$ (15)
where $`\stackrel{}{V}_0=\stackrel{}{p}_0/E_0`$. It is the standard wave packet representation. The function $`C(\stackrel{}{p}\stackrel{}{p}_0)`$ is assumed to be sharply peaked around the central momentum $`\stackrel{}{p}_0`$ with dispersion $`\mathrm{\Delta }\stackrel{}{p}`$. The particle is, by construction, on-mass-shell, i.e. $`E^2=p^2+m^2`$. This is of course also true for the central values $`E_0`$ and $`p_0`$. As the last expression shows, the particle behaves as a plane wave with frequency and wave vector given respectively by $`E_0`$ and $`\stackrel{}{p}_0`$ and with the shape function (envelope) given by $`\stackrel{~}{C}(\stackrel{}{r}\stackrel{}{V}_0t)`$ which is the Fourier transform of $`C(\stackrel{}{q})`$. Evidently the envelope moves with the classical velocity $`\stackrel{}{V}_0`$. Characteristic size of the wave packet is $`l_{pack}1/\mathrm{\Delta }p`$.
Let us consider the pion decay, $`\pi \mu +\nu `$. One would naturally expect $`\delta \omega \delta k\delta m^2/E`$. If this is true the probability of oscillation would be
$`P_{osc}\mathrm{cos}\left[{\displaystyle \frac{x+b\left(xVt\right)}{l_{osc}}}\right],`$ (16)
where $`l_{osc}`$ is given by expression (1) and $`b`$ is a numerical coefficient relating $`\delta p`$ with $`l_{osc}`$. For simplicity the one-dimensional expression is presented.
Thus to obtain the probability of neutrino registration one should average the factor $`(xVt)`$ over the size of the wave packet and for large packets, if $`b`$ is non-negligible, a considerable suppression of oscillations should be expected. The size of the neutrino wave packet from the pion decay at rest is macroscopically large, $`l_{pack}c\tau _\pi 7.8`$ m, where $`c`$ is the speed of light and $`\tau _\pi =2.610^8`$ sec is the pion life-time. The oscillation length is $`l_{osc}=0.4\mathrm{m}(E/\mathrm{MeV}/(\delta m^2/\mathrm{eV}^2)`$, so $`l_{osc}`$ could be smaller or comparable to $`l_{pack}`$ and the effect of suppression of oscillations due to a finite size of the wave packet might be significant. It is indeed true but only for the decay of a moving pion, and this suppression is related to an uncertainty in the position of decay. To see that we have to abandon the naive approach described above and to work formally using the standard set of quantum mechanical rules.
Let us assume that neutrinos are produced by a source with the wave function $`\psi _s(\stackrel{}{x},t)`$. This source produces neutrinos together with some other particles. We assume first the following experimental conditions: neutrinos are detected at space-time point $`\stackrel{}{x}_\nu ,t_\nu `$, while the accompanying particles are not registered. The complete set of stationary states of these particles is given by the wave functions $`\psi _n\mathrm{exp}(iE_nt)`$. The amplitude of registration of propagating state of neutrino of type $`j`$ (mass eigenstate) accompanying by other particles in the state $`\psi _n`$ is given by
$`A_j^{(n)}={\displaystyle 𝑑\stackrel{}{r}_s𝑑t_s\psi _s(\stackrel{}{r}_s,t_s)\psi _n(\stackrel{}{r}_s,t_s)G_{\nu _j}(\stackrel{}{r}_\nu \stackrel{}{r}_s,t_\nu t_s)}`$ (17)
In principle one even does not need to know the concrete form of $`\psi _n`$, the only necessary property of these functions is the condition that they form a complete set:
$`{\displaystyle \underset{n}{}}\psi _n(\stackrel{}{r},t)\psi _n^{}(\stackrel{}{r}^{},t)=\delta \left(\stackrel{}{r}\stackrel{}{r}^{}\right)`$ (18)
However in what follows we will use for simplicity the eigenfunctions of momentum, $`\psi _n\mathrm{exp}(i\stackrel{}{p}\stackrel{}{r}iEt)`$.
For the subsequent calculations we need the following representation of the Green’s function which is obtained by the sequence of integration:
$`G(\stackrel{}{r},t)`$ $`=`$ $`{\displaystyle \frac{d^4p_4}{p_4^2m^2}e^{ip_4x}}=`$ (19)
$`2\pi {\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega e^{i\omega t}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dpp^2}{\omega ^2p^2m^2}}{\displaystyle _1^1}𝑑\zeta e^{ipr\zeta }=`$
$`{\displaystyle \frac{i\pi }{r}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega e^{i\omega t}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{dpp^2}{\omega ^2p^2m^2}}\left(e^{ipr}e^{ipr}\right)`$
Here $`\omega `$ and $`p`$ are respectively the fourth and space components of the 4-vector $`p_4`$. We have omitted spin matrices because the final result is essentially independent of that. The integration over $`dp`$ was extended over the whole axis (from $`\mathrm{}`$ to $`+\mathrm{}`$) because the integrand is an even function of $`p`$. This permit to calculate this integral by taking residues in the poles on mass shell: $`p=\pm \sqrt{\omega ^2m^2+iϵ}`$. Both poles give the same contribution, so skipping unnecessary numerical coefficients, we finally obtain:
$`G(\stackrel{}{r},t)={\displaystyle \frac{1}{r}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega e^{i\omega t+i\sqrt{\omega ^2m^2}r}`$ (20)
As a source function $`\psi _s`$ we will take essentially expression (15) but assume that the source is a decaying particle with the decay width $`\gamma `$, that was born at the moment $`t=0`$. It corresponds to multiplication of $`\psi _s`$ by $`\theta (t)\mathrm{exp}(\gamma t)`$. Taking all together we obtain the following expression for the amplitude:
$`A_j^{(n)}(\stackrel{}{r},t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t_s{\displaystyle \frac{d\stackrel{}{r}_s}{|\stackrel{}{r}\stackrel{}{r}_s|}𝑑\stackrel{}{p}C\left(\stackrel{}{p}\stackrel{}{p}_0\right)e^{iEt_si\stackrel{}{p}\stackrel{}{r}_s\gamma t/2}e^{iE_nt_s+i\stackrel{}{p}_n\stackrel{}{r}_s}}`$ (21)
$`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega _je^{i\omega _j(tt_s)ik_j|\stackrel{}{r}\stackrel{}{r}_s|}`$
Integrals over $`d\stackrel{}{r}_s`$ and $`d\stackrel{}{p}`$ are taken over all infinite space. It is worthwhile to remind here that all the momenta are on mass shell, $`E^2=p^2+m_\pi ^2`$ (we assumed that the source is a decaying pion) and $`\omega _j^2=k_j^2+m_j^2`$, where $`m_j`$ is the mass of $`j`$-th neutrino mass eigenstate.
The integration over $`dt_s`$ is trivial and gives the factor $`\left(EE_n\omega _j+i\gamma /2\right)^1`$. The integration over $`d\stackrel{}{r}_s`$ can be easily done if the registration point is far from the source. In this case it is accurate enough to take $`1/|\stackrel{}{r}\stackrel{}{r}_s|1/r`$, while the same quantity in the exponent should be expanded up to the first order:
$`|\stackrel{}{r}\stackrel{}{r}_s|r\stackrel{}{\xi }\stackrel{}{r}_s`$ (22)
where $`\stackrel{}{\xi }=\stackrel{}{r}/r`$ is a unit vector directed from the center of the source taken at the initial moment $`t=0`$ to the detector at the point $`\stackrel{}{r}`$. In this limit the integral over $`d\stackrel{}{r}_s`$ gives $`\delta \left(\stackrel{}{p}\stackrel{}{p}_n\stackrel{}{k}_j\right)`$, ensuring momentum conservation:
$`\stackrel{}{p}=\stackrel{}{p}_n+\stackrel{}{k}_j\stackrel{}{p}_{\pi ,j}`$ (23)
The vector of neutrino momentum is formally defined as
$`\stackrel{}{k}_j=\stackrel{}{\xi }k_j=\stackrel{}{\xi }\sqrt{\omega _j^2m_j^2}`$ (24)
Ultimately we are left with the integral:
$`A_j^{(n)}={\displaystyle \frac{1}{r}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\omega _jC\left(\stackrel{}{p}_n+\stackrel{}{k}_j\stackrel{}{p}_0\right){\displaystyle \frac{e^{i\omega _jtik_jr}}{E_{\pi ,j}E_n\omega _j+i\gamma /2}}`$ (25)
where $`E_{\pi ,j}=\sqrt{(\stackrel{}{p}_n+\stackrel{}{k}_j)^2+m_\pi ^2}`$. This integral can be taken in the ’pole approximation’ and to do that we need to expand the integrand near the energy conservation law, see below eq. (28), as follows. The neutrino energy is presented as $`\omega _j=\omega _j^{(0)}+\mathrm{\Delta }\omega _j`$. To avoid confusion one should distinguish between the deviation of neutrino energy from the central value given by the conservation law, $`\mathrm{\Delta }\omega _j`$, from the difference of energies of different neutrino mass eigenstates, $`\delta \omega =\omega _1\omega _2`$. The neutrino momentum is expanded up to the first order in $`\mathrm{\Delta }\omega _j`$:
$`k_j=\sqrt{\omega _j^2m_j^2}k_j^{(0)}+\mathrm{\Delta }\omega _j/V_j^{(\nu )}`$ (26)
where $`V_j^{(\nu )}=k_j^{(0)}/\omega _j^{(0)}`$ is the velocity of $`j`$-th neutrino. The pion energy is determined by the momentum conservation (23) and is given by
$`E_{\pi ,j}=\sqrt{\left(\stackrel{}{p}_n+\stackrel{}{k}_j^{(0)}+\stackrel{}{\mathrm{\Delta }}k_j\right)^2+m_\pi ^2}E_{\pi ,j}^{(0)}+\stackrel{}{V}_{\pi ,j}\stackrel{}{\xi }\mathrm{\Delta }\omega _j`$ (27)
where the pion velocity is $`\stackrel{}{V}_j^{(\pi )}=(\stackrel{}{p}_n+\stackrel{}{k}_j^{(0)})/E_{\pi ,j}^{(0)}`$. The neutrino energy, $`\omega _j^{(0)}`$, satisfying the conservation law is defined from the equation:
$`E_{\pi ,j}^{(0)}E_n\omega _j^{(0)}=0`$ (28)
Now the integral over $`\omega _j`$ is reduced to
$`A_j^{(n)}={\displaystyle \frac{e^{i\omega _j^{(0)}tik_j^{(0)}r}}{r}}C\left(\stackrel{}{p}_n+\stackrel{}{k}_j^{(0)}\stackrel{}{p}_0\right)`$
$`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\mathrm{\Delta }\omega _j{\displaystyle \frac{e^{i\mathrm{\Delta }\omega _j\left(tr/V_j^\nu \right)}}{\left(\stackrel{}{V}_j^{(\pi )}\stackrel{}{\xi }/V_j^{(\nu )}\right)\mathrm{\Delta }\omega _j\mathrm{\Delta }\omega _j+i\gamma /2}}`$ (29)
The last integral vanishes if $`t<V_j^{(\nu )}r`$, while in the opposite case it can be taken as the residue in the pole and we finally obtain:
$`A_j^{(n)}={\displaystyle \frac{C\left(\stackrel{}{p}_{\pi ,j}\stackrel{}{p}_0\right)}{r}}\theta \left(rV_j^{(\nu )}t\right)\mathrm{exp}\left(i\omega ^{(0)}tik_j^{(0)}r{\displaystyle \frac{\gamma }{2}}{\displaystyle \frac{V_j^{(\nu )}tr}{V_j^{(\nu )}\stackrel{}{V}_j^{(\pi )}\stackrel{}{\xi }}}\right)`$ (30)
We obtained a neutrino wave packet moving with the velocity $`V_j^{(\nu )}`$ with a well defined front (given by the theta-function) and decaying with time in accordance with the decay law of the source. The similar wave packet, but moving with a slightly different velocity describes another oscillating state $`\nu _i`$. It is evident from this expressions that phenomenon of coherent oscillations takes place only if the packets overlap, as was noticed long ago .
The probability of registration of oscillating neutrinos at the space-time point $`(\stackrel{}{r},t)`$ is given by the density matrix
$`\rho _{ij}={\displaystyle 𝑑\stackrel{}{p}_nA_i^{(n)}(\stackrel{}{r}_\nu ,t_\nu )A_j^{(n)}(\stackrel{}{r}_\nu ,t_\nu )}`$ (31)
The oscillating part of the probability is determined by the phase difference (11) but now the quantities $`\delta \omega `$ and $`\delta k`$ are unambiguously defined. To this end we will use the conservation laws (23,28). They give:
$`\delta \omega ^{(0)}=\delta E_\pi \mathrm{and}\delta \stackrel{}{k}^{(0)}=\delta \stackrel{}{p}_\pi `$ (32)
The variation of neutrino energy is given by
$`\delta \omega =V^{(\nu )}\delta k+\delta m^2/2\omega `$ (33)
while the variation of the pion energy can be found from expression (27):
$`\delta E^{(\pi )}=\stackrel{}{V}^{(\pi )}\delta \stackrel{}{k}`$ (34)
From these equations follows
$`\delta \omega ={\displaystyle \frac{\delta m^2}{2\omega }}{\displaystyle \frac{\stackrel{}{V}^{(\pi )}\stackrel{}{\xi }}{V^{(\nu )}\stackrel{}{V}^{(\pi )}\stackrel{}{\xi }}}\mathrm{and}\delta k={\displaystyle \frac{\delta m^2}{2\omega }}{\displaystyle \frac{1}{V^{(\nu )}\stackrel{}{V}^{(\pi )}\stackrel{}{\xi }}}`$ (35)
One sees that generally both $`\delta \omega `$ and $`\delta k`$ are non-vanishing. Only in the case of pion decay at rest, $`\delta \omega =0`$ but $`\delta k`$ is non-zero in any case. Substituting the obtained results into expression (11) for the phase difference we come to the standard expression (14) if $`V_\pi =0`$. This result shows a remarkable stability with respect to assumptions made in its derivation. However if the pion is moving, then the oscillation phase contains an extra term
$`\delta \mathrm{\Phi }={\displaystyle \frac{\delta m^2}{2\omega }}{\displaystyle \frac{\stackrel{}{\xi }\left(\stackrel{}{r}\stackrel{}{V}^{(\pi )}t\right)}{V^{(\nu )}\stackrel{}{\xi }\stackrel{}{V}^{(\pi )}}}={\displaystyle \frac{r\delta m^2}{2k}}+{\displaystyle \frac{(\stackrel{}{\xi }\stackrel{}{V}^{(\pi )})(rV^{(\nu )}t)}{V^{(\nu )}\stackrel{}{\xi }\stackrel{}{V}^{(\pi )}}}`$ (36)
This extra term would lead to a suppression of oscillation after averaging over time. This suppression is related to the motion of the source and reflects the uncertainty in the position of pion at the moment of decay. So this result can be understood in the frameworks of the standard naive approach.
Similar expression can be derived for the case when both neutrino and muon from the decay $`\pi \mu +\nu _j`$ are registered in the space-time points $`\stackrel{}{r}_\nu ,t_\nu `$ and $`\stackrel{}{r}_\mu ,t_\mu `$ respectively. This case was considered in refs. . Here we will use the same approach as described above when the muon is not registered. The only difference is that in eq. (17) for the oscillation amplitude we have to substitute the Green’s function of muon $`G_\mu (\stackrel{}{r}_\mu r_s,t_\mu t_s)`$ instead of $`\psi _n(\stackrel{}{r}_s,t_s)`$. The calculations are essentially the same and after some algebra the following expression for the oscillation amplitude is obtained:
$`A_{\mu ,\nu }`$ $``$ $`{\displaystyle \frac{V_\mu V_\nu }{r_\mu r_\nu }}\mathrm{\Theta }(L_\mu +L_\nu )\mathrm{exp}\left[{\displaystyle \frac{\gamma (L_\mu +L_\nu )}{2(V_\mu +V_\nu )}}\right]\stackrel{~}{C}(V_\mu L_\nu V_\nu L_\mu )`$ (37)
$`\mathrm{exp}\left[i\left(k_\mu ^{(0)}r_\mu +k_\nu ^{(0)}r_\nu E_\mu ^{(0)}t_\mu E_\nu ^{(0)}t_\nu \right)\right]`$
where $`L=Vtr`$. Each kinematic variable depends upon the neutrino state $`j`$, so they should contain sub-index $`j`$. The upper indices “0” mean that these momenta and energies are the central values of the corresponding wave packets, so that the classical relation $`\stackrel{}{k}^{(0}=\stackrel{}{V}E^{(0)}`$ holds for them. Here the direction of momenta are defined as above along the vector indicated to the observation point. However, the kinematics in this case is different from the previous one and the change of the energy and momentum of each particle for reactions with different sorts of neutrinos are related through the equations:
$`\delta E_\mu +\delta E_\nu =0\mathrm{and}\delta \stackrel{}{k}_\mu +\delta \stackrel{}{k}_\nu =0.`$ (38)
For the central values of momenta the following relations are evidently true, $`\delta k_\nu =\delta E_\nu \delta m^2/2k_\nu `$, and the similar one for the muon without the last term proportional to the mass difference. Correspondingly the phase of oscillation is given by the expression
$`\delta \mathrm{\Phi }=\delta k_\mu r_\mu \delta E_\mu t_\mu +\delta k_\nu r_\nu \delta E_\nu t_\nu =`$
$`\delta E_\nu \left({\displaystyle \frac{r_\mu V_\mu t_\mu }{V_\mu }}{\displaystyle \frac{r_\nu V_\nu t_\nu }{V_\nu }}\right){\displaystyle \frac{\delta m^2}{2k_\nu }}r_\nu `$ (39)
The first two terms in the phase are proportional to the argument of $`\stackrel{~}{C}`$ in eq. (37) and thus they give the contribution equal to the size of the the source, i.e. to the wave packet of the initial pion. If the latter is small (as usually the case) we obtain again the standard expression for the oscillation phase.
## 2 Matter effects
Despite extremely weak interactions of neutrinos, matter may have a significant influence on the oscillations if/because the mass difference between the propagating eigenstates is very small. Description of neutrino oscillations in matter was first done in ref. . Somewhat later a very important effect of resonance neutrino conversion was discovered , when even with a very small vacuum mixing angle, mixing in medium could reach the maximal value.
Hamiltonian of free neutrinos in the mass eigenstate basis has the form:
$`_m^{(1,2)}=\left(\begin{array}{cc}E_1& 0\\ 0& E_2\end{array}\right)`$ (42)
where $`E_j=\sqrt{p^2+m_j^2}`$. In the interaction basis $`_m`$ is rotated by the matrix (5):
$`_m^{(a,b)}=U_m^{(1,2)}U^1=\left(\begin{array}{cc}\mathrm{cos}^2\theta E_1+\mathrm{sin}^2\theta E_2& g\mathrm{sin}\theta \mathrm{cos}\theta \\ g\mathrm{sin}\theta \mathrm{cos}\theta & \mathrm{sin}^2\theta E_1+\mathrm{cos}^2\theta E_2\end{array}\right)`$ (45)
Here $`g=\delta m^2/2E`$ and we returned to the more usual notation $`E`$ for neutrino energy.
The interaction Hamiltonian is diagonal in the interaction basis and if only first order effects in the Fermi coupling constant, $`G_F`$, are taken into account, then the Hamiltonian can be expressed through refraction index, $`n_a`$, of flavor $`a`$-neutrino in the medium (recall that the deviation of refraction index from unity is proportional to the forward scattering amplitude and thus contains $`G_F`$ in the first power):
$`H_{int}^{(a,b)}=\left(\begin{array}{cc}E(n_a1)& 0\\ 0& E(n_b1)\end{array}\right)`$ (48)
where a small difference between $`E_1`$ and $`E_2`$ in front of small factors $`(n1)`$ was neglected.
Thus, up to a unit matrix, the total Hamiltonian in the interaction basis can be written as
$`H_{tot}^{(a,b)}=\left(\begin{array}{cc}f& g\mathrm{sin}2\theta /2\\ g\mathrm{sin}2\theta /2& 0\end{array}\right)`$ (51)
where $`f=g\mathrm{cos}2\theta +E\delta n`$ and $`\delta n=n_an_b`$. This matrix is easy to diagonalize. Its eigenvalues are
$`\lambda _{1,2}={\displaystyle \frac{f\pm \sqrt{f^2+g^2\mathrm{sin}^22\theta }}{2}}`$ (52)
and the eigenstates in matter (up to normalization factor) are
$`|\nu _{1,2}=|\nu _a+{\displaystyle \frac{g\mathrm{sin}2\theta }{f\pm \sqrt{f^2+g^2\mathrm{sin}^22\theta }}}|\nu _b`$ (53)
Refraction index may change with time, as happens in cosmology, or with space point, if neutrinos propagate in inhomogeneous medium, for example from the center of the Sun to its surface. If somewhere (or sometime) $`f`$ vanishes then the resonance transition of one neutrino species to another is possible . Indeed let us assume that $`\nu _e`$ and $`\nu _\mu `$ are mixed with a small vacuum mixing angle $`\theta `$ and that initially an electronic neutrino was produced in vacuum. So the initial propagating state would be mostly $`\nu _e`$:
$`|\nu _1_{in}=|\nu _e+(1/2)\mathrm{tan}2\theta |\nu _\mu `$ (54)
After propagation in the media where the function $`f`$ changes sign passing through zero, the propagating state would become mostly $`\nu _\mu `$:
$`|\nu _1_{fin}=|\nu _e(2/\mathrm{tan}2\theta )|\nu _\mu `$ (55)
This effect of resonance conversion may play an important role in the resolution of the solar neutrino problem and in cosmology.
## 3 Neutrino oscillations in cosmology
### 3.1 A brief (and non-complete) review
Neutrino oscillations in the primeval plasma is significantly different from e.g. solar neutrino oscillations in the following two important aspects. First, cosmological plasma is almost charge symmetric. A relative excess of any particles over antiparticles is believed to be at the level $`10^910^{10}`$, while in stellar material the asymmetry is of order unity. Neutrino oscillations in the early universe may change the magnitude of asymmetry in the sector of active neutrinos. On the other hand, the asymmetry has a strong influence on the oscillations through the refraction index of the primeval plasma (see below). It leads to a strong non-linearity of the problem and makes calculations quite complicated. Second important point is that neutrino mean free path in the early universe is quite small at high temperatures and hence breaking of coherence becomes essential. Because of that one cannot use wave functions for description of oscillations and should turn to the density matrix formalism. It leads to a great complexity of equations. Kinetic equations for density matrix with the account of neutrino scattering and annihilation were derived in the papers . In ref. , where the impact of neutrino oscillations on big bang nucleosynthesis (BBN) was first considered, only the second order effects, proportional to $`G_F^2`$, were taken into account, while the deviation of refraction index from unity was neglected. This approximation is valid for a sufficiently high $`\delta m^2`$. In a subsequent paper implications for BBN of possible CP-violating effects in oscillations were discussed. Earlier works on neutrino oscillations also include refs. . The calculations of refraction index in cosmological plasma were performed in ref. . The results of this work permitted to make more accurate calculations of the role played by neutrino oscillations in BBN <sup>-</sup>.
It was noticed in ref. that the oscillations between an active and sterile neutrinos could generate an exponential rise of lepton asymmetry in the sector of active neutrinos. The origin of this instability is the following. Since lepton asymmetry comes with the opposite sign to refraction indices of neutrinos and antineutrinos (see below section 3.2), it may happen that the transformation of antineutrinos to their sterile partners would proceed faster than similar transformation of neutrinos, especially if resonance conditions are fulfilled. It would lead to an increase of the asymmetry and through the refraction index to more favorable conditions for its rise. However it was concluded that the rise is not significant and the effect of the generated asymmetry on BBN is small. This conclusion was reconsidered in ref. where it was argued that asymmetry generated by this mechanism could reach very large values, close to unity, and this effect, in accordance to the earlier paper of the same group , would have a significant influence on primordial abundances. This result attracted a great attention and was confirmed in several subsequent publications <sup>-</sup>. Moreover, some works showed not only rising and large asymmetry but also a chaotic behavior of its sign <sup>-</sup>.
Due to complexity of equations some simplifying approximations were made in their solutions. First accurate numerical solution of (almost) exact equations were done in refs. <sup>-</sup>. The most essential approximation was that the lost of coherence was described by the term $`\gamma (\rho _{eq}\rho )`$ instead of the exact two-dimensional collision integral taken from the square of the scattering amplitude over momenta. In fact this approximation was used in all the papers. The calculations of refs. <sup>-</sup> were performed for rather small mass difference, $`\delta m^2<10^7\mathrm{eV}^2`$. Neither chaoticity, nor a considerable rise of the asymmetry were found. For a larger mass difference a strong numerical instability was observed. However this result is not in a contradiction with other papers because the latter found the above mentioned effects for much larger mass differences. In our work we tried to extend the range of validity of direct numerical calculations to large $`\delta m^2`$. We were able to proceed only up to $`\delta m^210^6`$. For higher values we did not find a way to avoid numerical instability. To overcome the problem we analytically transformed kinetic equations to a simpler form that permitted a stable numerical integration. To this end an expansion in terms of a large parameter that corresponds to a high frequency of oscillations was used. This method is rather similar to the well known separation of fast and slow variables in differential equations. According to our results the asymmetry may rise, but only by 5-6 orders of magnitude (from initial $`10^{10}`$), i.e. 4-5 orders of magnitude weaker than was obtained in other papers . However in agreement with these papers no chaoticity was found. Our results were criticized in ref. where it was argued that lepton asymmetry generated by oscillations must be chaotic. However the author misunderstood our calculations and criticized us for the things that we never did and, second, the approximation used in that paper , namely calculations in terms of one fixed value of neutrino momentum that was chosen to be equal to the thermally averaged one, is intrinsically inappropriate for the solution of the problem of chaoticity. In that approximation even neutrino oscillations in vacuum would result in chaotic lepton asymmetry.
Thus, most probably a chaotic amplification of the lepton asymmetry does not take place, however the exact magnitude of the amplification remains uncertain. As discussed in ref. their recent calculations are exact and do not suffer from numerical instability. On the other hand, we do not see any shortcomings of our semi-analytical approach , all approximations are well under control and numerical part of calculations is quite simple and stable. More work is necessary to resolve the contradictions both in magnitude and possible chaoticity.
### 3.2 Refraction index
In this section we derive the Schroedinger equation for neutrino wave function in the primeval plasma. We will start with the neutrino quantum operator $`\nu _a(x)`$ of flavor $`a`$ that satisfies the usual Heisenberg equation of motion:
$`(i/)\nu _a(x)+{\displaystyle \frac{g}{2\sqrt{2}}}\delta _{ae}W_\alpha (x)O_\alpha ^{(+)}e(x)+{\displaystyle \frac{g}{4\mathrm{cos}\theta _W}}Z_\alpha (x)O_\alpha ^{(+)}\nu _a(x)=0`$ (56)
where $``$ is the neutrino mass matrix, $`W(x)`$, $`Z(x)`$ and $`e(x)`$ are respectively the quantum operators of intermediate bosons and electrons, and $`O_\alpha ^\pm =\gamma _\alpha \left(1\pm \gamma _5\right)`$. We assumed that the temperature of the plasma is in MeV range and thus only electrons are present in the plasma.
Equations of motion for the field operators of $`W`$ and $`Z`$ bosons have the form
$`G_{W,\alpha \beta }^1W_\beta (x)={\displaystyle \frac{g}{2\sqrt{2}}}\overline{\nu }_a(x)O_\alpha ^{(+)}\nu _a(x),`$ (57)
$`G_{Z,\alpha \beta }^1Z_\beta (x)={\displaystyle \frac{g}{4\mathrm{cos}\theta _W}}[\overline{\nu }_a(x)O_\alpha ^{(+)}\nu _a(x)+`$
$`(2\mathrm{sin}^2\theta _W1)\overline{e}(x)O_\alpha ^{(+)}e(x)+2\mathrm{sin}^2\theta _W\overline{e}(x)O_\alpha ^{()}e(x)]`$ (58)
where the differential operators $`G_{W,Z}^1`$ are inverse Green’s functions of $`W`$ and $`Z`$. In momentum representation they can be written as
$`G_{\alpha \beta }={\displaystyle \frac{g_{\alpha \beta }q_\alpha q_\beta /m^2}{m^2q^2}}`$ (59)
It can be be shown that the term $`q_\alpha q_\beta /m^2`$ gives contribution proportional to lepton masses and can be neglected.
In the limit of small momenta, $`qm_{W,Z}`$, equation (57) can be solved as
$`W_\alpha (x)={\displaystyle \frac{g}{2\sqrt{2}m_W^2}}\left(1{\displaystyle \frac{^2}{m_W^2}}\right)\left(\overline{e}(x)O_\alpha ^{(+)}\nu _e(x)\right)`$ (60)
A similar expression with an evident substitution for the r.h.s. can be obtained for $`Z_\alpha (x)`$. These expressions should be substituted into eq. (56) to obtain equation that contains only field operators of leptons, $`\nu _a(x)`$ and $`e(x)`$:
$`(i/)\nu _a(x)={\displaystyle \frac{G_F}{\sqrt{2}}}\{\delta _{ae}\left[(1{\displaystyle \frac{^2}{m_W^2}})\left(\overline{e}(x)O_\alpha ^{(+)}\nu _e(x)\right)\right]O_\alpha ^{(+)}e(x)+`$
$`{\displaystyle \frac{1}{2}}[(1{\displaystyle \frac{^2}{m_Z^2}})(\overline{\nu }_b(x)O_\alpha ^{(+)}\nu _b(x)+(2\mathrm{sin}^2\theta _W1)\overline{e}(x)O_\alpha ^{(+)}e(x)+`$
$`2\mathrm{sin}^2\theta _W\overline{e}(x)O_\alpha ^{()}e(x))]O_\alpha ^{(+)}\nu _a(x)\}`$ (61)
Neutrino wave function in the medium is defined as
$`\mathrm{\Psi }_a(x)=A|\nu _a(x)|A+\nu ^{(k)}`$ (62)
where $`A`$ describes the state of the medium and $`\nu ^{(k)}`$ is a certain one-neutrino state, specified by quantum numbers $`k`$, e.g. neutrino with momentum $`\stackrel{}{k}`$. Equation of motion for this wave function can be found from eq. (61) after averaging over medium. The theory is quantized perturbatively in the standard way. We define the free neutrino operator $`\nu _a^{(0)}`$ that satisfies the equation of motion:
$`(i/)\nu _a^{(0)}(x)=0`$ (63)
This operator is expanded as usually in terms of creation-annihilation operators:
$`\nu ^{(0)}(x)={\displaystyle \frac{d^3k}{(2\pi )^3\sqrt{2E_k}}\underset{s}{}\left(a_k^su^s(k)e^{ikx}+b_k^sv^s(k)e^{ikx}\right)}`$ (64)
and one-particle state is defined as $`|\nu ^{(k)}=a_k^{}|\mathrm{vac}`$.
The equation of motion for the neutrino wave function $`\mathrm{\Psi }_a(x)`$ can be obtained from expression (62) perturbatively by applying the operator $`(i/)`$ and using eq. (61) with free neutrino operators $`\nu _a^{(0)}`$ in the r.h.s. After some algebra which mostly consists in using equations of motion for the free fermion operators and (anti)commutation relations between the creation/annihilation operators, one would obtain the equation of the form:
$`i_t\mathrm{\Psi }(t)=\left(_m+V_{eff}\right)\mathrm{\Psi }`$ (65)
where $`_m`$ is the free Hamiltonian; in the mass eigenstate basis it has the form $`_0=\mathrm{diag}\left[\sqrt{p^2+m_j^2}\right]`$. The matrix-potential $`V_{eff}`$ describes interactions of neutrinos with media and is diagonal in the flavor basis. Up to the factor $`E`$ (i.e. neutrino energy) it is essentially the refraction index of neutrino in the medium. The potential contains two terms. The first one comes from the averaging of the external current $`J\overline{l}O_\alpha l`$. Due to homogeneity and isotropy of the plasma only its time component is non-vanishing and proportional to the charge asymmetry (i.e. to the excess of particles over antiparticles) in the plasma. This term has different signs for neutrinos and antineutrinos. However the interactions of neutrinos with the medium is not always of the (current)$`\times `$(current) form due to non-locality related to the exchange of $`W`$ or $`Z`$ bosons. If incoming and outgoing neutrinos interact in different space-time points, the interaction with the medium cannot be written as an interaction with the external current. The contribution of such terms is inversely proportional to $`m_{W,Z}^2`$ but formally it is of the first order in $`G_F`$. With these two types of contributions the diagonal matrix elements of the effective potential for the neutrino of flavor $`a`$ has the form:
$`V_{eff}^a=\pm C_1\eta G_FT^3+C_2^a{\displaystyle \frac{G_F^2T^4E}{\alpha }},`$ (66)
where $`E`$ is the neutrino energy, $`T`$ is the temperature of the plasma, $`G_F=1.16610^5`$ GeV<sup>-2</sup> is the Fermi coupling constant, $`\alpha =1/137`$ is the fine structure constant, and the signs “$`\pm `$” refer to anti-neutrinos and neutrinos respectively (this choice of sign describes the helicity state, negative for $`\nu `$ and positive for $`\overline{\nu }`$). According to ref. the coefficients $`C_j`$ are: $`C_10.95`$, $`C_2^e0.61`$ and $`C_2^{\mu ,\tau }0.17`$. These values are true in the limit of thermal equilibrium, but otherwise these coefficients are some integrals from the distribution functions over momenta. For oscillating neutrinos deviations from thermal equilibrium could be significant and in this case expression (66) should be modified. However it is technically rather difficult to take this effect into account in numerical calculations and the simplified version (66) is used.
### 3.3 Loss of Coherence and Density Matrix
Breaking of coherence appears in the second order in the Fermi coupling constant $`G_F`$, so that equations of motion for the operators of all leptonic fields (including electrons) should be solved up to the second order in $`G_F`$. Since the calculations are quite lengthy, we only sketch the derivation here. In the considered approximation the lepton operators $`l(x)`$, where $`l`$ stands for neutrino or electron, in the r.h.s. of eq. (61) should be expanded up to the first order in $`G_F`$. The corresponding expressions can be obtained from the formal solution of eq. (61) up to first order in $`G_F`$. Their typical form is the following:
$`l=l_0+G_l(\mathrm{r}.\mathrm{h}.\mathrm{s}._0)`$ (67)
where the matrix (in neutrino space) $`G_l`$ is the Green’s function of the corresponding lepton and $`\mathrm{r}.\mathrm{h}.\mathrm{s}._0`$ is the right hand side of eq. (61) taken in the lowest order in $`G_F`$, i.e with lepton operators taken in the zeroth order, $`l=l_0`$. Expression (67) should be substituted back into eq. (61) and this defines the r.h.s. up to the second order in $`G_F`$ in terms of free lepton operators $`l_0`$. Of course in the second approximation we neglect the non-local terms, $`1/m_{W,Z}^2`$.
Now we can derive kinetic equation for the density matrix of neutrinos, $`\widehat{\rho }_j^i=\nu ^i\nu _j^{}`$, where over-hut indicates that $`\widehat{\rho }`$ is a quantum operator. The $`C`$-valued density matrix is obtained from it by taking matrix element over the medium, $`\rho =\widehat{\rho }`$. We should apply to it the differential operator $`(i/)`$ and use eq. (61). The calculations of matrix elements of the free lepton operators $`l_0`$ are straightforward and can be achieved by using the standard commutation relations. There is an important difference between equations for the density matrix and the wave function. The latter contains only terms proportional to the wave function, $`i_t\mathrm{\Psi }=\mathrm{\Psi }`$, while equation for density matrix contains source term that does not vanish when $`\rho =0`$. Neutrino production or destruction is described by the imaginary part of the effective Hamiltonian. The latter is not hermitian because the system is not closed. By the optical theorem the imaginary part of the Hamiltonian is expressed through the cross-section of neutrino creation or annihilation. Such terms in kinetic equation for the density matrix are similar to the “normal” kinetic equation for the distribution functions:
$`{\displaystyle \frac{df_1}{dt}}=I_{coll}`$ (68)
where the collision integral is the integral over the proper phase space from the following combination of the distribution functions
$`F=f_1f_2(1f_3)(1f_4)+f_4f_3(1f_1)(1f_2)`$ (69)
A similar combination appears for the case of oscillating neutrinos but with a rather complicated matrix structure. For example there can be the contribution of the form of the anti-commutator:
$`\{\rho _1,g(1\rho _3)g\rho _2(1\rho _4)g\}+\left(\mathrm{inverse}\mathrm{reaction}\right)`$ (70)
and a few more of different structure that results in a very lengthy expression. The matrix $`g`$ describes neutrino interactions. It is diagonal in the flavor basis and has entries proportional to matrix elements squared of the relevant reactions. The Fermi-blocking factors $`(1\rho )`$ appear when one takes matrix elements of neutrino operators over the medium in which neutrino occupation numbers may be non-zero.
We will not present here the complete form of the equation. It can be found e.g. in the paper . Moreover in all the applications a “poor man” substitution has been done: all terms describing neutrino production or destruction were mimicked by
$`\mathrm{\Gamma }\left(\rho \rho _{eq}\right)`$ (71)
where $`\rho _{eq}`$ is the equilibrium value of the density matrix, i.e. the unit matrix multiplied by the equilibrium distribution function
$`f_{eq}=\left[\mathrm{exp}(E/T\xi )+1\right]`$ (72)
and the coefficient $`\mathrm{\Gamma }`$ is the reaction rate (see below).
### 3.4 Oscillations and lepton asymmetry
As we have already mentioned oscillations between active and sterile neutrinos may induce a significant lepton asymmetry in the sector of active neutrinos. Now we will consider this phenomenon in some more detail. Basic equations governing the evolution of the density matrix are:
$`i(_tHp_p)\rho _{aa}`$ $`=`$ $`F_0(\rho _{sa}\rho _{as})/2i\mathrm{\Gamma }_0(\rho _{aa}f_{eq}),`$ (73)
$`i(_tHp_p)\rho _{ss}`$ $`=`$ $`F_0(\rho _{sa}\rho _{as})/2,`$ (74)
$`i(_tHp_p)\rho _{as}`$ $`=`$ $`W_0\rho _{as}+F_0(\rho _{ss}\rho _{aa})/2i\mathrm{\Gamma }_1\rho _{as},`$ (75)
$`i(_tHp_p)\rho _{sa}`$ $`=`$ $`W_0\rho _{sa}F_0(\rho _{ss}\rho _{aa})/2i\mathrm{\Gamma }_1\rho _{sa},`$ (76)
where $`a`$ and $`s`$ mean “active” and “sterile” respectively, $`F_0=\delta m^2\mathrm{sin}2\theta /2E`$, $`W_0=\delta m^2\mathrm{cos}2\theta /2E+V_{eff}^a`$, $`H=\sqrt{8\pi \rho _{tot}/3M_p^2}`$ is the Hubble parameter, $`p`$ is the neutrino momentum. The antineutrino density matrix satisfies the similar set of equations with the opposite sign of the antisymmetric term in $`V_{eff}^a`$ and with a slightly different damping factor $`\overline{\gamma }`$ (this difference is proportional to the lepton asymmetry in the primeval plasma).
Equations (73-76) account exactly for the first order terms described by the refraction index, while the second order terms describing breaking of coherence are approximately modeled by the damping coefficients $`\mathrm{\Gamma }_j`$ in accordance with eq. (71). The latter are equal to :
$`\mathrm{\Gamma }_0=2\mathrm{\Gamma }_1=g_a{\displaystyle \frac{180\zeta (3)}{7\pi ^4}}G_F^2T^4p.`$ (77)
In general the coefficient $`g_a(p)`$ is a momentum-dependent function, but in the approximation of neglecting $`[1f]`$ factors in the collision integral it becomes a constant equal respectively to $`g_{\nu _e}4`$ and $`g_{\nu _\mu ,\mu _\tau }2.9`$ . In the following we will use more accurate values found from the thermal average of the complete electro-weak rates (with factors $`[1f]`$ included), which we calculated numerically from our Standard Model code . This gives us $`g_{\nu _e}3.56`$ and $`g_{\nu _\mu ,\mu _\tau }2.5`$.
It is convenient to introduce new variables
$`x=m_0R(t)\mathrm{and}y=pR(t)`$ (78)
where $`R(t)`$ is the cosmological scale factor so that $`H=\dot{R}/R`$ and $`m_0`$ is an arbitrary mass (just normalization), we choose $`m_0=1`$ MeV. In the approximation that we will work, we assume that $`\dot{T}=HT`$, so that we can take $`R=1/T`$. In terms of these variables the differential operator $`(_tHp_p)`$ transforms to $`Hx_x`$. We will normalize the density matrix elements to the equilibrium function $`f_{eq}`$:
$`\rho _{aa}`$ $`=`$ $`f_{eq}(y)[1+a(x,y)],\rho _{ss}=f_{eq}(y)[1+s(x,y)],`$ (79)
$`\rho _{as}`$ $`=`$ $`\rho _{sa}^{}=f_{eq}(y)[h(x,y)+il(x,y)],`$ (80)
and express the neutrino mass difference $`\delta m^2`$ in eV<sup>2</sup>.
As the next step we will take the sum and difference of eqs. (73)-(76) for $`\nu `$ and $`\overline{\nu }`$. The corresponding equations have the following form:
$`s_\pm ^{}`$ $`=`$ $`Fl_\pm ,`$ (81)
$`a_\pm ^{}`$ $`=`$ $`Fl_\pm 2\gamma _+a_\pm 2\gamma _{}a_{},`$ (82)
$`h_\pm ^{}`$ $`=`$ $`Ul_\pm VZl_{}\gamma _+h_\pm \gamma _{}h_{},`$ (83)
$`l_\pm ^{}`$ $`=`$ $`{\displaystyle \frac{F}{2}}(a_\pm s_\pm )Uh_\pm +VZh_{}\gamma _+l_\pm \gamma _{}l_{},`$ (84)
where $`a_\pm =(a\pm \overline{a})/2`$ etc, and the prime means differentiation with respect to $`x`$. We have used $`W=U\pm VZ`$, $`\gamma =\mathrm{\Gamma }_1/Hx`$, and $`\gamma _\pm =(\gamma \pm \overline{\gamma })/2`$, where $`\gamma _{}`$ parameterizes the difference of interaction rates between neutrino and anti-neutrinos, which is proportional to the neutrino asymmetry. With the approximation $`\rho _{tot}10.75\pi ^2T^4/30`$, the expressions for $`U`$, $`V`$, and $`Z`$ become:
$`U`$ $`=`$ $`1.1210^9\mathrm{cos}2\theta \delta m^2{\displaystyle \frac{x^2}{y}}+26.2{\displaystyle \frac{y}{x^4}},`$ (85)
$`V`$ $`=`$ $`{\displaystyle \frac{29.6}{x^2}},`$ (86)
$`Z`$ $`=`$ $`10^{10}\left(\eta _o{\displaystyle \frac{dy}{4\pi ^2}y^2f_{eq}a_{}}\right),`$ (87)
where $`\eta _o`$ is the asymmetry of the other particle species:
$`\eta _o^e`$ $`=`$ $`2\eta _{\nu _e}+\eta _{\nu _\mu }+\eta _{\nu _\tau }+\eta _e\eta _n/2(\mathrm{for}\nu _e),`$ (88)
$`\eta _o^\mu `$ $`=`$ $`2\eta _{\nu _\mu }+\eta _{\nu _e}+\eta _{\nu _\tau }\eta _n/2(\mathrm{for}\nu _\mu ),`$ (89)
and $`\eta `$ for $`\nu _\tau `$ is obtained from eq. (89) by the interchange $`\mu \tau `$. The asymmetry is normalized in the same way as the neutrino asymmetry (the second term in (87)). Here we have implicitly assumed that $`\nu _a=\nu _e`$.
Up to this point our equations are essentially the same as those used by other groups. The equations look rather innocent and at first sight one does not anticipate any problem with their numerical solution. However the contribution from $`Z`$ could be quite large with the increasing magnitude of the asymmetry. The exact value of the latter is determined by a delicate cancellation of the contributions from all energy spectrum of neutrinos. The function $`a_{}`$ under momentum integral is quickly oscillating and very good precision is necessary to calculate the integral with a desired accuracy. Even a small numerical error results in a large instability. To avoid this difficulty we analytically separated fast and slow variables in the problem and reduced this set of equations to a single differential equation for the asymmetry that can be easily numerically integrated. The corresponding algebra is somewhat complicated and we will not discuss it here. One can find the details in ref. .
We have found that asymmetry practically does not rise for a large mixing angles, $`\mathrm{sin}\theta >0.01`$. For smaller mixings some rise of the asymmetry is observed, though much weaker than that obtained in refs. . For example for $`\delta m^2=1`$ eV<sup>2</sup> we found that asymmetry rises by approximately 5 orders of magnitude reaching the value around $`10^5`$ for $`\mathrm{sin}2\theta `$ in the interval $`10^510^3`$. For $`\delta m^2=10^6`$ eV<sup>2</sup> the symmetry could rise up to 0.01 for $`\mathrm{sin}2\theta =310^6310^5`$, and only for huge mass difference $`\delta m^2=10^9`$ the asymmetry may reach unity in a rather narrow range of mixing angles, $`\mathrm{sin}2\theta =210^6410^6`$. For $`\mathrm{sin}2\theta `$ outside the indicated limits the asymmetry does not rise. For even larger mass difference, $`\delta m^2=10^{12}`$ eV<sup>2</sup>, the asymmetry practically does not rise.
Resolution of the contradiction between different groups is very important for the derivation of the constraints on the parameters of the oscillations from big bang nucleosynthesis (BBN). There are several effects by which the oscillations may influence abundances of light elements:
1. If sterile neutrinos are excited by the oscillations then the effective number of neutrino species at nucleosynthesis would be larger than 3. This effect, as is well known, results in an increase of mass fraction of helium-4 and deuterium.
2. Oscillations may distort the spectrum of neutrinos and, in particular, of electronic neutrinos. The sign of the effect is different, depending on the form of spectral distortion. A deficit of electronic neutrinos at high energy results in a smaller mass fraction of helium-4, while a deficit of $`\nu _e`$ at low energy works in the opposite direction. A decrease of total number/energy density of $`\nu _e`$ would result in an earlier freezing of neutrino-to-proton ratio and in a larger fraction of helium-4.
3. Oscillations may create an asymmetry between $`\nu _e`$ and anti-$`\nu _e`$. The $`n/p`$ ratio in this case would change as $`n/p\mathrm{exp}(\mu /T)`$, where $`\mu `$ is the chemical potential corresponding to the asymmetry. With the present day accuracy of the data the asymmetry in the sector of electronic neutrinos could be at the level of a few per cent, i.e. much larger than the standard $`10^{10}`$. Even if asymmetry is strongly amplified but still remains below 0.01 its direct influence on BBN would be negligible. It may however have an impact on the nucleosynthesis in an indirect way. Namely, the rise of the asymmetry by several orders of magnitude could suppress neutrino oscillations through refraction index so that new neutrino species corresponding to sterile neutrinos are not efficiently excited.
Thus one sees that the effects of oscillations may result both in a reduction or in an increase the effective number of neutrino species. In particular, oscillations may open room for additional particles at BBN. This conclusion strongly depends upon the magnitude of lepton asymmetry generated by the oscillations. Thus a resolution of the controversies between different theoretical calculations would be very important.
## Acknowledgments
This work was partly supported by Danmarks Grundforskningsfond through its funding of the Theoretical Astrophysics Center.
## References
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# Exact Solution of a Repulsive Fermi Model With Enhanced Superconducting Correlations
## Abstract
We present the exact solution of a model of interacting fermions in any dimension with a pure repulsive interaction projecting out a given Cooper channel. The solution rests upon the infinite ranged character of the interaction in real space, leading to a functional integral that is dominated by a Gaussian term. The solution produces strong superconducting enhancements and quasi long ranged order in a channel that is not present in the Hamiltonian explicitly, but of the form given by arguments from order by projection.
There is considerable recent interest in the possibility of models displaying a superconducting behaviour driven solely by repulsive interactions. The search is motivated by the High Tc problem, where no obvious,known attractive interaction can account for the phenomena; so one believes that repulsive interactions, originating in the Coulomb repulsion expressed within the Wannier basis of a few tightbinding bands, are ultimately responsible. In one dimension, the usual kind of repulsive interacions generically lead to enhanced SDW order, via powerlaw correlations with small decay exponents rather than superconducting enhancements. However recent interesting work shows that under certain conditions, superconducting correlations of non trivial symmetry are enhanced. In the physically important case of two dimensions the situation is not completely clear in relation to popular models, such as the Hubbard or the t-J model. While a treatment within the Random Phase Approximation near an Antiferromagnetic instability leads to d-wave superconductivity , one may worry whether the conclusions based on the (weak coupling) approximation are valid for strong repulsions. Similar worries exist regarding various versions of the Gauge theories that are in vogue currently. In this context, repulsive models which can be solved exactly have an important role to play. A new set of models was introduced recently that demonstrate the possibility of enhanced superconducting correlations quite explicitly. These models are generalizations of the Hubbard model, and include a term that is best interpreted as a projection operator that excludes a certain Cooper pairing channel from the problem. We present here an exact solution of the basic model involving the kinetic energy and the projection operator. Our solution is obtained by exploiting a certain feature of the interaction within the framework of functional integrals: namely that the projection operator is an infinite ranged repulsive operator, and its Hubbard-Stratanovic (HS) linearizing field is a single spatially uniform dynamical mode that admits only Gaussian fluctuations in the thermodynamic limit .This is in contrast to models with attractive interactions, where a new saddle point value of the HS field gets stabilized at low temperatures, and its fluctuations can be ignored. The repulsive case is much more subtle, and has some similarities to the situation found in Fermionic models in infinite dimensions where the relevant HS field is a Grassman variable that factorises the kinetic energy, and has only Gaussian fluctuations . We find that as a consequence the fermi gas develops enhanced correlations in a “compromise” pairing channel which is not explicitly present in the Hamiltonian. These are precisely of the sort that is expected from arguments from order by projection . The earlier treatments have used a variety of non perturbative techniques, such as rigorous inequalities and variational approaches. The exact solution presented here is consistent with these, and give us in addition, a clear understanding of the origin of these enhancements and of the quasi LRO in terms of the singularities of the two particle scattering amplitude. The model is defined by the Hamiltonian
$$H=T+U_sB^{}B.$$
(1)
Here $`B\zeta (k)b_k`$ is a Cooper pair operator, $`b_k=c_kc_k`$ are the pair destruction operators, $`T=ϵ_kn_{k\sigma }`$ is the kinetic energy corresponding to a band dispersion $`ϵ_k[=2_{\alpha =1}^d\mathrm{cos}(k_\alpha )`$ in the nearest neighbour problem \]. $`\zeta (k)`$ may be chosen arbitrarily. The two cases of interest in 2-dimensions are (i) $`\zeta (k)=1`$ giving rise to extended s-wave order, and (ii) $`\zeta (k)=\mathrm{cos}(k_x)\mathrm{cos}(k_y)`$ giving rise to second-neighbour d-wave order. For simplicity of presentation we initially focus on $`\zeta =1`$ and return to the other case later.
With the above choice of the model, we note that $`B`$ may be rewritten as $`_rc_{}(r)c_{}(r)`$, and hence the interaction term may be viewed as an infinite ranged hopping term for doubly occupied sites.
Using the H-S linearization within the standard Grassman variable formulation, we write the partition function for this model as the functional integral $`Z=Dc^{}DcD\varphi ^{}D\varphi \mathrm{exp}\beta \mathrm{\Psi }`$. The free energy functional $`\mathrm{\Psi }`$ is given in terms of the fermi fields $`c,c^{}`$ and the auxiliary bose fields $`\varphi ,\varphi ^{}`$ as
$`\beta \mathrm{\Psi }`$ $`=`$ $`{\displaystyle _0^\beta }\varphi ^{}(\tau )\varphi (\tau )𝑑\tau {\displaystyle \underset{k,\sigma }{}}{\displaystyle _0^\beta }c_{k,\sigma }^{}(\tau )(_\tau \xi _k)c_{k,\sigma }(\tau )𝑑\tau `$ (3)
$`i\sqrt{U_s}{\displaystyle _0^\beta }(\varphi ^{}(\tau )B(\tau )+B^{}(\tau )\varphi (\tau ))𝑑\tau ,`$
where $`\xi _k=ϵ_k\mu `$. We use a Fourier series expansion $`c_{k,\sigma }(\tau )=_n\mathrm{exp}(i\omega _n\tau )\widehat{c}_{k,\sigma }(n)`$, and $`\varphi (\tau )=_n\mathrm{exp}(i\mathrm{\Omega }_n\tau )\widehat{\varphi }_n`$, where the fermionic frequencies $`\omega _n=(2n+1)\pi \beta ^1`$ and the bosonic frequencies $`\mathrm{\Omega }_n=2n\pi \beta ^1`$. We denote the Fourier components $`\widehat{b}_k(n)=_m\widehat{c}_{k,}(m)\widehat{c}_{k,}(nm)`$ and so define $`\widehat{B}_n=_k\widehat{b}_k(n1)`$. In terms of these we can rewrite
$$\mathrm{\Psi }=\widehat{\varphi }_{}^{}{}_{n}{}^{}\widehat{\varphi }_n+\underset{k,n}{}(\xi _ki\omega _n)\widehat{c}_{k,\sigma }^{}(n)\widehat{c}_{k,\sigma }(n)+i\sqrt{U_s}\underset{n}{}(\widehat{\varphi }_{}^{}{}_{n}{}^{}\widehat{B}_n+\widehat{\varphi }_n\widehat{B}_{}^{}{}_{n}{}^{}).$$
We now trace out the Fermi degrees of freedom and find the reduced free energy functional $`\mathrm{\Psi }_\varphi =\mathrm{\Psi }_0+\widehat{\varphi }_{}^{}{}_{n}{}^{}\widehat{\varphi }_nk_bT_kTr\mathrm{log}\{1C(k)\}`$. Here $`\mathrm{\Psi }_0`$ is the noninteracting free energy, $`C`$ is an infinite dimentional matrix defined by its (frequency space) matrix elements $`C_{n,m}(k)=U_s_lG_0(k,n)G_0(k,m)\widehat{\varphi }_{}^{}{}_{m+l+1}{}^{}\widehat{\varphi }_{n+l+1}`$ in terms of the free Green’s function $`G_0(k,n)=1/(i\omega _n\xi _k)`$ , and the trace is in the frequency space.
We now analyse $`\mathrm{\Psi }_\varphi `$ in detail. We can expand
$`\mathrm{\Psi }_\varphi `$ $`=`$ $`\mathrm{\Psi }_0+{\displaystyle \underset{n}{}}\{1+U_s\mathrm{\Pi }_0(i\mathrm{\Omega }_n)\}\widehat{\varphi }_{}^{}{}_{n}{}^{}\widehat{\varphi }_n`$ (5)
$`+{\displaystyle \frac{U_s^2}{2}}{\displaystyle \mathrm{\Gamma }_{(m_1,m_2,n_1,n_2)}\widehat{\varphi }_{}^{}{}_{m_1}{}^{}\widehat{\varphi }_{}^{}{}_{m_2}{}^{}\widehat{\varphi }_{n_1}\widehat{\varphi }_{n_2}}+O(U_s^3).`$
Here $`\mathrm{\Pi }_0(i\mathrm{\Omega }_n)\frac{1}{}_k\pi _{k,n}`$ where $`\pi _{k,n}`$ are the polarizations
$$\pi _{k,n}=(2f(k)1)/(i\mathrm{\Omega }_n2\xi _k),$$
(6)
and $`f(k)`$ is the usual noninteracting Fermi function. The fourth order term is given explicitly as
$$\mathrm{\Gamma }_{(m_1,m_2,n_1,n_2)}=\frac{1}{}\underset{k}{}\frac{(2f(k)1)(D_{m_1,k}+D_{m_2,k})\delta _{m_1+m_2,n_1+n_2}}{D_{m_1,k}D_{m_2,k}D_{n_1,k}D_{n_2,k}},$$
where $`D_{m,k}(i\mathrm{\Omega }_m2\xi _k)`$ .
We note that both $`\mathrm{\Pi }_0`$ and $`\mathrm{\Gamma }`$ in the above equations are of $`O(1)`$ since these are normalized sums over momenta. Indeed every term in the expansion has a similar structure and is of the same order, namely $`O()`$. Hence one has the remarkable exact result that the Gaussian term dominates the rest of the terms in the thermodynamic limit . Roughly speaking, the Gaussian piece gives us the typical size $`\widehat{\varphi }_n\frac{1}{\sqrt{}}`$, and so the quadratic ( in $`U_s`$) piece is of the order $`1/`$, and likewise the $`m^{th}`$ term is of the order $`1/^{(m1)}`$. Thus in the thermodynamic limit, it suffices to keep the Gaussian term and to drop the remaining terms. This leads to the following remarkably simple result
$$<\widehat{\varphi }_{}^{}{}_{n}{}^{}\widehat{\varphi }_n>=\frac{\beta ^1}{1+U_s\mathrm{\Pi }_0(i\mathrm{\Omega }_n)}.$$
(7)
From the same arguments, the correlation function of the Cooper pair operators is given by
$$<b_k^{}(n)b_k^{}(n)>=\delta _{k,k^{}}\pi _{k,n}\frac{1}{}\pi _{k,n}V_{eff}(i\mathrm{\Omega }_n)\pi _{k^{},n}$$
(8)
where the effective interaction
$$V_{eff}(n)=\frac{V}{1+V\mathrm{\Pi }_0(i\mathrm{\Omega }_n)}$$
(9)
in terms of the (very large) coupling constant $`V=U_s`$.
We next present an alternate derivation of the above results starting from the equations of motion, which gives some more insight into them.
Define (as in ref ) the set of operators $`I_l(ϵ_k)^lb_k`$, and $`T_l(ϵ_k)^l(n_k+n_k1)`$ . Clearly $`I_0=B`$ and $`T_1=T`$. It is easy to ascertain that
$$[I_l,T]=2I_{l+1},[I_l,\widehat{N}]=2I_l,[I_l,I_{m}^{}{}_{}{}^{}]=T_{l+m},$$
(10)
whence,
$$[I_l,H\mu \widehat{N}]=2\mu I_l+2I_{l+1}U_sT_lI_0.$$
(11)
Now, we invoke the law of large numbers and argue that in the present problem the operator product $`T_lI_0`$ can be replaced by $`\mu _lI_0`$ where $`\mu _l=<T_l>`$ , the thermodynamic average of $`T_l`$, which is clearly of order $``$ . Then the equations of motion for the usual time ordered Green functions $`<<I_l;I_{m}^{}{}_{}{}^{}>>`$ reduce to the closed set
$$(i\mathrm{\Omega }_n+2\mu )<<I_l;I_{m}^{}{}_{}{}^{}>>=\mu _{l+m}+2<<I_{l+1};I_{m}^{}{}_{}{}^{}>>U_s\mu _l<<I_0;I_{m}^{}{}_{}{}^{}>>$$
(12)
which can be solved exactly. The solution is given by
$$<<I_l;I_{m}^{}{}_{}{}^{}>>=\mathrm{\Pi }_{l+m}+\mathrm{\Pi }_l\frac{U_s}{1+U_s\mathrm{\Pi }_0}\mathrm{\Pi }_m$$
(13)
where
$$\mathrm{\Pi }_l(i\mathrm{\Omega }_n)\underset{l^{}=0}{\overset{\mathrm{}}{}}\frac{2^l^{}\mu _{l+l^{}}}{(i\mathrm{\Omega }_n+2\mu )^{l^{}+1}}=(1/)\underset{k}{}\frac{(ϵ_k)^l(2<n_k>1)}{(i\mathrm{\Omega }_n2\xi _k)}.$$
This may be verified using the easily derived recursion relation $`\mathrm{\Pi }_{l+1}=[(i\mathrm{\Omega }_n+2\mu )\mathrm{\Pi }_l\mu _l]/2`$. Furthermore, one of the key features of the models being discussed is that the one particle propagators are unrenormalised , so that $`<n_k>=f(k)`$ whence, $`\mathrm{\Pi }_l(i\mathrm{\Omega }_n)=(1/)_k(ϵ_k)^l\pi _{k,n}`$. Then it is easily seen that the results in Eq. (13) are basically the same as in Eq. (8). Starting from the latter, multiplying by $`(ϵ_k)^l(ϵ_k^{})^m`$ and summing over $`k`$ and $`k^{}`$, we get the former.
An analysis of the detailed properties of the function $`V_{eff}`$ is of crucial importance for the rest of our discussion. For a simple model bandstructure, with a constant density of states $`g(ϵ)=1/2`$ for $`1<ϵ<1`$, we can compute it exactly as
$$V_{eff}(i\mathrm{\Omega })V_{eff}(\omega +i\eta )=\frac{V}{1+V\mathrm{\Pi }_0(\omega )},$$
(14)
with
$$\mathrm{\Pi }_0(\omega )=\frac{1}{4}\mathrm{log}(\frac{4|(\omega /2+\mu )^21|}{\omega ^2})+i\frac{\pi }{4}\theta (2|2\mu +\omega |)sgn(\omega )$$
The band extends from $`2+2\delta `$ to $`2+2\delta `$, where the hole filling $`\delta =1\rho =\mu `$. The schematic behaviour of $`V_{eff}`$ is as follows. At very high frequencies, $`V_{eff}V`$. It has two poles at certain large frequencies that essentially dominate its physics . In between these poles lies the intermediate frequency range where $`V_{eff}`$ is of order unity, which contains the branch cut corresponding to the one electron band of states. The location of the poles can be found by using a large frequency expansion for the function $`\mathrm{\Pi }_0(i\mathrm{\Omega }_n)`$, which for $`|\mathrm{\Omega }_n|>>W`$ ( $`W`$ is the band width) behaves as
$$\mathrm{\Pi }_0(i\mathrm{\Omega }_n)=\frac{\delta }{(i\mathrm{\Omega }_n)}\frac{2(|\mu _1|\mu \delta )}{(i\mathrm{\Omega }_n)^2}+O(\frac{1}{|\mathrm{\Omega }_n|^3}).$$
(15)
The higher order terms can be verified to be negligible. The poles of the $`V_{eff}`$ can then be found by solving for the zeroes of the denominator, a quadratic in $`\zeta 1/i\mathrm{\Omega }_n`$ , given as $`0=V^1\delta \zeta 2(|\mu _1|\mu \delta )\zeta ^2`$. The roots are always real corresponding to real frequency poles of $`V_{eff}`$. Hence,(a) at half filling, $`\delta =0`$, the poles are at $`\pm \omega _0=\pm \sqrt{2V|\mu _1|}`$, and near the poles $`V_{eff}\pm \frac{\omega _0V}{2(\omega \omega _0)}`$. The pole frequencies are infinite in the thermodynamic limit, but leave behind consequences in the groundstate as we see later. (b) away from but close to half filling , $`\delta <<1`$, one pole is at $`\omega _2V\delta `$ with residue $`V^2\delta `$, and the other at a large negative frequency $`\omega _1=(2|\mu _1|/\delta 2\mu )`$ with residue $`\omega _1^2/\delta `$. As $`\delta 0`$, the latter poles smoothly go over into those of the first case. The contribution from the branch cut does not have any particularly simple form, but is not important in the most interesting region of the problem, namely $`\delta 0`$.
Using the above properties of $`V_{eff}`$ we can compute exactly the instantaneous (expectation) values of the the extended s-wave and the s-wave correlation functions :
$`\alpha \frac{1}{4}<A^{}A>`$ with $`A=[T,B]=2_kϵ_kb_k`$, and $`ß\frac{1}{}<B^{}B>`$.
Using the main result Eq(8), we find
$`\alpha `$ $`=`$ $`k_BT{\displaystyle \underset{n}{}}[\mathrm{\Pi }_2(i\mathrm{\Omega }_n)V_{eff}(i\mathrm{\Omega }_n)\mathrm{\Pi }_1^2(i\mathrm{\Omega }_n)]e^{(i\mathrm{\Omega }_n0^+)}`$ (16)
$`ß`$ $`=`$ $`{\displaystyle \frac{k_BT}{V^2}}{\displaystyle \underset{n}{}}[VV_{eff}(i\mathrm{\Omega }_n)]e^{(i\mathrm{\Omega }_n0^+)}.`$ (17)
The main contributions to the resulting frequency sums are quite easily seen to be tied to the pole contributions of $`V_{eff}`$, the branch cut part giving an uninteresting subleading contribution. We find at half filling
$$\alpha _{pole}=\frac{\sqrt{U_s|\mu _1|^3}}{2\sqrt{2}};ß_{pole}=\frac{\sqrt{|\mu _1|}}{\sqrt{2U_s}}.$$
(18)
Near half filling we find
$$\alpha _{pole}=\frac{\mu _1^2}{\delta };ß_{pole}=\frac{4\mu _1^2}{\delta ^3U_s^2^2}.$$
(19)
At half filling Eq.(18) gives us the quasi long ranged order as well as the correction to ground state energy along with their appropriate coeffecients, the latter by an integration over the coupling constant $`U_s`$. These answers are in good numerical agreement with exact numerics on a finite sized system . Away from half filling we find that the uncertainty principle lower bounds are off by a factor of 2, i.e., from Eq.(19) $`\alpha _{pole}=2\alpha _{LowerBound}`$.
The results given above enable us to compute several other response functions exactly. For example we find that the charge stiffness or Meissner fraction is non- zero, and in fact unchanged from the non interacting value at half filling. We also find at half filling that the appropriate order parameter density $`\mathrm{\Delta }(r)`$ has correlations that are novel : $`<\mathrm{\Delta }^{}(r)\mathrm{\Delta }(0)>\frac{c}{\sqrt{}}`$. Thus although we do not have LRO of the usual sort, one has a divergent “structure function” $`𝑑\stackrel{}{r}<\mathrm{\Delta }^{}(r)\mathrm{\Delta }(0)>`$.
We next discuss the important case of $`\zeta (k)=\mathrm{cos}(k_x)\mathrm{cos}(k_y)`$ in two dimensions. This corresponds to suppressing d-wave order at the length scale of nearest neighbours.From the uncertainty principle argument of Refs , it follows that this would lead to enhancement again in the d-wave channel, but at a longer length scale, i.e. the resulting $`\widehat{A}=2ϵ_k\zeta (k)b_k`$ should have enhanced correlations. The functional integral solution sketched here bears this out exactly. We recover the results in Eq.s(18,19) with the replacements : $`\delta \widehat{\delta }=^1_k\zeta ^2(k)(12f(k))`$ and $`\mu \widehat{\mu }_1=^1_kϵ_k\zeta ^2(k)(12f(k))`$. The enhanced correlation function of $`\widehat{A}`$ then diverges at the point where $`\widehat{\delta }`$ vanishes. So long as the one electron dispersion has the bipartite symmetry, one can see that $`\widehat{\delta }`$ vanishes exactly at half filling. However, if the dispersion does not have this symmetry, e.g. by having a second neighbouring hopping $`t^{}`$, then $`\widehat{\delta }`$ vanishes at some other density determined by $`t^{}`$, as illustrated in Fig.1. The case of $`t^{}.4`$ is popular in literature, since it leads to a Fermi surface that is consistent with that seen in the photo-emission experiments in High $`T_c`$ systems , and it is an amusing coincidence that the filling $`\delta .18`$ is close to the optimum doping .
Finally, we can show that many of the above features of the infinite range model retain their relevance for more realistic models with finite range repulsion. To see this, consider (in the s wave case) a more general model of the form $`H=T+\frac{1}{}U(p)B^{}(p)B(p)`$, where $`B(p)=_kc_{k,}c_{k+p,}`$ is a Cooper pair operator with total momentum $`p`$ , with $`U(p)=U_s(l_c)^d`$ for a small set of $`/(l_c)^d`$ points surrounding the centre of the Brillouin zone. This would correspond to repulsive interactions with a long but finite range $`l_c`$ . Then it is not hard to see that the results we have discussed above would retain much of their validity (but for some differences in details) to leading order in $`(1/l_c)`$, with the replacement $`U_sU_s(l_c)^d`$, leading to enhanced pairing correlations over a finite range in momentum space. The quasi Long Ranged Order would be replaced by an enhancement of the $`O(1)`$, like that in the single mode model away from half filling. In this case, we see that all Cooper pairs with a finite ( but small) center of mass momentum are also influenced by the interaction, and thus the model is more realistic by way of helping current carrying states.
In summary, we have found an exact solution for an interesting model of Fermions with purely repulsive interactions with infinite range, which may be regarded as a meanfield repulsive model. The resulting solution has quasi long ranged order at half filling, as well as large unbounded enhancements as one approaches half filling, in the equal time pairing correlations. We have also argued that the above methods and results retain their relevance even for generalised models where the repulsive interactions have a large but finite range, and are therefore more realistic. Finally, it is remarkable that the enhanced pairing correlations in these models arise from very high energy poles in the scattering amplitude, not unlike the physics of Mott Hubbard systems, where the upper Hubbard band influences the properties of carriers in the lower band.
One of us (BSS) is grateful to the Aspen Center of Physics for support, where a part of this paper was written. He thanks P W Anderson and B Shraiman for stimulating discussions and constructive suggestions.
*Also at the JNCASR, Jakkur, Bangalore
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# Parametrizing the mixing matrix : A unified approach
## I Introduction
In the standard $`SU(3)\times SU(2)\times U(1)`$ model of strong , weak and electromagnetic interactions, all aspects of the charged weak interactions among quarks can be described in terms of a $`3\times 3`$ unitary matrix
$$V=\left(\begin{array}{ccc}V_{ud}& V_{us}& V_{ub}\\ V_{cd}& V_{cs}& V_{cb}\\ V_{td}& V_{ts}& V_{tb}\end{array}\right)$$
(1)
specified by four real parameters: three generalized Cabibbo angles and one Kobayashi-Maskawa phase. After the pioneering work of Kobayashi and Maskawa , this matrix, which describes the mixing between quark mass eigenstates and the charged weak current eigenstates, has been parametrized in a number of phenomenologically useful ways \[2-6\]. Generalizations to $`N3`$ generations of quarks, where the mixing matrix is characterized by $`N(N1)/2`$ angles and $`(N1)(N2)/2`$ phases, have also been proposed \[6-14\]. Analogues of the mixing matrix also arise in the lepton sector if the neutrinos are taken as massive Dirac particles. In most of the parametrizations hitherto proposed, the mixing matrix is expressed as an ordered product of $`N(N1)/2`$ factors each of which carries an angle. Of these $`N(N1)/2`$ factors, a prescribed set of $`(N1)(N2)/2`$ factors carry phases as well. Different parametrizations differ from each other in the ordering prescription and the location of the phase factors within the matrices carrying them. In this work, we present a parametrization of the mixing matrix based on a decomposition, involving, in the $`N=3`$ case, just two factors. This parametrization, apart from having a clear geometrical picture underlying it, also enables us to recover and relate other parametrizations and to generate new ones in a unified manner.
## II Parametrizing SU(n) elements by a sequence of complex unit vectors
The proposed parametrization of the mixing matrix is based on the observation that a generic matrix $`gSU(N)`$ can be parametrized by a sequence of complex unit vectors $`𝝃,\mathrm{},𝜸,𝜷,𝜶`$ of dimensions $`n,n1,\mathrm{},3,2`$ . This can be seen as follows. Let
$$\mathrm{\Sigma }_n=\{𝝍=\left(\begin{array}{c}\psi _1\\ \psi _2\\ \\ \\ \psi _n\end{array}\right)𝝍^{}𝝍=1\}$$
(2)
denote the set of unit vectors in complex n-dimensional Hilbert space i.e. a set of real dimension $`(2n1)`$. Any $`𝝍\mathrm{\Sigma }_n`$ can be mapped to the vector
$$𝐞_n=\left(\begin{array}{c}0\\ 0\\ \\ \\ 1\end{array}\right)$$
(3)
via a suitable $`SU(n)`$ element. ( Note that we are really using $`SU(n)`$ here , not $`U(n)`$) Therefore, $`SU(n)`$ acts transitively on $`\mathrm{\Sigma }_n`$. The subgroup of $`SU(n)`$ that leaves $`𝐞_n`$ invariant is $`SU(n1)`$ on the first $`(n1)`$ dimensions and hence
$$\mathrm{\Sigma }_ncosetspaceSU(n)/SU(n1)$$
(4)
Therefore, we expect that, apart from global matching problems or ambiguities on a subset of measure zero, any element in $`SU(n)`$ is uniquely specified by a pair consisting of an element in $`SU(n1)`$ and a unit vector $`𝝍\mathrm{\Sigma }_n`$. Therefore, recursively, we see that an element $`gSU(n)`$ can be parametrized as $`g=g(𝝃,\mathrm{},𝜸,𝜷,𝜶)`$ by a string of complex unit vectors $`𝝃,\mathrm{},𝜸,𝜷,𝜶`$ of dimensions $`n,n1,\mathrm{},4,3,2`$.
As a convention, we will let the above unit vectors stand for the last column. in the relevant $`SU(n)`$ matrices. This is because when a matrix of $`SU(n)`$ is multiplied on the right by a matrix of $`SU(n1)`$ (leaving $`𝐞_n`$ invariant), it is the last column in the former matrix that remains unchanged. For elements of $`SU(2)`$ we will thus write
$$ASU(2);A=A(𝜶)=\left(\begin{array}{cc}\alpha _2^{}& \alpha _1\\ \alpha _1^{}& \alpha _2\end{array}\right);𝜶^{}𝜶=1$$
(5)
This is globally well defined.
## III SU(3) and the Kobayashi-Maskawa phase
Let $`𝜷`$ denote a three component complex unit vector, $`𝜷^{}𝜷=1`$. Then for $`|\beta _1|<1`$, the matrix
$$B(𝜷)=\left(\begin{array}{ccc}P^1& 0& \beta _1\\ P\beta _1^{}\beta _2& P\beta _3^{}& \beta _2\\ P\beta _1^{}\beta _3& P\beta _2^{}& \beta _3\end{array}\right),P=(1|\beta _1|^2)^{1/2}$$
(6)
is in $`SU(3)`$. the unit vector $`𝜷`$ is a label for right $`SU(2)`$ cosets in $`SU(3)`$, and $`B(𝜷)`$ is a coset representative. So any $`BSU(3)`$, $`|B_{13}|<1`$, can be uniquely written as
$$B=B(𝜷,𝜶)=B(𝜷)\left(\begin{array}{cc}A(𝜶)& 0\\ 0& 1\end{array}\right).$$
(7)
On multiplying out the two matrices on the rhs one obtains
$$B(𝜷,𝜶)=\left(\begin{array}{ccccc}P^1\alpha _2^{}& & P^1\alpha _1& & \beta _1\\ P\beta _1^{}\beta _2\alpha _2^{}P\beta _3^{}\alpha _1^{}& & P\beta _1^{}\beta _2\alpha _1+P\beta _3^{}\alpha _2& & \beta _2\\ P\beta _1^{}\beta _3\alpha _2^{}+P\beta _2^{}\alpha _1^{}& & P\beta _1^{}\beta _3\alpha _1P\beta _2^{}\alpha _2& & \beta _3\end{array}\right)$$
(8)
Now we examine how $`B(𝜷,𝜶)`$ transforms under rephasing transformations i.e. we ask how $`𝜷`$ and $`𝜶`$ change when we multiply $`B(𝜷,𝜶)`$ on the left and on the right by independent diagonal elements of $`SU(3)`$:
$$B^{}=D(\theta ^{})BD(\theta )$$
(9)
where $`D(\theta )=diag(e^{i\theta _1+i\theta _2},e^{i\theta _1+i\theta _2},e^{2i\theta _2})`$ and $`D(\theta ^{})`$ is defined similarly. Then we find
$$D(\theta ^{})B(𝜷,𝜶)D(\theta )=B(𝜷^{},𝜶^{})$$
(10)
where
$`\alpha _1^{}`$ $`=`$ $`\alpha _1e^{i(\theta _1^{}+\theta _2^{}\theta _1+\theta _2)};\alpha _2^{}=\alpha _2e^{i(\theta _1^{}+\theta _2^{}+\theta _1+\theta _2)}`$ (11)
$`\beta _1^{}`$ $`=`$ $`\beta _1e^{i(\theta _1^{}+\theta _2^{}2\theta _2)};\beta _2^{}=\beta _2e^{i(\theta _1^{}+\theta _2^{}2\theta _2)};\beta _3^{}=\beta _3e^{2i(\theta _2^{}+\theta _2)}`$ (12)
These transformation laws can easily be written down from the locations of $`\alpha _1,\alpha _2,\beta _1,\beta _2,\beta _3`$ in the matrix $`(\text{8})`$.
As the dimension of $`SU(3)`$ is eight and we have four independent phases here, there should be four independent real invariants. Three of them, essentially the generalized Cabibbo angles may be chosen to be, say, $`|\alpha _1|,|\beta _1|,|\beta _2|`$. The fourth one can be found systematically as follows.
The phase $`\theta _1`$ does not occur in $`𝜷^{}`$. So we form the combination $`\alpha _1\alpha _2^{}`$ whose transformation law is $`\theta _1`$ independent:
$$\alpha _1^{}\alpha _2^{}=\alpha _1\alpha _2^{}e^{2i(\theta _1^{}+\theta _2^{}+\theta _2)}$$
(13)
Among the $`\beta `$’s, $`\theta _1^{}`$ occurs only in $`\beta _1^{}`$ and $`\beta _2^{}`$, so we form a combination which can cancel $`e^{2i\theta _1^{}}`$ on the rhs of $`(\text{13})`$
$$\beta _1^{}\beta _2\beta _1^{}\beta _2e^{2i\theta _1^{}}$$
(14)
From $`(\text{13})`$ and $`(\text{14})`$ we find
$$\alpha _1\alpha _2^{}\beta _1^{}\beta _2\alpha _1\alpha _2^{}\beta _1^{}\beta _2e^{2i(\theta _2^{}+\theta _2)}$$
(15)
Comparing this with $`\beta _3^{}`$ we see that $`arg(\alpha _1\alpha _2^{}\beta _1^{}\beta _2\beta _3)`$ is invariant under rephasing.
## IV Comparison with some well known parametrizations of the mixing matrix for N=3
Before we show how some well known parametrizations of the mixing matrix can be recovered from the considerations given above, it is useful to note that from the standard form$`(\text{8})`$ we can generate others by permutation of rows and columns and by taking transpose. The expressions for the invariants remain unchanged under these opertations as will become clear in section VII. This being the case, various parametrizations of the mixing matrix can be generated by choosing any one from $`\alpha _1,\alpha _2,\beta _1,\beta _2,\beta _3`$ which appear in the invariant $`arg(\alpha _1\alpha _2^{}\beta _1^{}\beta _2\beta _3)`$ to be complex and all others real in the matrix $`(\text{8})`$ or in the matrices obtained by permuting rows and columns or by taking transpose. Thus, for instance, choosing $`\beta _2`$ to be complex, all others real, and putting
$$\alpha _1=S_\theta ,\alpha _2=C_\theta ,\beta _1=S_\beta ,\beta _2=S_\gamma C_\beta e^{i\delta },\beta _3=C_\gamma C_\beta $$
(16)
in $`(\text{8})`$ one obtains the Maiani parametrization
$$\left(\begin{array}{ccc}C_\beta C_\theta & C_\beta S_\theta & S_\beta \\ S_\gamma C_\theta S_\theta e^{i\delta }S_\theta C_\gamma & C_\gamma C_\theta S_\gamma S_\beta S_\theta e^{i\delta }& S_\gamma C_\beta e^{i\delta }\\ S_\beta C_\gamma C_\theta +S_\gamma S_\theta e^{i\delta }& C_\gamma S_\beta S_\theta S_\gamma C_\theta e^{i\delta }& C_\gamma C_\beta \end{array}\right)$$
(17)
The Chau-Keung parametrization corresponds to choosing $`\beta _1`$ complex.
$$\alpha _1=S_{12},\alpha _2=C_{12},\beta _1=S_{13}e^{i\delta },\beta _2=S_{23}C_{13},\beta _3=C_{23}C_{13}$$
(18)
The mixing matrix, for this choice, is given by
$$\left(\begin{array}{ccc}C_{12}C_{13}& S_{12}C_{13}& S_{13}e^{i\delta _{13}}\\ S_{12}C_{23}C_{12}S_{23}S_{13}e^{i\delta _{13}}& C_{12}C_{23}S_{12}S_{23}S_{13}e^{i\delta _{13}}& S_{23}C_{13}\\ S_{12}S_{23}C_{12}C_{23}S_{13}e^{i\delta _{13}}& C_{12}S_{23}S_{12}C_{23}S_{13}e^{i\delta _{13}}& C_{23}C_{13}\end{array}\right)$$
(19)
The Kobayashi-Maskawa form corresponds to taking $`\beta _2`$ complex and putting
$$\alpha _1=S_3,\alpha _2=C_3,\beta _1=C_1,\beta _2=S_1C_2e^{i\delta },\beta _3=S_1S_2$$
(20)
in
$$B(𝜷,𝜶)=\left(\begin{array}{ccccc}\beta _1& & P^1\alpha _2^{}& & P^1\alpha _1\\ \beta _2& & P\beta _1^{}\beta _2\alpha _2^{}P\beta _3^{}\alpha _1^{}& & P\beta _1^{}\beta _2\alpha _1+P\beta _3^{}\alpha _2\\ \beta _3& & P\beta _1^{}\beta _3\alpha _2^{}+P\beta _2^{}\alpha _1^{}& & P\beta _1^{}\beta _3\alpha _1P\beta _2^{}\alpha _2\end{array}\right)$$
(21)
and leads to
$$\left(\begin{array}{ccc}C_1& S_1C_3& S_1S_3\\ S_1C_2e^{i\delta }& C_1C_2C_3e^{i\delta }+S_2S_3& C_1C_2S_3e^{i\delta }S_2C_3\\ S_1S_2& C_1S_2C_3+C_2S_3e^{i\delta }& C_1S_2S_3C_2C_3e^{i\delta }\end{array}\right)$$
(22)
which, on multiplying the second row by the phase factor $`(e^{i\delta })`$ gives precisely the mixing matrix originally given by Kobayashi and Maskawa.
Similarly, taking $`\beta _2`$ to be complex and putting
$$\alpha _1=C_{13},\alpha _2=S_{13},\beta _1=S_{12},\beta _2=C_{12}C_{23}e^{i\alpha },\beta _3=C_{12}S_{23}$$
(23)
in
$$B(𝜷,𝜶)=\left(\begin{array}{ccccc}P^1\alpha _1& & \beta _1& & P^1\alpha _2^{}\\ P\beta _1^{}\beta _2\alpha _1+P\beta _3^{}\alpha _2& & \beta _2& & P\beta _1^{}\beta _2\alpha _2^{}P\beta _3^{}\alpha _1^{}\\ P\beta _1^{}\beta _3\alpha _1P\beta _2^{}\alpha _2& & \beta _3& & P\beta _1^{}\beta _3\alpha _2^{}+P\beta _2^{}\alpha _1^{}\end{array}\right)$$
(24)
yields the parametrization due to Anselm et al :
$$\left(\begin{array}{ccc}C_{12}C_{13}& S_{12}& C_{12}S_{13}\\ S_{12}C_{13}C_{23}e^{i\alpha }S_{13}S_{23}& C_{12}C_{23}e^{i\alpha }& S_{12}S_{13}C_{23}e^{i\alpha }C_{13}S_{23}\\ S_{12}C_{13}S_{23}+S_{13}C_{23}e^{i\alpha }& C_{12}S_{23}& S_{12}S_{13}S_{23}+C_{13}C_{23}e^{i\alpha }\end{array}\right)$$
(25)
## V SU(4) and three Kobayashi-Maskawa phases
We now consider the $`N=4`$ case. Let $`𝜸`$ denote a four dimensional complex unit vector. Then, for $`|\gamma _1|^2+|\gamma _2|^2<1`$,
$$C(𝜸)=\left(\begin{array}{cccc}Q^1& 0& 0& \gamma _1\\ Q\gamma _1^{}\gamma _2& QR^1& 0& \gamma _2\\ Q\gamma _1^{}\gamma _3& QR\gamma _2^{}\gamma _3& R\gamma _4^{}& \gamma _3\\ Q\gamma _1^{}\gamma _4& QR\gamma _2^{}\gamma _4& R\gamma _3^{}& \gamma _4\end{array}\right)$$
(26)
where $`Q=(1|\gamma _1|^2)^{1/2}`$, $`R=(1|\gamma _1|^2|\gamma _2|^2)^{1/2}`$, is in $`SU(4)`$. The unit vector $`𝜸`$ is a label for right $`SU(3)`$ cosets in $`SU(4)`$ and $`C(𝜸)`$ is a coset representative. So, except on a subset of measure zero, for a $`CSU(4),|C_{14}|^2+|C_{24}|^2<1`$, there is a unique sequence of complex unit vectors $`𝜸,𝜷,𝜶`$ of dimensions $`4,3,2`$ respectively, such that
$$C=C(𝜸,𝜷,𝜶)=C(𝜸)\left(\begin{array}{cc}B(𝜷,𝜶)& 0\\ 0& 1\end{array}\right)$$
(27)
Now we multiply $`C`$ on the left and right by independent diagonal $`SU(4)`$ matrices, get the transformation laws for $`𝜸,𝜷,𝜶`$, and then construct the invariants.
$$C^{}=D(\theta ^{})C(𝜸,𝜷,𝜶)D(\theta )=C(𝜸^{},𝜷^{},𝜶^{})$$
(28)
where $`D(\theta )=diag(e^{i\theta _1+i\theta _2+i\theta _3},e^{i\theta _1+i\theta _2+i\theta _3},e^{2i\theta _2+i\theta _3},e^{3i\theta _3})`$ and $`D(\theta ^{})`$ is defined similarly. For simplicity, let $`B(𝜷,𝜶)`$ also denote the $`4\times 4`$ matrix obtained by an appropriate bordering. Then, because of the way we parametrized $`D(\theta )`$ and $`D(\theta ^{})`$, we find
$`C^{}`$ $`=`$ $`C(𝜸^{})B(𝜷^{},𝜶^{})`$ (29)
$`=`$ $`D(\theta ^{})C(𝜸)B(𝜷,𝜶)D(\theta )`$ (30)
$`=`$ $`D(\theta ^{})C(𝜸)diag(e^{i\theta _3},e^{i\theta _3},e^{i\theta _3},e^{3i\theta _3})`$ (32)
$`\times B(𝜷,𝜶)diag(e^{i\theta _1+i\theta _2},e^{i\theta _1+i\theta _2},e^{2i\theta _2},1)`$
The expressions for $`𝜸^{}`$ are easy to read off
$`\gamma _1^{}`$ $`=`$ $`\gamma _1e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}3\theta _3)};\gamma _2^{}=\gamma _2e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}3\theta _3)}`$ (33)
$`\gamma _3^{}`$ $`=`$ $`\gamma _3e^{i(2\theta _2^{}+\theta _3^{}3\theta _2)};\gamma _4^{}=\gamma _4e^{3i(\theta _3^{}+\theta _3)}`$ (34)
A little algebra shows that
$`D(\theta ^{})C(𝜸)diag(e^{i\theta _3},e^{i\theta _3},e^{i\theta _3},e^{3i\theta _3})`$ (35)
$`=`$ $`C(𝜸^{})diag(e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _3)},e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _3)},e^{2i(\theta _2^{}+\theta _3^{}+\theta _3)},1)`$ (36)
so that the rest reduces to an $`SU(3)`$ problem in $`3\times 3`$ matrix form
$`B(𝜷^{},𝜶^{})=`$ (37)
$`diag(e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _3)},e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _3)},e^{2i(\theta _2^{}+\theta _3^{}+\theta _3)})`$ (38)
$`\times B(𝜷,𝜶)diag(e^{i\theta _1+i\theta _2},e^{i\theta _1+i\theta _2},e^{2i\theta _2})`$ (39)
which is just the same as in $`(\text{10})`$ with the replacements $`\theta _1\theta _1,\theta _2\theta _2,\theta _1^{}\theta _1^{},\theta _2^{}\theta _2^{}+\theta _3^{}+\theta _3`$. Making these changes in $`(\text{11})`$ and $`(\text{12})`$ we see that for the $`SU(4)`$ problem to accompany $`(\text{34})`$, we have,
$`\alpha _1^{}`$ $`=`$ $`\alpha _1e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}\theta _1+\theta _2+\theta _3)};\alpha _2^{}=\alpha _2e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _1+\theta _2+\theta _3)}`$ (40)
$`\beta _1^{}`$ $`=`$ $`\beta _1e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}2\theta _2+\theta _3)};\beta _2^{}=\beta _2e^{i(\theta _1^{}+\theta _2^{}+\theta _3^{}2\theta _2+\theta _3)}`$ (41)
$`\beta _3^{}`$ $`=`$ $`\beta _3e^{2i(\theta _2^{}+\theta _3^{}+\theta _2+\theta _3)}`$ (42)
From $`(\text{34})`$, $`(\text{40})`$ and $`(\text{42})`$ we need to construct the invariants. The six ‘Cabibbo’ angles may be taken to be given by $`|\alpha _1|,|\beta _1|,|\beta _2|,|\gamma _1|,|\gamma _2|,|\gamma _3|`$. The three KM phases can be obtained systematically as follows. Since $`\theta _1`$ is involved only in $`𝜶^{}`$ , not in $`𝜷^{}`$ and $`\gamma ^{}`$, we see that the $`\alpha `$’s must enter only in the form $`\alpha _1\alpha _2^{}`$ which obeys
$$\alpha _1\alpha _2^{}\alpha _1\alpha _2^{}e^{2i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _2+\theta _3)}$$
(43)
Next, we see that $`\theta _2`$ is not involved in the $`\gamma ^{}`$’s at all, so we form independent expressions in $`\alpha _1\alpha _2^{}`$ and the $`\beta `$’s in which $`\theta _2`$ goes away.
$`\alpha _1\alpha _2^{}\beta _1`$ $`\alpha _1\alpha _2^{}\beta _1e^{3i(\theta _1^{}+\theta _2^{}+\theta _3^{}+\theta _3)}`$ (44)
$`\alpha _1\alpha _2^{}\beta _2`$ $`\alpha _1\alpha _2^{}\beta _2e^{i\theta _1^{}+3i(\theta _2^{}+\theta _3^{}+\theta _3)}`$ (45)
$`\alpha _1\alpha _2^{}\beta _3`$ $`\alpha _1\alpha _2^{}\beta _3e^{2i\theta _1^{}}`$ (46)
These three quantities and the four $`\gamma `$’s involve $`\theta _1^{},\theta _2^{},\theta _3^{},\theta _3`$. We now form independent combinations in which $`\theta _3`$ drops out. They are
$`\gamma _1\gamma _2^{}`$ $`\gamma _1\gamma _2^{}e^{2i\theta _1^{}}`$ (47)
$`\gamma _1\gamma _3^{}`$ $`\gamma _1\gamma _3^{}e^{i(\theta _1^{}+3\theta _2^{})}`$ (48)
$`\gamma _1\gamma _4^{}`$ $`\gamma _1\gamma _4^{}e^{i(\theta _1^{}+\theta _2^{}+4\theta _3^{})}`$ (49)
$`\alpha _1\alpha _2^{}\beta _1\gamma _4`$ $`\alpha _1\alpha _2^{}\beta _1\gamma _4e^{3i(\theta _1^{}+\theta _2^{})}`$ (50)
$`\alpha _1\alpha _2^{}\beta _2\gamma _4`$ $`\alpha _1\alpha _2^{}\beta _2\gamma _4e^{i(\theta _1^{}+3\theta _2^{})}`$ (51)
$`\alpha _1\alpha _2^{}\beta _3`$ $`\alpha _1\alpha _2^{}\beta _3e^{2i\theta _1^{}}`$ (52)
Here $`\theta _3^{}`$ appears only in the rule for $`\gamma _1\gamma _4^{}`$ so we just drop it. Then we quickly find a choice of three independent invariants:
$$arg(\alpha _1\alpha _2^{}\beta _1^{}\beta _2\beta _3);arg(\beta _1\beta _2^{}\gamma _1^{}\gamma _2);arg(\beta _2\beta _3^{}\gamma _2^{}\gamma _3\gamma _4);$$
(53)
The first of the three $`SU(4)`$ invariants is the same as the single $`SU(3)`$ invariant. This is explained by the observation that after the $`\gamma ^{}`$’s were determined in $`(\text{34})`$, the determination of the $`\beta ^{}`$’s and $`\alpha ^{}`$’s was reduced to the $`SU(3)`$ level problem - the $`SU(4)`$ expressions for the $`\beta ^{}`$’s and $`\alpha ^{}`$’s arise from those for $`SU(3)`$ in $`(\text{11})`$ and $`(\text{12})`$ by the replacements $`\theta _1\theta _1,\theta _2\theta _2,\theta _1^{}\theta _1^{},\theta _2^{}\theta _2^{}+\theta _3^{}+\theta _3`$.
The recursive procedure given above can easily be extended to $`N`$ generations.
## VI Comparison with some existing parametrizations of the mixing matrix for N=4
The parametrization due to Barger et al and Oakes correponds to choosing $`\beta _2,\gamma _2,\gamma _3`$ complex and all others real. Thus, on putting
$`\gamma _1`$ $`=`$ $`C_1;\gamma _2=S_1C_2e^{i(\delta _1+\delta _3)};\gamma _3=S_1S_2C_4e^{i\delta _2};\gamma _4=S_1S_2S_4`$ (54)
$`\beta _1`$ $`=`$ $`C_3;\beta _2=S_3C_6e^{i\delta _3};\beta _3=S_3S_6`$ (55)
$`\alpha _1`$ $`=`$ $`C_5;\alpha _2=S_5`$ (56)
in $`(\text{27})`$ and interchanging the first and the fourth columns and the second and the third we recover the parametrization in and after multiplying the second and the third row by phase factors $`e^{i(\delta _1+\delta _3)}`$ and $`e^{i\delta _2}`$ respectively.
The parametrization of the mixing matrix for $`N=4`$ due to Anselm et al is less economical. It corresponds to distributing the three phases four quantitities $`\beta _3,\gamma _2,\gamma _3,\gamma _4`$ with all others real:
$`\gamma _1`$ $`=`$ $`S_{12};\gamma _2=C_{12}C_{23}C_{24}e^{i\alpha };\gamma _3=C_{12}(S_{23}C_{24}C_{34}S_{24}S_{34})e^{i\gamma };\gamma _4=C_{12}(S_{23}C_{24}S_{34}e^{i(\beta +\gamma )}+S_{24}C_{34})e^{i\beta }`$ (57)
$`\beta _1`$ $`=`$ $`S_{13};\beta _2=C_{13}S_{23}/\sqrt{(1C_{23}^2C_{24}^2)};\beta _3=C_{13}C_{23}S_{24}e^{i(\alpha \beta )}/\sqrt{(1C_{23}^2C_{24}^2)}`$ (58)
$`\alpha _1`$ $`=`$ $`C_{14};\alpha _2=S_{14}`$ (59)
Substituting these expressions in $`(\text{27})`$ one obtains the results of Anselm et al after suitable permutation of the columns and multiplication of second and fourth row by factors $`e^{i\alpha }`$ and $`e^{i(\beta +\gamma )}`$.
The parametrization due to Harari and Leurer corresponds to choosing $`\beta _1,\gamma _1,\gamma _2`$ complex with all others real. Thus on putting
$$\alpha _1=S_{12},\alpha _2=C_{12},\beta _1=S_{13}e^{i\delta },\beta _2=S_{23}C_{13},\beta _3=C_{23}C_{13}$$
(60)
and
$$\gamma _1=S_{14}e^{i\delta _{14}};\gamma _2=C_{14}S_{24}e^{i\delta _{24}}\gamma _3=C_{14}C_{24}S_{34};\gamma _4=C_{14}C_{24}C_{34}$$
(61)
in $`(\text{27})`$ we recover their results. In fact in the Harari-Leurer parametrization, to go from $`3(orN1)`$ generations to $`4(orN)`$ generations, one needs to multiply the (appropriately augmented) mixing matrix at $`3(orN1)`$ generations, on the left, by a matrix consisting of $`3(orN1)`$ factors. In the present case, it is easily seen that the three matrices when multiplied out have precisely the same structure as in $`(\text{26})`$. In general, the $`N1`$ factors when multiplied out precisely correspond to the coset representative of $`SU(N)/SU(N1)`$ characterized by an $`N`$-dimensional complex unit vector with its first $`N2`$ components complex and the rest real.
## VII Phases in the mixing matrix and the Bargmann invariants
It is known that, under rephasing, apart from the obvious invariants $`|V_{\alpha i}|`$, the magnitudes of the matrix elements of the mixing matrix, the following quantities, quartic in $`V`$’s,
$$t_{\alpha i\beta j}V_{\alpha i}V_{\beta j}V_{\alpha j}^{}V_{\beta i}^{}$$
(62)
are invariant under the rephasing transformations
$$V_{\alpha i}e^{i\theta _\alpha ^{}}V_{\alpha i}e^{i\theta _i}$$
(63)
It is also evident that this set of invariants remains unchanged under row and column permutations.
In the present context, these invariants were first discussed by Jarlskog and by Greenberg for the case of three generations (for which there is only one independent invariant) and were later generalized to N-generations by Nieves and Pal who showed that of these the following $`(N1)(N2)/2`$ quantities can be taken as independent
$$t_{\alpha i1N};\alpha i,\alpha 1,iN$$
(64)
It can easily be verified by explicit calculations that the invariant phases given earlier for $`N=3,4`$ precisely coincide with $`arg(t_{\alpha i1N});\alpha i,\alpha 1,`$.
We would like to bring out the connection between these and the Bargmann invariants introduced by Bargmann in the context of Wigner’s unitary-antiunitary theorem. If $`\psi _1,\psi _2\mathrm{},\psi _n`$ are any $`n`$ vectors in a Hilbert space, with no two consecutive ones being orthogonal, the n-vertex Bargmann invariant is
$$\mathrm{\Delta }_n(\psi _1,\psi _2,\mathrm{},\psi _n)=<\psi _1|\psi _2><\psi _2|\psi _3>\mathrm{}<\psi _n|\psi _1>$$
(65)
It is easily seen that, under a common unitary transformation applied to all the $`\psi `$’s, and also, under independent phase changes of the $`\psi `$’s, $`\mathrm{\Delta }_n`$ remains unchanged. As an aside, we would like to remark here that there exists a deep connection between Bargmann invariants and the geometric phase as has been lucidly brought out by Mukunda and Simon .
To see the relevance of Bargmann invariants in the present context, notice that $`V`$ being a unitary matrix can be thought of as effecting a change of basis from one set of orthonormal vectors $`|f_i>`$ to another $`|e_\alpha >`$ so that $`V_{\alpha i}=<f_i|e_\alpha >`$ and one can express $`V_{\alpha i}V_{\beta j}V_{\alpha j}^{}V_{\beta i}^{}`$ as a Bargmann invariant
$$V_{\alpha i}V_{\beta j}V_{\alpha j}^{}V_{\beta i}^{}=<e_\alpha |f_i><f_i|e_\beta ><e_\beta |f_j><f_j|e_\alpha >$$
(66)
## VIII Summary
To summarize, the parametrization proposed here has the following special features:
* Introduction of $`N^{th}`$ generation requires one new $`N\times N`$ matrix determined by one $`N`$-dimensional complex unit vector, a $`SU(N)/SU(N1)`$ coset representative, multiplying the complete matrix at previous generation level after augmenting its dimension by one through bordering the last column and row suitably.
* All the invariants for $`N1`$ generations remain invariants for $`N`$-generations as well.
* One matrix of ours determined by an $`N`$-dimensional unit vector corresponds to a product of $`N1`$ factors of Harari and Leurer.
* The existing parametrizations are easily read off from our general expressions.
* Opens up new possibilities for alternatives parametrizations which may be phenomenologically useful, particularly for $`N4`$.
* The connection between the rephasing invariants and the Bargmann invariants is brought out.
We hope that the unified approach to parametrization of the mixing matrix developed here will prove to be phenomenologically useful as well. In particular, the connection between the phases and the Bargmann invariants brought out here may provide a new perspective on their origin.
One of us (SC) is grateful to Prof V. Gupta for asking a question which initiated this work. We are also grateful to Prof R. Simon for numerous discussions.
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# Instantons and the monopole-like equations in eight dimensions
## 1 Introduction
Yang-Mills instantons are among the simplest class of BPS states in the low-energy limit of superstring theory. When strings are compactified, it is often important to consider instantons on some special manifolds of dimension other than four. Mathematically, such instantons arise naturally as solutions to the eigenequations of a certain star operator acting on two-forms, and just as in the 4-dimensional case, the Yang-Mills action will reach its minimal values at these solutions. The present paper will be devoted to a study of instantons in eight dimensions.
The notion of Yang-Mills instantons in dimension greater than four is rather old and it may date back to the middle of the 1980’s . This problem has raised some renewed interest in recent years as a meanings of generalizing the Donaldson-Witten theory to higher dimensions . In particular, it is quite interesting to see whether Donaldson invariants have the holomorphic extension to Calabi-Yau four folds. Motivated by this as well as by the potential relevance to M-theory and D-brane physics, various aspects of higher dimensional cohomological Yang-Mills theories have been investigated, see e.g. . An extensive study of the relevant moduli geometry and its relations to certain calibrated submanifolds can be found in ref.. It is expected that the instanton configurations should correspond to supersymmetric D-branes embedded in some manifolds of special holonomies . In a more recent paper, Mariño, Minasian, Moore and Strominger explicitly found that a nonlinear deformation of the higher dimensional instanton equations can be derived from D-branes wrapping around supersymetric cycles, with the deformation parameter characterized by the $`B`$-field.
Presumably, the field theoretic approach to the instanton moduli problem based on BRST cohomology is perturbative in nature. The quantum degrees of freedom consist mainly of the nonabelian gauge fields $`A`$, which should be considered as fundamental fields when we try to develop a perturbative expansion in terms of the gauge coupling constant. One may ask whether there exists a nonperturbative theory within which one can use collective field variables to explore the underlying strong coupling physics . Inspired by the work of Seiberg and Witten in four dimensions, we tentatively expect that such a theory, if exists, should be closely related to a kind of $`S`$-duality. Moreover, in the dual description the collective variables should consist of an abelian gauge field together with a complex spinor satisfying certain “master equations” .
In this paper we take a modest step toward an $`S`$-dual description for the Yang-Mills instantons on some eight dimensional manifolds with special holonomy groups. Our description mimics the Seiberg-Witten theory , in which the nonabelian (anti) self-duality equation $`F^+=0`$ will be replaced by an abelian one, $`F^+=𝒬(\psi ^{},\psi )`$, with $`\psi `$ being a spinor field obeying the massless Dirac equation and $`𝒬`$ a suitable quadratic form. In writing down the explicit monopole-like equations, we shall consider two types of manifolds: Joyce manifolds of holonomy $`Spin(7)`$ as well as Calabi-Yau four-folds of holonomy $`SU(4)`$. We will compare our equations with the monopole equations constructed in 4-dimensions, and point out some problems yet to be resolved.
As a physical motivation of this investigation, we will also discuss the free abelian sector embedded in 8-dimensional Yang-Mills theory and provide a naive path-integral test of the $`S`$-duality in that sector. This discussion is an eight dimensional generalization of the usual electric-magnetic duality in four dimensions. The duality structure in eight dimensions may be alternatively understood as the existence of different gauge-fixings of a topological symmetry .
The paper is organized as follows. In Section 2 we recall some known facts about Yang-Mills instantons on manifolds of holonomy groups $`Spin(7)`$ and $`SU(4)`$. In Section 3, we give an explicit construction of the monopole-like equations. In Section 4 we turn to a discussion of the generalized $`S`$-duality in 8-dimensional abelian gauge theory. Finally we provide an Appendix where some useful properties of the 8-dimensional Clifford algebra are presented.
## 2 Yang-Mills Instantons in Eight Dimensions
Yang-Mills instantons in eight dimensions originate from a generalization of the usual concept of (anti) self-duality. Suppose that we have an eight-dimensional Riemannian manifold $`X`$ on which a closed 4-form $`\mathrm{\Omega }`$ is defined. One can use this $`\mathrm{\Omega }`$ to construct a star operator $`_\mathrm{\Omega }:\mathrm{\Lambda }^2\mathrm{\Lambda }^2`$, $`_\mathrm{\Omega }F(\mathrm{\Omega }F)`$ acting on the space of two-forms. The (anti) self-duality equations are then formulated as the eigenequation $`_\mathrm{\Omega }F=\lambda F`$ of the star operator.
In terms of components, the action of $`_\mathrm{\Omega }`$ is given by
$$(_\mathrm{\Omega }F)_{\mu \nu }=\frac{1}{2}\mathrm{\Omega }_{\mu \nu \alpha \beta }F^{\alpha \beta }.$$
(1)
Thus, if $`\mathrm{\Omega }`$ obeys an identity of the form
$$\mathrm{\Omega }_{\mu \nu \alpha \beta }\mathrm{\Omega }^{\alpha \beta \sigma \tau }=A(\delta _\mu ^\sigma \delta _\nu ^\tau \delta _\mu ^\tau \delta _\nu ^\sigma )+B\mathrm{\Omega }_{\mu \nu }^{}{}_{}{}^{\sigma \tau }$$
(2)
with some real constants $`A`$ and $`B`$ (where $`A>0`$ depends on the normalization of $`\mathrm{\Omega }`$), then the eigenequation has two solutions $`\lambda =\lambda ^\pm `$, $`F=F^\pm `$, determined by:
$$\begin{array}{ccc}\lambda ^\pm \hfill & =\hfill & \frac{B\sqrt{B^2+8A}}{4},\hfill \\ F_{\mu \nu }^\pm \hfill & =\hfill & \pm \frac{2}{\sqrt{B^2+8A}}\left(\lambda ^{}F_{\mu \nu }\frac{1}{2}\mathrm{\Omega }_{\mu \nu \alpha \beta }F^{\alpha \beta }\right).\hfill \end{array}$$
(3)
One may easily verify that $`F_{\mu \nu }=F_{\mu \nu }^++F_{\mu \nu }^{}`$. Accordingly, the space $`\mathrm{\Lambda }^2`$ of two-forms is decomposed into a direct sum $`\mathrm{\Lambda }_+^2\mathrm{\Lambda }_{}^2`$ of the eigenspaces of $`_\mathrm{\Omega }`$. We will call $`F`$ to be a self-dual form (resp. anti self-dual form) if it belongs to $`\mathrm{\Lambda }_{}^2`$ (resp. $`\mathrm{\Lambda }_+^2`$). The condition that $`F`$ is self-dual can be simply written as $`F_{\mu \nu }^+=0`$.
Actually we shall forcus on some $`G`$-bundle $`EX`$ and consider its connections $`A`$. In this context, $`A`$ is called a self-dual instanton (upto gauge transformations) if the corresponding curvature two-form $`F(A)`$ obeys the self-duality equation $`F^+(A)=0`$. An instanton will minimize the Yang-Mills action functional
$$S_{YM}[A]=\frac{1}{2g^2}_Xd^8x\sqrt{g}\mathrm{Tr}F_{\mu \nu }F^{\mu \nu }\frac{1}{2g^2}F^2.$$
(4)
In fact, if we decompose $`F`$ into components $`F^\pm \mathrm{\Lambda }_\pm ^2`$, then (3) gives:
$$S_{YM}[A]=\frac{1}{g^2(B+\sqrt{B^2+8A})}\left(2_X\mathrm{\Omega }\mathrm{Tr}(FF)+\sqrt{B^2+8A}F^+^2\right).$$
(5)
So the Yang-Mills action can be written as a non-negative term proportional to $`F^+^2`$, plus a topological invariant. Clearly, such an action will reach its minimal values at $`F^+=0`$. Manifolds with $`Spin(7)`$ holonomy. Now we briefly discuss the case when $`X`$ has the holonomy group $`Spin(7)`$. This means that $`X`$ is spin, and there is a real, non-zero parallel spinor $`\zeta S^+`$ on $`X`$ invariant under the action of $`Spin(7)Spin(8)`$. We will normalize such a spinor by imposing the condition
$$\zeta ^T\zeta =1.$$
(6)
According to the standard isomorphism $`S^+_SS^+\mathrm{\Lambda }^0\mathrm{\Lambda }_+^4`$, $`S^+S^+\mathrm{\Lambda }^2`$ between the space of forms and tensor product of the Clifford module , the 4-form $`\mathrm{\Omega }`$ considered above can be constructed as a “bispinor”
$$\mathrm{\Omega }_{\mu \nu \alpha \beta }=\zeta ^T\mathrm{\Gamma }_{\mu \nu \alpha \beta }\zeta ,$$
(7)
where $`\mathrm{\Gamma }_{\mu \nu \alpha \beta }`$ denotes the anti-symmetrized 4-fold product of the $`\gamma `$-matrices $`\mathrm{\Gamma }_\mu `$ in 8 dimensions, with a prefactor $`1/4!`$ included. Our convention of choosing the $`\gamma `$-matrices is given in the Appendix.
Eq.(7) obviously defines a $`Spin(7)`$ invariant rank-4 tensor. This tensor enjoys a couple of useful properties: First, it is covariantly constant, so that $`\mathrm{\Omega }`$ gives rise to a closed form. Second, using the $`\gamma `$-matrix identity $`\mathrm{\Gamma }_{\mu \nu \alpha \beta }\mathrm{\Gamma }_9=\frac{1}{4!}ϵ_{\mu \nu \alpha \beta \lambda \rho \sigma \tau }\mathrm{\Gamma }^{\lambda \rho \sigma \tau }`$, one easily sees that $`\mathrm{\Omega }`$ is self-dual with respect to the usual Hodge star operator, namely $`\mathrm{\Omega }=\mathrm{\Omega }`$, – this agrees with the fact that the symmetric tensor product of $`S^+`$ contains $`\mathrm{\Lambda }_+^4`$. The final property which we shall use is that (7) obeys an identity of the form (2):
$$\mathrm{\Omega }_{\mu \nu \alpha \beta }\mathrm{\Omega }^{\alpha \beta \sigma \tau }=6(\delta _\mu ^\sigma \delta _\nu ^\tau \delta _\mu ^\tau \delta _\nu ^\sigma )4\mathrm{\Omega }_{\mu \nu }^{}{}_{}{}^{\sigma \tau }.$$
(8)
In particular $`\mathrm{\Omega }`$ is normalized to be $`\mathrm{\Omega }^2\frac{1}{4!}\mathrm{\Omega }_{\mu \nu \alpha \beta }\mathrm{\Omega }^{\mu \nu \alpha \beta }=14`$.
Hence, on a manifold $`X`$ of holonomy $`Spin(7)`$, the (anti) self-duality equation $`_\mathrm{\Omega }F=\lambda F`$ has eigenvalues $`\lambda ^+=3`$, $`\lambda ^{}=1`$; the curvature two-form $`F_{\mu \nu }=F_{\mu \nu }^++F_{\mu \nu }^{}`$ is decomposed orthogonally into an anti self-dual part $`F_{\mu \nu }^+`$ and a self-dual part $`F_{\mu \nu }^{}`$, with
$$\begin{array}{ccc}F_{\mu \nu }^+\hfill & =\hfill & \frac{1}{4}\left(F_{\mu \nu }\frac{1}{2}\mathrm{\Omega }_{\mu \nu \alpha \beta }F^{\alpha \beta }\right)\mathrm{\Lambda }_+^2,\hfill \\ F_{\mu \nu }^{}\hfill & =\hfill & \frac{1}{4}\left(3F_{\mu \nu }+\frac{1}{2}\mathrm{\Omega }_{\mu \nu \alpha \beta }F^{\alpha \beta }\right)\mathrm{\Lambda }_{}^2.\hfill \end{array}$$
(9)
The Yang-Mills instantons are thus described by the equation $`F_{\mu \nu }^+=0`$. The dimensions of $`\mathrm{\Lambda }_\pm ^2`$ can be determined by a group-theoretic consideration , and the result turns out to be
$$dim\mathrm{\Lambda }_+^2=7,dim\mathrm{\Lambda }_{}^2=21.$$
(10)
For an alternative derivation of this result, note that $`\mathrm{Tr}(_\mathrm{\Omega })=0`$, so $`(3)dim\mathrm{\Lambda }_+^2+1dim\mathrm{\Lambda }_{}^2=0`$. This together with $`dim\mathrm{\Lambda }_+^2+dim\mathrm{\Lambda }_{}^2=dim\mathrm{\Lambda }^2=28`$ gives (10). Manifolds with $`SU(4)`$ holonomy. As just mentioned, $`X`$ has holonomy $`Spin(7)`$ if there is a generic parallel spinor $`\zeta 0`$ defined on it. However, this holonomy group may reduce to a subgroup of $`Spin(7)`$ when the parallel spinor obeys certain particular conditions . For example, one has the holonomy reduction $`Spin(7)SU(4)`$ provided there exists a parallel pure spinor $`\zeta `$ on $`X`$. Here we shall describe in some detail what a pure spinor is and explain why the existence of such a spinor will cause the manifold to have holonomy $`SU(4)`$ . We will then discuss a holomorphic version of the Yang-Mills instanton equations .
Let $`Cl_c(8)=Cl(8)\text{}`$ be the 8-dimensional Clifford algebra over , and let $`S_c`$ be a complex spinor space, on which an irreducible representation $`\rho ^c`$ of $`Cl_c(8)`$ is defined. Each spinor $`\psi S_c`$ can be associated to a -linear map
$$f_\psi :\text{}^8S_c,f_\psi (u)=\rho ^c(u)\psi ,$$
(11)
where $`u`$ is a complex linear combination of the Clifford generators $`e_\mu Cl(8)`$ and we have identified the space of all such linear combinations with $`\text{}^8`$. Let us consider the kernel of this map in the case $`\psi 0`$. If $`u=a^\mu e_\mu `$, $`v=b^\mu e_\mu \mathrm{Ker}f_\psi `$ (with complex coefficients $`a^\mu `$, $`b^\mu `$), then $`\rho ^c(u)\psi =\rho ^c(v)\psi =00=\rho ^c(\{u,v\})\psi =2g_{\mu \nu }a^\mu b^\nu \psi `$. It follows that the space $`\mathrm{Ker}f_\psi `$ is orthogonal to its complex conjugate $`\overline{\mathrm{Ker}f_\psi }`$ with respect to the standard hermitian inner product $`(a^\mu e_\mu ,b^\nu e_\nu )a^\mu e_\mu ,\overline{b}^\nu e_\nu =g_{\mu \nu }a^\mu \overline{b}^\nu `$ on $`\text{}^8`$. Thus, since $`\mathrm{Ker}f_\psi \overline{\mathrm{Ker}f_\psi }\text{}^8`$, we see that the complex dimensions of $`\mathrm{Ker}f_\psi `$ should not exceed 4. By definition, we say that $`\psi `$ is a pure spinor if $`dim_c\mathrm{Ker}f_\psi `$ reaches its maximally allowed value 4.
We will take the complex spinor space to be the complexification of the real $`Cl(8)`$-module $`S\text{}^{16}`$: $`S_c=S\text{}`$. Spinors in such a space can be written as linear combinations of a basis of $`S`$ with complex coefficients. Also, one takes $`\rho ^c`$ to be the -linear extension of the $`\gamma `$-matrix representation of $`Cl(8)`$. This allows us to choose $`\mathrm{\Gamma }_\mu \rho ^c(e_\mu )`$ as given in the Appendix.
Since $`\mathrm{\Gamma }_9=I(I)`$, $`S_c`$ is decomposed into subspaces $`S_c^\pm `$ of positive and negative chiralities. By construction, $`\psi S_c^+`$ means that $`\psi `$ is a linear combination of some real spinors in $`S^+`$ with complex coefficients, so its complex conjugate $`\overline{\psi }`$ also has positive chirality. One can show that if $`\psi `$ is a pure spinor, then either $`\psi `$ will be entirely in $`S_c^+`$ or it will be entirely in $`S_c^{}`$; namely, $`\psi `$ has a definite chirality. To see this, note that a change in the orthonormal basis $`\{e_\mu \}`$ of $`\text{}^8Cl(8)`$ will leave the matrix $`\mathrm{\Gamma }_9=\rho ^c(e_1\mathrm{}e_8)`$ invariant, upto a factor $`\pm 1`$ depending on the relative ordering. Also note that if $`\psi `$ is a pure spinor, then $`\mathrm{Ker}f_\psi \overline{\mathrm{Ker}f_\psi }=\text{}^8`$, so any basis $`\{h_{\overline{i}},1\overline{i}4\}`$ of $`\mathrm{Ker}f_\psi `$ along with its complex conjugate $`\{h_i\}\overline{\mathrm{Ker}f_\psi }`$ provides a basis of $`\text{}^8`$. We can take $`\{h_{\overline{i}}\}`$ to be orthonormal with respect to the natural hermitian metric on $`\text{}^8`$. Then the following vectors
$$\begin{array}{cccc}e_1=\frac{1}{\sqrt{2}}(h_1+h_{\overline{1}}),\hfill & e_2=\frac{1}{\sqrt{2}i}(h_1h_{\overline{1}}),\hfill & e_3=\frac{1}{\sqrt{2}}(h_2+h_{\overline{2}}),\hfill & e_4=\frac{1}{\sqrt{2}i}(h_2h_{\overline{2}})\hfill \\ e_5=\frac{1}{\sqrt{2}}(h_3+h_{\overline{3}}),\hfill & e_6=\frac{1}{\sqrt{2}i}(h_3h_{\overline{3}}),\hfill & e_7=\frac{1}{\sqrt{2}}(h_4+h_{\overline{4}}),\hfill & e_8=\frac{1}{\sqrt{2}i}(h_4h_{\overline{4}})\hfill \end{array}$$
(12)
form an orthonormal basis of $`\text{}^8`$, as they are all invariant under complex conjugation. Now as $`h_{\overline{1}}\mathrm{Ker}f_\psi `$, we have $`\rho ^c(e_1ie_2)\psi =0\rho ^c(e_1^2ie_1e_2)\psi =0\rho ^c(e_1e_2)\psi =i\psi `$. Similar arguments lead to $`\rho ^c(e_3e_4)\psi =\rho ^c(e_5e_6)\psi =\rho ^c(e_7e_8)\psi =i\psi `$. Thus we find $`\mathrm{\Gamma }_9\psi =\psi `$, indicating that $`\psi `$ has the positive chirality. A choice of the basis with different ordering will give $`\mathrm{\Gamma }_9\psi =\psi `$, but in any case the chirality of a pure spinor is definite.
Given now a pure spinor $`\zeta S_c^+`$, let us consider the maximal subgroup of $`Spin(8)`$ that keeps $`\zeta `$ invariant. An element of $`Spin(8)`$ can act adjointly on the real vector space spanned by $`\{e_\mu \}`$,
$$e_\mu e_\mu ^{}=g^1e_\mu g\rho _v(g)_{\mu }^{}{}_{}{}^{\nu }e_\nu ,\rho _v(g)SO(8),$$
(13)
which defines the representation $`\mathrm{𝟖}_v`$ of $`Spin(8)`$. The action (13) may be viewed as a change in the basis and it clearly preserves both the orthonormal property and the ordering of the basis. The -linear extension of $`\rho _v`$, which we will denote by $`\rho _v^c`$, is defined naturally on the space $`\text{}^8=\mathrm{Ker}f_\zeta \overline{\mathrm{Ker}f_\zeta }`$. For generic $`gSpin(8)`$, neither the subspace $`\mathrm{Ker}f_\zeta `$ nor $`\overline{\mathrm{Ker}f_\zeta }`$ is invariant under the action of $`\rho _v^c(g)`$. Elements of $`SO(8)`$ that leaves these subspaces invariant will map one orthonormal basis $`\{h_i\}\overline{\mathrm{Ker}f_\zeta }`$ (and $`\{h_{\overline{i}}\}\mathrm{Ker}f_\zeta `$) into another, thus forming the subgroup $`SU(4)`$. Such elements arise from those $`gSpin(8)`$ keeping $`\zeta `$ invariant. Indeed, for any $`h\mathrm{Ker}f_\zeta `$ and $`\rho ^c(g)\zeta =\zeta `$, we have $`\rho ^c(g^1hg)\zeta =\rho ^c(g)^1\rho ^c(h)\rho ^c(g)\zeta =\rho ^c(g)^1\rho ^c(h)\zeta =0g^1hg\mathrm{Ker}f_\zeta `$. We thus conclude that the isotropy group of a pure spinor $`\zeta S_c^+`$ is $`SU(4)`$. A globalized version of this discussion leads to the statement :
There exists a parallel pure spinor on $`X`$ $`X`$ has the holonomy group $`SU(4)`$ (or its subgroup).
Now we take a pure spinor $`\zeta S_c^+`$ and fix the almost complex structure on $`\text{}^8`$ as in (12), so that the basis $`h_i`$ of $`\overline{\mathrm{Ker}f_\zeta }`$ and the basis $`h_{\overline{i}}`$ of $`\mathrm{Ker}f_\zeta `$ have the $`\gamma `$-matrix representation:
$$\gamma _1=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_1+i\mathrm{\Gamma }_2),\gamma _2=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_3+i\mathrm{\Gamma }_3),\gamma _3=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_5+i\mathrm{\Gamma }_6),\gamma _4=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_7+i\mathrm{\Gamma }_8)$$
(14)
$$\gamma _{\overline{1}}=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_1i\mathrm{\Gamma }_2),\gamma _{\overline{2}}=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_3i\mathrm{\Gamma }_3),\gamma _{\overline{3}}=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_5i\mathrm{\Gamma }_6),\gamma _{\overline{4}}=\frac{1}{\sqrt{2}}(\mathrm{\Gamma }_7i\mathrm{\Gamma }_8).$$
(15)
The complex Clifford algebra is determined by the relations
$$\gamma _i\gamma _j+\gamma _j\gamma _i=\gamma _{\overline{i}}\gamma _{\overline{j}}+\gamma _{\overline{j}}\gamma _{\overline{i}}=0,\gamma _i\gamma _{\overline{j}}+\gamma _{\overline{j}}\gamma _i=2g_{i\overline{j}}.$$
(16)
We also need the dual basis
$$\gamma ^i=g^{i\overline{j}}\gamma _{\overline{j}},\gamma ^{\overline{i}}=\gamma _jg^{j\overline{i}}$$
(17)
as well as their anti-symmetrized products $`\gamma ^{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}`$. With this notation, a $`(p,q)`$-form $`t\mathrm{\Lambda }^{p,q}`$ has a natural representation in terms of the $`\gamma `$-matrices:
$$t\frac{1}{p!q!}t_{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}\gamma ^{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}.$$
(18)
Moreover, each such form should be associated to a bispinor $`\varphi ^{}\gamma _{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}\psi S_cS_c`$ as in the real case. Note that the isomorphism between the tensor product of spinors and forms has a -bilinear extension to the complex case
$$\rho ^c\rho ^c2(1+\rho _v^c+^2\rho _v^c+^3\rho _v^c)+^4\rho _v^c,$$
(19)
which shows that $`(S_c^+S_c^{})(S_c^+S_c^{})`$ can be identified with forms in $`^{}\text{}^8`$.
To warm up the complex Clifford calculus, let us establish an isomorphism between $`S_c^\pm \zeta ^{}`$ and certain particular forms. Since $`\gamma _{\overline{i}}`$ is in $`\mathrm{Ker}f_\zeta `$, we have
$$\gamma _{\overline{i}}\zeta =\gamma ^i\zeta =0\zeta ^{}\gamma _i=\zeta ^{}\gamma ^{\overline{i}}=0,$$
(20)
and this gives to $`\zeta ^{}\gamma _i\gamma _{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}=0`$. One may use this and the $`\gamma `$-matrix identity
$$\gamma _i\gamma _{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}=\gamma _{ii_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}+\underset{k=1}{\overset{q}{}}(1)^{k+p}g_{i\overline{j}_k}\gamma _{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\widehat{\overline{j}_k}\mathrm{}\overline{j}_q}$$
(21)
to deduce that the $`(p+1,q)`$ type bispinor $`\zeta ^{}\gamma _{ii_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}\psi `$ is in fact a linear combination of some $`(p,q1)`$ forms. This process can be proceeded inductively and we find that the $`(p+1,q)`$-bispinor finally becomes a linear combination of $`\zeta ^{}\gamma _{\overline{j}_1\mathrm{}\overline{j}_{qp1}}\psi `$ if $`qp+1`$ or a linear combination of $`\zeta ^{}\gamma _{i_1\mathrm{}i_{p+1q}}\psi 0`$ if $`q<p+1`$. Accordingly, for arbitrary $`\psi S_c=S_c^+S_c^{}`$, the tensor product $`\zeta ^{}\psi `$ can be identified to a form in $`\mathrm{\Lambda }^{0,}`$. Note that with our convention of the $`\gamma `$-matrices, $`\gamma _{\overline{j}_1\mathrm{}\overline{j}_q}`$ is block diagonal for $`q=`$ even and off-diagonal for $`q=`$ odd. It follows that
$$S_c^+\text{}\mathrm{\Lambda }^{0,\mathrm{even}},S_c^{}\text{}\mathrm{\Lambda }^{0,\mathrm{odd}}$$
(22)
here is the complex 1-dimensional space generated by $`\zeta ^{}`$. Similarly, tensor products $`\psi _\pm \zeta `$ for $`\psi _\pm S_c^\pm `$ should be identified with a form in $`\mathrm{\Lambda }^{\mathrm{even},0}`$ and in $`\mathrm{\Lambda }^{\mathrm{odd},0}`$, respectively.
Now we give a suitable normalization of $`\zeta `$. In the $`Spin(7)`$ case we have simply imposed the condition $`\zeta ^T\zeta =1`$. However, this normalization condition cannot be adopted here for a pure spinor $`\zeta `$. In fact from (20)-(21) we see that $`\gamma _{i\overline{j}}\zeta =g_{i\overline{j}}\zeta `$, so that $`\zeta ^T\gamma _{i\overline{j}}\zeta =g_{i\overline{j}}(\zeta ^T\zeta )`$, which together with the anti-symmetric property of the matrix $`\gamma _{i\overline{j}}`$ implies $`\zeta ^T\zeta =0`$. Nevertheless, one can still impose another normalization condition
$$\zeta ^{}\zeta =1,$$
(23)
and this looks more natural when we work in complex spaces. Using this normlization, we define an $`SU(4)`$ invariant closed (4,0)-form $`\mathrm{\Omega }`$ with the components
$$\mathrm{\Omega }_{ijkl}=\zeta ^T\gamma _{ijkl}\zeta .$$
(24)
Some properties of (24) can be explored using a complex version of the Fierz rearrangement formula:
$$\zeta \zeta ^T=\frac{1}{164!}\mathrm{\Omega }_{ijkl}\gamma ^{ijkl},$$
(25)
$$\zeta \zeta ^{}=\frac{1}{16}(1+g^{i\overline{j}}\gamma _{i\overline{j}})(1+\mathrm{\Gamma }_9)+\frac{1}{32}g^{i\overline{l}}g^{j\overline{k}}\gamma _{ij\overline{k}\overline{l}}.$$
(26)
For example one may apply (25) to a quick computation of the norm $`\mathrm{\Omega }^2`$. By definition, $`\mathrm{\Omega }^2\frac{1}{4!}\mathrm{\Omega }_{ijkl}\overline{\mathrm{\Omega }}^{ijkl}`$, $`\overline{\mathrm{\Omega }}^{ijkl}g^{i\overline{i}}g^{j\overline{j}}g^{k\overline{k}}g^{l\overline{l}}\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}`$, where $`\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}\mathrm{\Lambda }^{0,4}`$ is the complex conjugate of $`\mathrm{\Omega }`$. Since $`\overline{\mathrm{\Omega }}^{ijkl}=\zeta ^{}\gamma ^{ijkl}\overline{\zeta }`$, we see that (25) gives $`\mathrm{\Omega }^2=16`$. Notice that this normalization of $`\mathrm{\Omega }`$ is different from the $`Spin(7)`$ case, where $`\mathrm{\Omega }^2=14`$. As an application of (26), one can establish a more useful identity
$$\mathrm{\Omega }_{ijkl}\overline{\mathrm{\Omega }}^{mnkl}=32(\delta _i^m\delta _j^n\delta _i^n\delta _j^m),$$
(27)
which takes a form similar to (2).
We turn now to the self-duality equations. Given the $`SU(4)`$ invariant (4,0)-form $`\mathrm{\Omega }_{ijkl}`$ defined as above, its complex conjugate $`\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}`$, a (0,4)-form, may be used to construct an anti-linear star operator $`_\mathrm{\Omega }:\mathrm{\Lambda }^{0,2}\mathrm{\Lambda }^{0,2}`$ by means of
$$(_\mathrm{\Omega }\beta )_{\overline{i}\overline{j}}=\frac{1}{2}\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}\overline{\beta }^{\overline{k}\overline{l}},\beta _{\overline{i}\overline{j}}\mathrm{\Lambda }^{0,2},$$
(28)
where $`\overline{\beta }_{ij}\mathrm{\Lambda }^{2,0}`$ denotes the complex conjugate of $`\beta _{\overline{i}\overline{j}}`$ and $`\overline{\beta }^{\overline{k}\overline{l}}g^{i\overline{k}}g^{j\overline{l}}\overline{\beta }_{ij}`$. Thus, if $`F\mathrm{\Lambda }^2`$ is a curvature 2-form, we can decompose it into $`F=F^{(2,0)}+F^{(1,1)}+F^{(0,2)}`$ with $`F^{(2,0)}=\overline{F^{(0,2)}}`$ (assuming that the connection is unitary), and define
$$(_\mathrm{\Omega }F^{(0,2)})_{\overline{i}\overline{j}}=\frac{1}{2}\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}F^{(2,0)\overline{k}\overline{l}},(_\mathrm{\Omega }F^{(2,0)})_{ij}=\frac{1}{2}\mathrm{\Omega }_{ijkl}F^{(0,2)kl}.$$
(29)
Note that the (1,1)-component of $`F`$ is intact under the action of $`_\mathrm{\Omega }`$.
Just as in the $`Spin(7)`$ case, the (anti) self-duality equations should be formulated as the eigenvalue equation of $`_\mathrm{\Omega }`$. Hence, in terms of components, we call $`F_{\mu \nu }`$ to be (anti) self-dual if they satisfy the conditions
$$\frac{1}{2}\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}F^{(2,0)\overline{k}\overline{l}}=\lambda F_{\overline{i}\overline{j}}^{(0,2)},\frac{1}{2}\mathrm{\Omega }_{ijkl}F^{(0,2)kl}=\lambda F_{ij}^{(2,0)}.$$
(30)
Here the eigenvalues $`\lambda \text{}`$ are determined by
$$\lambda =\lambda _\pm ,\lambda _+=4,\lambda _{}=4.$$
(31)
Accordingly, the space of (0,2)-forms gets decomposed into the two eigenspaces of $`_\mathrm{\Omega }`$, $`\mathrm{\Lambda }^{0,2}=\mathrm{\Lambda }_+^{0,2}\mathrm{\Lambda }_{}^{0,2}`$, where $`\mathrm{\Lambda }_\pm ^{0,2}`$ correspond to the eigenvalues $`\lambda _\pm `$, respectively. The (0,2) component of $`F\mathrm{\Lambda }^2`$ then decomposes into an anti self-dual part $`F_+^{(0,2)}\mathrm{\Lambda }_+^{0,2}`$ and a self-dual part $`F_{}^{(0,2)}\mathrm{\Lambda }_{}^{0,2}`$, with
$$F_{\pm \overline{i}\overline{j}}^{(0,2)}=\frac{1}{2}\left(F_{\overline{i}\overline{j}}^{(0,2)}\pm \frac{1}{8}\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}F^{(2,0)\overline{k}\overline{l}}\right).$$
(32)
Holomorphic Yang-Mills instantons are thus characterized by the self-duality equation
$$F_+^{(0,2)}(A)=0\frac{1}{2}\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}F^{(2,0)\overline{k}\overline{l}}=4F_{\overline{i}\overline{j}}^{(0,2)}.$$
(33)
Sometimes it is useful to have a more compact description for the $`_\mathrm{\Omega }`$ operator, without reference to the unitary basis given in (12). To give such a description, note that there is a natural hermitian inner product on $`\mathrm{\Lambda }^{0,q}`$: for two arbitrary $`(0,q)`$ forms $`\alpha _{\overline{i}_1\mathrm{}\overline{i}_q}`$ and $`\beta _{\overline{i}_1\mathrm{}\overline{i}_q}`$, we can define an $`SU(4)`$ invariant paring
$$\alpha ,\beta \frac{1}{q!}\alpha _{\overline{i}_1\mathrm{}\overline{i}_q}\overline{\beta }^{\overline{i}_1\mathrm{}\overline{i}_q}=\frac{1}{q!}g^{i_1\overline{j}_1}\mathrm{}g^{i_q\overline{j}_q}\overline{\beta }_{i_1\mathrm{}i_q}\alpha _{\overline{j}_1\mathrm{}\overline{j}_q},$$
(34)
which is linear in $`\alpha `$ and anti-linear in $`\beta `$. In terms of this inner product, one then introduces an operator $`_\mathrm{\Omega }:\mathrm{\Lambda }^{0,q}\mathrm{\Lambda }^{0,4q}`$ through
$$\alpha _\mathrm{\Omega }\beta =\alpha ,\beta \overline{\mathrm{\Omega }}.$$
(35)
Clearly, this description manifests the $`SU(4)`$ invariance and does not depend on a particular choice of the basis of $`\mathrm{Ker}f_\zeta `$. One may see that this definition agrees with the previous one for $`q=2`$ .
Actually, it is possible to consider a slightly generalized case where we have an $`SU(n)`$ invariant $`(n,0)`$-form $`\mathrm{\Omega }_{i_1\mathrm{}i_n}`$ defined on some $`2n`$-dimensional space $`X`$. In that case, the star operator constructed by (35) should map $`\beta \mathrm{\Lambda }^{0,q}`$ into $`_\mathrm{\Omega }\beta \mathrm{\Lambda }^{0,nq}`$, so that $`\alpha _\mathrm{\Omega }\beta `$ is a $`(0,n)`$-form. The component of the left hand side of (35) is $`\frac{1}{q!(nq)!}\alpha _{[\overline{i}_1\mathrm{}\overline{i}_q}(_\mathrm{\Omega }\beta )_{\overline{i}_{q+1}\mathrm{}\overline{i}_n]}`$, while the component of the right hand side of (35) is $`\frac{1}{n!q!}\alpha _{\overline{j}_1\mathrm{}\overline{j}_q}\overline{\beta }^{\overline{j}_1\mathrm{}\overline{j}_q}\overline{\mathrm{\Omega }}_{\overline{i}_1\mathrm{}\overline{i}_n}`$; making them equal to each other for arbitrary $`\alpha \mathrm{\Lambda }^{0,q}`$ leads to
$$\delta _{[\overline{i}_1}^{\overline{j}_1}\mathrm{}\delta _{\overline{i}_n]}^{\overline{j}_n}(_\mathrm{\Omega }\beta )_{\overline{j}_{q+1}\mathrm{}\overline{j}_n}=\frac{(nq)!}{n!}\overline{\beta }^{\overline{j}_1\mathrm{}\overline{j}_q}\overline{\mathrm{\Omega }}_{\overline{i}_1\mathrm{}\overline{i}_n}.$$
By contracting the $`q`$ pairs $`(\overline{i}_1,\overline{j}_1),\mathrm{},(\overline{i}_q,\overline{j}_q)`$ of the tensor indices in this equation, and then using the identity
$$\delta _{[\overline{i}_1}^{\overline{i}_1}\mathrm{}\delta _{\overline{i}_q}^{\overline{i}_q}\delta _{\overline{i}_{q+1}}^{\overline{j}_{q+1}}\mathrm{}\delta _{\overline{i}_n]}^{\overline{j}_n}=\frac{q!(nq)!}{n!}\delta _{[\overline{i}_{q+1}}^{\overline{j}_{q+1}}\mathrm{}\delta _{\overline{i}_n]}^{\overline{j}_n},$$
we see that
$$(_\mathrm{\Omega }\beta )_{\overline{i}_{q+1}\mathrm{}\overline{i}_n}=\frac{1}{q!}\overline{\beta }^{\overline{i}_1\mathrm{}\overline{i}_q}\overline{\mathrm{\Omega }}_{\overline{i}_1\mathrm{}\overline{i}_q\overline{i}_{q+1}\mathrm{}\overline{i}_n}=\frac{(1)^{q(nq)}}{q!}\overline{\mathrm{\Omega }}_{\overline{i}_{q+1}\mathrm{}\overline{i}_n\overline{i}_1\mathrm{}\overline{i}_q}\overline{\beta }^{\overline{i}_1\mathrm{}\overline{i}_q}.$$
In particular for $`n=4`$ and $`q=2`$, this reduces to our earlier definition (28).
## 3 The Monopole-like Equations
Non-abelian instantons constitute a moduli problem. In 4-dimensions, this problem can be transformed into a simpler problem, where the gauge fields $`A`$ are taken to be abelian and one introduces certain new degrees of freedom – a spinor $`\psi `$, which satisfies the massless Dirac equations $`\mathrm{\Gamma }_\mu D_A^\mu \psi =0`$. The couplings between $`A`$ and $`\psi `$ are described by, in addtion to the Dirac eqautions, a non-linear relation $`F^+(A)=𝒬(\psi ,\overline{\psi })`$, where $`𝒬`$ is some quadratic form in $`\psi `$, taking values in the anti self-dual part $`\mathrm{\Lambda }_+^2`$ of two-forms. This is the basic setup of the Seiberg-Witten theory . Now a natural question arises as whether we can find an 8-dimensional analog of such a theory. Manifolds with $`Spin(7)`$ Holonomy. On 8-dimensional manifold $`X`$ with $`Spin(7)`$ holonomy, there also exists a natural quadratic form $`𝒬(\psi ,\overline{\psi })`$ valued in $`\mathrm{\Lambda }_+^2`$. Indeed, given a complex line bundle $``$ and a spinor field $`\psi S^+`$, one can construct a two-form $`\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi =\psi ^T\mathrm{\Gamma }_{\mu \nu }\zeta `$ and, according to , it takes values in $`\mathrm{\Lambda }_+^2`$. One can also form the inner product $`\overline{\psi }^T\zeta \mathrm{\Lambda }^0^1`$. It follows that the quadratic form $`𝒬_{\mu \nu }(\psi ,\overline{\psi })=(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )`$ belongs to $`(\mathrm{\Lambda }^0^1)(\mathrm{\Lambda }_+^2)\mathrm{\Lambda }_+^2\text{}`$. Thus, by choosing a unitary connection $`A`$ of $``$, it is possible to write down an 8-dimensional analog of the Seiberg-Witten equations
$$\begin{array}{ccc}F_{\mu \nu }^+(A)\hfill & =\hfill & ia\mathrm{}\left[(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )\right]+ib\mathrm{}\left[(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )\right],\hfill \\ \mathrm{\Gamma }_\mu D_A^\mu \psi \hfill & =\hfill & 0\hfill \end{array}$$
(36)
where $`a`$, $`b`$ are real constants.
In order to see that both of the real part and the imaginary part of $`𝒬`$ are not necessarily vanishing for generic $`\psi S^+`$, one may work out $`(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )`$ in a fully explicit form. Using the $`\gamma `$-matrices given in the Appendix we find
$$\zeta ^T\mathrm{\Gamma }_{12}\psi =\zeta _1\psi _2+\zeta _2\psi _1+\zeta _3\psi _4\zeta _4\psi _3\zeta _5\psi _6+\zeta _6\psi _5+\zeta _7\psi _8\zeta _8\psi _7,$$
$$\zeta ^T\mathrm{\Gamma }_{13}\psi =\zeta _1\psi _4\zeta _2\psi _3+\zeta _3\psi _2\zeta _4\psi _1+\zeta _5\psi _8\zeta _6\psi _7+\zeta _7\psi _6\zeta _8\psi _5,$$
$$\zeta ^T\mathrm{\Gamma }_{14}\psi =\zeta _1\psi _5+\zeta _2\psi _6\zeta _3\psi _7+\zeta _4\psi _8+\zeta _5\psi _1\zeta _6\psi _2+\zeta _7\psi _3\zeta _8\psi _4,$$
$$\zeta ^T\mathrm{\Gamma }_{15}\psi =\zeta _1\psi _6+\zeta _2\psi _5+\zeta _3\psi _8+\zeta _4\psi _7\zeta _5\psi _2\zeta _6\psi _1\zeta _7\psi _4\zeta _8\psi _3,$$
$$\zeta ^T\mathrm{\Gamma }_{16}\psi =\zeta _1\psi _3\zeta _2\psi _4+\zeta _3\psi _1+\zeta _4\psi _2+\zeta _5\psi _7+\zeta _6\psi _8\zeta _7\psi _5\zeta _8\psi _6,$$
$$\zeta ^T\mathrm{\Gamma }_{17}\psi =\zeta _1\psi _7+\zeta _2\psi _8\zeta _3\psi _5\zeta _4\psi _6+\zeta _5\psi _3+\zeta _6\psi _4\zeta _7\psi _1\zeta _8\psi _2,$$
$$\zeta ^T\mathrm{\Gamma }_{18}\psi =\zeta _1\psi _8+\zeta _2\psi _7+\zeta _3\psi _6\zeta _4\psi _5+\zeta _5\psi _4\zeta _6\psi _3\zeta _7\psi _2+\zeta _8\psi _1,$$
$$\zeta ^T\mathrm{\Gamma }_{23}\psi =\zeta _1\psi _3+\zeta _2\psi _4\zeta _3\psi _1\zeta _4\psi _2+\zeta _5\psi _7+\zeta _6\psi _8\zeta _7\psi _5\zeta _8\psi _6,$$
$$\zeta ^T\mathrm{\Gamma }_{24}\psi =\zeta _1\psi _6\zeta _2\psi _5+\zeta _3\psi _8+\zeta _4\psi _7+\zeta _5\psi _2+\zeta _6\psi _1\zeta _7\psi _4\zeta _8\psi _3,$$
$$\zeta ^T\mathrm{\Gamma }_{25}\psi =\zeta _1\psi _5+\zeta _2\psi _6+\zeta _3\psi _7\zeta _4\psi _8+\zeta _5\psi _1\zeta _6\psi _2\zeta _7\psi _3+\zeta _8\psi _4,$$
$$\zeta ^T\mathrm{\Gamma }_{26}\psi =\zeta _1\psi _4\zeta _2\psi _3+\zeta _3\psi _2\zeta _4\psi _1\zeta _5\psi _8+\zeta _6\psi _7\zeta _7\psi _6+\zeta _8\psi _5,$$
$$\zeta ^T\mathrm{\Gamma }_{27}\psi =\zeta _1\psi _8+\zeta _2\psi _7\zeta _3\psi _6+\zeta _4\psi _5\zeta _5\psi _4+\zeta _6\psi _3\zeta _7\psi _2+\zeta _8\psi _1,$$
$$\zeta ^T\mathrm{\Gamma }_{28}\psi =\zeta _1\psi _7\zeta _2\psi _8+\zeta _3\psi _5\zeta _4\psi _6+\zeta _5\psi _3+\zeta _6\psi _4+\zeta _7\psi _1+\zeta _8\psi _2,$$
$$\zeta ^T\mathrm{\Gamma }_{34}\psi =\zeta _1\psi _8+\zeta _2\psi _7+\zeta _3\psi _6+\zeta _4\psi _5\zeta _5\psi _4\zeta _6\psi _3\zeta _7\psi _2\zeta _8\psi _1,$$
$$\zeta ^T\mathrm{\Gamma }_{35}\psi =\zeta _1\psi _7\zeta _2\psi _8+\zeta _3\psi _5\zeta _4\psi _6\zeta _5\psi _3+\zeta _6\psi _4\zeta _7\psi _1+\zeta _8\psi _2,$$
$$\zeta ^T\mathrm{\Gamma }_{36}\psi =\zeta _1\psi _2\zeta _2\psi _1\zeta _3\psi _4+\zeta _4\psi _3\zeta _5\psi _6+\zeta _6\psi _5+\zeta _7\psi _8\zeta _8\psi _7,$$
$$\zeta ^T\mathrm{\Gamma }_{37}\psi =\zeta _1\psi _6+\zeta _2\psi _5+\zeta _3\psi _8\zeta _4\psi _7\zeta _5\psi _2+\zeta _6\psi _1+\zeta _7\psi _4\zeta _8\psi _3,$$
$$\zeta ^T\mathrm{\Gamma }_{38}\psi =\zeta _1\psi _5\zeta _2\psi _6+\zeta _3\psi _7+\zeta _4\psi _8+\zeta _5\psi _1+\zeta _6\psi _2\zeta _7\psi _3\zeta _8\psi _4,$$
$$\zeta ^T\mathrm{\Gamma }_{45}\psi =\zeta _1\psi _2\zeta _2\psi _1+\zeta _3\psi _4\zeta _4\psi _3+\zeta _5\psi _6\zeta _6\psi _5+\zeta _7\psi _8\zeta _8\psi _7,$$
$$\zeta ^T\mathrm{\Gamma }_{46}\psi =\zeta _1\psi _7+\zeta _2\psi _8+\zeta _3\psi _5\zeta _4\psi _6\zeta _5\psi _3+\zeta _6\psi _4+\zeta _7\psi _1\zeta _8\psi _2,$$
$$\zeta ^T\mathrm{\Gamma }_{47}\psi =\zeta _1\psi _3+\zeta _2\psi _4+\zeta _3\psi _1\zeta _4\psi _2+\zeta _5\psi _7\zeta _6\psi _8\zeta _7\psi _5+\zeta _8\psi _6,$$
$$\zeta ^T\mathrm{\Gamma }_{48}\psi =\zeta _1\psi _4\zeta _2\psi _3+\zeta _3\psi _2+\zeta _4\psi _1\zeta _5\psi _8\zeta _6\psi _7+\zeta _7\psi _6+\zeta _8\psi _5,$$
$$\zeta ^T\mathrm{\Gamma }_{56}\psi =\zeta _1\psi _8+\zeta _2\psi _7\zeta _3\psi _6\zeta _4\psi _5+\zeta _5\psi _4+\zeta _6\psi _3\zeta _7\psi _2\zeta _8\psi _1,$$
$$\zeta ^T\mathrm{\Gamma }_{57}\psi =\zeta _1\psi _4+\zeta _2\psi _3\zeta _3\psi _2\zeta _4\psi _1\zeta _5\psi _8\zeta _6\psi _7+\zeta _7\psi _6+\zeta _8\psi _5,$$
$$\zeta ^T\mathrm{\Gamma }_{58}\psi =\zeta _1\psi _3+\zeta _2\psi _4+\zeta _3\psi _1\zeta _4\psi _2\zeta _5\psi _7+\zeta _6\psi _8+\zeta _7\psi _5\zeta _8\psi _6,$$
$$\zeta ^T\mathrm{\Gamma }_{67}\psi =\zeta _1\psi _5+\zeta _2\psi _6+\zeta _3\psi _7+\zeta _4\psi _8\zeta _5\psi _1\zeta _6\psi _2\zeta _7\psi _3\zeta _8\psi _4,$$
$$\zeta ^T\mathrm{\Gamma }_{68}\psi =\zeta _1\psi _6+\zeta _2\psi _5\zeta _3\psi _8+\zeta _4\psi _7\zeta _5\psi _2+\zeta _6\psi _1\zeta _7\psi _4+\zeta _8\psi _3,$$
$$\zeta ^T\mathrm{\Gamma }_{78}\psi =\zeta _1\psi _2+\zeta _2\psi _1\zeta _3\psi _4+\zeta _4\psi _3+\zeta _5\psi _6\zeta _6\psi _5+\zeta _7\psi _8\zeta _8\psi _7.$$
So we get, for example,
$$(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{12}\psi )=\left(\underset{A=1}{\overset{8}{}}\overline{\psi }_A\zeta _A\right)\underset{a=1}{\overset{4}{}}(1)^a\left(\zeta _{2a1}\psi _{2a}\zeta _{2a}\psi _{2a1}\right).$$
(37)
If we write $`\psi =\chi +i\eta `$, then the real and imaginary parts of (37) read
$$\begin{array}{ccc}\mathrm{}\left[(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{12}\psi )\right]\hfill & =\hfill & \left(\underset{A=1}{\overset{8}{}}\chi _A\zeta _A\right)\underset{a=1}{\overset{4}{}}(1)^a\left(\zeta _{2a1}\chi _{2a}\zeta _{2a}\chi _{2a1}\right)\hfill \\ & +\hfill & \left(\underset{A=1}{\overset{8}{}}\eta _A\zeta _A\right)\underset{a=1}{\overset{4}{}}(1)^a\left(\zeta _{2a1}\eta _{2a}\zeta _{2a}\eta _{2a1}\right),\hfill \end{array}$$
(38)
$$\begin{array}{ccc}\mathrm{}\left[(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{12}\psi )\right]\hfill & =\hfill & \left(\underset{A=1}{\overset{8}{}}\chi _A\zeta _A\right)\underset{a=1}{\overset{4}{}}(1)^a\left(\zeta _{2a1}\eta _{2a}\zeta _{2a}\eta _{2a1}\right)\hfill \\ & \hfill & \left(\underset{A=1}{\overset{8}{}}\eta _A\zeta _A\right)\underset{a=1}{\overset{4}{}}(1)^a\left(\zeta _{2a1}\chi _{2a}\zeta _{2a}\chi _{2a1}\right).\hfill \end{array}$$
(39)
Other $`(\mu ,\nu )`$-components can be written down similarly.
The above explicit result shows that for generic spinors $`\psi `$, both the real part and the imaginary part of $`𝒬_{\mu \nu }(\psi ,\overline{\psi })`$ are indeed not zero. This is different from the 4-dimensional Seiberg-Witten theory, where the quadratic form $`𝒬_{\mu \nu }(\psi ,\overline{\psi })=\overline{\psi }^T\mathrm{\Gamma }_{\mu \nu }\psi `$ is essentially purely imaginary, as we can choose the $`Spin(4)`$ Lie algebra generators $`\mathrm{\Gamma }_{\mu \nu }`$ to be anti-hermitian. The difference stems from the fact that in 4-dimensional theory the quadratic takes the “diagonal form” $`\overline{\psi }^T\mathrm{\Gamma }_{\mu \nu }\psi \mathrm{\Lambda }_+^2`$ while in eight dimensions, such a diagonal form does not belong to $`\mathrm{\Lambda }_+^2`$ (though it is still purely imaginary). In order to define a reasonable $`𝒬\mathrm{\Lambda }_+^2`$ in 8 dimensions, we have to decompose the spinor $`\psi S^+\mathrm{𝟏}\mathrm{𝟕}`$ into two parts $`\psi =\psi _\mathrm{𝟏}+\psi _\mathrm{𝟕}`$, one of which, $`\psi _\mathrm{𝟏}(\zeta ^T\psi )\zeta `$, is in $`\mathrm{𝟏}`$, i.e. the trivial module of $`Spin(7)`$, and the other of which, $`\psi _\mathrm{𝟕}`$, belongs to $`\mathrm{𝟕}`$, namely the seven-dimensional irreducible module of $`Spin(7)`$. Since this decomposition is orthogonal and since $`\mathrm{\Gamma }_{\mu \nu }\zeta \mathrm{𝟕}`$, we have $`\overline{\psi }^T\zeta =\overline{\psi }_\mathrm{𝟏}^T\zeta `$ and $`\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi =\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi _\mathrm{𝟕}`$. Thus, the quadratic form $`𝒬_{\mu \nu }=(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )=(\overline{\psi }_\mathrm{𝟏}^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi _\mathrm{𝟕})`$ we have just constructed is really an “off-diagonal” product between the independent degrees of freedom $`\psi _\mathrm{𝟏}`$ and $`\psi _\mathrm{𝟕}`$. Such a product cannot be automatically real or purely imaginary. This explains why in the first equation of (36), we have splitted the quadratic form into its real and imaginary parts, and introduced two real coefficients $`a`$ and $`b`$.
There is a more compact way to write down the real and imaginary parts of $`𝒬`$:
$$\begin{array}{ccc}\mathrm{}\left[(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )\right]\hfill & =\hfill & \frac{1}{24!}\mathrm{\Omega }_{}^{\lambda \rho \sigma }{}_{[\mu }{}^{}(\psi ^{}\mathrm{\Gamma }_{\nu ]\lambda \rho \sigma }\psi ),\hfill \\ i\mathrm{}\left[(\overline{\psi }^T\zeta )(\zeta ^T\mathrm{\Gamma }_{\mu \nu }\psi )\right]\hfill & =\hfill & \frac{1}{8}(\psi ^{}\mathrm{\Gamma }_{\mu \nu }\psi )\frac{1}{16}\mathrm{\Omega }_{\mu \nu \alpha \beta }(\psi ^{}\mathrm{\Gamma }^{\alpha \beta }\psi )\hfill \\ & \hfill & \frac{1}{2}(P^+)_{\mu \nu }^{}{}_{}{}^{\alpha \beta }(\psi ^{}\mathrm{\Gamma }_{\alpha \beta }\psi )\hfill \end{array}$$
(40)
where $`P^+:\mathrm{\Lambda }^2\mathrm{\Lambda }_+^2`$ is the orthorgonal projection of two-forms onto $`\mathrm{\Lambda }_+^2`$. Note that the imaginary part of $`𝒬`$ resembles the term $`\psi ^{}\mathrm{\Gamma }_{\mu \nu }\psi `$ in the Seiberg-Witten theory, but in eight dimensions there are additional ingredients in the construction of a general quadratic form valued in $`\mathrm{\Lambda }_+^2`$: we have terms involving $`\mathrm{\Omega }_{}^{\lambda \rho \sigma }{}_{[\mu }{}^{}\mathrm{\Gamma }_{\nu ]\lambda \rho \sigma }`$. Such terms are forbidden in 4 dimensions since there $`\mathrm{\Omega }_{}^{\lambda \rho \sigma }{}_{\mu }{}^{}ϵ_{}^{\lambda \rho \sigma }{}_{\mu }{}^{}`$, $`\mathrm{\Gamma }_{\nu \lambda \rho \sigma }ϵ_{\nu \lambda \rho \sigma }`$ and $`ϵ_{}^{\lambda \rho \sigma }{}_{[\mu }{}^{}ϵ_{\nu ]\lambda \rho \sigma }=0`$.
Let us discuss another difference between the 8-dimensional and 4-dimensional theories. Writing down the equations in such theories requires to fix certain geometrical data on the underlying manifold. For example, in order to construct the anti self-dual part $`F^+`$ of the curvature tensor in the 4-dimensional theory, one has to pick up a Hodge star operator, whose definition depends on the conforml structure of the manifold. Thus, the geometrical data – a conformal structure of the 4-manifold – enters natually in the first Seiberg-Witten equation $`F_{\mu \nu }^+𝒬_{\mu \nu }`$. Such geometrical data also enters in the the second Seiberg-Witten equation, i.e. the massless Dirac equation in 4 dimensions, as that equation is conformally invariant and it also depends on the choice of a conformal structure. In the 8-dimensional theory, the construction of the first equation involves another data, $`\mathrm{\Omega }`$, which is the $`Spin(7)`$-invariant 4-form calibrating the underlying geometry. This can be expected, since as long as the self-duality structures are concerned $`\mathrm{\Omega }`$ will play a role similar to the Hodge star operator in 4 dimensions. What makes the 8-dimensional theory different from that in 4 dimensions is that the geometrical data $`\mathrm{\Omega }`$ does not enter in the Dirac equation. Thus, it should not be very suprising when we find that the functional formalism of (36) in general does not allow the delicate cancellations as in the 4-dimensional theory. In particular, we do not know at present how to handle the uncancelled terms involving $`F^{}`$, arsing from the functional $`\mathrm{\Gamma }_\mu D_A^\mu \psi ^2`$ of the Dirac equation. One possible resolusion is to modify the second equation in (36) so that it depends on the form $`\mathrm{\Omega }`$ (through the $`Spin(7)`$-invariant spinor $`\zeta `$). Manifolds with $`SU(4)`$ Holonomy. Now we try to formulate an eight dimensional analog of the Seiberg-Witten equations on manifolds with the $`SU(4)`$ holonomy group. The starting point will be similar to that in the $`Spin(7)`$ case: One wishes to replace the nonabelian instanton equation $`F_{+\overline{i}\overline{j}}^{(0,2)}=0`$ by an abelian, monopole-like equation $`F_{+\overline{i}\overline{j}}^{(0,2)}=𝒬_{\overline{i}\overline{j}}(\psi ,\overline{\psi })`$, where $`\psi S_c^+`$ is a spinor field twisted by some complex line bundle $``$, and $`𝒬_{\overline{i}\overline{j}}(\psi ,\overline{\psi })`$ denotes a certain quadratic form valued in $`\mathrm{\Lambda }_+^{0,2}`$. Our first task is thus to find out such a quadratic form.
The condition for a (0,2)-form $`\beta _{\overline{i}\overline{j}}`$ to be valued in $`\mathrm{\Lambda }_+^{0,2}`$ is that it obeys the eigenvalue equation $`_\mathrm{\Omega }\beta =4\beta `$. So according to (28), $`\beta \mathrm{\Lambda }_+^{0,2}`$ is characterized by the equations
$$\overline{\mathrm{\Omega }}_{\overline{i}\overline{j}\overline{k}\overline{l}}\overline{\beta }^{\overline{k}\overline{l}}=8\beta _{\overline{i}\overline{j}}\mathrm{\Omega }_{ijkl}\beta ^{kl}=8\overline{\beta }_{ij}.$$
(41)
Our key observation here is that the spinor $`\gamma ^{ij}\overline{\zeta }`$ satisfies an equation with the same structure as the second one in (41). To see this, multiplying (26) by $`\zeta ^T\gamma _{ijkl}`$ from the left, we find
$$\mathrm{\Omega }_{ijkl}\overline{\zeta }=\frac{1}{8}\zeta ^T(\gamma _{ijkl}+g^{m\overline{n}}\gamma _{ijkl}\gamma _{m\overline{n}})+\frac{1}{32}\zeta ^Tg^{m\overline{s}}g^{n\overline{r}}\gamma _{ijkl}\gamma _{mn\overline{r}\overline{s}}.$$
Notice that $`\gamma _{i_1\mathrm{}i_p}=\gamma _{i_1}\mathrm{}\gamma _{i_p}`$ and $`\gamma _{i_1\mathrm{}i_p\overline{j}_1\mathrm{}\overline{j}_q}=0`$ for $`p>4`$. So one can use Eq.(21) repeatedly to compute $`g^{m\overline{n}}\gamma _{ijkl}\gamma _{m\overline{n}}`$ as well as $`g^{m\overline{s}}g^{n\overline{r}}\gamma _{ijkl}\gamma _{mn\overline{r}\overline{s}}`$, and the result simply reads
$$g^{m\overline{n}}\gamma _{ijkl}\gamma _{m\overline{n}}=4\gamma _{ijkl},g^{m\overline{s}}g^{n\overline{r}}\gamma _{ijkl}\gamma _{mn\overline{r}\overline{s}}=12\gamma _{ijkl}.$$
Consequently, we have
$$\mathrm{\Omega }_{ijkl}\overline{\zeta }=\zeta ^T\gamma _{ijkl}=\gamma _{ijkl}\zeta $$
(42)
(where the last identity comes from the symmetric property of the matrix $`\gamma _{ijkl}`$). Now with the help of (42) and (21), we can do some further computations:
$$\mathrm{\Omega }_{ijkl}\gamma ^{kl}\overline{\zeta }=\gamma ^{kl}(\mathrm{\Omega }_{ijkl}\overline{\zeta })=\gamma ^{kl}\gamma _{ijkl}\zeta =g^{k\overline{m}}g^{l\overline{n}}\gamma _{\overline{m}\overline{n}}\gamma _{ijkl}\zeta $$
$$=g^{k\overline{m}}g^{l\overline{n}}\gamma _{\overline{m}}\gamma _{ijkl\overline{n}}\zeta +g^{k\overline{m}}\gamma _{\overline{m}}\gamma _{ijk}\zeta $$
$$=g^{k\overline{m}}g^{l\overline{n}}\gamma _{ijkl\overline{m}\overline{n}}\zeta 2g^{k\overline{m}}\gamma _{ijk\overline{m}}\zeta 2\gamma _{ij}\zeta =8\gamma _{ij}\zeta .$$
So finally we arrive at
$$\mathrm{\Omega }_{ijkl}\gamma ^{kl}\overline{\zeta }=8\gamma _{ij}\zeta ,$$
(43)
which shows that $`\gamma ^{ij}\overline{\zeta }`$ has the same tensor properties as an anti self-dual two-form $`\beta ^{ij}\mathrm{\Lambda }_+^{0,2}`$.
Thus, given any spinor $`\psi S_c^+`$, one can use the isomorphism (22) to construct a form $`\beta ^{ij}=\psi ^T\gamma ^{ij}\overline{\zeta }=\zeta ^{}\gamma ^{ij}\psi (S_c^+)\text{}\mathrm{\Lambda }^{0,\mathrm{even}}`$. Naively, the identity (43) indicates that such a form should obey the anti self-duality equation (41), and thus it would belong to the subspace $`\mathrm{\Lambda }_+^{0,2}`$:
$$\beta ^{ij}\zeta ^{}\gamma ^{ij}\psi \mathrm{\Lambda }_+^{0,2}.$$
(44)
To construct a quadratic form $`𝒬_{\overline{i}\overline{j}}(\psi ,\overline{\psi })\mathrm{\Lambda }_+^{0,2}`$, one still needs another form $`\alpha `$, which should be anti-linear in $`\psi `$ and valued in $`\mathrm{\Lambda }^{0,0}^1`$, so that the factor $``$ could be cancelled when forming the product $`\alpha \beta _{\overline{i}\overline{j}}`$. The simplest choice of such a form would be
$$\alpha \zeta ^T\overline{\psi }=\psi ^{}\zeta \mathrm{\Lambda }^{0,0}^1.$$
(45)
So at first sight we expect that the quadratic form we are seeking should look like<sup>1</sup><sup>1</sup>1(46) is very similar to the quadratic form $`𝒬_{\mu \nu }`$ constructed in the real case. One may decompose $`\psi `$ into $`\psi =\psi _{}+\psi _{}`$, where $`\psi _{}`$ is valued in the $`SU(4)`$ invariant subspace of $`S_c^+`$ spanned by $`\zeta `$ and $`\overline{\zeta }`$, and $`\psi _{}`$ lives in the subspace orthogonal to it. Using the facts that $`\zeta ^T\zeta =0`$, $`\zeta ^{}\gamma _{\overline{i}\overline{j}}\zeta =\zeta ^{}\gamma _{\overline{i}\overline{j}}\overline{\zeta }=0`$, we find that (46) can be represented as an “off-diagonal” product of the two independent degrees of freedom $`\psi _{}`$ and $`\psi _{}`$, namely $`𝒬_{\overline{i}\overline{j}}=(\psi _{}^{}\zeta )(\zeta ^{}\gamma _{\overline{i}\overline{j}}\psi _{})`$. This resembles the $`Spin(7)`$ case, where the quadratic form is also an off-diagonal product of two independent degrees of freedom.:
$$𝒬_{\overline{i}\overline{j}}(\psi ,\overline{\psi })\alpha \beta _{\overline{i}\overline{j}}=(\psi ^{}\zeta )(\zeta ^{}\gamma _{\overline{i}\overline{j}}\psi ).$$
(46)
However, there is a subtlety in the above construction, which appears only in the complex case. In our definition of self-duality, the star operator $`_\mathrm{\Omega }`$ given in (28) is conjugate-linear rather than linear. Thus, even if $`\beta _{\overline{i}\overline{j}}`$ is anti self-dual, namely it obeys the condition (41), the quantity $`\alpha \beta _{\overline{i}\overline{j}}`$ needs not to be such a form for complex $`\alpha \mathrm{\Lambda }^{0,0}^1`$. We cannot simply take $`\alpha `$ to be real as $``$ should be a nontrivial complex line bundle. Moreover, since $`\psi `$ is also a complex spinor, the (0,2)-form $`\beta ^{ij}=\psi ^T\gamma ^{ij}\overline{\zeta }`$ defined by (44) is not really valued in $`\mathrm{\Lambda }_+^{0,2}`$, even though $`\gamma ^{ij}\overline{\zeta }`$ behaves as an anti self-dual tensor. To solve this problem, let us introduce a pair $`\alpha `$, $`\alpha ^{}`$ of (0,0)-forms as well as a pair $`\beta _{\overline{i}\overline{j}}`$, $`\beta _{\overline{i}\overline{j}}^{}`$ of (0,2)-forms, specified as
$$\begin{array}{ccc}\alpha \hfill & =\hfill & \psi ^{}\zeta \mathrm{\Lambda }^{0,0}^1,\alpha ^{}=\overline{\alpha }=\zeta ^{}\psi \mathrm{\Lambda }^{0,0}\hfill \\ \beta ^{ij}\hfill & =\hfill & \zeta ^{}\gamma ^{ij}\psi \mathrm{\Lambda }^{0,2},\beta ^{ij}=\zeta ^{}\gamma ^{ij}\overline{\psi }\mathrm{\Lambda }^{0,2}^1,\hfill \end{array}$$
(47)
and construct the product
$$𝒬_{\overline{i}\overline{j}}(\psi ,\overline{\psi })c\alpha \beta _{\overline{i}\overline{j}}+\overline{c}\alpha ^{}\beta _{\overline{i}\overline{j}}^{}=c(\psi ^{}\zeta )(\zeta ^{}\gamma _{\overline{i}\overline{j}}\psi )+\overline{c}(\psi ^T\overline{\zeta })(\zeta ^{}\gamma _{\overline{i}\overline{j}}\overline{\psi })$$
(48)
with $`c\text{}`$ being an arbitrarily fixed complex number (similar to the real numbers $`a`$, $`b`$ in the $`Spin(7)`$ case). One then uses (43) to derive
$$\mathrm{\Omega }_{ijkl}𝒬^{kl}=c(\psi ^{}\zeta )(\psi ^T\mathrm{\Omega }_{ijkl}\gamma ^{kl}\overline{\zeta })\overline{c}(\psi ^T\overline{\zeta })(\psi ^{}\mathrm{\Omega }_{ijkl}\gamma ^{kl}\overline{\zeta })$$
$$=8c(\psi ^{}\zeta )(\psi ^T\gamma _{ij}\zeta )+8\overline{c}(\psi ^T\overline{\zeta })(\psi ^{}\gamma _{ij}\zeta )=8\overline{𝒬}_{ij},$$
which indicates that $`𝒬_{\overline{i}\overline{j}}`$ is now in $`\mathrm{\Lambda }_+^{0,2}`$.
Having constructed a quadratic form $`𝒬_{\overline{i}\overline{j}}`$ with the right properties, we can immediately write down the first equation analogous to Seiberg and Witten’s:
$$F_{+\overline{i}\overline{j}}^{(0,2)}=c\alpha \beta _{\overline{i}\overline{j}}+\overline{c}\alpha ^{}\beta _{\overline{i}\overline{j}}^{}=c(\psi ^{}\zeta )(\zeta ^{}\gamma _{\overline{i}\overline{j}}\psi )+\overline{c}(\psi ^{}\gamma _{\overline{i}\overline{j}}\overline{\zeta })(\zeta ^{}\psi ).$$
(49)
One may also write down a similar equation for $`F_{+ij}^{(2,0)}`$ by taking the complex conjugate of (49). It should be pointed out, however, that at this stage we have not yet established another kind of equation (something like $`F_\omega ^{(1,1)}\frac{1}{2}\omega (\alpha ^2\beta ^2)`$ as in the 4-dimensional theory), which governs the (1,1)-component of $`F`$ in the direction along the Käller form $`\omega _{i\overline{j}}=ig_{i\overline{j}}`$. To obtain such an equation, one should use a new star operator $`_\mathrm{\Theta }`$ with $`\mathrm{\Theta }\mathrm{}(\mathrm{\Omega })+\omega ^2`$ to define self-duality .
Next we consider the Dirac equation $`D_A\psi =0`$. With the isomorphism $`(S_c^+,S_c^{})\text{}(\mathrm{\Lambda }^{0,\mathrm{even}},\mathrm{\Lambda }^{0,\mathrm{odd}})`$, the twisted Dirac operator $`D_A:S_c^+S_c^{}`$ becomes $`D_A=\overline{}_A+\overline{}_A^{}:\mathrm{\Lambda }^{0,\mathrm{even}}\mathrm{\Lambda }^{0,\mathrm{odd}}`$. Let us restrict this operator to $`(\mathrm{\Lambda }^{0,0}\mathrm{\Lambda }^{0,2})`$. Under this restriction, the spinor $`\psi \mathrm{\Lambda }^{0,\mathrm{even}}`$ has the components $`\alpha ^{}=\overline{\alpha }\mathrm{\Lambda }^{0,0}`$ and $`\beta \mathrm{\Lambda }^{0,2}`$, and the Dirac equation is reduced simply to
$$\overline{}_A\overline{\alpha }+\overline{}_A^{}\beta =0.$$
(50)
This constitutes the second equation in our theory.
Although Eq.(49)-(50) resemble the four-dimensional Seiberg-Witten equations on Kähler manifolds , it should be pointed out that here the Dirac equation (50) in general does not allow a simple decomposition into $`\overline{}_A\overline{\alpha }=\overline{}_A^{}\beta =0`$. This makes a computation of the relevant invariants quite difficult. This difficulty is related to a problem appeared in the $`Spin(7)`$ case, where we mentioned that there is an uncancelled term involving $`F^{}`$ in the functional formalism.
## 4 $`S`$-duality in Abelian Gauge Theory
In this section we turn to the abelian gauge theory in eight dimensions. For simplicity, we will consider only the case when $`X`$ has the holonomy group $`Spin(7)`$. Classically we have a $`U(1)`$ gauge field $`A_\mu `$ and its field strength $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$, together with the following action functional
$$\begin{array}{ccc}S[A]\hfill & =\hfill & \frac{1}{2g^2}_XF_{\mu \nu }F^{\mu \nu }+\frac{i\theta }{8\pi ^2}_X\mathrm{\Omega }FF\hfill \\ & =\hfill & \frac{1}{2g^2}_XF_{\mu \nu }F^{\mu \nu }+\frac{i\theta }{32\pi ^2}_X\mathrm{\Omega }_{\mu \nu \alpha \beta }F^{\mu \nu }F^{\alpha \beta }.\hfill \end{array}$$
(51)
Given such data, the partition function $`Z(g,\theta )`$ can be formally defined as the Euclidean path-integral
$$Z(g,\theta )=[dA]e^{S[A]}.$$
(52)
Let us analyze the partition function in some detail. Usually, it is convenient to change the integration variables $`AF`$. The routine is quite standard: Just as in four dimensions, $`F`$ is not an independent variable, and it must be subject to the Bianchi identity $`dF=0`$. If we write $`dF=\frac{1}{2}_\lambda F_{\mu \nu }dx^\lambda dx^\mu dx^\nu `$, then one easily deduces from the closeness and self-duality of $`\mathrm{\Omega }`$ that
$$dF\mathrm{\Omega }dx^\nu =\frac{1}{2}_\mu [(_\mathrm{\Omega }F)^{\mu \nu }]d(\mathrm{vol}).$$
This implies that the Bianchi identity $`dF=0`$ can be replaced by a constraint
$$_\nu [(_\mathrm{\Omega }F)^{\mu \nu }]=0$$
(53)
on the field strength. Consequently, the partition function (52) has a path-integral representation over $`F`$, with the delta function $`\delta (_\nu [(_\mathrm{\Omega }F)^{\mu \nu }])`$ inserted. Such a delta function can be written as another path-integral over some auxiliary field $`A_\mu ^D`$. Thus, one may write
$$\begin{array}{ccc}Z(g,\theta )\hfill & =\hfill & [dF][dA^D]e^{S+i_XA_\mu ^D_\nu (_\mathrm{\Omega }F)^{\mu \nu }}\hfill \\ & =\hfill & [dF][dA^D]e^{S+\frac{i}{2}_XF_{\mu \nu }^D(_\mathrm{\Omega }F)^{\mu \nu }},\hfill \end{array}$$
(54)
where $`F_{\mu \nu }^D=_\mu A_\nu ^D_\nu A_\mu ^D`$ is the field strength of $`A_\mu ^D`$. If we first integrate out the auxiliary field $`A^D`$, then the resulting expressonis is nothing but $`[dF]\delta (_\nu [(_\mathrm{\Omega }F)^{\mu \nu }])e^S`$, which is equivalent to the original partition function (52).
Alternatively, one can integrate out $`F`$ in (54) first, leaving an effective action for the auxiliary field, which governs the dynamics of the collective variable $`A_\mu ^D`$. To achieve this, let us decompose $`F=F^++F^{}`$ into two independent components $`F^+`$, $`F^{}`$ and write $`[dF]=[dF^+][dF^{}]`$. In terms of these components, we have
$$\begin{array}{ccc}S\hfill & =\hfill & \left(\frac{1}{2g^2}i\frac{3\theta }{16\pi ^2}\right)_XF_{\mu \nu }^+F^{+\mu \nu }+\left(\frac{1}{2g^2}+i\frac{\theta }{16\pi ^2}\right)_XF_{\mu \nu }^{}F^{\mu \nu }\hfill \\ & \hfill & \frac{3}{4e_+^2}_XF_{\mu \nu }^+F^{+\mu \nu }+\frac{1}{4e_{}^2}_XF_{\mu \nu }^{}F^{\mu \nu },\hfill \end{array}$$
(55)
$$\frac{i}{2}_XF_{\mu \nu }^D(_\mathrm{\Omega }F)^{\mu \nu }=\frac{3i}{2}_XF_{\mu \nu }^{D+}F^{+\mu \nu }+\frac{i}{2}_XF_{\mu \nu }^DF^{\mu \nu }.$$
(56)
So substituting (55)-(56) into (54) yields a product of two gaussian integrals over $`F^\pm `$. An explicit evaluation of these integrals gives
$$\begin{array}{ccc}Z(g,\theta )\hfill & =\hfill & [dA^D]e^{\stackrel{~}{S}[A^D]},\hfill \\ \stackrel{~}{S}[A^D]\hfill & =\hfill & \frac{3e_+^2}{4}_XF_{\mu \nu }^{D+}F^{D+\mu \nu }+\frac{e_{}^2}{4}_XF_{\mu \nu }^DF^{D\mu \nu }.\hfill \end{array}$$
(57)
We thus obtain a dual description of the original theory using the collective field $`A^D`$, in which the coupling constants get transformed:
$$e_\pm ^2\stackrel{~}{e}_\pm ^2=\frac{1}{e_\pm ^2}.$$
(58)
Eq.(58) characterizes a generalized $`S`$-duality in eight dimensions.
The above discussion is somewhat rough and we ignored several subtleties arising from regularization. In four dimensions, a more careful study shows that the partition function transforms as a modular form, and this provides a precise test of the $`S`$-duality. When entering in eight dimensions, however, one sees from (55) that the action does not takes the form $`Si(\tau (F^+)^2\overline{\tau }(F^{})^2)`$, so the partition function $`Z(g,\theta )`$ will not be parametrized neatly by a single complex coupling $`\tau `$ along with its conjugate $`\overline{\tau }`$; more naturally, $`Z(g,\theta )`$ should be parametrized by $`(e^+,e^{})`$, and $`e^\pm `$ are not complex conjugate to each other. It seems rather difficult to write down a simple modular form expression for the partition function of the eight-dimensional theory. Without such a modular form our understanding of the generalized $`S`$-duality is quite incomplete.
###### Acknowledgments.
We would like to thank I. Singer for some helpful conversations.
## Appendix A $`\gamma `$-Matrices and Clifford Calculus
In the text we used $`Cl(8)`$ to denote the 8-dimensional Clifford algebra. This algebra has a real, irreducible representation $`\rho :Cl(8)End(S)`$. According to the standard argument, $`\rho (Cl(8))`$ constitutes the algebra $`\text{}(16)`$ of $`16\times 16`$ real matrices, acting on the 16-dimensional vector space $`S\text{}^{16}`$. The following isomorphism between $`Cl(8)`$ and the wedge algebra $`^{}\text{}^8`$ is quite evident:
$$Cl(8)^{}\text{}^8.$$
(59)
In particular, if we introduce a set of orthogonal generators of $`Cl(8)`$, $`e_\mu ^1\text{}^8`$ ($`1\mu 8`$), with the rule of Clifford multiplications
$$e_\mu e_\nu +e_\nu e_\mu =2e_\mu ,e_\nu 2g_{\mu \nu },$$
(60)
then the $`p`$-“form” $`e_{\mu _1}e_{\mu _2}\mathrm{}e_{\mu _p}^p\text{}^8`$ canonically has the representation
$$e_{\mu _1}e_{\mu _2}\mathrm{}e_{\mu _p}\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p}\mathrm{\Gamma }_{[\mu _1}\mathrm{\Gamma }_{\mu _2}\mathrm{}\mathrm{\Gamma }_{\mu _p]}$$
(61)
where
$$\mathrm{\Gamma }_\mu =\rho (e_\mu )\text{}(16)$$
(62)
are known as “$`\gamma `$-matrices”, and the square bracket indicates anti-symmetrization of the indices, with a prefactor $`1/p!`$. We shall use the notation $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}=I_{16\times 16}`$ for $`p=0`$.
The Clifford multiplication “$``$” between $`u=_\mu C^\mu e_\mu ^1\text{}^8Cl(8)`$ and any element $`wCl(8)^{}\text{}^8`$ can be identified with an operation on the wedge algebra:
$$uwuwi_u(w),$$
(63)
where the interior product $`i_u(w)`$ is defined by the linear map $`i_u:^p\text{}^8^{p1}\text{}^8`$, via
$$i_u(u_1u_2\mathrm{}u_p)=\underset{i=1}{\overset{p}{}}(1)^{i+1}u_i,uu_1\mathrm{}\widehat{u_i}\mathrm{}u_p.$$
(64)
Applying this to the matrix representation $`\rho `$, (63) becomes an identity between $`\gamma `$-matrices
$$\mathrm{\Gamma }_\mu \mathrm{\Gamma }_{\nu _1\nu _2\mathrm{}\nu _p}=\mathrm{\Gamma }_{\mu \nu _1\nu _2\mathrm{}\nu _p}g_{\mu \nu _1}\mathrm{\Gamma }_{\nu _2\nu _3\mathrm{}\nu _p}+g_{\mu \nu _2}\mathrm{\Gamma }_{\nu _1\nu _3\mathrm{}\nu _p}+\mathrm{}+(1)^pg_{\mu \nu _p}\mathrm{\Gamma }_{\nu _1\mathrm{}\nu _{p1}}.$$
(65)
Sometimes we need to fix a particular basis and construct the $`\gamma `$-matrices $`\mathrm{\Gamma }_\mu `$ explicitly. Our convention of choosing such matrices is as follows. Since $`\text{}(16)\text{}(2)\text{}(2)\text{}(2)\text{}(2)`$, $`\mathrm{\Gamma }_\mu `$ can be expressed by a 4-fold tensor product of some basis in $`\text{}(2)`$. Thus, we take a basis of $`\text{}(2)`$ to be
$$I=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),ϵ=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
(66)
It is easy to checks that the $`16\times 16`$ matrices
$$\begin{array}{cc}\mathrm{\Gamma }_1=ϵϵϵ\sigma _1,\hfill & \mathrm{\Gamma }_2=I\sigma _1ϵ\sigma _1,\hfill \\ \mathrm{\Gamma }_3=I\sigma _3ϵ\sigma _1,\hfill & \mathrm{\Gamma }_4=\sigma _1ϵI\sigma _1,\hfill \\ \mathrm{\Gamma }_5=\sigma _3ϵI\sigma _1,\hfill & \mathrm{\Gamma }_6=ϵI\sigma _1\sigma _1,\hfill \\ \mathrm{\Gamma }_7=ϵI\sigma _3\sigma _1,\hfill & \mathrm{\Gamma }_8=IIIϵ\hfill \end{array}$$
(67)
obey the relations $`\{\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_\nu \}=2\delta _{\mu \nu }`$.
In our convention (67), the matrices $`\mathrm{\Gamma }_\mu `$ are all anti-symmetric. More generally we have
$$(\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p})^T=(1)^p\mathrm{\Gamma }_{\mu _p\mathrm{}\mu _2\mu _1}=(1)^{\frac{p(p+1)}{2}}\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p}.$$
(68)
Thus $`\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p}`$ is anti-symmetric when $`p1,2`$ (mod 4) and symmetric when $`p3,4`$ (mod 4). Consequently, the matrix representation of the “volume element” $`\omega e_1e_2\mathrm{}e_8Cl(8)`$, namely
$$\rho (\omega )=\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_8\mathrm{\Gamma }_9,$$
(69)
has the symmetric property $`(\mathrm{\Gamma }_9)^T=\mathrm{\Gamma }_9`$. In fact $`\mathrm{\Gamma }_9`$ is diagonal in our basis:
$$\mathrm{\Gamma }_9=III\sigma _3=\left(\begin{array}{cc}I_{8\times 8}& 0\\ 0& I_{8\times 8}\end{array}\right).$$
(70)
(70) shows that the irreducible $`Cl(8)`$ module $`S\text{}^{16}`$ has a decomposition $`S=S^+S^{}`$ into the $`\pm 1`$ eigenspaces $`S^\pm `$ of $`\mathrm{\Gamma }_9`$. Of course neither of these eight-dimensional eigenspaces are invariant under the action of $`Cl(8)`$. To be a little more explicit, notice that $`\mathrm{\Gamma }_9\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}=(1)^p\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}\mathrm{\Gamma }_9`$, we have
$$\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}\{\begin{array}{cc}\left(\begin{array}{cc}& 0\\ 0& \end{array}\right),\hfill & p=\mathrm{even}\hfill \\ \left(\begin{array}{cc}0& \\ & 0\end{array}\right),\hfill & p=\mathrm{odd}\hfill \end{array}\mathrm{on}S=\left(\begin{array}{c}S^+\\ S^{}\end{array}\right),$$
(71)
so for odd $`p`$ the matrix $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}`$ swaps $`S^\pm `$. Nevertheless, if one considers a subalgebra of $`Cl(8)`$ spanned by some even elements $`a=e_{\mu _1}\mathrm{}e_{\mu _{2k}}`$, then the matrix representation $`\rho (a)=\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _{2k}}`$ will keep both the subspaces $`S^\pm S`$ invariant.
At this point we consider a linear space $`^2\text{}^8`$ spanned by elements of the form $`L_{\mu \nu }=\frac{1}{2}e_\mu e_\nu =\frac{1}{4}(e_\mu e_\nu e_\nu e_\mu )`$. This spcae forms a Lie algebra under the bracket
$$[L_{\mu \nu },L_{\alpha \beta }]L_{\mu \nu }L_{\alpha \beta }L_{\alpha \beta }L_{\mu \nu },$$
(72)
where “$``$” again stands for the Clifford multiplication. In fact, a simple computation shows that
$$[L_{\mu \nu },L_{\alpha \beta }]=g_{\mu \alpha }L_{\nu \beta }+g_{\nu \beta }L_{\mu \alpha }g_{\nu \alpha }L_{\mu \beta }g_{\mu \beta }L_{\nu \alpha },$$
(73)
so $`\{L_{\mu \nu }\}`$ generates the Lie algebra of $`Spin(8)`$. According to the previous discussion, $`S^\pm `$ can be considered as $`Spin(8)`$-modules, and actually they are two inequivenlent irreducible modules of $`Spin(8)`$. The representation $`\rho (L_{\mu \nu })=\frac{1}{2}\mathrm{\Gamma }_{\mu \nu }`$ of (the Lie algebra of) $`Spin(8)`$ then decomposes into two irreducible ones: $`\rho =\rho ^+\rho ^{}`$. $`\rho ^+`$ is the spin representation with positive chirality and $`\rho ^{}`$ the spin representation with negative chirality. That $`\rho ^\pm `$ are inequivalent stems from the “central element” $`\mathrm{\Gamma }_9=\rho (\omega )=\pm 1`$ having different values on $`S^\pm `$.
So far we have only constructed two irreducible spin representations of $`Spin(8)`$, $`\rho ^\pm `$, acting on $`S^+=\mathrm{𝟖}_s`$ and $`S^{}=\mathrm{𝟖}_c`$, respectively. There is another inequivalent eight-dimentional irreducible representation of $`Spin(8)`$, the so-called “vector representation” $`\rho _v`$, which will act on the vector space $`\mathrm{𝟖}_vSpan\{e_\mu \}^1\text{}^8`$ adjointly:
$$\rho _v(L_{\mu \nu })(e_\alpha )[L_{\mu \nu },e_\alpha ].$$
(74)
The matrix elements of $`\rho _v(L_{\mu \nu })`$ are determined by the following commutative relations:
$$[L_{\mu \nu },e_\alpha ]=g_{\mu \alpha }e_\nu g_{\nu \alpha }e_\mu .$$
(75)
It follows that the image of $`Spin(8)`$ under $`\rho _v`$ is isomorphic to $`SO(8)`$. The existence of the three inequivalent irreducible 8-dimensional modules $`\mathrm{𝟖}_s`$, $`\mathrm{𝟖}_c`$ and $`\mathrm{𝟖}_v`$ is often summarized as the triality of the $`Spin(8)`$-representations.
As a natural extension of the vector representation $`\rho _v`$, it is possible to construct tensor representations of $`Spin(8)`$ on $`^p\text{}^8`$. One verifies by induction that (75) is extended to
$$[L_{\mu \nu },e_{\alpha _1}\mathrm{}e_{\alpha _p}]=M_{\alpha _1\mathrm{}\alpha _p}^{\beta _1\mathrm{}\beta _p}(L_{\mu \nu })e_{\beta _1}\mathrm{}e_{\beta _p}$$
(76)
with some adjoint matrices $`M(L_{\mu \nu })`$. This therefore defines a tensor representation $`^p\rho _v`$ of $`Spin(8)`$. Alternatively, $`^p\rho _v`$ may also be obtained by considering tensor product of the spin representations $`\rho ^\pm `$.
To see this, we first need to establish an isomorphism between the spaces $`(S^+S^{})(S^+S^{})`$ and $`^{}\text{}^8`$. Note that both spaces have the same dimensions: $`16\times 16=2^8`$. By choosing an orthogonal basis $`\{v_A\}_{1A16}`$ of $`S=S^+S^{}`$, we associate each $`v_Av_BSS`$ to an element of $`^{}\text{}^8`$ as follows:
$$v_Av_B\underset{p=0}{\overset{8}{}}v_A,\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}v_Be_{\mu _1}\mathrm{}e_{\mu _p}.$$
(77)
This correspondence is 1:1 since if $`v_Av_B`$, $`v_Cv_D`$ are associated to the same element of $`^{}\text{}^8`$, then we must have $`v_A,\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}v_Bv_C,\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}v_D`$ for all $`p`$ and, by irreducibility of the Clifford group<sup>2</sup><sup>2</sup>2Clifford group $`G_dCl(d)`$ in $`d`$-dimensions is a finite group whose generators can be presented by the abstract elements $`\{e_1,\mathrm{},e_d,1\}`$ subject to the relation that $`1`$ is central and that $`(1)^2=1,e_i^2=1`$ and $`e_ie_j=(1)e_je_i`$ for all $`ij`$. acting on $`S`$, the two matrix elements must be orthogonal and never identical to each other, unless $`v_A=v_C`$, $`v_B=v_D`$.
Now we can use the isomorphism $`SS^{}\text{}^8`$ specified by (77). We see that the action of any group element $`gSpin(8)`$ on $`v_A`$, i.e. $`v_A\stackrel{~}{v}_A(\rho ^+\rho ^{})(g)v_A`$, will induce two equivalent actions on $`v_Av_B`$ and on its image in $`^{}\text{}^8`$. The first action is simply $`(\rho ^+\rho ^{})(\rho ^+\rho ^{})(g):v_Av_B\stackrel{~}{v}_A\stackrel{~}{v}_B`$. The second action, when restricted to the components $`T^{\mu _1\mathrm{}\mu _p}v_A,\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}v_B^p\text{}^8`$, is determined by $`T^{\mu _1\mathrm{}\mu _p}\stackrel{~}{T}^{\mu _1\mathrm{}\mu _p}\stackrel{~}{v}_A,\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}\stackrel{~}{v}_B`$, which in turn gives rise to the adjoint action $`\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}\rho (g)^T\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}\rho (g)=\rho (g)^1\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}\rho (g)`$, leading to the tensor representation $`^p\rho _v`$. It follows that
$$(\rho ^+\rho ^{})(\rho ^+\rho ^{})\underset{p=0}{\overset{8}{}}^p\rho _v2(1\rho _v^2\rho _v^3\rho _v)^4\rho _v.$$
(78)
The isomorphism discussed above gives an identity known as the Fierz rearrangement formula. The vectors $`v_A`$, $`v_B`$ in (77) can be replaced by arbitrary spinors $`\varphi =\varphi ^Av_A,\psi =\psi ^Av_AS`$. On the left hand side of this correspondence, we have the tensor product $`\varphi \psi `$ with components $`\varphi ^A\psi ^B`$, which can be viewed as a $`16\times 16`$ matrix acting on $`S`$. The right hand side can also be considered as such a matrix if we replace the Clifford elements $`e_{\mu _1}\mathrm{}e_{\mu _p}`$ by their $`\gamma `$-matrix representation $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}`$. Since these two matrices are the same object, we must have
$$\varphi ^A\psi ^B=\frac{1}{16}\underset{p=0}{\overset{8}{}}\frac{1}{p!}(\varphi ^T\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}\psi )\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}^{}{}_{}{}^{AB}$$
(79)
here $`(\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p})^{AB}`$ denotes the matrix-element of $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}`$ in the basis $`v_A`$. The coefficients $`\frac{1}{p!}`$ in (79) are introduced so as to ensure that the sum runs over each of the basis elements of $`^{}\text{}^8`$ exactly once, and the factor $`\frac{1}{16}`$ comes from a group-theoretical consideration, which is nothing but the inverse of the dimension of the irreducible representation for the Clifford group. That this factor must be equal to $`\frac{1}{16}`$ may also be checked by taking the trace of (79): From (71) we recall that for odd $`p`$, the matrix $`\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}`$ is always off-diagonal, thus having a vanishing trace. For even $`p>0`$, the cyclic property of the trace $`\mathrm{Tr}(\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p})=\mathrm{Tr}(\mathrm{\Gamma }_{\mu _2\mathrm{}\mu _p\mu _1})`$ together with the $`\gamma `$-matrix identity $`\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p}=\mathrm{\Gamma }_{\mu _2\mathrm{}\mu _p\mu _1}`$ also gives $`\mathrm{Tr}(\mathrm{\Gamma }_{\mu _1\mu _2\mathrm{}\mu _p})=0`$. So when taking the trace, only the first term ($`p=0`$) in the right hand side of (79) survives and it takes the value $`\frac{1}{16}\varphi ^T\psi \mathrm{Tr}(I_{16\times 16})=\varphi ^T\psi `$, which agrees exactly with the trace of the left hand side. As an aside, note that this kind of argument allows us to write down a trace formula for the $`\gamma `$-matrices: Multiplying (79) by $`\mathrm{\Gamma }_{}^{\nu _1\mathrm{}\nu _q}{}_{AB}{}^{}=(1)^{[\frac{q+1}{2}]}\mathrm{\Gamma }_{}^{\nu _1\mathrm{}\nu _q}{}_{BA}{}^{}`$, the left hand side becomes $`\varphi ^T\mathrm{\Gamma }^{\nu _1\mathrm{}\nu _q}\psi `$, while the right hand side is $`_{p=0}^8\frac{(1)^{[\frac{q+1}{2}]}}{16p!}(\varphi ^T\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _p}\psi )\mathrm{Tr}(\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}\mathrm{\Gamma }^{\nu _1\mathrm{}\nu _q})`$, and this gives
$$\mathrm{Tr}(\mathrm{\Gamma }_{\mu _1\mathrm{}\mu _p}\mathrm{\Gamma }^{\nu _1\mathrm{}\nu _q})=16p!(1)^{[\frac{p+1}{2}]}\delta _{pq}\delta _{[\mu _1}^{\nu _1}\mathrm{}\delta _{\mu _p]}^{\nu _p}.$$
(80)
For example we have $`\mathrm{Tr}(\mathrm{\Gamma }_\mu \mathrm{\Gamma }^\nu )=16\delta _\mu ^\nu `$, $`\mathrm{Tr}(\mathrm{\Gamma }_{\mu \nu }\mathrm{\Gamma }^{\alpha \beta })=16(\delta _\mu ^\alpha \delta _\nu ^\beta \delta _\mu ^\beta \delta _\nu ^\alpha )`$, etc..
We end this appendix with a few remarks. If $`\varphi =\psi `$ has a definite chirality, then many terms in the sum (79) will vanish. Such terms correspond to $`p1,2`$ (mod 4) when the $`\gamma `$-matrices are anti-symmetric or $`p=`$ odd when the $`\gamma `$-matrices map $`\varphi `$ into some spinors with opposite chirality, which are orthogonal to $`\varphi ^T`$. In this case the Fierz reaarangement formula gets much simplified:
$$\varphi \varphi ^T=\frac{1}{16}\left((\varphi ^T\varphi )I_{16\times 16}+(\varphi ^T\mathrm{\Gamma }_9\varphi )\mathrm{\Gamma }_9+\frac{1}{4!}(\varphi ^T\mathrm{\Gamma }^{\mu \nu \alpha \beta }\varphi )\mathrm{\Gamma }_{\mu \nu \alpha \beta }\right).$$
(81)
Thus, since $`\mathrm{\Gamma }_9\varphi =\pm \varphi `$ for $`\varphi S^\pm `$, we have
$$\varphi \varphi ^T=\frac{1}{16}\left((\varphi ^T\varphi )(I_{16\times 16}\pm \mathrm{\Gamma }_9)+\frac{1}{4!}(\varphi ^T\mathrm{\Gamma }^{\mu \nu \alpha \beta }\varphi )\mathrm{\Gamma }_{\mu \nu \alpha \beta }\right),\varphi S^\pm .$$
(82)
Clearly (82) defines a projector from $`S`$ onto its one-dimensional subspace spaned by $`\varphi `$. Moreover, if $`\varphi \psi `$ but they still have the same definite chirality – say, both of them are in $`S^+`$ or in $`S^{}`$, then by symmetrizing (79) we get
$$\varphi ^A\psi ^B+\varphi ^B\psi ^A=\frac{1}{8}\left((\varphi ^T\psi )(\delta ^{AB}\pm \mathrm{\Gamma }_{9}^{}{}_{}{}^{AB})+\frac{1}{4!}(\varphi ^T\mathrm{\Gamma }^{\mu \nu \alpha \beta }\psi )\mathrm{\Gamma }_{\mu \nu \alpha \beta }^{}{}_{}{}^{AB}\right),$$
(83)
where “$`\pm `$” corresponds to $`\varphi `$, $`\psi S^\pm `$, respectively. We can also anti-symmetrize (79) to derive, for $`\varphi `$ and $`\psi `$ having the same chirality,
$$\varphi ^A\psi ^B\varphi ^B\psi ^A=\frac{1}{8}\left(\frac{1}{2!}(\varphi ^T\mathrm{\Gamma }^{\mu \nu }\psi )\mathrm{\Gamma }_{\mu \nu }^{}{}_{}{}^{AB}+\frac{1}{6!}(\varphi ^T\mathrm{\Gamma }^{\mu \nu \alpha \beta \lambda \rho }\psi )\mathrm{\Gamma }_{\mu \nu \alpha \beta \lambda \rho }^{}{}_{}{}^{AB}\right).$$
(84)
The two terms in the right hand side of (84) are in fact equal to each other upto a factor $`\pm \mathrm{\Gamma }_9`$ and we finally have $`\varphi \psi ^T\psi \varphi ^T=\frac{1}{16}(\varphi ^T\mathrm{\Gamma }^{\mu \nu }\psi )\mathrm{\Gamma }_{\mu \nu }(1\pm \mathrm{\Gamma }_9)`$, if both $`\varphi `$, $`\psi `$ are in $`S^\pm `$.
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# I Introduction
## I Introduction
The quark masses are some of basic parameters of the standard model. There are various hadronic phenomenologies relating to the light quark($`u,d,s`$) masses. For instance, they break the chiral symmetry of QCD explicitly. $`m_um_d`$ breaks isospin symmetry or charge symmetry, and $`(m_u+m_d)/2m_s`$ breaks SU(3) symmetry in hadron physics respectively. However, in QCD the masses of the light quarks are not directly measureable in inertial experiments, but enter the theory only indrectly as parameters in the fundamental lagrangian. The purpose of this paper is to study light quark masses at energy scale $`\mu =m_\rho `$ in a non-perturbative way.
In general, at low energy information about the light quark mass ratios is extracted in order by order by a rigorous, semiphenomenological method, chiral perturbative theory(ChPT). To the first order, the results, $`m_s/m_d=19`$ and $`m_u/m_d=0.556`$, are well-known. Many authors have studied the mass ratios up to next to leading order of the chiral expansion. Gasser and Leutwyler first obtained $`m_s/m_d=20.2`$ and $`m_u/m_d=0.554`$. Then Kaplan and Manohar extracted $`m_s/m_d=`$15 to 23 and $`m_u/m_d=`$0 to 0.8 with very larger error bar. These values have been improved to $`m_s/m_d=(20.5\pm 2.5),m_u/m_d=0.52\pm 0.13`$ by Leutwyler, to $`m_s/m_d=18,m_u/m_d=0.66`$ by Gerard, and to $`m_s/m_d=21,m_u/m_d=0.30\pm 0.07`$ by Donoghue et. al. respectively. Finally, Leutwyler analysis previous results and obtained $`m_s/m_d=18.9\pm 0.8,m_u/m_d=0.553\pm 0.043`$. So far, however, the study on light quark masses in framework of ChPT is limited by the following shortages: 1) In this framework, to obtain light quark mass ratio beyond the next to leading-order is very diffcult, since more and more free parameters are included with raising of perturbative order. 2) Due to Kaplan-Manohar ambiguity, the defination of the light quark masses in ChPT is not uniquely. The reason has been pointed out in ref. that, at the next to leading-order, QCD renormalization is mass-dependent, and the symmetry can not distinguish renormalized quark masses from “bare” quark masses. 3) In framework of ChPT, we can only obtain light quark masse ratios. For obtaining the individual quark masses, other approachs, such as QCD sum rules or lattice calculation are needed. These shortages are urgently wanted to be improved by theorectical studies of QCD and experiment.
In ref. we have the constructed chiral constituent quark model(ChCQM) including the lowest vector meson resonances following the spirit of Manohar-Georgi model. This model provides a formulation to perform rigorous field theory calculation at energy scale lying between ChPT($`\mu 0.5`$GeV) and chiral symmetry spontaneously broken(CSSB) scale($`\mu 1.2`$GeV), and a successful description on physics in this energy region. Up to $`O(p^4)`$, low energy limit of the model agree with ChPT well. Thus this model can be treated as an approach to extend ChPT investigation inspired by QCD. The most important advantage of this approach is that we can perform calculation beyond the low energy expansion in ChCQM, and only fewer free parameters are required. In general, there are three types of expansion working at low energy. They are momentum expansion, light quark mass expansion and $`N_c^1`$ expansion. In ref., we have provided a rigorous method to perform calculation to include all order terms of the momentum expansion and up to the next to leading order of $`N_c^1`$ expansion. In this the present paper, we will extend this method to reflect all order information of light quark mass expansion.
The Kaplan-Manohar ambiguity of ChPT has caused many debates. In particular, due to this ambiguity, the authors of ref argued that the observed mass spectrum is consistent with a broad range of quark mass ratios, which specially includes the possibility $`m_u=0`$. However, it is disagreed by other anthors. This problem can also be discussed in ChCQM. Since ChCQM is an effective approach with features of low energy QCD, light current quark masses are defined uniquely in ChCQM, that are just renormalized “physical” masses of u, d, and s quarks. In principle, therefore, there is no Kaplan-Manohar ambiguity in ChCQM. Light current quark masses can be determined uniquely via meson spectrum.
In ref., we studied the isospin breaking process $`\omega \pi ^+\pi ^{}`$ and determined isospin breaking parameter $`m_dm_u=(3.9\pm 0.22)`$MeV at vector meson energy scale. This result together with meson spectrum provide so much information that we can obtain not only light quark mass ratios but also individual quark masses.
In ref., we have shown that the chiral expansion at vector meson energy scale converge slowly. In particular, we have also pointed out that, if we neglect strange quark masses, the chiral expansion at $`\varphi (1020)`$ energy scale will be divergent! It implies that $`m_s`$ play very important role at $`\varphi `$-physics. In ref. we have successfully studies the chiral expansion at $`m_\rho `$ and $`m_\omega `$ energy scale. In order to extend this study to $`K^{}(892)`$ and $`\varphi (1020)`$, individual quark masses are neccessary. It has been recognized that the light quark masses obtained in different approaches are with larger difference. Thus we have to extract quark masses by this formalism itself. It is one of our goals.
In general, in ChCQM the light quark masses can be extracted not only by pseudoscalar meson spectrums, but also by the lowest vector meson resonance spectrums. However, one-loop effects of mesons will contribute to vector meson masses, and calculation on one-loop effects of mesons is related to the chiral expansion at vector meson energy scale. As shown in ref., this relation is very complicate. It makes that the relationship between vector meson spectrums and the light quark masses are also very complicate and indirect. Thus in this paper we will extract information about the light current quark masses from pseudoscalar meson spectrums and their decay constants. The vector meson spectrums will be predicted by this formalism in other paper.
The contents of the paper are organized as follows. In sect. 2 we review the basic notations of the chiral constituent quark model with the lowest vector meson resonances. In sect. 3, masses and decay constants of pseudoscalar meson octet are calculated. The results will including all order information of the light quark masses. In sect. 4, one-loop effects of pseudoscalar mesons and renormalization are discussed. The numerical results are given in sect. 5 and a brief summary is included in sect. 6.
## II Chiral Constituent Quark Model
The simplest version of chiral quark model which was originated by Weinberg, and developed by Manohar and Georgi provides a QCD-inspired description on the simple constituent quark model. In view of this model, in the energy region between the CSSB scale and the confinement scale ($`\mathrm{\Lambda }_{QCD}0.10.3GeV`$), the dynamical field degrees of freedom are constituent quarks(quasi-particle of quarks), gluons and Goldstone bosons associated with CSSB(these Goldstone bosons correspond to lowest pseudoscalar octet). In this quasiparticle description, the effective coupling between gluon and quarks is small and the important interaction is the coupling between quarks and Goldstone bosons. In I we have further included the lowest vector meson resonances into this formalism. At chiral limit, this model is parameterized by the following chiral constituent quark lagrangian
$`_\chi `$ $`=`$ $`i\overline{q}(/+/\mathrm{\Gamma }+g__A/\mathrm{\Delta }\gamma _5i/V)qm\overline{q}q+{\displaystyle \frac{F^2}{16}}<_\mu U^\mu U^{}>+{\displaystyle \frac{1}{4}}m_0^2<V_\mu V^\mu >.`$ (1)
Here $`<\mathrm{}>`$ denotes trace in SU(3) flavour space, $`\overline{q}=(\overline{q}_u,\overline{q}_d,\overline{q}_s)`$ are constituent quark fields. $`V_\mu `$ denotes vector meson octet and singlet. Since in this paper we only focus on pseudoscalar meson spectrums and decay constants, we will neglect vector meson fields in the following. The $`\mathrm{\Delta }_\mu `$ and $`\mathrm{\Gamma }_\mu `$ are defined as follows,
$`\mathrm{\Delta }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi ^{}(_\mu ir_\mu )\xi \xi (_\mu il_\mu )\xi ^{}\},`$ (2)
$`\mathrm{\Gamma }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi ^{}(_\mu ir_\mu )\xi +\xi (_\mu il_\mu )\xi ^{}\},`$ (3)
and covariant derivative are defined as follows
$`_\mu U`$ $`=`$ $`_\mu Uir_\mu U+iUl_\mu =2\xi \mathrm{\Delta }_\mu \xi ,`$ (4)
$`_\mu U^{}`$ $`=`$ $`_\mu U^{}il_\mu U^{}+iU^{}r_\mu =2\xi ^{}\mathrm{\Delta }_\mu \xi ^{},`$ (5)
where $`l_\mu =v_\mu +a_\mu `$ and $`r_\mu =v_\mu a_\mu `$ are linear combinations of external vector field $`v_\mu `$ and axial-vector field $`a_\mu `$, $`\xi `$ associates with non-linear realization of spontanoeusly broken global chiral symmetry introduced by Weinberg. This realization is obtained by specifying the action of global chiral group $`G=SU(3)_L\times SU(3)_R`$ on element $`\xi (\mathrm{\Phi })`$ of the coset space $`G/SU(3)__V`$:
$$\xi (\mathrm{\Phi })g_R\xi (\mathrm{\Phi })h^{}(\mathrm{\Phi })=h(\mathrm{\Phi })\xi (\mathrm{\Phi })g_L^{},g_L,g_RG,h(\mathrm{\Phi })H=SU(3)__V.$$
(6)
Explicit form of $`\xi (\mathrm{\Phi })`$ is usual taken
$$\xi (\mathrm{\Phi })=\mathrm{exp}\{i\lambda ^a\mathrm{\Phi }^a(x)/2\},U(\mathrm{\Phi })=\xi ^2(\mathrm{\Phi }),$$
(7)
where the Goldstone boson $`\mathrm{\Phi }^a`$ are treated as pseudoscalar meson octet. In ref. we have shown that the lagrangian( 1) is invariant under $`G_{\mathrm{global}}\times G_{\mathrm{local}}`$.
Distingushing from some versions of chiral quark models, there is a kinetic term of pseudoscalar mesons in lagrangian( 1). Therefore, the kinetic term of pseudoscalar mesons generated by one-loop effects of constituent quarks can be renormalized. Note that there is no mass term of pseudoscalar mesons in eq. (1). All of these are due to the basic assumption of the model. In the other words, in this energy region, the dynamical field degrees of freedom are constituent quarks and massless Goldstone bosons(pseudoscalar octet) associated with CSSB. Masses of pseudoscalar mesons will be genarated by quark loops as current quark mass parameters emerge in the dynamics (to see below). In ref. we have fitted the parameter $`g_A=0.75`$ via $`\beta `$-decay of neutron, and $`m=480`$MeV via low energy limit of the model. It has been also pointed out that the value of $`g_A`$ has included effects of intermediate axial-vector meson resonances exchanges at low energy.
The light current quark mass-dependent term has been introduced in ref. based on requirement of the chiral symmetry,
$$\frac{1}{2}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}+\xi \stackrel{~}{\chi }^{}\xi )q\frac{\kappa }{2}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}\xi \stackrel{~}{\chi }^{}\xi )\gamma _5q,$$
(8)
where $`\stackrel{~}{\chi }=s+ip`$, $`s=s_{\mathrm{ext}}+`$, $`=\mathrm{diag}\{m_u,m_d,m_s\}`$ is light current quark mass matrix, $`s_{\mathrm{ext}}`$ and $`p`$ are scalar and pseudoscalar external field respectively. Eq. (8) will return to standard quark mass term of QCD, $`\overline{\psi }\psi `$, in absence of pseudoscalar mesons at high energy for arbitrary $`\kappa `$. Therefore, although the light current quark masses are defined uniquely in this formalism, the symmetry and the constrains of underlying QCD still can not fixed the coupling between pseudoscalar mesons and constituent quarks. $`\kappa `$ will be treated as an initial parameter of the model and be fitted phenomenologically.
To conclude this section, the ChCQM lagrangian with light current quark masses is
$`_\chi `$ $`=`$ $`i\overline{q}(/+/\mathrm{\Gamma }+g__A/\mathrm{\Delta }\gamma _5)qm\overline{q}q{\displaystyle \frac{1}{2}}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}+\xi \stackrel{~}{\chi }^{}\xi )q{\displaystyle \frac{\kappa }{2}}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}\xi \stackrel{~}{\chi }^{}\xi )\gamma _5q`$ (10)
$`+{\displaystyle \frac{F^2}{16}}<_\mu U^\mu U^{}>,`$
where vector meson fields have been omitted. The effects of isospin breaking due to inequality of light quark masses will be generated via constituent quark loops. For example, it generated masses of pseudoscalar mesons and splits their decay constants.
## III Quark Loops
In this section, we will calculate pseudoscalar meson masses and decay constants induced by one-loop effects of constituent quarks. They are the leading order in $`N_c^1`$ expansion.
In this framework, the effective action describing meson interaction can be obtained via integrating over degrees of freedom of fermions
$$e^{iS_{\mathrm{eff}}}𝒟\overline{q}𝒟qe^{i{\scriptscriptstyle d^4x_\chi (x)}}=<vac,out|in,vac>_{\mathrm{Ext}},$$
(11)
where $`<vac,out|in,vac>_{\mathrm{Ext}}`$ is vacuum expectation value in presence external sources. The above path integral can be performed formally, and many methods, such as heat kernel manner, have been used to regulate the bilinear operator yielded via path integral. By using those formal integral methods, however, the explicit calculations of high order contributions to the chiral expansion are extremely tedious. In ref. we have provided a convenient method to evaluate effective action via calculating one-loop diagrams of constituent quarks directly. This method can capture all high order contributions of the chiral expansion.
In interaction picture, the equation( 11) is rewritten as follow
$`e^{iS_{\mathrm{eff}}}`$ $`=`$ $`<0|𝒯_qe^{i{\scriptscriptstyle d^4x_\chi ^\mathrm{I}(x)}}|0>`$ (12)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}i{\displaystyle d^4p_1\frac{d^4p_2}{(2\pi )^4}\mathrm{}\frac{d^4p_n}{(2\pi )^4}\stackrel{~}{\mathrm{\Pi }}_n(p_1,\mathrm{},p_n)\delta ^4(p_1p_2\mathrm{}p_n)}`$ (13)
$``$ $`i\mathrm{\Pi }_1(0)+{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}i{\displaystyle \frac{d^4p_1}{(2\pi )^4}\mathrm{}\frac{d^4p_{n1}}{(2\pi )^4}\mathrm{\Pi }_n(p_1,\mathrm{},p_{n1})},`$ (14)
where $`𝒯_q`$ is time-order product of constituent quark fields, $`_\chi ^\mathrm{I}`$ is quark-meson interaction part of lagrangian( 10), $`\stackrel{~}{\mathrm{\Pi }}_n(p_1,\mathrm{},p_n)`$ is one-loop effects of constituent quarks with $`n`$ external sources, $`p_1,p_2,\mathrm{},p_n`$ are four-momentas of $`n`$ external sources respectively and
$$\mathrm{\Pi }_n(p_1,\mathrm{},p_{n1})=d^4p_n\stackrel{~}{\mathrm{\Pi }}_n(p_1,\mathrm{},p_n)\delta ^4(p_1p_2\mathrm{}p_n).$$
(15)
To get rid of all disconnected diagrams, we have
$`S_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}S_n,`$ (16)
$`S_1`$ $`=`$ $`\mathrm{\Pi }_1(0),`$ (17)
$`S_n`$ $`=`$ $`{\displaystyle \frac{d^4p_1}{(2\pi )^4}\mathrm{}\frac{d^4p_{n1}}{(2\pi )^4}\mathrm{\Pi }_n(p_1,\mathrm{},p_{n1})},(n2).`$ (18)
Hereafter we will call $`S_n`$ as $`n`$-point effective action.
For the purpose of this paper, the lagrangian( 10) can be rewritten as follow
$`_\chi (x)=_q(x)+_2^{(0)}(x)+[\pi (x),\eta _8(x)]+[K^\pm (x)]+[K^0(x)],`$ (19)
where $`_q`$ and $`_2^{(0)}`$ are free field lagrangian of constituent quarks of pseudoscalar meson respectively,
$`_q`$ $`=`$ $`{\displaystyle \underset{i=u,d,s}{}}\overline{q}_i(i/\overline{m}_i)q_i,\overline{m}_i=m+m_i,`$ (20)
$`_2^{(0)}`$ $`=`$ $`{\displaystyle \frac{F^2}{16}}<_\mu U^\mu U^{}>`$ (21)
$`=`$ $`{\displaystyle \frac{F^2}{8}}_\mu \mathrm{\Phi }^a^\mu \mathrm{\Phi }^a+4a_\mu ^a^\mu \mathrm{\Phi }^a+\mathrm{},a=1,2,\mathrm{},8,`$ (22)
where
$`\lambda ^a\mathrm{\Phi }^a(x)`$ $`=`$ $`\sqrt{2}\left(\begin{array}{ccc}\frac{\pi _3}{\sqrt{2}}+\frac{\eta _8}{\sqrt{6}}& \pi ^+& K^+\\ \pi ^{}& \frac{\pi _3}{\sqrt{2}}+\frac{\eta _8}{\sqrt{6}}& K^0\\ K^{}& \overline{K}^0& \frac{2}{\sqrt{3}}\eta _8\end{array}\right),`$
$`\lambda ^aa_\mu ^a`$ $`=`$ $`\sqrt{2}\left(\begin{array}{ccc}A_\mu ^{(u)}& a_\mu ^+& A_\mu ^+\\ a_\mu ^{}& A_\mu ^{(d)}& A_\mu ^0\\ A_\mu ^{}& \overline{A}_\mu ^0& A_\mu ^{(s)}\end{array}\right).`$
$`[\mathrm{\Phi }(x)]`$ denotes quark-meson coupling lagrangian,
$`[\pi ,\eta _8]`$ $`=`$ $`{\displaystyle \frac{g_A}{\sqrt{2}}}[(_\mu \pi ^++2a_\mu ^+)\overline{u}\gamma ^\mu \gamma _5d+c.c.]+{\displaystyle \frac{g_A}{2}}{\displaystyle \underset{i=u,d,s}{}}(_\mu P_i+2A_\mu ^{(i)})\overline{q}_i\gamma ^\mu \gamma _5q_i`$ (27)
$`+{\displaystyle \frac{i}{\sqrt{2}}}\kappa (m_u+m_d)(\pi ^+\overline{u}\gamma _5d+c.c.)+i\kappa {\displaystyle \underset{i=u,d,s}{}}m_iP_i\overline{q}_i\gamma _5q_i`$
$`+{\displaystyle \frac{1}{2}}(m_u+m_d)\pi ^+\pi ^{}(\overline{u}u+\overline{d}d)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=u,d,s}{}}m_iP_i^2\overline{q}_iq_i,`$
$`[K^\pm ]`$ $`=`$ $`{\displaystyle \frac{g_A}{\sqrt{2}}}[(_\mu K^++2A_\mu ^+)\overline{u}\gamma ^\mu \gamma _5s+c.c.]+{\displaystyle \frac{i}{\sqrt{2}}}\kappa (m_u+m_s)(K^+\overline{u}\gamma _5s+c.c.)`$ (29)
$`+{\displaystyle \frac{1}{2}}(m_u+m_s)K^+K^{}(\overline{u}u+\overline{s}s),`$
$`[K^0]`$ $`=`$ $`{\displaystyle \frac{g_A}{\sqrt{2}}}[(_\mu K^0+2A_\mu ^0)\overline{d}\gamma ^\mu \gamma _5s+(_\mu \overline{K}^0+2\overline{A}_\mu ^0)d\gamma ^\mu \gamma _5\overline{s}]`$ (31)
$`+{\displaystyle \frac{i}{\sqrt{2}}}\kappa (m_d+m_s)(K^0\overline{d}\gamma _5s+\overline{K}^0d\gamma _5\overline{s})+{\displaystyle \frac{1}{2}}(m_d+m_s)K^0\overline{K}^0(\overline{d}d+\overline{s}s),`$
where $`c.c.`$ denotes charge conjugate term of pervious term,
$`P_u=\pi _3+{\displaystyle \frac{1}{\sqrt{3}}}\eta _8,P_d=\pi _3+{\displaystyle \frac{1}{\sqrt{3}}}\eta _8,P_s={\displaystyle \frac{2}{\sqrt{3}}}\eta _8,`$ (32)
and $`A_\mu ^{(u)},A_\mu ^{(d)};A_\mu ^{(s)}`$ etc. are axial-vector external fields correponding these pseudoscalar meson fields. From eq.( 27) we can see that the $`K^\pm `$-quark coupling and $`K^0`$-quark coupling are similar to $`\pi ^\pm `$-quark coupling. Thus we only need to calculate masses and decay constants of $`\pi ^\pm ,\pi ^0`$ and $`\eta ^8`$. Then masses and decay constants of $`K^\pm `$ can be obtained via replacing $`m_d`$ by $`m_s`$ in one of $`\pi ^\pm `$, and masses and decay constants of $`K^0`$ can be obtained via replacing $`m_u`$ by $`m_s`$ in one of $`\pi ^\pm `$.
The one-point effective action is generated by tadpole-loop of constituent quarks. Calculation about this tadpole-loop contribution is simple.
$`iS_1[\pi ,\eta _8]`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle }d^4x\{(m_u+m_d)\pi ^+(x)\pi ^{}(x)<0|T(\overline{u}(x)u(x)+\overline{d}(x)d(x))|0>`$ (34)
$`+{\displaystyle \underset{i=u,d,s}{}}m_iP_i^2(x)<0|T\overline{q}_i(x)q_i(x)|0>`$
$`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle }d^4x{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}\{[(m_u+m_d)\pi ^+(x)\pi ^{}(x)+m_uP_1^2(x)]Tr[S_F^{(u)}(k)]`$ (36)
$`+[(m_u+m_d)\pi ^+(x)\pi ^{}(x)+m_dP_2^2(x)]Tr[S_F^{(d)}(k)]+m_sP_3^2(x)Tr[S_F^{(s)}(k)]\},`$
where $`Tr`$ denotes trace taking over color and Lorentz space, $`S_F^{(q)}`$ is propagator of constituent quark fields, $`S_F^{(q)}(k)=i(/k\overline{m}_q+iϵ)^1`$. In terms of dimensional regularization, we can integrate over internal line momenta $`k`$ in the above equation. The result is
$`iS_1[\pi ,\eta _8]`$ $`=`$ $`{\displaystyle \frac{2N_c}{(4\pi )^{D/2}}}\mathrm{\Gamma }(1{\displaystyle \frac{D}{2}})i{\displaystyle }d^4x\{({\displaystyle \frac{\mu ^2}{\overline{m}_u^2}})^{ϵ/2}\overline{m}_u^3[(m_u+m_d)\pi ^+\pi ^{}+m_uP_1^2]`$ (38)
$`+({\displaystyle \frac{\mu ^2}{\overline{m}_d^2}})^{ϵ/2}\overline{m}_d^3[(m_u+m_d)\pi ^+\pi ^{}+m_dP_2^2]+({\displaystyle \frac{\mu ^2}{\overline{m}_s^2}})^{ϵ/2}\overline{m}_s^3m_sP_3^2\}.`$
Defining a constant $`B_0`$ to absorbe the quadratic divergence from loop integral,
$$\frac{F_0^2}{16}B_0=\frac{N_c}{(4\pi )^{D/2}}(\frac{\mu ^2}{m^2})^{ϵ/2}\mathrm{\Gamma }(1\frac{D}{2})m^3,$$
(39)
we have
$`S_1[\pi ,\eta _8]`$ $`=`$ $`{\displaystyle d^4x_1[\pi (x),\eta _8(x)]}`$ (40)
$`_1[\pi ,\eta _8]`$ $`=`$ $`{\displaystyle \frac{F_0^2}{8}}B_0\{(x_u^3+x_d^3)(m_u+m_d)\pi ^+\pi ^{}+{\displaystyle \underset{i=u,d,s}{}}x_i^3m_iP_i^2\}`$ (42)
$`{\displaystyle \frac{N_c}{8\pi ^2}}m^3\{(x_u^3\mathrm{ln}x_u^2+x_d^3\mathrm{ln}x_d^2)(m_u+m_d)\pi ^+\pi ^{}+{\displaystyle \underset{i=u,d,s}{}}m_ix_i^3\mathrm{ln}x_i^2P_i^2\}`$
where $`x_q=\overline{m}_q/m`$.
The two-point effective action concerning to masses and decay constants of charge pion can be obtained as follow,
$`iS_2[\pi ^\pm ]={\displaystyle \frac{g_A^2}{2}}{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}(iq_\mu \pi ^+(q)+2a_\mu ^+(q))(iq_\nu \pi ^{}(q)+2a_\nu ^{}(q))Tr[\gamma ^\mu \gamma _5S_F^{(d)}(kq)\gamma ^\nu \gamma _5S_F^{(u)}(k)]}`$ (45)
$`+{\displaystyle \frac{i}{2}}g_A\kappa (m_u+m_d){\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}{\displaystyle \frac{d^4k}{(2\pi )^4}}\{(iq_\mu \pi ^+(q)+2a_\mu ^+(q))\pi ^{}(q)Tr[\gamma ^\mu \gamma _5S_F^{(d)}(kq)\gamma _5S_F^{(u)}(k)]+c.c.\}`$
$`{\displaystyle \frac{\kappa ^2}{2}}(m_u+m_d)^2{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\pi ^+(q)\pi ^{}(q)Tr[\gamma _5S_F^{(d)}(kq)\gamma _5S_F^{(u)}(k)]}`$
$`=`$ $`{\displaystyle \frac{2N_c}{(4\pi )^2}}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})g_A^2i{\displaystyle d^4x\frac{d^4q}{(2\pi )^4}e^{iqx}(iq_\mu \pi ^+(q)+2a_\mu ^+(q))(^\mu \pi ^{}(x)+2a^\mu (x))}`$ (51)
$`\times {\displaystyle _0^1}dt[\overline{m_u}^2+\overline{m}_u\overline{m}_d+t(\overline{m}_d^2\overline{m}_u^2)]\left({\displaystyle \frac{\overline{m_u}^2+t(\overline{m}_d^2\overline{m}_u^2)t(1t)q^2}{4\pi \mu ^2}}\right)^{\frac{ϵ}{2}}`$
$`+{\displaystyle \frac{2N_c}{(4\pi )^2}}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})\kappa g_A(m_u+m_d)i{\displaystyle }d^4x{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}e^{iqx}[(iq_\mu \pi ^+(q)+2a_\mu ^+(q))^\mu \pi ^{}(x)+c.c.]`$
$`\times {\displaystyle _0^1}dt[\overline{m_u}+t(m_dm_u)]\left({\displaystyle \frac{\overline{m_u}^2+t(\overline{m}_d^2\overline{m}_u^2)t(1t)q^2}{4\pi \mu ^2}}\right)^{\frac{ϵ}{2}}`$
$`+{\displaystyle \frac{2N_c}{(4\pi )^2}}\kappa ^2(m_u+m_d)^2i{\displaystyle d^4x\frac{d^4q}{(2\pi )^4}e^{iqx}\pi ^+(q)\pi ^{}(x)_0^1𝑑t\left(\frac{\overline{m_u}^2+t(\overline{m}_d^2\overline{m}_u^2)t(1t)q^2}{4\pi \mu ^2}\right)^{\frac{ϵ}{2}}}`$
$`\times \{\mathrm{\Gamma }(1{\displaystyle \frac{D}{2}})[\overline{m_u}^2+t(\overline{m}_d^2\overline{m}_u^2)t(1t)q^2]\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})[\overline{m_u}^2\overline{m}_u\overline{m}_d+t(\overline{m}_d^2\overline{m}_u^2)2t(1t)q^2]\}.`$
There are both quadratic divergence and logarithmic divergence in the above effective action. The quadratic divergence can be canceled by constant $`B_0`$ defined in eq.( 39), and the logarithmic divergence can be conceled via defining
$`g^2`$ $`=`$ $`{\displaystyle \frac{8N_c}{3(4\pi )^2}}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})({\displaystyle \frac{4\pi \mu ^2}{m^2}})^{\frac{ϵ}{2}},`$ (52)
$`{\displaystyle \frac{F_0^2}{16}}`$ $`=`$ $`{\displaystyle \frac{F^2}{16}}+{\displaystyle \frac{N_c}{(4\pi )^2}}g_A^2m^2\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})({\displaystyle \frac{4\pi \mu ^2}{m^2}})^{\frac{ϵ}{2}}.`$ (53)
$`g`$ is an universal coupling constant of this model. In ref., it has been determined as $`g^2=\frac{N_c}{3\pi ^2}`$ by the first KSRF sum rule.
Then eq.( 45) together with eq.( 40) give $`O(N_c)`$ effective lagrangian containing the terms linear or quadratic in the charge pion fields as follow
$`_2[\pi ^\pm (x)]`$ $`=`$ $`{\displaystyle \frac{F^2(m_u,m_d)}{4}}_\mu \pi ^+^\mu \pi ^{}+{\displaystyle \frac{\overline{f}^2(m_u,m_d)}{2}}(a_\mu ^+^\mu \pi ^{}+c.c.){\displaystyle \frac{F_0^2}{4}}\overline{M}^2(m_u,m_d)\pi ^+\pi ^{}`$ (55)
$`+{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}e^{iqx}\{{\displaystyle \frac{\alpha (q^2;m_u,m_d)}{2}}iq^\mu (a_\mu ^+(x)\pi ^{}(q)+c.c.){\displaystyle \frac{F_0^2}{4}}\beta (q^2;m_u,m_d))\pi ^+(q)\pi ^{}(x)\},`$
where
$`F^2(m_u,m_d)=F_0^2+{\displaystyle \frac{3}{2}}g^2g_A^2(m_u+m_d)(4m+m_u+m_d)+3g^2\kappa g_A(m_u+m_d)(\overline{m}_u+\overline{m}_d)`$ (56)
$`{\displaystyle \frac{N_c}{2\pi ^2}}g_A[g_A(\overline{m}_u+\overline{m}_d)+2\kappa (m_u+m_d)]{\displaystyle _0^1}𝑑t[\overline{m}_u+t(m_dm_u)]\mathrm{ln}(x_u^2+t(x_d^2x_u^2))`$ (57)
$`+4\kappa ^2(m_u+m_d)^2\{({\displaystyle \frac{g^2}{4}}{\displaystyle \frac{F_0^2B_0}{48m^3}}{\displaystyle \frac{N_c}{24\pi ^2}}){\displaystyle \frac{N_c}{8\pi ^2}}{\displaystyle _0^1}𝑑tt(1t)[3\mathrm{ln}(x_u^2+t(x_d^2x_u^2)){\displaystyle \frac{\overline{m}_u\overline{m}_d}{\overline{m}_u^2+t(\overline{m}_d^2\overline{m}_u^2)}}]\},`$ (58)
$`\overline{f}^2(m_u,m_d)=F_0^2+{\displaystyle \frac{3}{2}}g^2g_A^2(m_u+m_d)(4m+m_u+m_d)+{\displaystyle \frac{3}{2}}g^2\kappa g_A(m_u+m_d)(\overline{m}_u+\overline{m}_d)`$ (59)
$`{\displaystyle \frac{N_c}{2\pi ^2}}g_A[g_A(\overline{m}_u+\overline{m}_d)+\kappa (m_u+m_d)]{\displaystyle _0^1}𝑑t[\overline{m}_u+t(m_dm_u)]\mathrm{ln}(x_u^2+t(x_d^2x_u^2)),`$ (60)
$`\overline{M}^2(m_u,m_d)=B_0(m_u+m_d)\{{\displaystyle \frac{1}{2}}(x_u^3+x_d^3)+{\displaystyle \frac{N_c}{2\pi ^2}}{\displaystyle \frac{m^3}{F_0^2B_0}}(x_u^3\mathrm{ln}x_u^2+x_d^3\mathrm{ln}x_d^3){\displaystyle \frac{\kappa ^2}{4}}{\displaystyle \frac{m_u+m_d}{m}}(x_u^2+x_d^2)\}`$ (61)
$`+{\displaystyle \frac{3}{2}}g^2\kappa ^2{\displaystyle \frac{(m_d^2m_u^2)^2}{F_0^2}}{\displaystyle \frac{N_c}{2\pi ^2}}\kappa ^2{\displaystyle \frac{(m_u+m_d)^2}{F_0^2}}{\displaystyle _0^1}𝑑t[2\overline{m}_u^2\overline{m}_u\overline{m}_d+2t(\overline{m}_d^2\overline{m}_u^2)]\mathrm{ln}(x_u^2+t(x_d^2x_u^2)),`$ (62)
$`\alpha (q^2;m_u,m_d)={\displaystyle \frac{N_c}{2\pi ^2}}g_A[g_A(\overline{m}_u+\overline{m}_d)+\kappa (m_u+m_d)]{\displaystyle _0^1}𝑑t[\overline{m}_u+t(m_dm_u)]\mathrm{ln}\left(1{\displaystyle \frac{t(1t)q^2}{\overline{m}_u^2+t(\overline{m}_d^2\overline{m}_u^2)}}\right),`$ (63)
$`\beta (q^2;m_u,m_d)={\displaystyle \frac{N_c}{2\pi ^2}}g_A{\displaystyle \frac{q^2}{F_0^2}}[g_A(\overline{m}_u+\overline{m}_d)+2\kappa (m_u+m_d)]{\displaystyle _0^1}𝑑t[\overline{m}_u+t(m_dm_u)]\mathrm{ln}\left(1{\displaystyle \frac{t(1t)q^2}{\overline{m}_u^2+t(\overline{m}_d^2\overline{m}_u^2)}}\right)`$ (64)
$`{\displaystyle \frac{N_c}{2\pi ^2}}\kappa ^2{\displaystyle \frac{(m_u+m_d)^2}{F_0^2}}{\displaystyle _0^1}dt\{[2\overline{m}_u^2\overline{m}_u\overline{m}_d+2t(\overline{m}_d^2\overline{m}_u^2)3t(1t)q^2]\mathrm{ln}(1{\displaystyle \frac{t(1t)q^2}{\overline{m}_u^2+t(\overline{m}_d^2\overline{m}_u^2)}})`$ (65)
$`+t(1t)q^2(2{\displaystyle \frac{\overline{m}_u\overline{m}_d}{\overline{m}_u^2+t(\overline{m}_d^2\overline{m}_u^2)}})\}.`$ (66)
It should be pointed out that $`\alpha (q^2;m_u,m_d)`$ is order $`q^2`$ at least and $`\beta (q^2;m_u,m_d)`$ is order $`q^4`$ at least. Since in this paper we focus on pseudoscalar meson spectrums, these high order derivative terms should obey motion equation of pseudoscalar mesons. In momentum space, the motion equation of physical pseudoscalar mesons is generally written
$`(q^2m_\phi ^2)\phi (q)=if_\phi q^\mu A_\mu ^{(\phi )}(q),`$ (67)
where $`m_\phi `$ and $`f_\phi `$ are physical mass and decay constants of pseudoscalar, e.g., $`m_\pi =135`$MeV and $`f_\pi =185.2`$MeV. Due to this motion equation, we have
$`\alpha (q^2;m_u,m_d)a_\mu ^+(x)\pi ^{}(q)`$ $`=`$ $`\alpha (m_\pi ^2;m_u,m_d)a_\mu ^+(x)\pi ^{}(q),`$ (68)
$`\beta (q^2;m_u,m_d)\pi ^+(q)\pi ^{}(q)`$ $`=`$ $`\beta (m_\pi ^2;m_u,m_d)\pi ^+(q)\pi ^{}(q){\displaystyle \frac{i}{2}}q^\mu f_\pi \beta ^{}(m_\pi ^2;m_u,m_d)(a_\mu ^+(x)\pi ^{}(q)+c.c.),`$ (69)
where
$`\beta ^{}(m_\pi ^2;m_u,m_d)={\displaystyle \frac{d}{dq^2}}\beta ^{}(q^2;m_u,m_d)|_{q^2=m_\pi ^2}.`$ (70)
Thus eq.( 55) can be rewritten
$`_2[\pi ^\pm (x)]`$ $`=`$ $`{\displaystyle \frac{F^2(m_u,m_d)}{4}}_\mu \pi ^+^\mu \pi ^{}{\displaystyle \frac{F_0^2}{4}}\{\overline{M}^2(m_u,m_d)+\beta (m_\pi ^2;m_u,m_d)\}\pi ^+\pi ^{}`$ (72)
$`+\{{\displaystyle \frac{\overline{f}^2(m_u,m_d)}{2}}+{\displaystyle \frac{\alpha (m_\pi ^2;m_u,m_d)}{2}}+{\displaystyle \frac{F_0^2}{4F(m_u,m_d)}}f_\pi \beta ^{}(m_\pi ^2;m_u,m_d)\}(a_\mu ^+^\mu \pi ^{}+c.c.).`$
The two-point effective action concerning to masses and decay constants of neutral pion and $`\eta _8`$ can be evaluated similarly. However the decay constants for neutral mesons cannot be extracted directly from the data. It means that the decay constants for neutral mesons cannot be used to determine light current quark masses. Therefore, in this paper we do not need to evaluate the decay constants for neutral pion and $`\eta _8`$. The effective action containing the quadratic terms of neutral pion and $`\eta _8`$ is then
$`iS_2[\pi ^0,\eta _8]`$ (73)
$`=`$ $`{\displaystyle \frac{g_A^2}{8}}{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\underset{j=u,d,s}{}q^2P_j(q)P_j(q)Tr[\gamma ^\mu \gamma _5S_F^{(j)}(kq)\gamma ^\nu \gamma _5S_F^{(j)}(k)]}`$ (76)
$`+{\displaystyle \frac{i}{2}}g_A\kappa {\displaystyle \underset{j=u,d,s}{}}m_j{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}(iq_\mu )P_j(q)P_j(q)Tr[\gamma ^\mu \gamma _5S_F^{(j)}(kq)\gamma _5S_F^{(j)}(k)]}`$
$`{\displaystyle \frac{\kappa ^2}{2}}{\displaystyle \underset{j=u,d,s}{}}m_j^2{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}P_j(q)P_j(q)Tr[\gamma _5S_F^{(j)}(kq)\gamma _5S_F^{(j)}(k)]}`$
$`=`$ $`{\displaystyle \frac{N_c}{(4\pi )^2}}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})g_A^2i{\displaystyle \frac{d^4q}{(2\pi )^4}q^2\underset{j=u,d,s}{}\overline{m}_j^2q^2P_j(q)P_j(q)_0^1𝑑t\left(\frac{\mu ^2}{\overline{m}_u^2t(1t)q^2}\right)^{\frac{ϵ}{2}}}`$ (80)
$`+{\displaystyle \frac{2N_c}{(4\pi )^2}}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})\kappa g_Ai{\displaystyle \frac{d^4q}{(2\pi )^4}q^2\underset{j=u,d,s}{}m_j\overline{m}_jq^2P_j(q)P_j(q)_0^1𝑑t\left(\frac{\mu ^2}{\overline{m}_u^2t(1t)q^2}\right)^{\frac{ϵ}{2}}}`$
$`+{\displaystyle \frac{2N_c}{(4\pi )^2}}\kappa ^2i{\displaystyle \frac{d^4q}{(2\pi )^4}\underset{j=u,d,s}{}m_j^2P_j(q)P_j(q)_0^1𝑑t\left(\frac{\mu ^2}{\overline{m}_u^2t(1t)q^2}\right)^{\frac{ϵ}{2}}}`$
$`\times \{\mathrm{\Gamma }(1{\displaystyle \frac{D}{2}})[\overline{m}_j^2t(1t)q^2]+2t(1t)q^2\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}})\}.`$
The divergences in the above effective action can be canceled by eqs.( 39) and ( 52). Then from eqs.( 40) and ( 73) we can obtain $`O(N_c)`$ effective lagrangian describing two-point vertex of $`\pi ^0`$ and $`\eta _8`$ as follow
$`_2[\pi ^0(x),\eta _8(x)]={\displaystyle \underset{i=u,d,s}{}}\{{\displaystyle \frac{F_i^2(m_i)}{8}}_\mu P_i^\mu P_i{\displaystyle \frac{F_0^2}{8}}\overline{M}_i^2(m_i)P_i^2{\displaystyle \frac{F_0^2}{8}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\beta _i(q^2;m_i)P_i(q)P_i(x)}\},`$ (81)
where
$`F_i^2(m_i)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{2}}+3g^2g_A^2m_i(2m+m_i)+6g^2\kappa g_Am_i\overline{m}_i+8\kappa ^2({\displaystyle \frac{g^2}{4}}{\displaystyle \frac{F_0^2B_0}{48m^3}}{\displaystyle \frac{N_c}{48\pi ^2}})m_i^2`$ (83)
$`{\displaystyle \frac{N_c}{2\pi ^2}}(g_A^2\overline{m}_i^2+2g_A\kappa m_i\overline{m}_i+\kappa ^2m_i^2)\mathrm{ln}x_i^2,`$
$`\overline{M}_i^2(m_i)`$ $`=`$ $`(B_0+{\displaystyle \frac{N_c}{\pi ^2}}{\displaystyle \frac{m^3}{F_0^2}}\mathrm{ln}x_i^2)m_ix_i^2(x_i{\displaystyle \frac{m_i}{m}}\kappa ^2),`$ (84)
$`\beta _i(q^2;m_i)`$ $`=`$ $`{\displaystyle \frac{N_c}{2\pi ^2F_0^2}}q^2\{{\displaystyle \frac{\kappa ^2}{3}}m_i^2+\left(g_A^2\overline{m}_i^2+2\kappa m_i\overline{m}_i+2\kappa ^2m_i^2[\overline{m}_i^2q^23t(1t)]\right){\displaystyle _0^1}𝑑t\mathrm{ln}(1{\displaystyle \frac{t(1t)q^2}{\overline{m}_i^2}})\}.`$ (85)
Since auxiliary fields $`P_i`$ do not lie in physical hadron spectrums, the equation of motion (67) can not be used in (81) simply. In section 5, we will use propagator method to deal with the terms with high power momenta in (81) and diagonalize $`\pi _3\eta _8`$ mixing.
## IV Meson Loops and Renormaliztion
Purpose of this section is to evaluate one-loop effects of pseudoscalar mesons. Due to parity conservation, there are only tadpole diagrams of pseudoscalar mesons contributing to masses of decay constants of $`0^{}`$ mesons(fig. 1). Since ChCQM is a nonrenormalizable effective theory, it is very difficlut to calculate meson loop effects completely. However, we can expect that, in mass spectrums and decay constants of pseudoscalar mesons, the dominant one-loop effects is generated by the lowest order effective lagrangian, because neglect is $`O(m_q^3)`$ and suppressed by $`N_c^1`$ expansion. Fortunately, the formalism is renormilzed up to this order.
The lowest order effective lagrangian is well-known
$`_2`$ $`=`$ $`{\displaystyle \frac{F_0^2}{16}}<_\mu U^\mu U^{}+\chi U^{}+U\chi ^{}>`$ (86)
$`=`$ $`{\displaystyle \frac{F_0^2}{48}}<[\lambda ^a,\mathrm{\Delta }_\mu ][\lambda ^a,\mathrm{\Delta }^\mu ]>+{\displaystyle \frac{3F_0^2}{256}}<\lambda ^a\lambda ^a(\xi \chi ^{}\xi +\xi ^{}\chi \xi ^{})>,`$ (87)
where $`\lambda ^a(a=1,2,\mathrm{},8)`$ are Gell-Mann matrices, $`\chi =2B_0\stackrel{~}{\chi }`$ and the following SU(N) completeness relation have been used
$`{\displaystyle \underset{a=1}{\overset{N^21}{}}}<\lambda ^aA\lambda ^aB>`$ $`=`$ $`{\displaystyle \frac{2}{N}}<AB>+2<A><B>,`$ (88)
$`{\displaystyle \underset{a=1}{\overset{N^21}{}}}<\lambda ^aA><\lambda ^aB>`$ $`=`$ $`\mathrm{\hspace{0.33em}2}<AB>{\displaystyle \frac{2}{N}}<A><B>.`$ (89)
To evaluate the one-loop graphs generated by this lagrangian, we consider the quantum fluctuation $`\phi (x)=\phi ^a(x)\lambda ^a`$ around the solution $`\overline{U}(x)=\overline{\xi }^2(x)`$ to the classical equations of motion,
$`U=\overline{\xi }e^{i\phi }\overline{\xi }.`$ (90)
Substituting expansion( 90) into $`_2`$ and retaining terms up to and including $`\phi ^2`$ one obtains
$`_2\overline{}_2+{\displaystyle \frac{F_0^2}{8}}(_\mu \phi ^a^\mu \phi ^am_\phi ^2\phi ^a\phi ^a){\displaystyle \frac{F_0^2}{16}}<[\phi ,\mathrm{\Delta }_\mu ][\phi ,\mathrm{\Delta }^\mu ]+{\displaystyle \frac{1}{4}}\{\phi ,\phi \}(\overline{\xi }\chi ^{}\overline{\xi }+\overline{\xi }^{}\chi \overline{\xi }^{}2)>,`$ (91)
where we have omitted some terms which do not contribute to masses and decay constants via one-loop graphs. The contribution of tadpole graphs can be calculated easily
$`_2^{(\mathrm{tad})}`$ $`=`$ $`{\displaystyle \frac{1}{4}}<[\lambda ^a,\mathrm{\Delta }_\mu ][\lambda ^a,\mathrm{\Delta }^\mu ]+{\displaystyle \frac{1}{4}}\{\lambda ^a,\lambda ^a\}(\overline{\xi }\chi ^{}\overline{\xi }+\overline{\xi }^{}\chi \overline{\xi }^{}2)>{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{i}{k^2m_\phi ^2+iϵ}}`$ (92)
$`=`$ $`{\displaystyle \frac{1}{4}}(m_\phi ^2\lambda {\displaystyle \frac{m_\phi ^2}{16\pi ^2}}\mathrm{ln}{\displaystyle \frac{m_\phi ^2}{\mu ^2}})<[\lambda ^a,\mathrm{\Delta }_\mu ][\lambda ^a,\mathrm{\Delta }^\mu ]+{\displaystyle \frac{1}{4}}\{\lambda ^a,\lambda ^a\}(\overline{\xi }\chi ^{}\overline{\xi }+\overline{\xi }^{}\chi \overline{\xi }^{}2)>,`$ (93)
where
$$\lambda =\frac{1}{16\pi ^2}\{\frac{2}{ϵ}+\mathrm{ln}(4\pi )+\gamma +1\}.$$
Comparing eq.( 92) and eq.( 86) we can see that the divergence $`\lambda `$ can be absorbed by free parameters $`F_0`$ and $`B_0`$. Thus the sum of tree graphs and tadpole contribution is
$`_2^{(t)}={\displaystyle \frac{F_0^2}{48}}(13\mu _\phi )v<[\lambda ^a,\mathrm{\Delta }_\mu ][\lambda ^a,\mathrm{\Delta }^\mu ]>+{\displaystyle \frac{3F_0^2}{256}}(1{\displaystyle \frac{8}{3}}\mu _\phi )<\lambda ^a\lambda ^a(\xi \chi ^{}\xi +\xi ^{}\chi \xi ^{})>,`$ (94)
where
$`\mu _\phi ={\displaystyle \frac{m_\phi ^2}{4\pi ^2F_0^2}}\mathrm{ln}{\displaystyle \frac{m_\phi ^2}{\mu ^2}}.`$ (95)
In lagrangian( 94) we appoint that $`m_\phi =m_\pi `$ when $`a=1,2,3`$, $`m_\phi =m__K`$ when $`a=4,5,6,7`$ and $`m_\phi =m_{\eta _8}`$ when $`a=8`$.
## V Light Quark Mass Determination Beyond the Chiral Perturbation Expansion
For extracting $`(m_{_{K^+}})_{\mathrm{QCD}}`$ from experimental data, the electromagnetic mass splitting of $`K`$-meson is required. The prediction of Dashen theorem, $`(m_{_{K^+}}m_{_{K^0}})_{e.m.}=1.3`$MeV, has been corrected in serveral recent analysis with considering contribution from vector meson exchange. A larger correction is first obtained by Donoghue, Holstein and Wyler, who find $`(m_{_{K^+}}m_{_{K^0}})_{e.m.}=2.3`$MeV. Then Bijnens and Prades, who evaluated both long-distance contribution using ENJL model and short-distance contribution using perturbative QCD and factorization, find $`(m_{_{K^+}}m_{_{K^0}})_{e.m.}=2.4\pm 0.3`$MeV at $`\mu =m_\rho `$. Gao et.al. also gave $`(m_{_{K^+}}m_{_{K^0}})_{e.m.}=2.5`$MeV. Baur and Urech, however, obtained a smaller correction, $`(m_{_{K^+}}m_{_{K^0}})_{e.m.}=1.6`$MeV at $`\mu =m_\rho `$. In addition, calculation of lattice QCD found $`(m_{_{K^+}}m_{_{K^0}})_{e.m.}=1.9`$MeV. These estimates indicate that the corrections to Dashen theorem are indeed substantial. In this the present paper, we average the above results and take $`m_{_{K^+}}m_{_{K^0}})_{e.m.}=2.1\pm 0.1`$MeV at energy scale $`\mu =m_\rho `$.
From eqs.( 72) and (94), the masses and decay constants of koan and chrage pion can be obtained via solve the following equations
$`m_{\pi ^+}^2`$ $`=`$ $`{\displaystyle \frac{F_0^2}{F__R^2(m_u,m_d)}}\{\overline{M}__R^2(m_u,m_d)+\beta (m_\pi ^2;m_u,m_d)\}`$ (96)
$`m_{_{K^+}}^2`$ $`=`$ $`{\displaystyle \frac{F_0^2}{F__R^2(m_u,m_s)}}\{\overline{M}__R^2(m_u,m_s)+\beta (m__K^2;m_u,m_s)\},`$ (97)
$`m_{_{K^0}}^2`$ $`=`$ $`{\displaystyle \frac{F_0^2}{F__R^2(m_d,m_s)}}\{\overline{M}__R^2(m_d,m_s)+\beta (m__K^2;m_d,m_s)\},`$ (98)
$`f_{\pi ^+}`$ $`=`$ $`{\displaystyle \frac{\overline{f}__R^2(m_u,m_d)}{F__R(m_u,m_d)}}+{\displaystyle \frac{\alpha (m_\pi ^2;m_u,m_d)}{F__R(m_u,m_d)}}+{\displaystyle \frac{F_0^2}{2F__R^2(m_u,m_d)}}f_{\pi ^+}\beta ^{}(m_\pi ^2;m_u,m_d),`$ (99)
$`f_{_{K^+}}`$ $`=`$ $`{\displaystyle \frac{\overline{f}__R^2(m_u,m_s)}{F__R(m_u,m_s)}}+{\displaystyle \frac{\alpha (m__K^2;m_u,m_s)}{F__R(m_u,m_s)}}+{\displaystyle \frac{F_0^2}{2F__R^2(m_u,m_s)}}f__K\beta ^{}(m__K^2;m_u,m_s),`$ (100)
$`f_{_{K^0}}`$ $`=`$ $`{\displaystyle \frac{\overline{f}__R^2(m_d,m_s)}{F__R(m_d,m_s)}}+{\displaystyle \frac{\alpha (m__K^2;m_d,m_s)}{F__R(m_d,m_s)}}+{\displaystyle \frac{F_0^2}{2F__R^2(m_d,m_s)}}f__K\beta ^{}(m__K^2;m_d,m_s),`$ (101)
where subscript “R” denotes renormalized quantity,
$`F__R^2(m_u,m_d)`$ $`=`$ $`F^2(m_u,m_d)F_0^2(2\mu _\pi +\mu __K),`$ (102)
$`F__R^2(m_i,m_s)`$ $`=`$ $`F^2(m_i,m_s){\displaystyle \frac{3}{4}}F_0^2(\mu _\pi +2\mu __K+\mu _{\eta _8}),(i=u,d)`$ (103)
$`\overline{M}__R^2(m_u,m_d)`$ $`=`$ $`\overline{M}^2(m_u,m_d)B_0(m_u+m_d)({\displaystyle \frac{3}{2}}\mu _\pi +\mu __K+{\displaystyle \frac{1}{6}}\mu _{\eta _8}),`$ (104)
$`\overline{M}__R^2(m_i,m_s)`$ $`=`$ $`\overline{M}^2(m_i,m_s)B_0(m_i+m_s)({\displaystyle \frac{3}{4}}\mu _\pi +{\displaystyle \frac{3}{2}}\mu __K+{\displaystyle \frac{5}{12}}\mu _{\eta _8}),`$ (105)
$`\overline{f}__R^2(m_u,m_d)`$ $`=`$ $`\overline{f}^2(m_u,m_d)F_0^2(2\mu _\pi +\mu __K),`$ (106)
$`\overline{f}__R^2(m_i,m_s)`$ $`=`$ $`\overline{f}^2(m_i,m_s){\displaystyle \frac{3}{4}}F_0^2(\mu _\pi +2\mu __K+\mu _{\eta _8}),`$ (107)
Here the quantity $`\mu _\phi `$ is defined in (95) and depend on the renormalization scale $`\mu `$. It has been recognized that the scale plays an important role in low-energy QCD. In this formalism, many parameters, such as light current quark masses, constituent quark mass $`m`$, axial-vector coupling constant $`g_A`$, are scale-dependent. The characteristic scale of the model is described by the universal coupling constant $`g`$ which is determined by the first KSRF sum rule at energy scale $`\mu =m_\rho `$. Hence we take $`\mu =m_\rho =770`$MeV in $`\mu _\phi `$, and the physical masses and decay constants, however, should be independent of the renormalization scale.
In order to obtain the masses of $`\pi ^0`$ and $`\eta ^8`$, the $`\pi _3\eta _8`$ mixing in eq.( 81) must be diagonalized. Eq.(81) together with eq.(94) lead to the quadratic terms for the $`\pi _3`$ and $`\eta _8`$ are of the form
$`S_2[\pi _3,\eta _8]={\displaystyle \frac{d^4q}{(2\pi )^4}\{\frac{1}{2}(q^2M_3^2(q^2))\pi _3^2+\frac{1}{2}(q^2M_8^2(q^2))\eta _8^2M_{38}^2(q^2)\pi _3\eta _8\}}`$ (108)
where
$`M_3^2(q^2)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{F_3^2}}\{\overline{M}_u^2(m_u)+\overline{M}_d^2(m_d)+\beta _u(q^2;m_u)+\beta _d(q^2;m_d)2B_0\widehat{m}({\displaystyle \frac{3}{2}}\mu _\pi +\mu __K+{\displaystyle \frac{1}{6}}\mu _{\eta _8})\},`$ (109)
$`M_8^2(q^2)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{3F_8^2}}\{\overline{M}_u^2(m_u)+\overline{M}_d^2(m_d)+4\overline{M}_s^2(m_s)+\beta _u(q^2;m_u)+\beta _d(q^2;m_d)+4\beta _s(q^2;m_s)`$ (111)
$`2B_0(\widehat{m}+2m_s)(2\mu __K+{\displaystyle \frac{2}{3}}\mu _{\eta _8})2B_0\widehat{m}({\displaystyle \frac{3}{2}}\mu _\pi +\mu __K+{\displaystyle \frac{1}{2}}\mu _{\eta _8})\},`$
$`M_{38}^2(q^2)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{\sqrt{3}F_3F_8}}\{\overline{M}_u^2(m_u)\overline{M}_d^2(m_d)+\beta _u(q^2;m_u)\beta _d(q^2;m_d)B_0(m_um_d)({\displaystyle \frac{3}{2}}\mu _\pi +\mu __K+{\displaystyle \frac{1}{6}}\mu _{\eta _8})\}`$ (113)
$`+q^2{\displaystyle \frac{F_u^2(m_u)F_d^2(m_d)}{\sqrt{3}F_3F_8}},`$
with $`\widehat{m}=(m_u+m_d)/2`$ and
$`F_3^2`$ $`=`$ $`F_u^2(m_u)+F_d^2(m_d)F_0^2(2\mu _\pi +\mu __K),`$ (114)
$`F_8^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}\{F_u^2(m_u)+F_d^2(m_d)+4F_s^2(m_s)\}3F_0^2\mu __K.`$ (115)
Due to $`\pi _3\eta _8`$ mixing, the “physical” propagators of $`\pi ^0`$ and $`\eta _8`$ are obtained via the chain approximation in momentum space
$`{\displaystyle \frac{i}{q^2m_{\pi ^0}^2+iϵ}}`$ $`=`$ $`{\displaystyle \frac{i}{q^2M_3^2(q^2)+iϵ}}+{\displaystyle \frac{iM_{38}^4(q^2)}{(q^2M_8^2(q^2)+iϵ)(q^2M_3^2(q^2)+iϵ)^2}}+\mathrm{}`$ (116)
$`=`$ $`{\displaystyle \frac{i}{q^2M_3^2(q^2)\frac{M_{38}^4(q^2)}{q^2M_8^2(q^2)}+iϵ}},`$ (117)
$`{\displaystyle \frac{i}{q^2m_{\eta _8}^2+iϵ}}`$ $`=`$ $`{\displaystyle \frac{i}{q^2M_8^2(q^2)+iϵ}}+{\displaystyle \frac{iM_{38}^4(q^2)}{(q^2M_3^2(q^2)+iϵ)(q^2M_8^2(q^2)+iϵ)^2}}+\mathrm{}`$ (118)
$`=`$ $`{\displaystyle \frac{i}{q^2M_8^2(q^2)\frac{M_{38}^4(q^2)}{q^2M_3^2(q^2)}+iϵ}}.`$ (119)
Then the masses of $`\pi ^0`$ and $`\eta _8`$ are just solutions of the following equations,
$`m_{\pi ^0}^2`$ $`=`$ $`M_3^2(m_{\pi ^0}^2)+{\displaystyle \frac{M_{38}^2(m_{\pi ^0}^2)}{m_{\pi ^0}^2M_8^2(m_{\pi ^0}^2)}},`$ (120)
$`m_{\eta ^8}^2`$ $`=`$ $`M_8^2(m_{\eta _8}^2)+{\displaystyle \frac{M_{38}^2(m_{\eta _8}^2)}{m_{\eta _8}^2M_8^2(m_{\eta _8}^2)}}.`$ (121)
In eqs.( 56) and ( 83), the parameters $`\kappa `$, $`F_0`$ and $`B_0`$ are still not determined. In order to determine them and three light quark masses, six inputs are required. In this paper we choose $`f_{\pi ^+}=185.2\pm 0.5`$MeV, $`f_{_{K^+}}=226.0\pm 2.5`$MeV, $`m_{\pi ^0}=134.98`$MeV, $`m_{_{K^0}}=497.67`$MeV and $`(m_{_{K^+}})_{\mathrm{QCD}}=491.6\pm 0.1`$MeV. Another input is $`m_dm_u=3.9\pm 0.22`$MeV, which is extracted from $`\omega \pi ^+\pi ^{}`$ decay at energy scale $`\mu =m_\rho `$ in ref.. Recalling $`m=480`$MeV, $`g_A=0.75`$ and $`g=\pi ^1`$ for $`N_c=3`$, we can fit light quark masses as in table 1.
In table 1, the errors of results are from uncertainties in decay constant of $`K^+`$. electromagnetic mass splitting of $`K`$-mesons and isospin violation parameter $`m_dm_u`$ respectively. The first column in table 1 we show the results for $`f_{_{K^+}}=223.5`$MeV. The third column corresponds the center value of $`f_{_{K^+}}`$, $`226.0`$MeV and the fifth column corresponds $`f_{_{K^+}}=228.5`$MeV. Then from table 1 we have
$`m_s`$ $`=`$ $`160\pm 15.5\mathrm{MeV},m_d=7.9\pm 2.7\mathrm{MeV},m_u=4.1\pm 1.5\mathrm{MeV},`$ (122)
$`{\displaystyle \frac{m_s}{m_d}}`$ $`=`$ $`20.2\pm 3.0,{\displaystyle \frac{m_u}{m_d}}=0.5\pm 0.09.`$ (123)
Here the large errors are from the uncertainty of $`f_{_{K^+}}`$. From table 1 we also have $`(m_{\pi ^+}m_{\pi ^0})_{\mathrm{QCD}}=0.25`$MeV(in this paper the contribution from $`\pi _3\eta ^{}`$ mixing is neglected). This result allows that the electromagnetic mass splitting of pion is $`(m_{\pi ^+}m_{\pi ^0})_{\mathrm{e}.\mathrm{m}.}=4.8`$MeV.
## VI Discussion and Summary
In this paper we study information on light quark masses at energy scale $`\mu =m_\rho `$ in the framework of chiral constituent quark model. The analysis of quark masses beyond the next to leading order in the chiral expansion is a challenging subject. The attempt is stoped in ChPT due to the difficulties mentioned in Introduction. In approach of quark model, however, it is not necessary to treat the light current quark masses as the quantities used to construct a perturbative expansion. In addtion, since the light current quark masses can be defined uniquely, the Kaplan-Manohar ambiguity is avoided in this formalism. From eq. (122) $`m_u0`$ is also confirmed in our results. This conclusion supports viewpoint of ref. but disagree with one of ref..
Our results are yielded by a non-perturbative method and contain complete information on light quark masses. Not only quark mass ratios but also individual light quark masses are obtained. The results agree with one obtained by other approachs(e.g., QCD sum rule or ChPT) well.
Finally, it is interesting to expand non-perturtation results (96) up to the next to leading order of the chiral expansion and compare with one of ChPT. If we neglect the mass difference $`m_dm_u`$, up to this order the masses of decay constants of koan and pion read
$`m_\pi ^2`$ $`=`$ $`B_0(m_u+m_d)\{1+{\displaystyle \frac{1}{2}}\mu _\pi {\displaystyle \frac{1}{6}}\mu _{\eta _8}+[{\displaystyle \frac{1}{2}}(3\kappa ^2)+{\displaystyle \frac{3m^2}{\pi ^2F_0^2}}({\displaystyle \frac{m}{B_0}}\kappa g_A{\displaystyle \frac{g_A^2}{2}}{\displaystyle \frac{B_0}{6m}}g_A^2)]{\displaystyle \frac{m_u+m_d}{m}}\},`$ (124)
$`m__K^2`$ $`=`$ $`B_0(\widehat{m}+m_s)\{1+{\displaystyle \frac{1}{3}}\mu _{\eta _8}+[{\displaystyle \frac{1}{2}}(3\kappa ^2)+{\displaystyle \frac{3m^2}{\pi ^2F_0^2}}({\displaystyle \frac{m}{B_0}}\kappa g_A{\displaystyle \frac{g_A^2}{2}}{\displaystyle \frac{B_0}{6m}}g_A^2)]{\displaystyle \frac{\widehat{m}+m_s}{m}}\},`$ (125)
$`f_\pi `$ $`=`$ $`F_0\{1\mu _\pi {\displaystyle \frac{1}{2}}\mu __K+{\displaystyle \frac{3}{2\pi ^2}}g_A^2{\displaystyle \frac{m(m_u+m_d)}{F_0^2}}\},`$ (126)
$`f__K`$ $`=`$ $`F_0\{1{\displaystyle \frac{3}{8}}\mu _\pi {\displaystyle \frac{3}{4}}\mu __K{\displaystyle \frac{3}{8}}\mu _{\eta _8}+{\displaystyle \frac{3}{2\pi ^2}}g_A^2{\displaystyle \frac{m(\widehat{m}+m_s)}{F_0^2}}\}.`$ (127)
Comparing the above equation with one of ChPT, we predict the $`O(p^4)`$ chiral coupling constants, $`L_4,L_5L_6`$ and $`L_8`$, as follow
$`L_4`$ $`=`$ $`L_6=0,L_5={\displaystyle \frac{3m}{32\pi ^2B_0}}g_A^2,`$ (128)
$`L_8`$ $`=`$ $`{\displaystyle \frac{F_0^2}{128B_0m}}(3\kappa ^2)+{\displaystyle \frac{3m}{64\pi ^2B_0}}({\displaystyle \frac{m}{B_0}}\kappa g_A{\displaystyle \frac{g_A^2}{2}}{\displaystyle \frac{B_0}{6m}}g_A^2)+{\displaystyle \frac{L_5}{2}}.`$ (129)
Numerically, inuptting $`\kappa =0.35\pm 0.15`$, $`F_0=0.156`$GeV and $`B_0=1.6\pm 0.4`$GeV, we obtain $`L_5=(1.6\pm 0.3)\times 10^3`$ and $`L_8=(0.7\pm 0.5)\times 10^3`$. These value well agree with one of ChPT, $`L_5=(1.4\pm 0.5)\times 10^3`$ and $`L_8=(0.9\pm 0.3)\times 10^3`$ at energy scale $`\mu =m_\rho `$. This fact can interpret why the light quark mass ratio in eq. (122) are close to results extracted from ChPT by Leutwyler. The above results indicate that the contribution from scalar meson resonance exchange is small. In fact, in hardon spectrum, there is no scalar meson otect or singnlet which belong to composited fields of $`q\overline{q}`$. Thus it is a ad hoc assumption to agrue some low energy coupling constants, such as $`L_5`$ and $`L_8`$, receiving large contribution from scalar meson exchange.
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